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warning/0003/hep-ph0003261.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The increasing precision of high-energy scattering experiments has reached the level that next-to-next-to-leading theoretical results for $`2\mathrm{\hspace{0.17em}2}`$ scattering processes of massless particles are demanded . Several steps have already been taken towards this goal: in Refs. the tensor reduction of two-loop planar boxes and the computation of the relevant master integrals have been given. Reference dealt with the tensor reduction of the pentabox, while other simpler two-loop diagrams obtained by one-loop insertion into one-loop four-point functions have been treated in Ref. .
The final unresolved issue is the tensor reduction of two-loop crossed boxes of Fig. 1 and the evaluation of the corresponding master integrals. One of the master integrals, corresponding to the scalar integral with all powers of propagators equal to unity, was computed in Ref. as an analytic expansion in $`ϵ=(4D)/2`$, where $`D`$ is the space-time dimension.
A first result of this paper is to present a method to reduce tensor integrals of crossed two-loop boxes to two crossed-box master integrals, plus master integrals for simpler-topology diagrams. We choose as first master integral the already evaluated crossed box with all powers of propagators equal to one, and as second master integral the scalar integral where the power of the second propagator is equal to two, all the others being one.
As it is well known , tensor integrals can be related to scalar integrals with higher powers of the propagators in higher dimensions. We make use of recurrence relations obtained by integration-by-parts and Lorentz-invariance identities to connect integrals with different powers of the propagators and to derive a reduction algorithm that expresses the generic scalar integral as a function of the two master crossed boxes, plus simpler sub-topologies. We explicitly give the equations that connect the two master integrals in different dimensions (dimensional shift).
Quite in general, it has been shown that the master integrals for any Feynman graph topology satisfy a system of first-order differential equations on any of the Mandelstam variables on which they depend . We derive a system of two coupled differential equations for the crossed-box master integrals by two independent methods: first using the raising and lowering operators and second by taking the on-shell limit of the differential equations for the crossed box with one off-shell leg .
Inserting the analytic expression of the first master crossed box into the differential equation for the first master integral, we obtain an algebraic equation for the second master crossed box, that can be solved to give the analytic expansion in $`ϵ`$ of the second master integral. The differential equation for the second master integral provides an overall check of the calculation.
Our paper is organized as follows: in Section 2 we introduce the notation we are going to use. In Section 3 we briefly review how tensor integrals are related to scalar ones, and we give the reduction algorithm to express these integrals as a combination of the two master crossed boxes, plus simpler master integrals for topologies with fewer propagators. In Sections 4 and 5 we derive the differential equations that the two master integrals satisfy in arbitrary dimensions $`D`$, using the two different methods. The analytic expansion of the second master integral is given in Section 6. For completeness, in Appendix A we present the algorithm for the tensor reduction and for the dimensional shift of the crossed triangle, since this is one of the non-trivial sub-topologies produced by the reduction procedure. Finally we conclude with Section 7.
## 2 Notation
We denote the generic two-loop tensor crossed (or non-planar) four-point function in $`D`$ dimensions of Fig. 1 with seven propagators $`A_i`$ raised to arbitrary powers $`\nu _i`$ as
$`\mathrm{Xbox}^D(\nu _1,\nu _2,\nu _3,\nu _4,\nu _5,\nu _6,\nu _7;s,t)[1;k^\mu ;l^\mu ;k^\mu k^\nu ;k^\mu l^\nu ;\mathrm{}]`$
$`={\displaystyle \frac{d^Dk}{i\pi ^{D/2}}\frac{d^Dl}{i\pi ^{D/2}}\frac{[1;k^\mu ;l^\mu ;k^\mu k^\nu ;k^\mu l^\nu ;\mathrm{}]}{A_1^{\nu _1}A_2^{\nu _2}A_3^{\nu _3}A_4^{\nu _4}A_5^{\nu _5}A_6^{\nu _6}A_7^{\nu _7}}},`$ (2.1)
where the propagators are
$`A_1`$ $`=`$ $`(k+l+p_{34})^2+i0,`$
$`A_2`$ $`=`$ $`(k+l+p_{134})^2+i0,`$
$`A_3`$ $`=`$ $`(k+l)^2+i0,`$
$`A_4`$ $`=`$ $`l^2+i0,`$ (2.2)
$`A_5`$ $`=`$ $`(l+p_3)^2+i0,`$
$`A_6`$ $`=`$ $`k^2+i0,`$
$`A_7`$ $`=`$ $`(k+p_4)^2+i0.`$
The external momenta $`p_j`$ are in-going and light-like, $`p_j^2=0`$, $`j=1\mathrm{}4`$, so that the only momentum scales are the usual Mandelstam variables $`s=(p_1+p_2)^2`$ and $`t=(p_2+p_3)^2`$, together with $`u=st`$. For ease of notation, we define $`p_{ij}=p_i+p_j`$ and $`p_{ijk}=p_i+p_j+p_k`$. In the square brackets we keep trace of the tensor structure that may be present in the numerator of Eq. (2). When the numerator is unity we have the scalar integral
$$\mathrm{Xbox}^D(\nu _1,\nu _2,\nu _3,\nu _4,\nu _5,\nu _6,\nu _7;s,t)[1]\mathrm{Xbox}^D(\nu _1,\nu _2,\nu _3,\nu _4,\nu _5,\nu _6,\nu _7;s,t).$$
(2.3)
## 3 Tensor reduction using raising and lowering operators
Tensor integrals can be related to combinations of scalar integrals with higher powers of propagators and/or different values of $`D`$ . In fact, introducing the Schwinger parameters $`x_i`$, we can write the integrand of Eq. (2) in the form
$$\frac{1}{A_1^{\nu _1}\mathrm{}A_7^{\nu _7}}=𝒟x\mathrm{exp}\left(\underset{i=1}{\overset{7}{}}x_iA_i\right),$$
(3.1)
where
$`{\displaystyle 𝒟x}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{7}{}}}{\displaystyle \frac{(1)^{\nu _i}}{\mathrm{\Gamma }(\nu _i)}}{\displaystyle _0^{\mathrm{}}}𝑑x_ix_i^{\nu _i1},`$ (3.2)
$`{\displaystyle \underset{i=1}{\overset{7}{}}}x_iA_i`$ $`=`$ $`ak^2+bl^2+2ckl+2dk+2el+f,`$ (3.3)
and
$`a`$ $`=`$ $`x_1+x_2+x_3+x_6+x_7`$
$`b`$ $`=`$ $`x_1+x_2+x_3+x_5+x_4`$
$`c`$ $`=`$ $`x_1+x_2+x_3`$
$`d^\mu `$ $`=`$ $`x_1p_{34}^\mu +x_2p_{134}^\mu +x_7p_4^\mu `$
$`e^\mu `$ $`=`$ $`x_1p_{34}^\mu +x_2p_{134}^\mu +x_5p_3^\mu `$
$`f`$ $`=`$ $`x_1s.`$ (3.4)
We can diagonalize the exponent with the change of variables
$`k^\mu `$ $``$ $`\left(K{\displaystyle \frac{cL}{a}}+𝒳\right)^\mu ,`$ (3.5)
$`l^\mu `$ $``$ $`\left(L+𝒴\right)^\mu ,`$ (3.6)
where
$`𝒳^\mu `$ $`=`$ $`{\displaystyle \frac{1}{𝒫}}\{[(x_2+x_1)(x_7+x_5+x_4)+(x_5+x_4+x_3)x_7]p_4^\mu +[x_3x_5x_4(x_2+x_1)]p_3^\mu `$ (3.7)
$`x_2(x_5+x_4)p_1^\mu \},`$
$`𝒴^\mu `$ $`=`$ $`{\displaystyle \frac{1}{𝒫}}\{[x_3x_7x_6(x_2+x_1)]p_4^\mu [(x_2+x_1)(x_7+x_6+x_5)+(x_7+x_6+x_3)x_5)]p_3^\mu `$ (3.8)
$`x_2(x_7+x_6)p_1^\mu \},`$
and
$$𝒫=\left(x_7+x_6+x_5+x_4\right)\left(x_3+x_2+x_1\right)+\left(x_5+x_4\right)\left(x_7+x_6\right).$$
(3.9)
The generic tensor integral (dropping all the dependences on $`\nu _i`$ and on the external scales) becomes
$`\mathrm{Xbox}^D[k^{\mu _1}\mathrm{}k^{\mu _m}l^{\nu _1}\mathrm{}l^{\nu _n}]={\displaystyle 𝒟x\frac{d^DK}{i\pi ^{D/2}}\frac{d^DL}{i\pi ^{D/2}}}`$
$`\times \left(K{\displaystyle \frac{cL}{a}}+𝒳\right)^{\mu _1}\mathrm{}\left(K{\displaystyle \frac{cL}{a}}+𝒳\right)^{\mu _m}\left(L+𝒴\right)^{\nu _1}\mathrm{}\left(L+𝒴\right)^{\nu _n}`$
$`\times \mathrm{exp}\left(aK^2+{\displaystyle \frac{𝒫}{a}}L^2+{\displaystyle \frac{𝒬}{𝒫}}\right),`$ (3.10)
where
$$𝒬=x_2(x_5x_6x_4x_7)t+\left[x_1x_3(x_7+x_6+x_5+x_4)+x_3x_5x_7x_2x_4x_7+x_1x_4x_6\right]s.$$
(3.11)
The integration over the loop momenta $`K`$ and $`L`$ is now straightforward. For example, the scalar integral and the two one-index tensor integrals are given by
$`\mathrm{Xbox}^D(\nu _1,\nu _2,\nu _3,\nu _4,\nu _5,\nu _6,\nu _7;s,t)`$ $`=`$ $`{\displaystyle 𝒟x},`$ (3.12)
$`\mathrm{Xbox}^D(\nu _1,\nu _2,\nu _3,\nu _4,\nu _5,\nu _6,\nu _7;s,t)[k^\mu ]`$ $`=`$ $`{\displaystyle 𝒟x𝒳^\mu },`$ (3.13)
$`\mathrm{Xbox}^D(\nu _1,\nu _2,\nu _3,\nu _4,\nu _5,\nu _6,\nu _7;s,t)[l^\mu ]`$ $`=`$ $`{\displaystyle 𝒟x𝒴^\mu },`$ (3.14)
where
$$=\frac{1}{𝒫^{D/2}}\mathrm{exp}\left(\frac{𝒬}{𝒫}\right).$$
(3.15)
Recalling the definition (3.2), we see that we can absorb the factors $`x_i`$ of $`𝒳^\mu `$ and $`𝒴^\mu `$ into $`𝒟x`$, increasing the power of the $`i`$-th propagator by one
$$\frac{(1)^{\nu _i}x_i^{\nu _i1}}{\mathrm{\Gamma }(\nu _i)}x_i\nu _i\frac{(1)^{\nu _i+1}x_i^{\nu _i}}{\mathrm{\Gamma }(\nu _i+1)}\nu _i𝐢^\mathbf{+},$$
(3.16)
where
$$𝐢^\mathbf{\pm }\mathrm{Xbox}^D(\mathrm{},\nu _i,\mathrm{})=\mathrm{Xbox}^D(\mathrm{},\nu _i\pm 1,\mathrm{}),$$
(3.17)
while the factor $`𝒫`$ can be absorbed into $``$ (see Eq. (3.15))
$$\frac{1}{𝒫^{D/2}}\frac{1}{𝒫}\frac{1}{𝒫^{(D+2)/2}},$$
(3.18)
so that $`1/𝒫`$ acts as a dimension increaser
$$\frac{1}{𝒫}𝐝^\mathbf{+},𝒫𝐝^{\mathbf{}},$$
(3.19)
where
$$𝐝^\mathbf{\pm }\mathrm{Xbox}^D=\mathrm{Xbox}^{D\pm 2}.$$
(3.20)
For example, using the expression for $`𝒳^\mu `$ of Eq. (3.7), we can write
$`\mathrm{Xbox}^D[k^\mu ]`$ $`=`$ $`\{[\nu _3\mathrm{𝟑}^\mathbf{+}\nu _5\mathrm{𝟓}^\mathbf{+}\nu _4\mathrm{𝟒}^\mathbf{+}(\nu _2\mathrm{𝟐}^\mathbf{+}+\nu _1\mathrm{𝟏}^\mathbf{+})]p_3^\mu \nu _2\mathrm{𝟐}^\mathbf{+}(\nu _5\mathrm{𝟓}^\mathbf{+}+\nu _4\mathrm{𝟒}^\mathbf{+})p_1^\mu `$ (3.21)
$`[(\nu _2\mathrm{𝟐}^\mathbf{+}+\nu _1\mathrm{𝟏}^\mathbf{+})(\nu _7\mathrm{𝟕}^\mathbf{+}+\nu _5\mathrm{𝟓}^\mathbf{+}+\nu _4\mathrm{𝟒}^\mathbf{+})`$
$`+(\nu _5\mathrm{𝟓}^\mathbf{+}+\nu _4\mathrm{𝟒}^\mathbf{+}+\nu _3\mathrm{𝟑}^\mathbf{+})\nu _7\mathrm{𝟕}^\mathbf{+}]p_4^\mu \}\mathrm{Xbox}^{D+2}[1].`$
The task to compute tensor integrals has then been moved to the computation of scalar integrals with higher powers of the propagators in higher dimensions.
In the next sections, we will follow the procedure already used in Ref. :
* we first reduce the powers of the propagators of the generic scalar diagram and we express it in terms of a finite set (basis) of master diagrams;
* then we relate the master integrals in higher dimensions to the master integrals in $`D`$ dimensions (dimensional shift).
### 3.1 The scalar crossed-box reduction
The strategy to reduce the generic scalar integral to a linear combination of master ones is based on identities that relate scalar integrals with different powers of propagators. Some of these identities can be obtained using the integration-by-parts method . Since there are three external independent momenta and two loop momenta, we can build ten identities, imposing that
$$\frac{d^Dk}{i\pi ^{D/2}}\frac{d^Dl}{i\pi ^{D/2}}\frac{}{a^\mu }\left[b^\mu f(k,l,p_i)\right]=0,$$
(3.22)
where
$`a^\mu `$ $`=`$ $`k^\mu ,l^\mu `$ (3.23)
$`b^\mu `$ $`=`$ $`k^\mu ,l^\mu ,p_1^\mu ,p_3^\mu ,p_4^\mu `$ (3.24)
$`f(k,l,p_i)`$ $`=`$ $`{\displaystyle \frac{1}{A_1^{\nu _1}A_2^{\nu _2}A_3^{\nu _3}A_4^{\nu _4}A_5^{\nu _5}A_6^{\nu _6}A_7^{\nu _7}}}.`$ (3.25)
Not all the scalar products that appear in the application of the integration-by-parts identities can be written in terms of combinations of propagators: with the two loop momenta and with the three linearly independent external ones, we can form nine scalar products involving $`k`$ and $`l`$. Since we have seven linearly independent propagators, we are left with two irreducible scalar products in the numerator, that we choose to be $`\left(lp_1\right)`$ and $`\left(lp_4\right)`$.
Three more identities among the scalar integrals can be derived if we exploit the Lorentz invariance of the Feynman diagrams . In fact, since the Feynman integral is a function only of scalar products of the external momenta, it is invariant under the (infinitesimal) rotation
$$p_i^\mu =\mathrm{\Lambda }_\nu ^\mu p_i^\nu ,\mathrm{\Lambda }_{\mu \nu }=g_{\mu \nu }+\delta ϵ_{\mu \nu },ϵ_{\mu \nu }=ϵ_{\nu \mu }.$$
(3.26)
We can then write
$$\frac{d^Dk}{i\pi ^{D/2}}\frac{d^Dl}{i\pi ^{D/2}}f(k,l,p_i)=\frac{d^Dk}{i\pi ^{D/2}}\frac{d^Dl}{i\pi ^{D/2}}f(k,l,p_i^{}).$$
(3.27)
Expanding in a Taylor series in $`\delta `$ the right-hand side of Eq. (3.27), we obtain
$$\frac{d^Dk}{i\pi ^{D/2}}\frac{d^Dl}{i\pi ^{D/2}}\underset{j=1}{\overset{3}{}}\frac{f(k,l,p_i)}{p_j^\mu }ϵ_\nu ^\mu p_j^\nu =0.$$
(3.28)
Using the three independent external momenta, we can build three independent antisymmetric tensors
$`ϵ_1^{\mu \nu }`$ $`=`$ $`p_1^\mu p_2^\nu p_2^\mu p_1^\nu ,`$
$`ϵ_2^{\mu \nu }`$ $`=`$ $`p_1^\mu p_3^\nu p_3^\mu p_1^\nu ,`$
$`ϵ_3^{\mu \nu }`$ $`=`$ $`p_2^\mu p_3^\nu p_3^\mu p_2^\nu ,`$
that, once inserted into Eq. (3.28), give rise to three more identities.
Taking linear combinations of these thirteen equations, we can build the following system
$`s\nu _1\mathrm{𝟏}^\mathbf{+}\nu _7\mathrm{𝟕}^\mathbf{+}\mathrm{𝟔}^{\mathbf{}}\nu _5\mathrm{𝟓}^\mathbf{+}\mathrm{𝟒}^{\mathbf{}}\left(\nu _2\mathrm{𝟐}^\mathbf{+}+\nu _1\mathrm{𝟏}^\mathbf{+}\right)\mathrm{𝟑}^{\mathbf{}}\nu _{1257}2\nu _{346}+2D=0`$ (3.29)
$`s\nu _3\mathrm{𝟑}^\mathbf{+}\nu _6\mathrm{𝟔}^\mathbf{+}\mathrm{𝟕}^{\mathbf{}}\nu _4\mathrm{𝟒}^\mathbf{+}\mathrm{𝟓}^{\mathbf{}}\left(\nu _3\mathrm{𝟑}^\mathbf{+}+\nu _2\mathrm{𝟐}^\mathbf{+}\right)\mathrm{𝟏}^{\mathbf{}}\nu _{2346}2\nu _{157}+2D=0`$ (3.30)
$`2\left(lp_4\right)\nu _4\mathrm{𝟒}^\mathbf{+}\left(\nu _6\mathrm{𝟔}^\mathbf{+}+\nu _5\mathrm{𝟓}^\mathbf{+}\right)\mathrm{𝟕}^{\mathbf{}}+\nu _5\mathrm{𝟓}^\mathbf{+}\mathrm{𝟏}^{\mathbf{}}+\nu _4\mathrm{𝟒}^\mathbf{+}\left(\mathrm{𝟑}^{\mathbf{}}\mathrm{𝟔}^{\mathbf{}}\right)`$
$`\nu _{456}2\nu _7+D=0`$ (3.31)
$`2\left(lp_4\right)\nu _5\mathrm{𝟓}^\mathbf{+}+\nu _7\mathrm{𝟕}^\mathbf{+}\mathrm{𝟔}^{\mathbf{}}+\nu _5\mathrm{𝟓}^\mathbf{+}\left(\mathrm{𝟕}^{\mathbf{}}\mathrm{𝟏}^{\mathbf{}}+s\right)+\nu _4\mathrm{𝟒}^\mathbf{+}\left(\mathrm{𝟔}^{\mathbf{}}\mathrm{𝟑}^{\mathbf{}}\right)`$
$`+\nu _{457}+2\nu _6D=0`$ (3.32)
$`2\left(lp_4\right)\nu _6\mathrm{𝟔}^\mathbf{+}+\nu _7\mathrm{𝟕}^\mathbf{+}\left(\mathrm{𝟓}^{\mathbf{}}\mathrm{𝟏}^{\mathbf{}}\right)+\nu _6\mathrm{𝟔}^\mathbf{+}\left(\mathrm{𝟕}^{\mathbf{}}+\mathrm{𝟓}^{\mathbf{}}\mathrm{𝟏}^{\mathbf{}}+s\right)+\nu _4\mathrm{𝟒}^\mathbf{+}\mathrm{𝟓}^{\mathbf{}}`$
$`+\nu _{47}+2\nu _5D=0`$ (3.33)
$`2\left(lp_4\right)\nu _7\mathrm{𝟕}^\mathbf{+}+\nu _6\mathrm{𝟔}^\mathbf{+}\left(\mathrm{𝟑}^{\mathbf{}}\mathrm{𝟒}^{\mathbf{}}\right)\nu _7\mathrm{𝟕}^\mathbf{+}\left(\mathrm{𝟔}^{\mathbf{}}+\mathrm{𝟒}^{\mathbf{}}\mathrm{𝟑}^{\mathbf{}}\right)\nu _5\mathrm{𝟓}^\mathbf{+}\mathrm{𝟒}^{\mathbf{}}`$
$`\nu _{56}2\nu _4+D=0`$ (3.34)
$`2s\left(lp_4\right)\nu _2\mathrm{𝟐}^\mathbf{+}+4s\left(lp_4\right)\nu _3\mathrm{𝟑}^\mathbf{+}+\nu _2\mathrm{𝟐}^\mathbf{+}\left[\left(t+s\right)\left(s+\mathrm{𝟏}^{\mathbf{}}\right)t\mathrm{𝟑}^{\mathbf{}}\right]`$
$`+s\left(2\nu _6\mathrm{𝟔}^\mathbf{+}+2\nu _3\mathrm{𝟑}^\mathbf{+}+\nu _2\mathrm{𝟐}^\mathbf{+}\right)\mathrm{𝟕}^{\mathbf{}}s\left(2\nu _3\mathrm{𝟑}^\mathbf{+}+\nu _2\mathrm{𝟐}^\mathbf{+}\right)\mathrm{𝟔}^{\mathbf{}}`$
$`\left(2D2\nu _{13457}\nu _2\right)s=0`$ (3.35)
$`2s\left(lp_4\right)\nu _3\mathrm{𝟑}^\mathbf{+}2s\left(lp_4\right)\nu _1\mathrm{𝟏}^\mathbf{+}+(t+s)\nu _2\mathrm{𝟐}^\mathbf{+}\left(\mathrm{𝟏}^{\mathbf{}}\mathrm{𝟑}^{\mathbf{}}\right)+s\left(\nu _1\mathrm{𝟏}^\mathbf{+}\nu _3\mathrm{𝟑}^\mathbf{+}\right)\mathrm{𝟔}^{\mathbf{}}`$
$`+s\left(\nu _6\mathrm{𝟔}^\mathbf{+}+\nu _3\mathrm{𝟑}^\mathbf{+}\nu _1\mathrm{𝟏}^\mathbf{+}\right)\mathrm{𝟕}^{\mathbf{}}s\nu _5\mathrm{𝟓}^\mathbf{+}\mathrm{𝟒}^{\mathbf{}}s\nu _1\mathrm{𝟏}^\mathbf{+}\mathrm{𝟑}^{\mathbf{}}`$
$`\left(\nu _6\nu _{15}\right)s=0`$ (3.36)
$`2\left(lp_4\right)\nu _3\mathrm{𝟑}^\mathbf{+}2\left(lp_1\right)\nu _2\mathrm{𝟐}^\mathbf{+}+\left(t+s\right)\nu _2\mathrm{𝟐}^\mathbf{+}+\left(\nu _6\mathrm{𝟔}^\mathbf{+}+\nu _3\mathrm{𝟑}^\mathbf{+}+\nu _2\mathrm{𝟐}^\mathbf{+}+\nu _1\mathrm{𝟏}^\mathbf{+}\right)\mathrm{𝟕}^{\mathbf{}}`$
$`\left(\nu _2\mathrm{𝟐}^\mathbf{+}+\nu _1\mathrm{𝟏}^\mathbf{+}\right)\mathrm{𝟓}^{\mathbf{}}\nu _3\mathrm{𝟑}^\mathbf{+}\mathrm{𝟒}^{\mathbf{}}+\nu _{1236}+2\nu _7D=0`$ (3.37)
$`2s\left(lp_4\right)\nu _3\mathrm{𝟑}^\mathbf{+}2s\left(lp_1\right)\nu _7\mathrm{𝟕}^\mathbf{+}+\nu _7\mathrm{𝟕}^\mathbf{+}\left[\left(t+s\right)\left(s\mathrm{𝟏}^{\mathbf{}}\right)+t\mathrm{𝟓}^{\mathbf{}}+s\mathrm{𝟐}^{\mathbf{}}\right]`$
$`+\nu _6\mathrm{𝟔}^\mathbf{+}\left(t\mathrm{𝟒}^{\mathbf{}}t\mathrm{𝟑}^{\mathbf{}}+s\mathrm{𝟕}^{\mathbf{}}\right)+s\nu _3\mathrm{𝟑}^\mathbf{+}\left(\mathrm{𝟕}^{\mathbf{}}\mathrm{𝟔}^{\mathbf{}}\right)+(t+s)\nu _2\mathrm{𝟐}^\mathbf{+}\mathrm{𝟏}^{\mathbf{}}+s\nu _1\mathrm{𝟏}^\mathbf{+}\mathrm{𝟐}^{\mathbf{}}`$
$`+s\left(\nu _{126}+2\nu _{3457}2D\right)t\left(D\nu _{67}+\nu _22\nu _{45}\right)=0`$ (3.38)
$`2s\left(lp_4\right)\nu _3\mathrm{𝟑}^\mathbf{+}2s\left(lp_1\right)\nu _5\mathrm{𝟓}^\mathbf{+}+\left(t+s\right)\nu _5\mathrm{𝟓}^\mathbf{+}\left(\mathrm{𝟕}^{\mathbf{}}\mathrm{𝟏}^{\mathbf{}}+s\right)+(t+s)\nu _2\mathrm{𝟐}^\mathbf{+}\mathrm{𝟏}^{\mathbf{}}`$
$`+(t+s)\nu _4\mathrm{𝟒}^\mathbf{+}\left(\mathrm{𝟔}^{\mathbf{}}\mathrm{𝟑}^{\mathbf{}}\right)+s\left(\nu _6\mathrm{𝟔}^\mathbf{+}+\nu _3\mathrm{𝟑}^\mathbf{+}\right)\mathrm{𝟕}^{\mathbf{}}s\nu _3\mathrm{𝟑}^\mathbf{+}\mathrm{𝟔}^{\mathbf{}}s\nu _1\mathrm{𝟏}^\mathbf{+}\mathrm{𝟐}^{\mathbf{}}`$
$`+s\left(\nu _{1456}\nu _2+2\nu _7D\right)+t\left(2\nu _{67}+\nu _{45}\nu _2D\right)=0`$ (3.39)
$`2s\left(lp_4\right)\nu _3\mathrm{𝟑}^\mathbf{+}+2s\left(lp_1\right)\nu _6\mathrm{𝟔}^\mathbf{+}+\nu _6\mathrm{𝟔}^\mathbf{+}\left[t\left(\mathrm{𝟒}^{\mathbf{}}\mathrm{𝟑}^{\mathbf{}}\right)s\left(\mathrm{𝟐}^{\mathbf{}}\mathrm{𝟏}^{\mathbf{}}+s\right)\right]s\nu _4\mathrm{𝟒}^\mathbf{+}\mathrm{𝟓}^{\mathbf{}}`$
$`+s\nu _3\mathrm{𝟑}^\mathbf{+}\left(\mathrm{𝟕}^{\mathbf{}}\mathrm{𝟔}^{\mathbf{}}\mathrm{𝟐}^{\mathbf{}}\right)+t\nu _7\mathrm{𝟕}^\mathbf{+}\left(\mathrm{𝟓}^{\mathbf{}}\mathrm{𝟏}^{\mathbf{}}\right)+(t+s)\nu _2\mathrm{𝟐}^\mathbf{+}\mathrm{𝟏}^{\mathbf{}}`$
$`+s\left(\nu _{34}\nu _2\right)+t\left(\nu _{67}+2\nu _{45}\nu _2D\right)=0`$ (3.40)
$`2s\left(lp_4\right)\nu _3\mathrm{𝟑}^\mathbf{+}+2s\left(lp_1\right)\nu _4\mathrm{𝟒}^\mathbf{+}+\nu _4\mathrm{𝟒}^\mathbf{+}\left[(t+s)\left(\mathrm{𝟔}^{\mathbf{}}\mathrm{𝟑}^{\mathbf{}}\right)+s\mathrm{𝟓}^{\mathbf{}}\right]+2s\nu _6\mathrm{𝟔}^\mathbf{+}\mathrm{𝟕}^{\mathbf{}}`$
$`+(t+s)\nu _5\mathrm{𝟓}^\mathbf{+}\left(\mathrm{𝟕}^{\mathbf{}}\mathrm{𝟏}^{\mathbf{}}\right)+(t+s)\nu _2\mathrm{𝟐}^\mathbf{+}\mathrm{𝟏}^{\mathbf{}}+s\nu _3\mathrm{𝟑}^\mathbf{+}\left(\mathrm{𝟕}^{\mathbf{}}\mathrm{𝟔}^{\mathbf{}}+\mathrm{𝟐}^{\mathbf{}}\right)`$
$`+s\left(\nu _{23}+2\nu _{146}+3\nu _5+4\nu _73D\right)+t\left(\nu _{45}\nu _2+2\nu _{67}D\right)=0,`$ (3.41)
where we have introduced the shorthand $`\nu _{ij}=\nu _i+\nu _j`$, $`\nu _{ijk}=\nu _i+\nu _j+\nu _k`$, etc. Each equation acts on the integrand of the generic crossed box before any loop integration has taken place.
Equations (3.29) and (3.30) of the system, being independent of the two irreducible scalar products, need no further manipulation, and can be rewritten in the form
$`s\nu _1\mathrm{𝟏}^\mathbf{+}`$ $`=`$ $`\nu _7\mathrm{𝟕}^\mathbf{+}\mathrm{𝟔}^{\mathbf{}}+\nu _5\mathrm{𝟓}^\mathbf{+}\mathrm{𝟒}^{\mathbf{}}+\left(\nu _2\mathrm{𝟐}^\mathbf{+}+\nu _1\mathrm{𝟏}^\mathbf{+}\right)\mathrm{𝟑}^{\mathbf{}}+\nu _{1257}+2\nu _{346}2D,`$ (3.42)
$`s\nu _3\mathrm{𝟑}^\mathbf{+}`$ $`=`$ $`\nu _6\mathrm{𝟔}^\mathbf{+}\mathrm{𝟕}^{\mathbf{}}+\nu _4\mathrm{𝟒}^\mathbf{+}\mathrm{𝟓}^{\mathbf{}}+\left(\nu _3\mathrm{𝟑}^\mathbf{+}+\nu _2\mathrm{𝟐}^\mathbf{+}\right)\mathrm{𝟏}^{\mathbf{}}+\nu _{2346}+2\nu _{157}2D.`$ (3.43)
By repeated application of these two identities, we can reduce $`\nu _1`$ and $`\nu _3`$ to one. During this process, the generic scalar box is expressed as a linear combination of crossed-box diagrams with $`\nu _1=\nu _3=1`$ and diagrams belonging to simpler topologies, that originate when powers of propagators are reduced (pinched) to zero by the decreasing operators. We will deal with the pinched diagrams later, concentrating now on the reduction of the remaining propagators.
In order to use the other equations of the system, we have to eliminate the irreducible scalar products in the numerator.
For example, applying the operator $`\nu _7\mathrm{𝟕}^\mathbf{+}`$ to Eq. (3.1) and $`\nu _4\mathrm{𝟒}^\mathbf{+}`$ to Eq. (3.1), and taking the difference, we get
$`\left(D2\nu _7\nu _{56}2\right)\nu _7\mathrm{𝟕}^\mathbf{+}\left(D\nu _{56}2\nu _42\right)\nu _4\mathrm{𝟒}^\mathbf{+}\left(\nu _7\nu _4\right)\left(\nu _5\mathrm{𝟓}^\mathbf{+}+\nu _6\mathrm{𝟔}^\mathbf{+}\right)`$
$`+\nu _5\mathrm{𝟓}^\mathbf{+}\nu _7\mathrm{𝟕}^\mathbf{+}\mathrm{𝟏}^{\mathbf{}}\nu _4\mathrm{𝟒}^\mathbf{+}\nu _6\mathrm{𝟔}^\mathbf{+}\mathrm{𝟑}^{\mathbf{}}=0.`$ (3.44)
In the same way, we can apply $`\nu _6\mathrm{𝟔}^\mathbf{+}`$ to Eq. (3.1) and $`\nu _5\mathrm{𝟓}^\mathbf{+}`$ to Eq. (3.1) and take the difference, to obtain
$`\left(D\nu _{47}2\nu _52\right)\nu _5\mathrm{𝟓}^\mathbf{+}\left(D\nu _{47}2\nu _62\right)\nu _6\mathrm{𝟔}^\mathbf{+}+\left(\nu _6\nu _5\right)\left(\nu _4\mathrm{𝟒}^\mathbf{+}+\nu _7\mathrm{𝟕}^\mathbf{+}\right)`$
$`+\nu _5\mathrm{𝟓}^\mathbf{+}\nu _7\mathrm{𝟕}^\mathbf{+}\mathrm{𝟏}^{\mathbf{}}\nu _4\mathrm{𝟒}^\mathbf{+}\nu _6\mathrm{𝟔}^\mathbf{+}\mathrm{𝟑}^{\mathbf{}}=0.`$ (3.45)
Combining Eq. (3.1) and (3.1) to eliminate $`\nu _5\mathrm{𝟓}^\mathbf{+}`$, we have
$`\nu _4\mathrm{𝟒}^\mathbf{+}`$ $`=`$ $`{\displaystyle \frac{\left(D2\nu _{57}2\right)}{\left(D2\nu _{45}2\right)}}\nu _7\mathrm{𝟕}^\mathbf{+}2{\displaystyle \frac{\left(\nu _7\nu _4\right)}{\left(D2\nu _{45}2\right)}}\nu _6\mathrm{𝟔}^\mathbf{+}`$ (3.46)
$`+{\displaystyle \frac{1}{\left(D\nu _{4567}2\right)}}\left(\nu _5\mathrm{𝟓}^\mathbf{+}\nu _7\mathrm{𝟕}^\mathbf{+}\mathrm{𝟏}^{\mathbf{}}\nu _4\mathrm{𝟒}^\mathbf{+}\nu _6\mathrm{𝟔}^\mathbf{+}\mathrm{𝟑}^{\mathbf{}}\right),`$
that can be used to reduce $`\nu _4`$ to one, at the expense of increasing $`\nu _6`$ and $`\nu _7`$. If, on the other hand, we eliminate $`\nu _4\mathrm{𝟒}^\mathbf{+}`$, we obtain the symmetric equation that can reduce $`\nu _5`$ to one. At this point, all the powers of the propagators except $`\nu _2`$, $`\nu _6`$ and $`\nu _7`$ have been reduced to one.
In the same spirit we can derive
$`st\nu _2\mathrm{𝟐}^\mathbf{+}\nu _6\mathrm{𝟔}^\mathbf{+}`$ $`=`$ $`\left(2\nu _{1567}+\nu _2+22D\right)s\nu _6\mathrm{𝟔}^\mathbf{+}+\left(\nu _{467}+2\nu _5D\right)s\nu _2\mathrm{𝟐}^\mathbf{+}`$ (3.47)
$`2\left(D\nu _{467}2\nu _5\right)\left(\nu _6\mathrm{𝟔}^\mathbf{+}\mathrm{𝟕}^{\mathbf{}}+\nu _4\mathrm{𝟒}^\mathbf{+}\mathrm{𝟓}^{\mathbf{}}+\nu _3\mathrm{𝟑}^\mathbf{+}\mathrm{𝟏}^{\mathbf{}}+\nu _2\mathrm{𝟐}^\mathbf{+}\mathrm{𝟏}^{\mathbf{}}\right)`$
$`+s\left(2\nu _3\mathrm{𝟑}^\mathbf{+}+\nu _2\mathrm{𝟐}^\mathbf{+}\right)\left[\left(\nu _7\mathrm{𝟕}^\mathbf{+}+\nu _6\mathrm{𝟔}^\mathbf{+}+\nu _4\mathrm{𝟒}^\mathbf{+}\right)\mathrm{𝟓}^{\mathbf{}}\nu _7\mathrm{𝟕}^\mathbf{+}\mathrm{𝟏}^{\mathbf{}}\right]`$
$`+t\nu _2\mathrm{𝟐}^\mathbf{+}\nu _6\mathrm{𝟔}^\mathbf{+}\left(\mathrm{𝟑}^{\mathbf{}}\mathrm{𝟏}^{\mathbf{}}\right)+2s\nu _4\mathrm{𝟒}^\mathbf{+}\nu _6\mathrm{𝟔}^\mathbf{+}\mathrm{𝟓}^{\mathbf{}}`$
$`+2\left(D\nu _{467}2\nu _5\right)\left(2D2\nu _{157}\nu _{2346}\right),`$
that, together with the symmetric one for $`\nu _2\mathrm{𝟐}^\mathbf{+}\nu _7\mathrm{𝟕}^\mathbf{+}`$ and with
$`s\nu _6\mathrm{𝟔}^\mathbf{+}\nu _7\mathrm{𝟕}^\mathbf{+}`$ $`=`$ $`\left(D\nu _{567}2\nu _41\right)\nu _6\mathrm{𝟔}^\mathbf{+}+\left(D\nu _{467}2\nu _51\right)\nu _7\mathrm{𝟕}^\mathbf{+}`$ (3.48)
$`+\nu _6\mathrm{𝟔}^\mathbf{+}\nu _7\mathrm{𝟕}^\mathbf{+}\left(\mathrm{𝟑}^{\mathbf{}}+\mathrm{𝟏}^{\mathbf{}}\mathrm{𝟓}^{\mathbf{}}\mathrm{𝟒}^{\mathbf{}}\right)\nu _4\mathrm{𝟒}^\mathbf{+}\nu _7\mathrm{𝟕}^\mathbf{+}\mathrm{𝟓}^{\mathbf{}}\nu _5\mathrm{𝟓}^\mathbf{+}\nu _6\mathrm{𝟔}^\mathbf{+}\mathrm{𝟒}^{\mathbf{}}`$
$`+\nu _7(\nu _7+1)\mathrm{𝟕}^{\mathbf{+}\mathbf{+}}\left(\mathrm{𝟏}^{\mathbf{}}\mathrm{𝟓}^{\mathbf{}}\right)+\nu _6(\nu _6+1)\mathrm{𝟔}^{\mathbf{+}\mathbf{+}}\left(\mathrm{𝟑}^{\mathbf{}}\mathrm{𝟒}^{\mathbf{}}\right),`$
reduces all powers except one ($`\nu _2`$ or $`\nu _6`$ or $`\nu _7`$) to unity.
We can decrease $`\nu _2`$ at the expense of increasing $`\nu _6`$ and $`\nu _7`$ using
$`\left[\left(\nu _4\nu _7+2\nu _{23}+2D\right)s+\left(\nu _{45}\nu _{67}\right)t\right]\nu _2\mathrm{𝟐}^\mathbf{+}=\left(\nu _5\nu _3\right)s\nu _4\mathrm{𝟒}^\mathbf{+}+\left(\nu _7\nu _3\right)s\nu _6\mathrm{𝟔}^\mathbf{+}`$
$`\left(D2\nu _7\nu _{16}2\right)s\nu _7\mathrm{𝟕}^\mathbf{+}\left(D2\nu _5\nu _{14}2\right)s\nu _5\mathrm{𝟓}^\mathbf{+}`$
$`+(t+s)\nu _2\mathrm{𝟐}^\mathbf{+}\left[\nu _7\mathrm{𝟕}^\mathbf{+}\left(\mathrm{𝟏}^{\mathbf{}}\mathrm{𝟓}^{\mathbf{}}\right)+\nu _4\mathrm{𝟒}^\mathbf{+}\left(\mathrm{𝟔}^{\mathbf{}}\mathrm{𝟑}^{\mathbf{}}\right)\right]`$
$`+t\nu _2\mathrm{𝟐}^\mathbf{+}\left[\nu _5\mathrm{𝟓}^\mathbf{+}\left(\mathrm{𝟕}^{\mathbf{}}\mathrm{𝟏}^{\mathbf{}}\right)+\nu _6\mathrm{𝟔}^\mathbf{+}\left(\mathrm{𝟑}^{\mathbf{}}\mathrm{𝟒}^{\mathbf{}}\right)\right]`$
$`s\nu _1\mathrm{𝟏}^\mathbf{+}\left(\nu _7\mathrm{𝟕}^\mathbf{+}\mathrm{𝟓}^{\mathbf{}}+\nu _5\mathrm{𝟓}^\mathbf{+}\mathrm{𝟕}^{\mathbf{}}\right)`$
$`+s\nu _3\mathrm{𝟑}^\mathbf{+}\left[\left(2\nu _7\mathrm{𝟕}^\mathbf{+}+\nu _4\mathrm{𝟒}^\mathbf{+}\right)\mathrm{𝟔}^{\mathbf{}}+\left(2\nu _5\mathrm{𝟓}^\mathbf{+}+\nu _6\mathrm{𝟔}^\mathbf{+}\right)\mathrm{𝟒}^{\mathbf{}}\right]`$
$`\left(2D2\nu _{57}3\nu _{46}\right)\left[\nu _6\mathrm{𝟔}^\mathbf{+}\mathrm{𝟕}^{\mathbf{}}+\nu _4\mathrm{𝟒}^\mathbf{+}\mathrm{𝟓}^{\mathbf{}}+\left(\nu _3\mathrm{𝟑}^\mathbf{+}+\nu _2\mathrm{𝟐}^\mathbf{+}\right)\mathrm{𝟏}^{\mathbf{}}\right]`$
$`+\left(\nu _{57}2\nu _2\right)\left[\nu _7\mathrm{𝟕}^\mathbf{+}\mathrm{𝟔}^{\mathbf{}}+\nu _5\mathrm{𝟓}^\mathbf{+}\mathrm{𝟒}^{\mathbf{}}+\left(\nu _2\mathrm{𝟐}^\mathbf{+}+\nu _1\mathrm{𝟏}^\mathbf{+}\right)\mathrm{𝟑}^{\mathbf{}}\right]`$
$`+4D^22\left(5\nu _{57}+4\nu _{46}+\nu _3\nu _2+2\nu _1\right)D+\nu _7\left(5\nu _{17}+10\nu _{456}+4\nu _3+\nu _2\right)`$
$`+\nu _6\left(3\nu _6+10\nu _5+6\nu _{14}+3\nu _3\nu _2\right)+\nu _5\left(5\nu _5+10\nu _4+4\nu _3+\nu _2+5\nu _1\right)`$
$`+\nu _4\left(3\nu _4+3\nu _3\nu _2+6\nu _1\right)2\nu _2\left(2\nu _3+\nu _{12}\right).`$ (3.49)
The power of the seventh propagator can be reduced with
$`s(t+s)\nu _7(\nu _7+1)\mathrm{𝟕}^{\mathbf{+}\mathbf{+}}=\sigma \nu _7\mathrm{𝟕}^\mathbf{+}(\nu _71){\displaystyle \frac{\left(D6\right)t+\left(5D2\nu _726\right)s}{D2\nu _76}}\mathrm{𝟔}^\mathbf{+}`$
$`+\rho {\displaystyle \frac{\left(2D\nu _711\right)t+\left(3D2\nu _715\right)s}{D\nu _75}}\mathrm{𝟓}^\mathbf{+}\nu _7\mathrm{𝟕}^\mathbf{+}\mathrm{𝟏}^{\mathbf{}}`$
$`+\rho {\displaystyle \frac{\left(D2\nu _74\right)t+\left(5D\right)s}{D\nu _75}}\mathrm{𝟒}^\mathbf{+}\mathrm{𝟔}^\mathbf{+}\mathrm{𝟑}^{\mathbf{}}+\rho (t+s)\{\nu _7(\nu _7+1)\mathrm{𝟕}^{\mathbf{+}\mathbf{+}}\mathrm{𝟔}^{\mathbf{}}`$
$`+\mathrm{𝟓}^\mathbf{+}\nu _7\mathrm{𝟕}^\mathbf{+}\left(\mathrm{𝟑}^{\mathbf{}}\mathrm{𝟒}^{\mathbf{}}\mathrm{𝟔}^{\mathbf{}}\right)+\mathrm{𝟒}^\mathbf{+}\nu _7\mathrm{𝟕}^\mathbf{+}\left[2\left(\mathrm{𝟑}^{\mathbf{}}\mathrm{𝟔}^{\mathbf{}}\right)+\mathrm{𝟏}^{\mathbf{}}\right]`$
$`\mathrm{𝟐}^\mathbf{+}(\nu _7\mathrm{𝟕}^\mathbf{+}+\mathrm{𝟒}^\mathbf{+})\mathrm{𝟏}^{\mathbf{}}+\mathrm{𝟓}^\mathbf{+}\mathrm{𝟔}^\mathbf{+}\mathrm{𝟑}^{\mathbf{}}2\mathbf{\hspace{0.17em}5}^{\mathbf{+}\mathbf{+}}\mathrm{𝟒}^{\mathbf{}}\}`$
$`2\nu _7{\displaystyle \frac{t+s}{D2\nu _76}}(\mathrm{𝟒}^\mathbf{+}\nu _7\mathrm{𝟕}^\mathbf{+}\mathrm{𝟓}^{\mathbf{}}+\mathrm{𝟓}^\mathbf{+}\mathrm{𝟔}^\mathbf{+}\mathrm{𝟒}^{\mathbf{}})+(t+s)[2\mathbf{\hspace{0.17em}6}^{\mathbf{+}\mathbf{+}}(\mathrm{𝟒}^{\mathbf{}}\mathrm{𝟑}^{\mathbf{}})`$
$`+\nu _7(\nu _7+1)\mathrm{𝟕}^{\mathbf{+}\mathbf{+}}(\mathrm{𝟓}^{\mathbf{}}\mathrm{𝟏}^{\mathbf{}})+\mathrm{𝟔}^\mathbf{+}\nu _7\mathrm{𝟕}^\mathbf{+}(\mathrm{𝟓}^{\mathbf{}}+\mathrm{𝟒}^{\mathbf{}}\mathrm{𝟑}^{\mathbf{}}\mathrm{𝟏}^{\mathbf{}})]`$
$`+\rho s[\mathrm{𝟑}^\mathbf{+}\mathrm{𝟓}^\mathbf{+}(\mathrm{𝟏}^{\mathbf{}}\mathrm{𝟒}^{\mathbf{}}\mathrm{𝟕}^{\mathbf{}})\mathrm{𝟑}^\mathbf{+}\mathrm{𝟔}^\mathbf{+}(\mathrm{𝟒}^{\mathbf{}}+\mathrm{𝟕}^{\mathbf{}})\mathrm{𝟑}^\mathbf{+}\nu _7\mathrm{𝟕}^\mathbf{+}(\mathrm{𝟒}^{\mathbf{}}+\mathrm{𝟐}^{\mathbf{}})`$
$`\mathrm{𝟒}^\mathbf{+}(\nu _7\mathrm{𝟕}^\mathbf{+}\mathrm{𝟐}^{\mathbf{}}+\mathrm{𝟔}^\mathbf{+}\mathrm{𝟕}^{\mathbf{}}+\mathrm{𝟏}^\mathbf{+}\mathrm{𝟐}^{\mathbf{}}+\mathrm{𝟑}^\mathbf{+}\mathrm{𝟕}^{\mathbf{}})]`$
$`+(2D3\nu _77)\rho \left[\mathrm{𝟔}^\mathbf{+}\mathrm{𝟕}^{\mathbf{}}+\mathrm{𝟒}^\mathbf{+}\mathrm{𝟓}^{\mathbf{}}+\left(\mathrm{𝟑}^\mathbf{+}+\mathrm{𝟐}^\mathbf{+}\right)\mathrm{𝟏}^{\mathbf{}}2(D\nu _74)\right],`$ (3.50)
where we have introduced the shorthands
$`\rho `$ $`=`$ $`{\displaystyle \frac{D6}{D2\nu _76}}`$
$`\sigma `$ $`=`$ $`{\displaystyle \frac{\left(5D^28\nu _7D50D+2\nu _7^2+42\nu _7+124\right)s+\left(2D^23\nu _7D21D+18\nu _7+54\right)t}{D2\nu _76}}.`$
Equation (3.1) is not as general as the previous ones since we have set all the powers of the other propagators to unity. In addition, since this equation contains $`\mathrm{𝟕}^{\mathbf{+}\mathbf{+}}`$, we cannot always reduce $`\nu _7`$ to one, but we also have integrals where $`\nu _7=2`$. A similar identity can be obtained by symmetry for $`\mathrm{𝟔}^{\mathbf{+}\mathbf{+}}`$, so that we are left with three integrals: $`\mathrm{Xbox}^D(1,1,1,1,1,1,1;s,t)`$, $`\mathrm{Xbox}^D(1,1,1,1,1,1,2;s,t)`$ and $`\mathrm{Xbox}^D(1,1,1,1,1,2,1;s,t)`$.
The last step is to write the integral with $`\nu _6=2`$ as a combination of the other two. This can be done with the identity that links $`\mathrm{𝟔}^\mathbf{+}`$ with $`\mathrm{𝟕}^\mathbf{+}`$. We derived such an identity, equating the expressions obtained by acting with $`\nu _7\mathrm{𝟕}^\mathbf{+}`$ on $`\nu _2\mathrm{𝟐}^\mathbf{+}\nu _4\mathrm{𝟒}^\mathbf{+}`$ and by acting with $`\nu _2\mathrm{𝟐}^\mathbf{+}`$ on $`\nu _4\mathrm{𝟒}^\mathbf{+}\nu _7\mathrm{𝟕}^\mathbf{+}`$
$`(D6)(D5){\displaystyle \frac{t}{t+s}}\mathrm{𝟔}^\mathbf{+}=(D6)(D5)\mathrm{𝟕}^\mathbf{+}4{\displaystyle \frac{(D5)^3}{t+s}}`$
$`\left(\mathrm{𝟓}^\mathbf{+}\mathrm{𝟕}^\mathbf{+}+\mathrm{𝟏}^\mathbf{+}\mathrm{𝟕}^\mathbf{+}+\mathrm{𝟏}^\mathbf{+}\mathrm{𝟓}^\mathbf{+}+\mathrm{𝟏}^\mathbf{+}\mathrm{𝟒}^\mathbf{+}\right)\mathrm{𝟔}^{\mathbf{}}\left(\mathrm{𝟑}^\mathbf{+}\mathrm{𝟕}^\mathbf{+}+\mathrm{𝟒}^\mathbf{+}\mathrm{𝟔}^\mathbf{+}+\mathrm{𝟑}^\mathbf{+}\mathrm{𝟔}^\mathbf{+}+\mathrm{𝟑}^\mathbf{+}\mathrm{𝟒}^\mathbf{+}\right)\mathrm{𝟓}^{\mathbf{}}`$
$`+\left(\mathrm{𝟑}^\mathbf{+}\mathrm{𝟒}^\mathbf{+}+\mathrm{𝟒}^\mathbf{+}\mathrm{𝟔}^\mathbf{+}+\mathrm{𝟓}^\mathbf{+}\mathrm{𝟔}^\mathbf{+}+\mathrm{𝟑}^\mathbf{+}\mathrm{𝟔}^\mathbf{+}\right)\mathrm{𝟏}^{\mathbf{}}(D7)\mathrm{𝟒}^\mathbf{+}\mathrm{𝟕}^\mathbf{+}\mathrm{𝟏}^{\mathbf{}}`$
$`{\displaystyle \frac{1}{2}}\left(\mathrm{𝟐}^\mathbf{+}\mathrm{𝟕}^\mathbf{+}\mathrm{𝟔}^{\mathbf{}}+\mathrm{𝟐}^\mathbf{+}\mathrm{𝟓}^\mathbf{+}\mathrm{𝟔}^{\mathbf{}}+\mathrm{𝟒}^\mathbf{+}\mathrm{𝟕}^\mathbf{+}\mathrm{𝟓}^{\mathbf{}}+\mathrm{𝟐}^\mathbf{+}\mathrm{𝟔}^\mathbf{+}\mathrm{𝟓}^{\mathbf{}}\mathrm{𝟏}^\mathbf{+}\mathrm{𝟒}^\mathbf{+}\mathrm{𝟑}^{\mathbf{}}\right)`$
$`{\displaystyle \frac{t+s}{2}}\left[\mathrm{𝟐}^\mathbf{+}\mathrm{𝟒}^\mathbf{+}\mathrm{𝟕}^\mathbf{+}\left(\mathrm{𝟔}^{\mathbf{}}+\mathrm{𝟓}^{\mathbf{}}\right)+2\left(\mathrm{𝟐}^{\mathbf{+}\mathbf{+}}\mathrm{𝟕}^\mathbf{+}\mathrm{𝟐}^\mathbf{+}\mathrm{𝟒}^\mathbf{+}\mathrm{𝟕}^\mathbf{+}+\mathrm{𝟐}^{\mathbf{+}\mathbf{+}}\mathrm{𝟒}^\mathbf{+}\right)\mathrm{𝟏}^{\mathbf{}}\right]`$
$`{\displaystyle \frac{s}{2}}\{\mathrm{𝟐}^\mathbf{+}[\mathrm{𝟑}^\mathbf{+}(\mathrm{𝟔}^\mathbf{+}+\mathrm{𝟓}^\mathbf{+}+\mathrm{𝟒}^\mathbf{+})\mathrm{𝟕}^{\mathbf{}}+\mathrm{𝟒}^\mathbf{+}\mathrm{𝟔}^\mathbf{+}\mathrm{𝟕}^{\mathbf{}}+\mathrm{𝟑}^\mathbf{+}(\mathrm{𝟕}^\mathbf{+}+\mathrm{𝟔}^\mathbf{+}+\mathrm{𝟓}^\mathbf{+})\mathrm{𝟒}^{\mathbf{}}]`$
$`+\mathrm{𝟑}^\mathbf{+}\mathrm{𝟓}^\mathbf{+}(2\mathbf{\hspace{0.17em}7}^\mathbf{+}\mathrm{𝟐}^\mathbf{+})\mathrm{𝟏}^{\mathbf{}}\}(D6){\displaystyle \frac{(t+2s)}{2(t+s)}}[\mathrm{𝟐}^\mathbf{+}(\mathrm{𝟕}^\mathbf{+}\mathrm{𝟒}^{\mathbf{}}+\mathrm{𝟒}^\mathbf{+}\mathrm{𝟕}^{\mathbf{}}+\mathrm{𝟓}^\mathbf{+}\mathrm{𝟒}^{\mathbf{}})`$
$`+2(\mathrm{𝟏}^\mathbf{+}\mathrm{𝟕}^\mathbf{+}+\mathrm{𝟓}^\mathbf{+}\mathrm{𝟕}^\mathbf{+}+\mathrm{𝟏}^\mathbf{+}\mathrm{𝟔}^\mathbf{+}+\mathrm{𝟏}^\mathbf{+}\mathrm{𝟓}^\mathbf{+})\mathrm{𝟒}^{\mathbf{}}2\mathbf{\hspace{0.17em}1}^\mathbf{+}\mathrm{𝟔}^\mathbf{+}\mathrm{𝟑}^{\mathbf{}}]{\displaystyle \frac{2(D5)t+s}{2s}}\mathrm{𝟐}^\mathbf{+}\mathrm{𝟕}^\mathbf{+}\mathrm{𝟓}^{\mathbf{}}`$
$`(D6)\left[\mathrm{𝟒}^\mathbf{+}\mathrm{𝟕}^\mathbf{+}\mathrm{𝟑}^{\mathbf{}}+\left(\mathrm{𝟑}^\mathbf{+}\mathrm{𝟒}^\mathbf{+}+\mathrm{𝟒}^\mathbf{+}\mathrm{𝟔}^\mathbf{+}+\mathrm{𝟑}^\mathbf{+}\mathrm{𝟔}^\mathbf{+}+\mathrm{𝟑}^\mathbf{+}\mathrm{𝟓}^\mathbf{+}\right)\mathrm{𝟕}^{\mathbf{}}\right]`$
$`+(D6){\displaystyle \frac{s}{t+s}}\left[\mathrm{𝟕}^\mathbf{+}\mathrm{𝟑}^\mathbf{+}\mathrm{𝟒}^{\mathbf{}}+\left(\mathrm{𝟕}^\mathbf{+}\mathrm{𝟑}^\mathbf{+}+\mathrm{𝟕}^\mathbf{+}\mathrm{𝟒}^\mathbf{+}+\mathrm{𝟏}^\mathbf{+}\mathrm{𝟒}^\mathbf{+}\right)\mathrm{𝟐}^{\mathbf{}}\right]+{\displaystyle \frac{2D13}{2}}\mathrm{𝟒}^\mathbf{+}\mathrm{𝟕}^\mathbf{+}\mathrm{𝟔}^{\mathbf{}}`$
$`+(D5){\displaystyle \frac{(D5)t+(2D11)s}{s(t+s)}}\left[\left(\mathrm{𝟐}^\mathbf{+}+\mathrm{𝟏}^\mathbf{+}\right)\mathrm{𝟑}^{\mathbf{}}+\mathrm{𝟕}^\mathbf{+}\mathrm{𝟔}^{\mathbf{}}\right]`$
$`+(D5){\displaystyle \frac{t}{t+s}}\left[2\mathbf{\hspace{0.17em}3}^\mathbf{+}\mathrm{𝟕}^\mathbf{+}\mathrm{𝟔}^{\mathbf{}}+\mathrm{𝟑}^\mathbf{+}\mathrm{𝟒}^\mathbf{+}\mathrm{𝟔}^{\mathbf{}}\mathrm{𝟏}^\mathbf{+}\mathrm{𝟕}^\mathbf{+}\mathrm{𝟓}^{\mathbf{}}\mathrm{𝟏}^\mathbf{+}\mathrm{𝟓}^\mathbf{+}\mathrm{𝟕}^{\mathbf{}}\right]+{\displaystyle \frac{t}{2s}}\mathrm{𝟐}^\mathbf{+}\mathrm{𝟓}^\mathbf{+}\mathrm{𝟑}^{\mathbf{}}`$
$`(D5){\displaystyle \frac{(D5)ts}{s(t+s)}}\left[\mathrm{𝟒}^\mathbf{+}\mathrm{𝟓}^{\mathbf{}}+\mathrm{𝟐}^\mathbf{+}\mathrm{𝟏}^{\mathbf{}}+\mathrm{𝟔}^\mathbf{+}\mathrm{𝟕}^{\mathbf{}}+\mathrm{𝟑}^\mathbf{+}\mathrm{𝟏}^{\mathbf{}}\right]{\displaystyle \frac{2D15}{2}}\mathrm{𝟑}^\mathbf{+}\mathrm{𝟕}^\mathbf{+}\mathrm{𝟏}^{\mathbf{}}`$
$`+(D5)\left[\mathrm{𝟏}^{\mathbf{}}\left(\mathrm{𝟑}^\mathbf{+}\mathrm{𝟓}^\mathbf{+}+2\mathbf{\hspace{0.17em}2}^{\mathbf{+}\mathbf{+}}+\mathrm{𝟐}^\mathbf{+}\mathrm{𝟑}^\mathbf{+}\right)\right]+{\displaystyle \frac{t}{2}}\left[\mathrm{𝟐}^\mathbf{+}\left(\mathrm{𝟓}^\mathbf{+}\mathrm{𝟕}^\mathbf{+}\mathrm{𝟏}^{\mathbf{}}+\mathrm{𝟒}^\mathbf{+}\mathrm{𝟔}^\mathbf{+}\mathrm{𝟑}^{\mathbf{}}\right)\right]`$
$`+{\displaystyle \frac{5t2(D7)s}{2(t+s)}}\mathrm{𝟓}^\mathbf{+}\mathrm{𝟕}^\mathbf{+}\mathrm{𝟏}^{\mathbf{}}{\displaystyle \frac{(2D9)t^2(D5)st2(D5)s^2}{2s(t+s)}}\mathrm{𝟐}^\mathbf{+}\mathrm{𝟓}^\mathbf{+}\mathrm{𝟏}^{\mathbf{}}`$
$`(D6){\displaystyle \frac{t+2s}{2s}}\mathrm{𝟐}^\mathbf{+}\mathrm{𝟕}^\mathbf{+}\mathrm{𝟑}^{\mathbf{}}{\displaystyle \frac{2(D5)t^2+(D6)st+2(D6)s^2}{2s(t+s)}}\mathrm{𝟐}^\mathbf{+}\mathrm{𝟔}^\mathbf{+}\mathrm{𝟒}^{\mathbf{}}`$
$`+{\displaystyle \frac{(3D16)t+2(D5)s}{2s}}\mathrm{𝟐}^\mathbf{+}\mathrm{𝟕}^\mathbf{+}\mathrm{𝟏}^{\mathbf{}}+{\displaystyle \frac{(D4)t+2s}{2(t+s)}}\mathrm{𝟐}^\mathbf{+}\mathrm{𝟔}^\mathbf{+}\mathrm{𝟕}^{\mathbf{}}`$
$`+{\displaystyle \frac{2(D5)t^2(D6)st2(D6)s^2}{2s(t+s)}}\mathrm{𝟐}^\mathbf{+}\mathrm{𝟓}^\mathbf{+}\mathrm{𝟕}^{\mathbf{}}(4D21){\displaystyle \frac{t}{2(t+s)}}\mathrm{𝟒}^\mathbf{+}\mathrm{𝟔}^\mathbf{+}\mathrm{𝟑}^{\mathbf{}}`$
$`+{\displaystyle \frac{(D5)t+(D6)s}{t+s}}\mathrm{𝟑}^\mathbf{+}\mathrm{𝟔}^\mathbf{+}\mathrm{𝟒}^{\mathbf{}}{\displaystyle \frac{(D4)t2(D6)s}{2s}}\mathrm{𝟐}^\mathbf{+}\mathrm{𝟒}^\mathbf{+}\mathrm{𝟑}^{\mathbf{}}`$
$`+{\displaystyle \frac{(2D11)t^2+(D7)st+2(D6)s^2}{2s(t+s)}}\mathrm{𝟐}^\mathbf{+}\mathrm{𝟔}^\mathbf{+}\mathrm{𝟑}^{\mathbf{}}+{\displaystyle \frac{t+2s}{2s}}\mathrm{𝟐}^\mathbf{+}\mathrm{𝟔}^\mathbf{+}\mathrm{𝟏}^{\mathbf{}}`$
$`+{\displaystyle \frac{2(D5)t+(D6)s}{t+s}}\mathrm{𝟑}^\mathbf{+}\mathrm{𝟓}^\mathbf{+}\mathrm{𝟒}^{\mathbf{}}+(D5){\displaystyle \frac{(D5)t+(2D11)s}{s(t+s)}}\mathrm{𝟓}^\mathbf{+}\mathrm{𝟒}^{\mathbf{}}`$
$`+{\displaystyle \frac{2D11}{2}}\mathrm{𝟐}^\mathbf{+}\mathrm{𝟒}^\mathbf{+}\mathrm{𝟓}^{\mathbf{}}{\displaystyle \frac{(D6)t2s}{2s}}\mathrm{𝟐}^\mathbf{+}\mathrm{𝟒}^\mathbf{+}\mathrm{𝟏}^{\mathbf{}}+{\displaystyle \frac{2(D5)ts}{2s}}\mathrm{𝟐}^\mathbf{+}\mathrm{𝟒}^\mathbf{+}\mathrm{𝟔}^{\mathbf{}},`$ (3.51)
where we have set all the powers of the propagators to unity.
At the end of this reduction program, we are left with the following two crossed-box integrals: $`\mathrm{Xbox}^D(1,1,1,1,1,1,1;s,t)`$ and $`\mathrm{Xbox}^D(1,1,1,1,1,1,2;s,t)`$, plus simpler diagrams that can always be expressed as a combination of:
* the master crossed triangle of Fig. 2 (a)
$$\mathrm{Xtri}^D\left(s\right)=\mathrm{Xbox}^D(1,0,1,1,1,1,1;s,t),$$
(3.52)
* the master diagonal box of Fig. 2 (b), produced by
$$\mathrm{Dbox}^D(s,t)=\mathrm{Xbox}^D(0,1,1,0,1,1,1;s,t)=\mathrm{Xbox}^D(1,1,0,1,1,1,0;s,t),$$
(3.53)
together with
$$\mathrm{Dbox}^D(s,u)=\mathrm{Xbox}^D(1,1,0,1,0,1,1;s,t)=\mathrm{Xbox}^D(0,1,1,1,1,0,1;s,t),$$
(3.54)
and
$$\mathrm{Dbox}^D(t,u)=\mathrm{Xbox}^D(0,1,0,1,1,1,1;s,t),$$
(3.55)
* the master box with a bubble insertion of Fig. 2 (c), produced by
$$\mathrm{Bbox}^D(s,t)=\mathrm{Xbox}^D(1,1,1,0,1,1,0;s,t),$$
(3.56)
together with
$$\mathrm{Bbox}^D(s,u)=\mathrm{Xbox}^D(1,1,1,1,0,0,1;s,t),$$
(3.57)
* the master triangle with a bubble insertion of Fig. 2 (d), produced by
$$\mathrm{Btri}^D\left(s\right)=\mathrm{Xbox}^D(1,0,1,0,1,1,0;s,t)=\mathrm{Xbox}^D(1,0,1,1,0,0,1;s,t),$$
(3.58)
* the master sunset diagram of Fig. 2 (e), produced by
$$\mathrm{Sset}^D\left(s\right)=\mathrm{Xbox}^D(0,0,1,0,1,0,1;s,t)=\mathrm{Xbox}^D(1,0,0,1,0,1,0;s,t),$$
(3.59)
together with
$$\mathrm{Sset}^D\left(t\right)=\mathrm{Xbox}^D(0,1,0,0,1,1,0;s,t),$$
(3.60)
and
$$\mathrm{Sset}^D\left(u\right)=\mathrm{Xbox}^D(0,1,0,1,0,0,1;s,t),$$
(3.61)
where we explicitly used the symmetry of the crossed box
$$\mathrm{Xbox}^D(\nu _1,\nu _2,\nu _3,\nu _4,\nu _5,\nu _6,\nu _7;s,t)=\mathrm{Xbox}^D(\nu _3,\nu _2,\nu _1,\nu _7,\nu _6,\nu _5,\nu _4;s,t).$$
(3.62)
The reduction and dimensional shift of the generic scalar diagonal box and of the box with a bubble insertion have been treated in Ref. . The tensor reduction and dimensional shift of the crossed triangle is treated in Appendix A. All the diagrams obtained from the pinchings of the other propagators can be easily reduced to the previous ones using the integration-by-parts identities for the crossed box and for the crossed triangle, with the corresponding powers of propagators set to zero.
Instead of keeping $`\mathrm{Xbox}^D(1,1,1,1,1,1,2;s,t)`$ as one of the two members of the basis, we prefer to switch to a more symmetric integral: $`\mathrm{Xbox}^D(1,2,1,1,1,1,1;s,t)`$. The expression for this one, in terms of the other two, can be easily obtained through the application of the reduction formalism outlined above
$`\mathrm{Xbox}^D(1,2,1,1,1,1,1;s,t)=c_1\mathrm{Xbox}^D(1,1,1,1,1,1,1;s,t)`$
$`+c_2\mathrm{Xbox}^D(1,1,1,1,1,1,2;s,t)+c_3\mathrm{Xtri}^D\left(s\right)+c_4\mathrm{Dbox}^D(s,t)`$
$`+c_5\mathrm{Dbox}^D(s,u)+c_6\mathrm{Dbox}^D(t,u)+c_7\mathrm{Bbox}^D(s,t)+c_8\mathrm{Bbox}^D(s,u)`$
$`+c_9\mathrm{Btri}^D\left(s\right)+c_{10}\mathrm{Sset}^D\left(s\right)+c_{11}\mathrm{Sset}^D\left(t\right)+c_{12}\mathrm{Sset}^D\left(u\right),`$ (3.63)
where the coefficients $`c_j`$ are collected in Appendix B. This allows us to define as master integrals
$`\mathrm{Xbox}_1^D(s,t)`$ $`=`$ $`\mathrm{Xbox}^D(1,1,1,1,1,1,1;s,t)`$ (3.64)
$`\mathrm{Xbox}_2^D(s,t)`$ $`=`$ $`\mathrm{Xbox}^D(1,2,1,1,1,1,1;s,t),`$ (3.65)
which are symmetric under the exchange $`tu`$. This will produce more compact results in the rest of the paper.
### 3.2 The dimensional shift for the two master integrals
Exploiting the equivalence between $`𝐝^{\mathbf{}}`$ and $`𝒫`$ of Eq. (3.19), we can rewrite Eq. (3.20) using (3.9)
$`\mathrm{Xbox}^D(\nu _1,\nu _2,\nu _3,\nu _4,\nu _5,\nu _6,\nu _7;s,t)`$
$`=[(\nu _7\mathrm{𝟕}^\mathbf{+}+\nu _6\mathrm{𝟔}^\mathbf{+}+\nu _5\mathrm{𝟓}^\mathbf{+}+\nu _4\mathrm{𝟒}^\mathbf{+})(\nu _3\mathrm{𝟑}^\mathbf{+}+\nu _2\mathrm{𝟐}^\mathbf{+}+\nu _1\mathrm{𝟏}^\mathbf{+})`$
$`+(\nu _5\mathrm{𝟓}^\mathbf{+}+\nu _4\mathrm{𝟒}^\mathbf{+})(\nu _7\mathrm{𝟕}^\mathbf{+}+\nu _6\mathrm{𝟔}^\mathbf{+})\left]\mathrm{Xbox}^{D+2}\right(\nu _1,\nu _2,\nu _3,\nu _4,\nu _5,\nu _6,\nu _7;s,t),`$
and more specifically
$`\mathrm{Xbox}_1^D(s,t)`$ $`=`$ $`[(\mathrm{𝟕}^\mathbf{+}+\mathrm{𝟔}^\mathbf{+}+\mathrm{𝟓}^\mathbf{+}+\mathrm{𝟒}^\mathbf{+})(\mathrm{𝟑}^\mathbf{+}+\mathrm{𝟐}^\mathbf{+}+\mathrm{𝟏}^\mathbf{+})`$ (3.67)
$`+(\mathrm{𝟓}^\mathbf{+}+\mathrm{𝟒}^\mathbf{+})(\mathrm{𝟕}^\mathbf{+}+\mathrm{𝟔}^\mathbf{+})\left]\mathrm{Xbox}_1^{D+2}\right(s,t),`$
$`\mathrm{Xbox}_2^D(s,t)`$ $`=`$ $`[(\mathrm{𝟕}^\mathbf{+}+\mathrm{𝟔}^\mathbf{+}+\mathrm{𝟓}^\mathbf{+}+\mathrm{𝟒}^\mathbf{+})(\mathrm{𝟑}^\mathbf{+}+2\mathbf{\hspace{0.17em}2}^\mathbf{+}+\mathrm{𝟏}^\mathbf{+})`$ (3.68)
$`+(\mathrm{𝟓}^\mathbf{+}+\mathrm{𝟒}^\mathbf{+})(\mathrm{𝟕}^\mathbf{+}+\mathrm{𝟔}^\mathbf{+})\left]\mathrm{Xbox}_2^{D+2}\right(s,t).`$
The right-hand side of the two equations can then be reduced to a linear combination of master integrals in $`D+2`$ dimensions, and the system can then be inverted to give $`\mathrm{Xbox}_1^{D+2}(s,t)`$ and $`\mathrm{Xbox}_2^{D+2}(s,t)`$ as a function of master integrals in $`D`$ dimensions
$`\mathrm{Xbox}_1^{D+2}(s,t)`$ $`=`$ $`A(t,u)+A(u,t),`$ (3.69)
$`\mathrm{Xbox}_2^{D+2}(s,t)`$ $`=`$ $`B(t,u)+B(u,t),`$ (3.70)
where
$`A(t,u)`$ $`=`$ $`a_1\mathrm{Xbox}_1^D(s,t)+a_2\mathrm{Xbox}_2^D(s,t)+a_3\mathrm{Xtri}^D\left(s\right)+a_4\mathrm{Dbox}^D(s,t)`$
$`+a_5\mathrm{Dbox}^D(t,u)+a_6\mathrm{Bbox}^D(s,t)+a_7\mathrm{Btri}^D\left(s\right)`$
$`+a_8\mathrm{Sset}^D\left(s\right)+a_9\mathrm{Sset}^D\left(t\right),`$
$`B(t,u)`$ $`=`$ $`b_1\mathrm{Xbox}_1^D(s,t)+b_2\mathrm{Xbox}_2^D(s,t)+b_3\mathrm{Xtri}^D\left(s\right)+b_4\mathrm{Dbox}^D(s,t)`$ (3.71)
$`+b_5\mathrm{Dbox}^D(t,u)+b_6\mathrm{Bbox}^D(s,t)+b_7\mathrm{Btri}^D\left(s\right)`$
$`+b_8\mathrm{Sset}^D\left(s\right)+b_9\mathrm{Sset}^D\left(t\right).`$
Some of the coefficients $`a_i`$ and $`b_i`$ of the system are collected in Appendix C.
## 4 Differential equations for the two master integrals
In the previous sections we showed how tensor integrals can be expressed in terms of the two master crossed boxes, $`\mathrm{Xbox}_1^D`$ and $`\mathrm{Xbox}_2^D`$ plus simpler master diagrams. The analytic expansion in $`ϵ=(4D)/2`$ of the first master integral was computed in Ref. . We can obtain the analytic form for the second one by writing the derivative of $`\mathrm{Xbox}_1^D`$ with respect to one of the two independent physical scales (that we choose to be $`t`$), as a combination of master integrals, and solving the equation for $`\mathrm{Xbox}_2^D`$.
Moreover we can verify the correctness of both the expressions of $`\mathrm{Xbox}_1^D(s,t)`$ and $`\mathrm{Xbox}_2^D(s,t)`$, by deriving an analogous differential equation for $`\mathrm{Xbox}_2^D`$, and checking that the obtained identity is satisfied.
Using the procedure outlined in Section 3, we differentiate with respect to $`t`$ the two master integrals written in the form of Eq. (3.12). In this way, the only $`t`$-dependence comes from $`𝒬`$ of Eq. (3.15) and the result is
$`{\displaystyle \frac{}{t}}\mathrm{Xbox}_1^D(s,t)`$ $`=`$ $`\mathrm{Xbox}^{D+2}(1,2,1,2,1,1,2;s,t)\mathrm{Xbox}^{D+2}(1,2,1,1,2,2,1;s,t),`$
$`{\displaystyle \frac{}{t}}\mathrm{Xbox}_2^D(s,t)`$ $`=`$ $`2\mathrm{Xbox}^{D+2}(1,3,1,2,1,1,2;s,t)2\mathrm{Xbox}^{D+2}(1,3,1,1,2,2,1;s,t).`$
Applying the reduction formulae of Section 3.1 and the dimensional-shift of Section 3.2 to the right-hand sides of the system, we can rewrite them as a combination of master crossed boxes and pinchings in $`D`$ dimensions.
The final system of differential equations in arbitrary $`D`$ is given by
$`{\displaystyle \frac{}{t}}\mathrm{Xbox}_1^D(s,t)`$ $`=`$ $`{\displaystyle \frac{1}{tu}}[H(t,u)+H(u,t),]`$ (4.1)
$`{\displaystyle \frac{}{t}}\mathrm{Xbox}_2^D(s,t)`$ $`=`$ $`{\displaystyle \frac{1}{tu}}[K(t,u)+K(u,t),]`$ (4.2)
where
$`H(t,u)`$ $`=`$ $`h_1\mathrm{Xbox}_1^D(s,t)+h_2\mathrm{Xbox}_2^D(s,t)+h_3\mathrm{Xtri}^D\left(s\right)+h_4\mathrm{Dbox}^D(s,t)`$ (4.3)
$`+h_5\mathrm{Dbox}^D(t,u)+h_6\mathrm{Bbox}^D(s,t)+h_7\mathrm{Btri}^D\left(s\right)`$
$`+h_8\mathrm{Sset}^D\left(s\right)+h_9\mathrm{Sset}^D\left(t\right),`$
$`K(t,u)`$ $`=`$ $`k_1\mathrm{Xbox}_1^D(s,t)+k_2\mathrm{Xbox}_2^D(s,t)+k_3\mathrm{Xtri}^D\left(s\right)+k_4\mathrm{Dbox}^D(s,t)`$ (4.4)
$`+k_5\mathrm{Dbox}^D(t,u)+k_6\mathrm{Bbox}^D(s,t)+k_7\mathrm{Btri}^D\left(s\right)`$
$`+k_8\mathrm{Sset}^D\left(s\right)+k_9\mathrm{Sset}^D\left(t\right),`$
and the coefficients $`h_i`$ and $`k_i`$ are listed in Appendix D.
From the symmetry $`tu`$ of the two master integrals, we expect the two derivatives to be anti-symmetric with respect to the exchange $`tu`$, and this is a further check of the correctness of the system.
## 5 The off-shell method
An alternative way to derive the system of equations (4.1) and (4.2) is based on the construction of a set of differential equations for the crossed-box integrals with the momentum $`p_1`$ taken off-shell. This is done using the algorithms and computer programs described in , by means of which the differential equations for any massless two-loop four-point function with three light-like and one off-shell leg can be obtained. We then take the on-shell limit $`p_1^20`$.
To derive the differential equations for the diagrams with one leg off-shell, we take $`p_2`$, $`p_3`$ and $`p_4`$, all massless and on-shell, as independent momenta, and form the derivatives with respect to the external invariants $`s_{23}`$, $`s_{24}`$, $`s_{34}`$, where $`s_{ij}=p_{ij}^2`$
$`s_{23}{\displaystyle \frac{}{s_{23}}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(+p_2^\mu {\displaystyle \frac{}{p_2^\mu }}+p_3^\mu {\displaystyle \frac{}{p_3^\mu }}p_4^\mu {\displaystyle \frac{}{p_4^\mu }}\right),`$
$`s_{24}{\displaystyle \frac{}{s_{24}}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(+p_2^\mu {\displaystyle \frac{}{p_2^\mu }}p_3^\mu {\displaystyle \frac{}{p_3^\mu }}+p_4^\mu {\displaystyle \frac{}{p_4^\mu }}\right),`$ (5.1)
$`s_{34}{\displaystyle \frac{}{s_{34}}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(p_2^\mu {\displaystyle \frac{}{p_2^\mu }}+p_3^\mu {\displaystyle \frac{}{p_3^\mu }}+p_4^\mu {\displaystyle \frac{}{p_4^\mu }}\right).`$
The derivatives with respect to $`p_i^\mu `$ act on the propagators in the Feynman integrand of Eq. (2) and give rise to a number of different integrals with increased powers of the propagators. Using integration-by-parts and Lorentz-invariance identities, all these different integrals can be reduced to a small number of master integrals. In particular, any crossed-box topology integral can be reduced to a linear combination of two master crossed-box integrals, which we can choose to be the off-shell versions of $`\mathrm{Xbox}_1^D`$ and $`\mathrm{Xbox}_2^D`$, plus a number of pinched master integrals. Thus, we obtain a system of two coupled linear first-order differential equations for the two master crossed-box integrals.
Since, in the on-shell limit, we are interested in the derivatives at fixed $`p_1^2=0`$, we rewrite the system of equations in terms of the variables $`p_1^2=s_{23}+s_{24}+s_{34}`$, $`s=s_{34}`$, $`t=s_{23}`$, and the corresponding derivative operators:
$`{\displaystyle \frac{}{s}}`$ $`=`$ $`{\displaystyle \frac{}{s_{34}}}{\displaystyle \frac{}{s_{24}}},`$
$`{\displaystyle \frac{}{t}}`$ $`=`$ $`{\displaystyle \frac{}{s_{23}}}{\displaystyle \frac{}{s_{24}}},`$ (5.2)
$`{\displaystyle \frac{}{p_1^2}}`$ $`=`$ $`{\displaystyle \frac{}{s_{24}}}.`$
The next step is to take the on-shell limit $`p_1^20`$. A complication here is the fact that factors of $`1/p_1^2`$ and $`1/(p_1^2)^2`$ appear in some of the coefficients of the pinched master integrals. It is also important to realize that several of the off-shell master integrals with less than seven propagators become reducible in the on-shell limit (an example of such a case will be given below). When these reductions are taken into account, all the terms proportional to $`1/p_1^2`$ and $`1/(p_1^2)^2`$ cancel out, so that the on-shell limit of the system of equations for the two master crossed-box integrals is indeed well defined.
However, because of these factors of $`p_1^2`$ in the denominator, it is not sufficient merely to replace all pinched master integrals by their limit $`p_1^20`$: subleading terms in $`p_1^2`$ also have to be included. The Taylor series around $`p_1^2=0`$ for the pinched diagrams can be obtained in a straightforward manner by using the differential equation in $`p_1^2`$. As an example, we consider the propagator diagram $`\mathrm{Sset}^D\left(p_1^2st\right)`$, which fulfills the homogeneous differential equation
$$\frac{}{p_1^2}\mathrm{Sset}^D\left(p_1^2st\right)=\frac{D3}{p_1^2st}\mathrm{Sset}^D\left(p_1^2st\right).$$
(5.3)
Iterating this equation, we obtain
$`\mathrm{Sset}^D\left(p_1^2st\right)`$ $`=`$ $`\mathrm{Sset}^D\left(st\right)+p_1^2{\displaystyle \frac{D3}{st}}\mathrm{Sset}^D\left(st\right)`$ (5.4)
$`+{\displaystyle \frac{1}{2}}p_1^4{\displaystyle \frac{(D3)(D4)}{(st)^2}}\mathrm{Sset}^D\left(st\right)+𝒪\left(p_1^6\right).`$
If a factor of $`1/p_1^2`$ appears in the homogeneous term of the differential equations for an off-shell master integral, this integral becomes reducible in the limit $`p_1^20`$.
An example of a reducible master integral is the vertex diagram $`\mathrm{Rtri}(p_1^2,s)`$ of Fig. 3, which fulfils the differential equation
$$\frac{}{p_1^2}\mathrm{Rtri}(p_1^2,s)=\frac{D4}{2}\frac{2p_1^2s}{p_1^2(p_1^2s)}\mathrm{Rtri}(p_1^2,s)\frac{3D8}{2}\frac{1}{p_1^2(p_1^2s)}\mathrm{Sset}^D\left(s\right).$$
(5.5)
Multiplying by $`p_1^2`$, one can immediately read off the limit
$$\mathrm{Rtri}\left(p_1^20,s\right)=\frac{3D8}{D4}\frac{1}{s}\mathrm{Sset}^D\left(s\right)+𝒪\left(p_1^2\right).$$
(5.6)
Subleading terms in $`p_1^2`$ are again obtained by successive iteration of the differential equation (5.5).
All the information required to derive the massless limit of the system of equations for the master crossed-box diagrams is contained in the differential equations for the pinched master integrals.
The differential equations obtained with this method are in perfect agreement with (4.1) and (4.2), giving an independent confirmation of the correctness of the two methods.
In addition, we checked that the two methods give the same results for the crossed boxes with the powers of propagators equal to unity and with one scalar product in the numerator, of the type $`p_ik`$ or $`p_il`$. This was done computing Eq.(3.21) with the method described in Section 3 and then contracting the result with one of the external momenta.
## 6 Analytic expansion of the second master integral
Inserting the $`ϵ`$ expansion of $`\mathrm{Xbox}_1^D(s,t)`$ computed in Ref. and the $`ϵ`$ expansions of the sub-topologies listed in Refs. into Eq. (4.1), and solving it with respect to $`\mathrm{Xbox}_2^D(s,t)`$, we obtain, in the physical region $`s>0`$, $`t,u<0`$,
$$\mathrm{Xbox}_2^D(s,t)=\mathrm{\Gamma }^2(1+ϵ)\left\{\frac{G_1(t,u)}{s^3t}+\frac{G_2(t,u)}{s^2t^2}+\frac{G_1(u,t)}{s^3u}+\frac{G_2(u,t)}{s^2u^2}\right\},$$
(6.1)
where
$`G_1(t,u)=s^{2ϵ}\{{\displaystyle \frac{6}{ϵ^3}}+{\displaystyle \frac{1}{ϵ^2}}(326T6U)`$
$`+{\displaystyle \frac{1}{ϵ}}\left(112\pi ^224T+T^224U+16TU+U^2\right)4318T+13T^2+{\displaystyle \frac{8}{3}}T^3`$
$`18U+16TU+11T^2U+13U^220TU^2+{\displaystyle \frac{8}{3}}U^3+\pi ^2\left(17T+17U{\displaystyle \frac{112}{3}}\right)`$
$`122\zeta (3)+62T\text{Li}_2\left({\displaystyle \frac{t}{s}}\right)62\text{Li}_3\left({\displaystyle \frac{t}{s}}\right)+62\text{S}_{1,2}\left({\displaystyle \frac{t}{s}}\right)`$
$`+i\pi [{\displaystyle \frac{1}{ϵ}}(16+6T+6U)349\pi ^26T10T^26U+14TU10U^2]\},`$
$`G_2(t,u)=s^{2ϵ}\{{\displaystyle \frac{2}{ϵ^4}}+{\displaystyle \frac{1}{ϵ^3}}(8+{\displaystyle \frac{5}{2}}T+{\displaystyle \frac{7}{2}}U)`$
$`+{\displaystyle \frac{1}{ϵ^2}}\left({\displaystyle \frac{29}{2}}{\displaystyle \frac{5}{12}}\pi ^2+7TT^2+20U4TUU^2\right)`$
$`+{\displaystyle \frac{1}{ϵ}}[{\displaystyle \frac{1}{2}}+17T+2T^2{\displaystyle \frac{T^3}{3}}+{\displaystyle \frac{\pi ^2}{6}}(14+5T29U)+13U28TU4U^2`$
$`+3TU^2U^3+{\displaystyle \frac{19}{2}}\zeta (3)2T\text{Li}_2({\displaystyle \frac{t}{s}})+2\text{Li}_3({\displaystyle \frac{t}{s}})2\text{S}_{1,2}({\displaystyle \frac{t}{s}})]`$
$`+{\displaystyle \frac{37}{2}}+{\displaystyle \frac{37}{40}}\pi ^4+7T5T^2{\displaystyle \frac{22}{3}}T^3+{\displaystyle \frac{2}{3}}T^4+5U20TU+{\displaystyle \frac{8}{3}}T^3U2U^2`$
$`+24TU^2T^2U^28U^3{\displaystyle \frac{4}{3}}TU^3+{\displaystyle \frac{4}{3}}U^4`$
$`+{\displaystyle \frac{\pi ^2}{6}}\left(7922T5T^2200U+76TU+25U^2\right)+\left(6813T33U\right)\zeta (3)`$
$`+\left(10\pi ^232T+17T^2+12TU\right)\text{Li}_2\left({\displaystyle \frac{t}{s}}\right)+\left(3260T12U\right)\text{Li}_3\left({\displaystyle \frac{t}{s}}\right)`$
$`+\left(28T6U32\right)\text{S}_{1,2}\left({\displaystyle \frac{t}{s}}\right)26\text{S}_{1,3}\left({\displaystyle \frac{t}{s}}\right)36\text{S}_{2,2}\left({\displaystyle \frac{t}{s}}\right)+86\text{Li}_4\left({\displaystyle \frac{t}{s}}\right)`$
$`+i\pi [{\displaystyle \frac{2}{ϵ^3}}+{\displaystyle \frac{1}{ϵ^2}}(11T+U)+{\displaystyle \frac{1}{ϵ}}(1{\displaystyle \frac{31}{6}}\pi ^210T2T^2+4U2TU2U^2)`$
$`+11+4T2T^2+{\displaystyle \frac{10}{3}}T^3+{\displaystyle \frac{\pi ^2}{3}}\left(65+28TU\right)+2U8TU8U^2`$
$`+2U^389\zeta (3)+(14T+18U)\text{Li}_2({\displaystyle \frac{t}{s}})32\text{Li}_3({\displaystyle \frac{t}{s}})+44\text{S}_{1,2}({\displaystyle \frac{t}{s}})]\},`$
and $`T=\mathrm{log}(t/s)`$, $`U=\mathrm{log}(u/s)`$. We used Nielsen’s generalized polylogarithms $`\text{S}_{n,p}`$ , defined by
$$\text{S}_{n,p}\left(x\right)=\frac{(1)^{n+p1}}{(n1)!p!}_0^1\text{d}t\frac{\mathrm{log}^{n1}(t)\mathrm{log}^p(1xt)}{t},n,p1,x1,$$
(6.4)
where the usual polylogarithms are given by
$$\mathrm{Li}_n(x)=\text{S}_{n1,1}\left(x\right).$$
(6.5)
The three kinematically accessible regions of the phase-space are depicted in Fig. 4.
* $`𝒔\mathbf{>}\mathrm{𝟎}\mathbf{,}𝒕\mathbf{,}𝒖\mathbf{<}\mathrm{𝟎}`$. All logarithms and polylogarithms occurring in Eqs. (6) and (6) are real in this region.
Formulae for the other two regions, (ii) and (iii), can be derived by analytic continuation, starting from region (i) and following the paths indicated in the figure.
The analytic continuation can be performed through a few simple steps.
* $`𝒕\mathbf{>}\mathrm{𝟎}\mathbf{,}𝒔\mathbf{,}𝒖\mathbf{<}\mathrm{𝟎}`$. Going from region (i) to region (ii), we have to pass through two branches: $`t=0`$ and $`s=0`$. We can then split the analytic continuation into two steps:
+ we first split the logarithm $`T=\mathrm{log}(t)\mathrm{log}(s)`$. At $`t=0`$, nothing happens to the polylogarithms $`\text{S}_{n,p}\left(t/s\right)`$, but $`\mathrm{log}(t)`$ gets an imaginary part: $`\mathrm{log}(t)\mathrm{log}(t)i\pi `$.
We are now in an unphysical region, where both $`s`$ and $`t`$ are positive and $`u`$ is negative. Using the transformation formulae for $`x\mathrm{\hspace{0.17em}1}/x`$ (see, eg. Refs. ), we can express $`\text{S}_{n,p}\left(t/s\right)`$ in terms of $`\text{S}_{n,p}\left(s/t\right)`$ and $`\mathrm{log}(t/s)`$.
+ To enter region (ii), we have to pass now the branch point at $`s=0`$. We split $`\mathrm{log}(t/s)=\mathrm{log}(t)\mathrm{log}(s)`$ and $`U=\mathrm{log}(s+t)\mathrm{log}(s)`$ and we analytically continue $`\mathrm{log}(s)\mathrm{log}(s)+i\pi `$.
In this way, for example, the logarithms in Eqs. (6) and (6) undergo the transformation
$`\mathrm{log}\left({\displaystyle \frac{t}{s}}\right)`$ $``$ $`\mathrm{log}\left({\displaystyle \frac{t}{s}}\right)2i\pi ,`$ (6.6)
$`\mathrm{log}\left({\displaystyle \frac{u}{s}}\right)`$ $``$ $`\mathrm{log}\left({\displaystyle \frac{u}{s}}\right)i\pi .`$ (6.7)
* $`𝒖\mathbf{>}\mathrm{𝟎}\mathbf{,}𝒔\mathbf{,}𝒕\mathbf{<}\mathrm{𝟎}`$. The procedure to go from region (i) to region (iii) is similar to the previous one, but it requires an additional step.
+ We rewrite $`\text{S}_{n,p}\left(t/s\right)`$ in terms of $`\text{S}_{n,p}\left((s+t)/s\right)`$, $`\mathrm{log}(t/s)`$ and $`\mathrm{log}((s+t)/s)`$, using the transformation $`x\mathrm{\hspace{0.17em}1}x`$, and we split the logarithms as before. In passing the first branch point at $`u=0`$, the polylogarithms are well defined while $`\mathrm{log}(s+t)\mathrm{log}(st)i\pi `$.
+ We invert now the argument of the polylogarithms, expressing $`\text{S}_{n,p}\left((s+t)/s\right)`$ in terms of $`\text{S}_{n,p}\left(s/(s+t)\right)`$ and $`\mathrm{log}((st)/s)=\mathrm{log}(st)\mathrm{log}(s)`$. Finally, $`\mathrm{log}(s)\mathrm{log}(s)+i\pi `$, as we pass the branch point at $`s=0`$ and enter region (iii).
The logarithms in Eqs. (6) and (6) undergo the transformation
$`\mathrm{log}\left({\displaystyle \frac{t}{s}}\right)`$ $``$ $`\mathrm{log}\left({\displaystyle \frac{t}{s}}\right)i\pi ,`$ (6.8)
$`\mathrm{log}\left({\displaystyle \frac{u}{s}}\right)`$ $``$ $`\mathrm{log}\left({\displaystyle \frac{u}{s}}\right)2i\pi .`$ (6.9)
The expression for $`G_1(t,u)`$ and $`G_2(t,u)`$ in this region can also be obtained directly from the expressions in region (ii), using the symmetry $`tu`$.
A non-trivial check of the correctness of the expressions of $`\mathrm{Xbox}_1^D(s,t)`$ and $`\mathrm{Xbox}_2^D(s,t)`$ comes from Eq. (4.2), that must be identically satisfied, once the respective $`ϵ`$ expansions are used.
## 7 Conclusions
High-energy scattering processes are one of the most important sources of information on short-distance physics. Recent improvements of experimental measurements demand nowadays the knowledge of $`2\mathrm{\hspace{0.17em}2}`$ scattering rates at two-loop order.
Two non-trivial topologies characterize these processes: planar two-loop boxes , and crossed two-loop boxes.
The crossed boxes constituted the last barrier towards the completion of this goal. In this paper we presented an algorithm to reduce tensor two-loop massless crossed boxes with light-like external legs to two master crossed boxes (the first one with all powers of propagators equal to one, and the second one with the power of the second propagator of Fig. 1 equal to two), plus simpler diagrams.
We derived the equations that connect the two master integrals in dimensions $`D`$ and $`D+2`$ and the system of first-order coupled differential equations satisfied by the two master integrals.
This last part was done following two different methods: using the raising and lowering operators, and using the on-shell limit of the differential equations for the crossed box with one off-shell leg . The agreement between the results of the two methods is a strong support for the validity of the two different procedures.
The differential equation for the first master integral allowed us to derive an $`ϵ`$ expansion for the second master integral, once the known expression of the first master integral is inserted. The differential equation for the second master integral then provided us a non-trivial way to check the obtained expansion.
### Acknowledgement
We thank E.W.N. Glover, M.E. Tejeda-Yeomans and J.J. van der Bij for assistance and useful suggestions. C.A. acknowledges the financial support of the Greek Government, C.O. acknowledges the financial support of the INFN, E.R. thanks the Alexander-von-Humboldt Stiftung for supporting his stay at the Institut für Theoretische Teilchenphysik of the University of Karlsruhe and J.B.T. acknowledges the financial support of the DFG-Forschergruppe “Quantenfeldtheorie, Computeralgebra und Monte-Carlo-Simulation”. We gratefully acknowledge the support of the British Council and German Academic Exchange Service under ARC project 1050.
## A Tensor reduction and dimensional shift of the crossed two-loop triangle
For completeness, we give here the reduction formulae for the crossed two-loop triangle that appears as a sub-topology in the scalar reduction of the crossed box. We make use of integration-by-parts and Lorentz-invariance identities to reduce the generic scalar triangle to one master integral which has all powers of propagators equal to one , plus pinchings. For an alternative solution to this reduction problem, exploiting a connection with massless three-loop propagator integrals, see Ref. .
The generic two-loop scalar crossed triangle of Fig. 5 is given by
$$\mathrm{Xtri}^D(\nu _1,\nu _2,\nu _3,\nu _4,\nu _5,\nu _6;s)=\frac{d^Dk}{i\pi ^{D/2}}\frac{d^Dl}{i\pi ^{D/2}}\frac{1}{A_1^{\nu _1}A_2^{\nu _2}A_3^{\nu _3}A_4^{\nu _4}A_5^{\nu _5}A_6^{\nu _6}},$$
(A.1)
where the propagators are
$`A_1`$ $`=`$ $`(k+l+p_3+p_4)^2+i0,`$ (A.2)
$`A_2`$ $`=`$ $`(k+l)^2+i0,`$
$`A_3`$ $`=`$ $`l^2+i0,`$
$`A_4`$ $`=`$ $`(l+p_3)^2+i0,`$
$`A_5`$ $`=`$ $`k^2+i0,`$
$`A_6`$ $`=`$ $`(k+p_4)^2+i0.`$
The external momenta are in-going, two of them being light-like, $`p_3^2=p_4^2=0`$, while $`p_{12}^2=s`$ is the only physical scale.
Repeating the reasoning of Section 3.1, we can build eight integration-by-parts identities and a single Lorentz-invariance identity. They depend upon one irreducible scalar product in the numerator, that we choose to be $`(lp_4)`$:
$`s\nu _1\mathrm{𝟏}^\mathbf{+}+\left(2D2\nu _{235}\nu _{146}\right)\nu _4\mathrm{𝟒}^\mathbf{+}\mathrm{𝟑}^{\mathbf{}}\nu _1\mathrm{𝟏}^\mathbf{+}\mathrm{𝟐}^{\mathbf{}}\nu _6\mathrm{𝟔}^\mathbf{+}\mathrm{𝟓}^{\mathbf{}}=0`$ (A.3)
$`s\nu _2\mathrm{𝟐}^\mathbf{+}+\left(2D2\nu _{146}\nu _{235}\right)\nu _3\mathrm{𝟑}^\mathbf{+}\mathrm{𝟒}^{\mathbf{}}\nu _2\mathrm{𝟐}^\mathbf{+}\mathrm{𝟏}^{\mathbf{}}\nu _5\mathrm{𝟓}^\mathbf{+}\mathrm{𝟔}^{\mathbf{}}=0`$ (A.11)
$`2\left(lp_4\right)\nu _1\mathrm{𝟏}^\mathbf{+}\left(D\nu _{24}2\nu _3\right)+\nu _1\mathrm{𝟏}^\mathbf{+}\left(\mathrm{𝟐}^{\mathbf{}}+\mathrm{𝟒}^{\mathbf{}}\mathrm{𝟓}^{\mathbf{}}\right)`$
$`+\nu _2\mathrm{𝟐}^\mathbf{+}\left(\mathrm{𝟑}^{\mathbf{}}\mathrm{𝟓}^{\mathbf{}}\right)+\nu _4\mathrm{𝟒}^\mathbf{+}\mathrm{𝟑}^{\mathbf{}}=0`$
$`2\left(lp_4\right)\nu _2\mathrm{𝟐}^\mathbf{+}\left(D\nu _{125}2\nu _6\right)+\nu _2\mathrm{𝟐}^\mathbf{+}\left(\mathrm{𝟔}^{\mathbf{}}\mathrm{𝟑}^{\mathbf{}}\right)`$
$`\nu _1\mathrm{𝟏}^\mathbf{+}\left(\mathrm{𝟔}^{\mathbf{}}\mathrm{𝟒}^{\mathbf{}}\right)\nu _5\mathrm{𝟓}^\mathbf{+}\mathrm{𝟔}^{\mathbf{}}=0`$
$`2\left(lp_4\right)\nu _3\mathrm{𝟑}^\mathbf{+}+\left(D\nu _{345}\nu _6\right)+\nu _3\mathrm{𝟑}^\mathbf{+}\left(\mathrm{𝟐}^{\mathbf{}}\mathrm{𝟓}^{\mathbf{}}\right)+\nu _4\mathrm{𝟒}^\mathbf{+}\left(\mathrm{𝟏}^{\mathbf{}}\mathrm{𝟔}^{\mathbf{}}\right)`$
$`\nu _5\mathrm{𝟓}^\mathbf{+}\mathrm{𝟔}^{\mathbf{}}=0`$
$`2\left(lp_4\right)\nu _4\mathrm{𝟒}^\mathbf{+}+s\nu _4\mathrm{𝟒}^\mathbf{+}(D\nu _{346}2\nu _5)+\nu _3\mathrm{𝟑}^\mathbf{+}\left(\mathrm{𝟓}^{\mathbf{}}\mathrm{𝟐}^{\mathbf{}}\right)`$
$`+\nu _4\mathrm{𝟒}^\mathbf{+}\left(\mathrm{𝟔}^{\mathbf{}}\mathrm{𝟏}^{\mathbf{}}\right)+\nu _6\mathrm{𝟔}^\mathbf{+}\mathrm{𝟓}^{\mathbf{}}=0`$
$`2\left(lp_4\right)\nu _5\mathrm{𝟓}^\mathbf{+}+s\nu _5\mathrm{𝟓}^\mathbf{+}\left(D\nu _{36}2\nu _4\right)+\nu _5\mathrm{𝟓}^\mathbf{+}\left(\mathrm{𝟒}^{\mathbf{}}+\mathrm{𝟔}^{\mathbf{}}\mathrm{𝟏}^{\mathbf{}}\right)`$
$`+\nu _6\mathrm{𝟔}^\mathbf{+}\left(\mathrm{𝟒}^{\mathbf{}}\mathrm{𝟏}^{\mathbf{}}\right)+\nu _3\mathrm{𝟑}^\mathbf{+}\mathrm{𝟒}^{\mathbf{}}=0`$
$`2\left(lp_4\right)\nu _6\mathrm{𝟔}^\mathbf{+}\left(D\nu _{45}2\nu _3\right)\nu _6\mathrm{𝟔}^\mathbf{+}\left(\mathrm{𝟑}^{\mathbf{}}+\mathrm{𝟓}^{\mathbf{}}\mathrm{𝟐}^{\mathbf{}}\right)`$
$`+\nu _5\mathrm{𝟓}^\mathbf{+}\left(\mathrm{𝟐}^{\mathbf{}}\mathrm{𝟑}^{\mathbf{}}\right)+\nu _4\mathrm{𝟒}^\mathbf{+}\mathrm{𝟑}^{\mathbf{}}=0`$
$`2\left(lp_4\right)\nu _1\mathrm{𝟏}^\mathbf{+}\left(D\nu _{2356}\right)+\nu _1\mathrm{𝟏}^\mathbf{+}\left(\mathrm{𝟐}^{\mathbf{}}+\mathrm{𝟔}^{\mathbf{}}\mathrm{𝟓}^{\mathbf{}}\right)+\nu _4\mathrm{𝟒}^\mathbf{+}\mathrm{𝟑}^{\mathbf{}}=0.`$
By eliminating the irreducible scalar product in the numerator, we obtain
$$s\nu _1\mathrm{𝟏}^\mathbf{+}=\left(2D2\nu _{235}\nu _{146}\right)+\nu _4\mathrm{𝟒}^\mathbf{+}\mathrm{𝟑}^{\mathbf{}}+\nu _1\mathrm{𝟏}^\mathbf{+}\mathrm{𝟐}^{\mathbf{}}+\nu _6\mathrm{𝟔}^\mathbf{+}\mathrm{𝟓}^{\mathbf{}},$$
(A.12)
which, together with the symmetric one for $`\nu _2\mathrm{𝟐}^\mathbf{+}`$, can reduce $`\nu _1`$ and $`\nu _2`$ to unity. Using
$`\nu _3\mathrm{𝟑}^\mathbf{+}`$ $`=`$ $`{\displaystyle \frac{1}{D2\nu _{3456}}}\left(\nu _4\mathrm{𝟒}^\mathbf{+}\nu _6\mathrm{𝟔}^\mathbf{+}\mathrm{𝟏}^{\mathbf{}}\nu _3\mathrm{𝟑}^\mathbf{+}\nu _5\mathrm{𝟓}^\mathbf{+}\mathrm{𝟐}^{\mathbf{}}\right)`$ (A.13)
$`+{\displaystyle \frac{1}{D22\nu _{34}}}\left[\left(D22\nu _{46}\right)\nu _6\mathrm{𝟔}^\mathbf{+}+2\left(\nu _3\nu _6\right)\nu _5\mathrm{𝟓}^\mathbf{+}\right],`$
and the symmetric one for $`\nu _4\mathrm{𝟒}^\mathbf{+}`$, we can reduce $`\nu _3`$ and $`\nu _4`$ to one. To complete the reduction, we use
$$\left(\nu _{56}\nu _{34}\right)=\nu _2\mathrm{𝟐}^\mathbf{+}\left(\mathrm{𝟑}^{\mathbf{}}\mathrm{𝟓}^{\mathbf{}}\right)+\nu _1\mathrm{𝟏}^\mathbf{+}\left(\mathrm{𝟒}^{\mathbf{}}\mathrm{𝟔}^{\mathbf{}}\right),$$
(A.14)
that can be re-iterated until $`\left(\nu _{56}\nu _{34}\right)=0`$. Since we are applying this identity to scalar integrals where $`\nu _3`$ and $`\nu _4`$ have already been reduced to one, the reduction procedure will stop when $`\nu _5=\nu _6=1`$. This integral cannot be reduced any further, and we choose the crossed master triangle to be
$$\mathrm{Xtri}^D(s)=\mathrm{Xtri}^D(1,1,1,1,1,1;s).$$
(A.15)
The application of the identity
$$𝐝^{\mathbf{}}=\left(\nu _6\mathrm{𝟔}^\mathbf{+}+\nu _5\mathrm{𝟓}^\mathbf{+}+\nu _4\mathrm{𝟒}^\mathbf{+}+\nu _3\mathrm{𝟑}^\mathbf{+}\right)\left(\nu _2\mathrm{𝟐}^\mathbf{+}+\nu _1\mathrm{𝟏}^\mathbf{+}\right)+\left(\nu _4\mathrm{𝟒}^\mathbf{+}+\nu _3\mathrm{𝟑}^\mathbf{+}\right)\left(\nu _6\mathrm{𝟔}^\mathbf{+}+\nu _5\mathrm{𝟓}^\mathbf{+}\right)$$
(A.16)
to $`\mathrm{Xtri}^{D+2}(s)`$ and its further reduction, gives rise to the dimensional-shift formula
$`\mathrm{Xtri}^{D+2}\left(s\right)`$ $`=`$ $`{\displaystyle \frac{(D4)s^2}{4(D2)(2D7)(2D5)}}\mathrm{Xtri}^D\left(s\right)`$ (A.17)
$`{\displaystyle \frac{37D^3313D^2+858D752}{2(D4)(D2)(2D7)(2D5)(3D8)}}\mathrm{Btri}^D\left(s\right)`$
$`+{\displaystyle \frac{43D^4478D^3+1963D^23530D+2352}{2(D4)^2(D3)(D2)(2D7)(2D5)s}}\mathrm{Sset}^D\left(s\right).`$
The expression of the master integral of Eq. (A.15) has been computed in Refs. .
## B Coefficients $`𝒄_𝒊`$
We collect here the coefficients $`c_i`$ of Eq. (3.1).
$`c_1`$ $`=`$ $`4{\displaystyle \frac{(D5)^2u}{(D6)st}}`$
$`c_2`$ $`=`$ $`{\displaystyle \frac{tu}{t}}`$
$`c_3`$ $`=`$ $`2{\displaystyle \frac{(D4)(2D9)}{(D6)st}}`$
$`c_4`$ $`=`$ $`3{\displaystyle \frac{(D4)(3D14)u^2}{(D6)^2(2D11)s^3t^3}}\left[2\left(7D^276D+206\right)t+(D5)(5D28)s\right]`$
$`c_5`$ $`=`$ $`3{\displaystyle \frac{(D4)(3D14)t}{(D6)^2(2D11)s^3u^2}}\left[2\left(7D^276D+206\right)t+(D5)(5D28)s\right]`$
$`c_6`$ $`=`$ $`3{\displaystyle \frac{(D5)(D4)(3D14)s\left[2(D5)t(D6)s\right]}{(D6)^2(2D11)t^3u^2}}`$
$`c_7`$ $`=`$ $`6{\displaystyle \frac{(D3)(3D14)\left[(5D26)t+2(D5)s\right]}{(D6)^2st^3}}`$
$`c_8`$ $`=`$ $`6{\displaystyle \frac{(D3)(3D14)\left[(5D26)t+2(D5)s\right]}{(D6)^2stu^2}}`$
$`c_9`$ $`=`$ $`{\displaystyle \frac{3}{2}}{\displaystyle \frac{(D3)(3D10)}{(D6)^2(D5)(D4)s^3t^3u^2}}[(41D^3620D^2+3124D5248)st^3`$
$`+2(D6)^2(D4)t^4\left(61D^3922D^2+4640D7776\right)s^2t^2`$
$`4(D5)(3D14)(4D21)s^3t4(D5)^2(3D14)s^4]`$
$`c_{10}`$ $`=`$ $`3{\displaystyle \frac{(D3)(3D10)(3D8)}{(D6)^2(D5)(D4)^2(2D11)s^4t^3u^2}}`$
$`\times [2(D6)(D4)(2D11)(5D22)t^4`$
$`2\left(64D^41305D^3+9978D^233902D+43180\right)st^3`$
$`\left(208D^44209D^3+31899D^2107314D+135216\right)s^2t^2`$
$`(D5)(3D14)(34D^2359D+944)s^3t(D5)^2(3D14)(5D28)s^4]`$
$`c_{11}`$ $`=`$ $`3{\displaystyle \frac{(D3)(3D14)(3D10)(3D8)}{(D6)^2(D5)(D4)^2(2D11)s^3t^4u^2}}`$
$`\times [2(D5)(7D^276D+206)t^4+(61D^3952D^2+4941D8528)st^3`$
$`+\left(92D^31417D^2+7253D12336\right)s^2t^2`$
$`+(56D^3853D^2+4317D7258)s^3t+2(D5)(2D11)(2D9)s^4]`$
$`c_{12}`$ $`=`$ $`3{\displaystyle \frac{(D3)(3D14)(3D10)(3D8)}{(D6)^2(D5)(D4)^2(2D11)s^3t^3u^3}}[2(D5)(7D^276D+206)t^4`$ (B.1)
$`\left(9D^3130D^2+615D948\right)st^33(D6)(D5)(D4)s^2t^2`$
$`(D5)(3D^230D+74)s^3t+(D6)(D5)^2s^4]`$
## C Coefficients of the dimensional-shift system
We give here only the first three coefficients of the system (3.71). The whole list can be obtained from the authors (C.A. and C.O.).
$`a_1`$ $`=`$ $`{\displaystyle \frac{3}{512}}{\displaystyle \frac{(D4)s^2\left[(3D14)(3D10)s^24(5D^239D+74)tu\right]}{(D3)^3(2D9)(2D7)tu}}`$
$`a_2`$ $`=`$ $`{\displaystyle \frac{(D6)(D4)s\left[4(D3)tu3(3D10)s^2\right]}{512(D5)(D3)^3(2D9)(2D7)}}`$
$`a_3`$ $`=`$ $`{\displaystyle \frac{(D4)s\left[4(D6)(D3)tu+3(3D14)(3D10)s^2\right]}{512(D5)(D3)^3(2D7)tu}}`$
$`b_1`$ $`=`$ $`{\displaystyle \frac{s\left[3(D4)(3D14)s^24(D6)(D5)tu\right]}{128(D3)(2D9)(2D7)tu}}`$
$`b_2`$ $`=`$ $`{\displaystyle \frac{(D6)\left[4(D3)tu+3(D4)s^2\right]}{128(D5)(D3)(2D9)(2D7)}}`$
$`b_3`$ $`=`$ $`{\displaystyle \frac{3(D4)(3D14)s^24(D6)(D3)tu}{128(D5)(D3)(2D7)tu}}`$ (C.1)
## D Coefficients of the differential equations for the two master <br>integrals
The coefficients of the system of differential equations (4.1) and (4.2) for the two master integrals are given by
$`h_1`$ $`=`$ $`{\displaystyle \frac{(D4)s^24tu}{4tu}}`$
$`h_2`$ $`=`$ $`{\displaystyle \frac{(D6)s}{4(D5)}}`$
$`h_3`$ $`=`$ $`{\displaystyle \frac{(D4)(2D9)s}{4(D5)tu}}`$
$`h_4`$ $`=`$ $`{\displaystyle \frac{3}{2}}{\displaystyle \frac{(D4)(3D14)u}{(D5)st^2}}`$
$`h_5`$ $`=`$ $`{\displaystyle \frac{3}{2}}{\displaystyle \frac{(D4)(3D14)s^2}{(D6)t^2u^2}}`$
$`h_6`$ $`=`$ $`3{\displaystyle \frac{(D3)(3D14)}{(D5)t^2}}`$
$`h_7`$ $`=`$ $`{\displaystyle \frac{3}{4}}{\displaystyle \frac{(D3)(3D10)\left[(3D14)\left(u^2+t^2\right)+2(D4)tu\right]}{(D5)(D4)st^2u^2}}`$
$`h_8`$ $`=`$ $`{\displaystyle \frac{3}{4}}{\displaystyle \frac{(D3)(3D10)(3D8)\left[(D5)(3D14)\left(u^2+t^2\right)(D6)(D4)tu\right]}{(D5)^2(D4)^2s^2t^2u^2}}`$
$`h_9`$ $`=`$ $`{\displaystyle \frac{3}{2}}{\displaystyle \frac{(D3)(3D14)(3D10)(3D8)}{(D6)(D5)^2(D4)^2st^3u^2}}[(2D9)(3D16)u^2`$ (D.1)
$`+(7D^268D+164)tu+2(D5)^2t^2],`$
and by
$`k_1`$ $`=`$ $`{\displaystyle \frac{(D5)^2s}{tu}}`$
$`k_2`$ $`=`$ $`{\displaystyle \frac{(D6)(u^2+t^2)}{2tu}}`$
$`k_3`$ $`=`$ $`{\displaystyle \frac{(D4)(2D9)}{tu}}`$
$`k_4`$ $`=`$ $`6{\displaystyle \frac{(D4)(3D14)u\left[(5D)u+(2D11)t\right]}{(D6)s^2t^3}}`$
$`k_5`$ $`=`$ $`{\displaystyle \frac{3}{2}}{\displaystyle \frac{(D5)(D4)(3D14)s^3}{(D6)t^3u^3}}`$
$`k_6`$ $`=`$ $`3{\displaystyle \frac{(D3)(3D14)}{(D6)st^3u}}\left[(5D28)tu+(D6)t^22(D5)u^2\right]`$
$`k_7`$ $`=`$ $`{\displaystyle \frac{3}{2}}{\displaystyle \frac{(D3)(3D10)}{(D6)(D4)s^2t^3u^3}}[2(D6)(3D14)tu(u^2+t^2)`$
$`(D5)(3D14)(u^4+t^4)+2(5D^249D+118)t^2u^2]`$
$`k_8`$ $`=`$ $`3{\displaystyle \frac{(D3)(3D10)(3D8)}{(D6)(D5)(D4)^2s^3t^3u^3}}[3(D5)^2(3D14)tu(u^2+t^2)`$
$`(D5)^2(3D14)(u^4+t^4)(D4)(7D^270D+176)t^2u^2]`$
$`k_9`$ $`=`$ $`3{\displaystyle \frac{(D3)(3D14)(3D10)(3D8)}{(D6)^2(D5)(D4)^2s^2t^4u^3}}[(D5)^2(D2)u^4`$ (D.2)
$`+(D6)\left(13D^2129D+318\right)tu^3+2\left(5D^380D^2+422D734\right)t^2u^2`$
$`+(D6)(D5)(5D24)t^3u+(D6)(D5)^2t^4].`$ |
warning/0003/hep-th0003087.html | ar5iv | text | # References
Introduction
The conclusions of this paper, that electrons are particles, will not surprise our experimental colleagues. What may come as a surprise to them is the fact that the hitherto accepted wisdom in the theoretical community was that there does not exist any relativistic description of the electron as a particle!
The root of the problem lies in the masslessness of the photon. This generates long range interactions and it is well known that these fall off so slowly that they cannot be neglected even for widely separated charges a long time before or after scattering processes. This is seen in $`S`$-matrix calculations as the lack of a pole for fermionic external legs.
This lack of a particle interpretation of an electron presents a radical departure from the usual view of particle physics which we find highly unsatisfactory. We want to show that it is also unnecessary. We note that an understanding of how particles can arise in gauge theories is important to improve our insight into the physical structures of gauge theories, to help with QCD phenomenology and it may well have spin offs for the vexed question of how to describe unstable particles.
In the context of QED this lack of any particle language has not hindered progress. Our knowledge of the classical limit of QED acts as a guide to the extraction of physical predictions from suitably defined cross-sections via the Bloch-Nordsieck framework. In theories such as QCD such intuition is still greatly lacking. In particular we do not understand hadronisation which relies on coloured particles metamorphosing into jets. The route which leads from partons to effective quarks and glue is essentially uncharted.
In this letter we will show that a relativistic particle description of the electron (or any other charged particle) immediately follows once a correct physical identification of the electron has been made. Having first recalled the standard statement of the problem , we will show that the effects which prevent the identification of a charged particle structure in QED disappear if the right fields are used. We stress that this identification only holds in the asymptotic region a long time before or after scattering occurs. We then demonstrate that a particle structure is also asymptotically present for photons in full QED. Finally we discuss how these results have been verified in perturbative calculations and present some conclusions.
The Interaction
We have already noted that in theories like QED we have long range interactions due to the masslessness of the photon. It has been shown by various authors that such interactions cannot be neglected even at large times before or after scattering. In particular, Kulish and Faddeev showed that the annihilation operator of the matter fields of QED which is defined as the large time limit of the operator
$$b(q,s,t):=d^3x\frac{1}{\sqrt{2E_q}}u^s(q)\psi (x)\mathrm{e}^{iqx},$$
(1)
does not become just the usual free particle mode, but rather takes on the form<sup>4</sup><sup>4</sup>4There is also a distortion contribution to the non-observable phase of $`S`$-matrix elements which we will ignore in this letter.
$$b(q,s,t)=D_{\mathrm{soft}}(q,t)b(q,s),$$
(2)
where
$$D_{\mathrm{soft}}(q,t)=\mathrm{exp}\left\{e\underset{\mathrm{soft}}{}\frac{d^3k}{(2\pi )^3}\frac{1}{2\omega _k}\left(\frac{qa(k)}{qk}\mathrm{e}^{itkq/E_q}\frac{qa^{}(k)}{qk}\mathrm{e}^{itkq/E_q}\right)\right\},$$
(3)
is called a *distortion operator* . The creation and annihilation operators for the photonic variables enter into this expression and, as long as $`D_{\mathrm{soft}}1`$, a particle mode for the electron will not be recovered. Of course, we should only expect to regain a particle description at asymptotic times and, for large $`t`$, the integral in (3) only receives contributions from soft photons, but it still does not reduce to the unit operator. This important observation has lead to the conclusion that it is *not* possible to describe the electron as a particle.
However, since the interaction does not switch off, the matter field $`\psi `$ *never* becomes gauge invariant and thus we should not expect to identify it with physical particles via Equation 1! We now construct the fields which do have a particle description at large times.
The Electron
We recall that physical fields must be invariant under the local gauge transformations, $`A_\mu (x)A_\mu (x)+_\mu \theta (x)`$ and $`\psi (x)\mathrm{e}^{ie\theta (x)}\psi (x)`$. A physical matter field must then be given by a product of the form $`h^1(x)\psi (x)`$ where, under a gauge transformation, $`h^1(x)h^1(x)\mathrm{e}^{ie\theta (x)}`$. We call such a field with this gauge transformation property a dressing.
Of course there are a multiplicity of such dressing fields which satisfy this minimal requirement . To describe a charged particle we need to further demand that the dressing satisfies the dressing equation
$$uh^1(x)=ieh^1(x)uA(x),$$
(4)
where $`u^\mu =\gamma (\eta +v)^\mu `$ is the four velocity of the charged particle, $`\eta `$ is the unit time-like vector, $`v=(0,𝒗)`$ is the velocity and $`\gamma =(1|𝒗|^2)^{1/2}`$. It is important to note the velocity dependence here: we can only expect to have a particle interpretation of a charge at asymptotic times and in that regime the velocity is a well defined quantum number. Our dressed charges will, therefore, only correspond to particles at the appropriate point on the mass shell characterised by the velocity in the dressing.
In QED we have been able to solve these two requirements and have found the following description of a charged field moving with a given velocity:
$$h^1(x)\psi (x)=\mathrm{e}^{ieK(x)}\mathrm{e}^{ie\chi (x)}\psi (x).$$
(5)
The $`K`$ term is separately gauge invariant and contributes to the unobservable phase. We ignore it in this letter. The $`\chi `$ part of the dressing may be written as
$$\chi (x)=\frac{𝒢A}{𝒢},$$
(6)
where $`𝒢^\mu =(\eta +v)^\mu (\eta v)^\mu `$.
We now define an annihilation operator for the gauge invariant, dressed field
$$b(q,s,v,t):=d^3x\frac{1}{\sqrt{2E_q}}u^s(q)\mathrm{e}^{ie\chi (x)}\psi (x)\mathrm{e}^{iqx}.$$
(7)
Note that this has an explicit velocity dependence coming from the form of the dressing. To find its asymptotic form we proceed in the same manner as Kulish and Faddeev . The annihilation operator now picks up two distortion factors: the original one associated with the unphysical matter field and a further correction from the dressing. We find
$$b(q,s,t,v)h_{\mathrm{soft}}^1(q,t,v)D_{\mathrm{soft}}(q,t)b(q,s),$$
(8)
where $`D_{\mathrm{soft}}(q,t)`$ is given in (3). The distortion from the soft part of the dressing is
$$h_{\mathrm{soft}}^1(q,t,v)=\mathrm{exp}\left\{e\underset{\mathrm{soft}}{}\frac{d^3k}{(2\pi )^3}\frac{1}{2\omega _k}\left(\frac{Va(k)}{Vk}\mathrm{e}^{itkq/E_q}\frac{Va^{}(k)}{Vk}\mathrm{e}^{itkq/E_q}\right)\right\},$$
(9)
where $`V^\mu =(\eta +v)^\mu (\eta v)kk^\mu `$, which is essentially the Fourier transform of $`𝒢^\mu `$. Combining these distortions we find the overall distortion factor
$`h_{\mathrm{soft}}^1(q,t,v)D_{\mathrm{soft}}(q,t)`$ $`=`$ $`\mathrm{exp}(e{\displaystyle \underset{\mathrm{soft}}{}}{\displaystyle \frac{d^3k}{(2\pi )^3}}{\displaystyle \frac{1}{2\omega _k}}[({\displaystyle \frac{Va(k)}{Vk}}{\displaystyle \frac{qa(k)}{qk}})\mathrm{e}^{itkq/E_q}`$ (10)
$`({\displaystyle \frac{Va^{}(k)}{Vk}}{\displaystyle \frac{qa^{}(k)}{qk}})\mathrm{e}^{itkq/E_q}]).`$
We now note that we can write (recall that $`k`$ is on-shell)
$`{\displaystyle \frac{V^\mu }{Vk}}{\displaystyle \frac{q^\mu }{qk}}`$ $`=`$ $`{\displaystyle \frac{(\eta +v)^\mu (\eta v)kk^\mu }{(\eta +v)k(\eta v)k}}{\displaystyle \frac{q^\mu }{qk}}`$ (11)
$`=`$ $`{\displaystyle \frac{(\eta +v)^\mu }{(\eta +v)k}}{\displaystyle \frac{q^\mu }{qk}}{\displaystyle \frac{k^\mu }{Vk}}.`$
Hence, at the point in the mass-shell where $`q^\mu =m\gamma (\eta +v)^\mu `$, the distortion operator becomes a trivial operator since the argument of the exponential becomes
$$e\frac{d^3k}{(2\pi )^3}\frac{1}{2\omega _k}\left(\frac{ka(k)}{Vk}\mathrm{e}^{it\omega _k}\frac{ka^{}(k)}{Vk}\mathrm{e}^{it\omega _k}\right),$$
(12)
which vanishes between physical states due to the Gupta-Bleuler subsidiary condition (it only contains unphysical photon degrees of freedom which correspond to the Nakanishi-Lautrup $`B`$ field). We thus see that the modes of the dressed field, at the appropriate place in the mass shell, are free particle modes at large times. There is no distortion and we have recovered a particle picture for charged matter in QED.
The Photon
In the usual framework it is assumed that the coupling switches off asymptotically. There are then many ways to show that the physical components of the non-interacting vector potential are the transverse degrees of freedom (see, e.g., Chap. 19 of ). However, we have seen that in the true asymptotic domain the coupling does not vanish due to the masslessness of the photon. This means that previous arguments are incomplete. We will now show that, although we still have an interaction, the physical photonic degrees of freedom do decouple and a particle description for the photon can be recovered in full QED.
The asymptotic form of the (interaction picture) vector boson can be readily obtained once the form of the asymptotic interaction Hamiltonian, $`H_{\mathrm{int}}^{\mathrm{as}}`$, has been identified . Using this result we can straightforwardly transform from the free vector boson, $`A^\mathrm{f}`$, to the asymptotic Heisenberg field, $`A^{\mathrm{as}}`$, as follows:
$`A_\mu ^{\mathrm{as}}(x)`$ $`=`$ $`\mathrm{exp}\left(i{\displaystyle _{\mathrm{}}^t}𝑑\tau H_{\mathrm{int}}^{\mathrm{as}}(\tau )\right)A_\mu ^\mathrm{f}(x)\mathrm{exp}\left(i{\displaystyle _{\mathrm{}}^t}𝑑\tau H_{\mathrm{int}}^{\mathrm{as}}(\tau )\right)`$ (13)
$`=`$ $`A_\mu ^\mathrm{f}(x)e{\displaystyle _{\mathrm{}}^t}𝑑\tau d^3yD(\tau t,𝒚𝒙)J_\mu ^{\mathrm{as}}(\tau ,𝒚),`$
where the asymptotic current is given by
$$J_{\mathrm{as}}^\mu (t,𝒙)=\frac{d^3p}{(2\pi )^3}\frac{p^\mu }{E_p}\rho (p)\delta ^3\left(𝒙t𝒑/E_p\right).$$
(14)
This shows, as expected, that the vector boson field is also not free at large times. In the spirit of our discussions above, we can, though, straightforwardly recover a particle picture for photons by showing that the second term in (13) vanishes for the transverse, physical components, $`(\delta _{ij}_i_j/^2)A_j^{\mathrm{as}}`$. This shows that in the far field domain, where the potential takes on the asymptotic form (13), a particle description emerges even in the interacting theory.
Conclusions
In this letter we have shown how a particle description for the electron and the photon emerges from QED. The essential observation in this construction is that a particle cannot be identified with the raw matter field that enters into the QED Lagrangian. Rather, a physical particle corresponds to an appropriately dressed, gauge invariant field. The dressing for a charged asymptotic particle with a specific velocity depends explicitly on that velocity. A particle description for the electron is only recovered at that point in the mass shell corresponding to the velocity.
For a photon it is not surprising that not all the components of the vector potential are physical. What is new in our discussion of the photon is that we have seen how the free photon emerges in the asymptotic regime even though there is still a residual interaction with matter for the non-physical components.
The formal arguments presented in this letter have been checked in a wide variety of detailed calculations. We have shown for both scalar and fermionic QED that the on-shell propagator and other Green’s functions of the dressed matter fields have, to all orders in perturbation theory, a good pole structure. This is in contrast to the Lagrangian matter fields whose on-shell Green’s functions are plagued by infra-red singularities. This means that for our fields the $`S`$-matrix elements can be constructed via the traditional LSZ-formalism.
We should also point out that these charged fields have good ultra-violet behaviour . Multiplicative renormalisation is possible and the dressed field operators do not mix under renormalisation . In addition, there is an ultra-violet logarithm in scattering processes which is just the universal Isgur-Wise (or equivalently the Wilson line kink) renormalisation constant. This logarithm structure is just the one that appears in the Bloch-Nordsieck formalism.
This programme is currently being extended to the non-abelian domain where the theoretical problems are more severe and the experimental situation in identifying physical quarks and gluons, e.g., in jets, is much more subtle. The essential new ingredient found there that further obstructs a particle interpretation is the gluonic self-interaction. Such massless charges spawn a new class of asymptotic dynamics which is no longer spin independent and generates collinear singularities. This can be studied in massless QED and early indications are that the appropriate dressed massless charges are indeed free. Furthermore, we have already seen that the minimal, non-abelian component of the dressing is responsible for the anti-screening interaction that drives asymptotic freedom in both three and four dimensional QCD. These results give us a great deal of confidence that we can regain a particle picture of quarks and gluons in the pre-hadronisation regime.
Acknowledgements: This work was supported by the British Council/Spanish Education Ministry Acciones Integradas grant no. 1801/HB1997-0141. We thank Robin Horan, Tom Steele, Shogo Tanimura and Izumi Tsutsui for discussions. |
warning/0003/gr-qc0003034.html | ar5iv | text | # Solar System Tests of Higher-Dimensional Gravity
## 1 Introduction
There is now a substantial literature on the higher-dimensional extension of Einstein’s general theory of relativity known as Kaluza-Klein gravity (Overduin and Wesson (1997); Wesson (1999)). There are several ways to test the theory, with perhaps the most straightforward involving the motion of test particles in the field of a static, spherically-symmetric mass like the Sun or the Earth. Birkhoff’s theorem in the usual sense does not hold in higher dimensions (Bronnikov and Melnikov (1995); Schmidt (1997)), so some question arises in identifying the appropriate metric to use for this problem. In the five-dimensional (5D) case (with one extra coordinate $`yx^4`$), most attention has focused on the soliton metric (Gross and Perry (1983); Sorkin (1983); Davidson and Owen (1985)), which satisfies the 5D vacuum field equations, reduces to the standard four-dimensional (4D) Schwarzschild solution on hypersurfaces $`y=`$ const, and contains no explicit $`y`$-dependence. The assumption of a vacuum in 5D is consistent with the spirit of Kaluza’s idea, that 4D matter and gauge fields appear as a manifestation of pure geometry in the higher-dimensional world. The soliton metric has been generalized in various ways to incorporate time-dependence (Liu et al. (1993)), $`y`$-dependence (Billyard and Wesson (1996)) and electric charge (Liu and Wesson (1997)), among other things (eg, Wesson and Liu (1998)); see for review Overduin and Wesson (1997). We confine ourselves here to the original (two-parameter) soliton metric.
The motion of test bodies in the gravitational field of the soliton can be studied using the familiar classical tests of general relativity (gravitational redshift, light deflection, perihelion advance and time delay), along with the geodetic precession test. Work done so far along these lines (Lim, Overduin and Wesson 1995; Kalligas, Wesson and Everitt 1995, hereafter “KWE”) has demonstrated the existence of small but potentially measurable departures from the standard 4D Einstein predictions. In the present paper, we extend these earlier calculations in several ways, clarifying the physical meaning of the light deflection and time delay results for massless test particles and presenting new generalizations of the perihelion shift and geodetic precession formulas for massive ones. We take special care to compare our results to the latest experimental data in each case, obtaining new numerical constraints on the small parameter $`b`$ associated with the extra part of the soliton metric.
## 2 The Soliton Metric
In what follows, lowercase Greek indices $`\mu ,\nu \mathrm{}`$ will be taken to run over $`0,1,2,3`$ as usual, while capital Latin indices $`A,B,C\mathrm{}`$ run over all five coordinates ($`0,1,2,3,4`$). Units are such that $`G=c=1`$ except where stated otherwise. It is important to distinguish between the 4D line element ($`ds`$) and its 5D counterpart ($`dS`$), the two being related by
$$dS^{\mathrm{\hspace{0.17em}2}}=ds^2+g_{44}dy^2.$$
(1)
To interpret an expression containing $`d/dS`$ physically, one can always make the conversion
$$\frac{d}{dS}=\frac{ds}{dS}\frac{d}{ds}=\sqrt{1g_{44}\left(\frac{dy}{dS}\right)^2}\frac{d}{ds}.$$
(2)
We emphasize in particular that $`d/dSd/ds`$ if $`dy/dS0`$.
The soliton metric may be written (following Gross and Perry (1983), but switching to nonisotropic form, and defining $`a1/\alpha `$, $`b\beta /\alpha `$ and $`M2m`$)
$`dS^{\mathrm{\hspace{0.17em}2}}`$ $`=`$ $`A^adt^2A^{ab}dr^2A^{1ab}\times `$ (3)
$`r^2\left(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2\right)A^bdy^2,`$
where $`A(r)12M/r`$, $`M`$ is a parameter related to the mass of the object at the center of the geometry, and the constants $`a,b`$ satisfy a consistency relation
$$a^2+ab+b^2=1,$$
(4)
so that any two of $`M,a,b`$ may be taken as independent metric parameters. We will treat $`b`$ as the primary free parameter of the theory in what follows, noting that the 4D Schwarzschild metric is recovered (on hypersurfaces $`y=`$ const) in the limit $`b0`$ (and $`a+1`$). In general, larger values of $`|b|`$ will give rise to increasing departures from Einstein’s theory, subject to the upper bound $`|b|2/\sqrt{3}1.15`$ imposed by equation (4). Possible theoretical expectations for this parameter in the solar system and elsewhere are discussed further in §8.
## 3 Equation of Motion
We proceed now with the analysis of experimental constraints. The Lagrangian for a test particle in the field described by the metric (3) is
$``$ $`=`$ $`[A^a\dot{t}^2A^{ab}\dot{r}^2A^{1ab}\times `$ (5)
$`r^2(\dot{\theta }^2+\mathrm{sin}^2\theta \dot{\varphi }^2)A^b\dot{y}^2]^{1/2},`$
where the overdot represents differentiation with respect to an affine parameter $`\lambda `$ along the geodesics.
The Euler-Lagrange equations read
$$\frac{d}{d\lambda }\left(\frac{}{\dot{x}^C}\right)\frac{}{x^C}=0.$$
(6)
We confine ourselves to orbits with $`\theta =\pi /2`$ and $`\dot{\theta }=0`$, so that $``$ becomes
$$=\left(A^a\dot{t}^2A^{ab}\dot{r}^2A^{1ab}r^2\dot{\varphi }^2A^b\dot{y}^2\right)^{1/2}.$$
(7)
We can identify three constants of motion
$`\mathrm{}`$ $``$ $`{\displaystyle \frac{1}{}}A^a\dot{t}=A^a{\displaystyle \frac{dt}{dS}},`$
$`h`$ $``$ $`{\displaystyle \frac{1}{}}A^{1ab}r^2\dot{\varphi }=A^{1ab}r^2{\displaystyle \frac{d\varphi }{dS}},`$
$`k`$ $``$ $`{\displaystyle \frac{1}{}}A^b\dot{y}=A^b{\displaystyle \frac{dy}{dS}},`$ (8)
where we have used the relation $`=dS/d\lambda `$. From these equations we find that
$`\left({\displaystyle \frac{dr}{d\varphi }}\right)^2`$ $`+`$ $`Ar^2({\displaystyle \frac{\mathrm{}^{\mathrm{\hspace{0.17em}2}}}{h^2}}A^{22ab}{\displaystyle \frac{k^2}{h^2}}A^{2a2b}`$ (9)
$``$ $`{\displaystyle \frac{1}{h^2}}A^{2ab})r^4=0.`$
The derivation here differs slightly from that of KWE, where $`(dS/d\lambda )^2`$. Although the two approaches are physically equivalent, we have found that results are obtained more simply if the three constants of motion $`\mathrm{},h,k`$ are defined in terms of $`d/dS`$ rather than $`d/d\lambda `$ (or $`d/ds`$).
## 4 Light deflection
Experimental upper limits on possible violations of local Lorentz invariance are extremely tight (Will (1993)), so that we are justified in assuming that photons follow 4D null geodesics, $`ds=0`$. The situation is not so clear with regard to the 5D line element. However, it is economical to follow KWE and suppose that all particles follow ND null geodesics in N-dimensional gravity, whether massive or not.<sup>1</sup><sup>1</sup>1This assumption is supported by various lines of argument. In one version of 5D gravity, for example, the fifth coordinate $`y`$ is related to rest mass $`m`$ (Wesson (1984)), so that one has $`dS^{\mathrm{\hspace{0.17em}2}}=ds^2+g_{44}(G/c^2)^2dm^2`$. If all particles move on 5D null geodesics, then $`ds^2=g_{44}(G/c^2)^2dm^2`$. It then follows that $`ds=0`$ for photons, which have $`m=`$ const$`=0`$. For massive particles with $`ds0`$, one expects variations in rest mass $`m`$, which are however below currently detectable levels, owing to the small size of the dimension-transposing constant $`G/c^2`$ (Overduin and Wesson 1997, 1998). Recent work on incorporating non-relativistic quantum theory into higher-dimensional gravity also strongly suggests that all test particles travel on ND null geodesics in the classical limit (Seahra (2000)). Proceeding on this assumption, and substituting $`ds=dS=0`$ into equation (1), we get $`dy=0`$ also, so that $`\mathrm{},h\mathrm{}`$ and $`k`$ is undefined. The ratios $`\mathrm{}/h`$ and $`k/h`$ are however well-behaved, and read
$`{\displaystyle \frac{\mathrm{}}{h}}`$ $`=`$ $`r^2A^{2a+b1}{\displaystyle \frac{dt}{d\varphi }}=\text{ finite},`$
$`{\displaystyle \frac{k}{h}}`$ $`=`$ $`r^2A^{a+2b1}{\displaystyle \frac{dy}{d\varphi }}=0.`$ (10)
For self-consistency, therefore, the terms in $`k/h`$ can be dropped from KWE equations (7),(8),(11) and (12). Equation (8) of that paper, in particular, reduces to
$$\left(\frac{du}{d\varphi }\right)^2+Au^2\frac{\mathrm{}^{\mathrm{\hspace{0.17em}2}}}{h^2}A^{22ab}=0,$$
(11)
and the definition of the parameter $`p`$, KWE equation (12), becomes just
$$p(22ab)\frac{\mathrm{}^{\mathrm{\hspace{0.17em}2}}}{h^2}.$$
(12)
The photon’s trajectory is deflected by an angle
$$\delta \varphi =\omega =\frac{4M}{r_\mathrm{o}}+2Mpr_\mathrm{o},$$
(13)
which agrees with KWE equation (18.1). At the closest approach to the central body, we have $`u=u_o=1/r_\mathrm{o}`$ and $`du/d\varphi =0`$, so that equation (11) gives
$$\frac{\mathrm{}^{\mathrm{\hspace{0.17em}2}}}{h^2}=A_\mathrm{o}^{2a+b1}u_o^2=\frac{1}{r_\mathrm{o}^2}+O(\epsilon ),$$
(14)
where $`\epsilon M`$ is a small parameter. Putting equations (12) and (14) into equation (13), we find for the final light deflection angle
$$\delta \varphi =(4a+2b)\frac{M}{r_\mathrm{o}}+O(\epsilon ^2),$$
(15)
as in KWE equation (18.2), where however it is presented as a special case $`k=0`$. We see here that equation (15) is in fact entirely general for light deflection, and does not depend on any choice of $`k`$, which is in any case undefined when $`ds=dS=0`$.
To obtain experimental constraints from the light deflection result, let us express equation (15) for the Sun in terms of the deviation $`\mathrm{\Delta }_{\mathrm{LD}}`$ from the general relativity prediction $`\delta \varphi _{\mathrm{GR}}`$, as follows
$$\delta \varphi =\delta \varphi _{\mathrm{GR}}(1+\mathrm{\Delta }_{\mathrm{LD}}),$$
(16)
where (to first order in $`\epsilon `$):
$`\delta \varphi _{\mathrm{GR}}`$ $``$ $`4M_{}/r_\mathrm{o},`$
$`\mathrm{\Delta }_{\mathrm{LD}}`$ $``$ $`a+b/21.`$ (17)
Using the consistency relation (4) we find
$$a=b/2\pm (13b^2/4)^{1/2}.$$
(18)
Theoretical and numerical work indicates that $`|b|1`$ in the solar system (§8), and our experimental limits bear this out. The negative roots of equation (18) may also be ignored, as they are inconsistent with the limiting Schwarzschild case, and also imply the possibility of negative gravitational and/or inertial soliton mass (Gross and Perry (1983); Lim et al. (1995); Overduin and Wesson (1997)). We therefore take
$$a=1b/23b^2/8+O(b^4),$$
(19)
in the solar system, whereupon equation (17) gives
$$\mathrm{\Delta }_{\mathrm{LD}}=3b^2/8+O(b^4).$$
(20)
The best available constraints on $`\mathrm{\Delta }_{\mathrm{LD}}`$ come from long-baseline radio interferometry, which implies that $`|\mathrm{\Delta }_{\mathrm{LD}}|0.0017`$ (Robertson et al. (1991); Lebach et al. (1995)). We therefore infer an upper limit
$$|b|0.07,$$
(21)
for the Sun. This could potentially be tightened by more than an order of magnitude using a proposed astrometric optical interferometer sensitive to departures from Einstein’s theory of as little as $`|\mathrm{\Delta }_{\mathrm{LD}}|10^5`$ (Reasenberg and Shapiro (1986)).
It is important to bear in mind, however, that the parameter $`b`$ characterizing the soliton metric (3) is not a universal constant of nature like $`G`$ or $`c`$, but may in principle vary from soliton to soliton. Kaluza-Klein gravity as an alternative to 4D general relativity is therefore best constrained by the application of two or more tests to the same system. With this in mind we can use a recent measurement of light deflection by Jupiter, for which $`|\mathrm{\Delta }_{\mathrm{LD}}|0.17`$ (Treuhaft and Lowe (1991)), to obtain
$$|b|0.7,$$
(22)
for that planet. It has also been proposed to measure light deflection by the Earth using the Hipparcos satellite, with an estimated precision of 12% (Gould (1993)). Such a test would be sensitive to values of $`|b|0.6`$ for the Earth. The Gravity Probe B satellite should also be able to detect this effect by means of its guide star telescope, though with a somewhat lower precision (Adler (2000)).
## 5 Time delay
The arguments in the previous section regarding the parameter $`k`$ also apply to the time delay (or radar ranging) test, and circular photon orbits as well. That is, terms in $`k/h`$ and $`k/\mathrm{}`$ may be dropped from KWE equations (20-24) for radar ranging, and KWE equations (28-30) for circular orbits. The final results given there, however — equations (25) and (31) respectively — are correct. In fact, they hold not only for the special case $`k=0`$, but quite generally.
In particular, the excess round-trip time delay $`\mathrm{\Delta }`$$`\tau `$ for signals emitted from Earth (at distance $`r_\mathrm{e}`$ from the Sun) which graze the Sun (at nearest distance $`r_\mathrm{o}`$) and bounce off another planet (at $`r_\mathrm{p}`$) may be calculated by setting $`k/\mathrm{}`$=0 in KWE equation (24) to obtain
$$\mathrm{\Delta }\tau =\mathrm{\Delta }\tau _{\mathrm{GR}}(1+\mathrm{\Delta }_{\mathrm{TD}}),$$
(23)
where (to first order in $`\epsilon `$)
$`\mathrm{\Delta }\tau _{\mathrm{GR}}`$ $``$ $`4M_{}[\mathrm{ln}\left({\displaystyle \frac{r_\mathrm{p}+\sqrt{r_\mathrm{p}^2r_\mathrm{o}^2}}{r_\mathrm{o}}}\right)`$
$`+\mathrm{ln}\left({\displaystyle \frac{r_\mathrm{e}+\sqrt{r_\mathrm{e}^2r_\mathrm{o}^2}}{r_\mathrm{o}}}\right)],`$
$`\mathrm{\Delta }_{\mathrm{TD}}`$ $``$ $`a+b/21`$ (24)
$`=`$ $`3b^2/8+O(b^4).`$
We note that departures from 4D general relativity for time delay have exactly the same form as they do for light deflection.
The best experimental constraint on time delay so far has come from the Viking lander on Mars, and gives $`|\mathrm{\Delta }_{\mathrm{TD}}|0.002`$ (Reasenberg et al. (1979)). This leads immediately to the upper bound
$$|b|0.07,$$
(25)
for the Sun, exactly the same as the limit obtained in the case of light deflection using long-baseline interferometry.
Keeping in mind that values of $`b`$ can differ from soliton to soliton, however, it is possible that different physical setups could provide new information. For instance, one could attempt to measure $`b`$ for the Earth by sending grazing signals from an orbiting satellite past our planet and bouncing them off the Moon; retroflectors left there by Apollo astronauts are routinely used for lunar laser ranging (Williams et al. (1996)). Substituting $`M_\mathrm{e}`$ for $`M_{}`$ and replacing $`r_\mathrm{e},r_\mathrm{p}`$ and $`r_\mathrm{o}`$ with the appropriate distances, we find an expected excess time delay of order 400 ps using a satellite in geostationary orbit. This is well above the currently available resolution of $`50`$ ps (Samain et al. (1998)). The feasibility of such a proposal would likely be limited by the weakness of the reflected signal. Better results might be obtained by active ranging between two orbiting satellites, or by statistical analysis of ranging data between two such satellites and an Earth station (the latter would however require excellent atmospheric modelling).
In the same vein, one might attempt to measure $`b`$ for the Moon by grazing it with signals from the Earth and bouncing them off the Viking lander on Mars. This might be done when Mars is at nearest approach (on the same side of the Sun as the Earth) to minimize signal contamination from the competing effect of the Sun. Substituting $`M_\mathrm{m}`$ for $`M_{}`$ in equation (23), however, and replacing $`r_\mathrm{e},r_\mathrm{p}`$ and $`r_\mathrm{o}`$ with the appropriate distances, we find that the Moon’s excess time delay (of order 10 ps) would be so short as to make this a daunting task at present.
## 6 Perihelion advance
We now switch our attention to massive test particles. In terms of a new variable $`u1/r`$, equation (9) becomes
$`\left({\displaystyle \frac{du}{d\varphi }}\right)^2`$ $`+`$ $`Au^2({\displaystyle \frac{\mathrm{}^{\mathrm{\hspace{0.17em}2}}}{h^2}}A^{22ab}{\displaystyle \frac{k^2}{h^2}}A^{2a2b}`$ (26)
$``$ $`{\displaystyle \frac{1}{h^2}}A^{2ab})=0.`$
Differentiating with respect to $`\varphi `$ (and letting primes denote $`d/d\varphi `$), we find that noncircular orbits ($`u^{}0`$) are governed by the following differential equation
$$u^{\prime \prime }+(1+\gamma ϵ)u=B+ϵB^1u^2+O\left(ϵ^2\right),$$
(27)
where five new quantities have been introduced
$`\gamma `$ $``$ $`{\displaystyle \frac{f}{3d}},ϵ3MB,B{\displaystyle \frac{Md}{h^2}},`$
$`d`$ $``$ $`(2ab)\mathrm{}^{\mathrm{\hspace{0.17em}2}}(22ab)`$
$`+k^2(2a2b),`$
$`f`$ $``$ $`2(2ab)(1+a+b)`$ (28)
$`+\mathrm{\hspace{0.17em}2}\mathrm{}^{\mathrm{\hspace{0.17em}2}}(2+2a+b)(1+2a+b)`$
$`+\mathrm{\hspace{0.17em}2}k^2(2a2b)(1+a+2b).`$
These expressions agree with KWE equations (32-36). (We have however chosen to relabel their $`e`$ as $`f`$, for reasons that will become clear shortly.)
The solution of the differential equation (27) is
$`u={\displaystyle \frac{1}{r}}=B`$ $`+`$ $`\left(1{\displaystyle \frac{\gamma }{2}}\right)C\mathrm{cos}\left\{\left[1ϵ\left(1{\displaystyle \frac{\gamma }{2}}\right)\right]\varphi \right\}`$ (29)
$`+`$ $`ϵ(1\gamma )B+ϵ{\displaystyle \frac{C^{\mathrm{\hspace{0.17em}2}}}{2B}}\left(1{\displaystyle \frac{\gamma }{2}}\right)^2`$
$``$ $`ϵ{\displaystyle \frac{C^{\mathrm{\hspace{0.17em}2}}}{6B}}\left(1{\displaystyle \frac{\gamma }{2}}\right)^2\mathrm{cos}2\varphi +O(ϵ^2),`$
where $`C`$ is an integration constant. \[This result differs slightly from KWE equation (37), where the factors of $`(1\gamma /2)^2`$ were omitted.\] Equation (29) can be written in a physically more transparent form by introducing two new quantities $`e`$ and $`\omega `$ via
$$\left(1\frac{\gamma }{2}\right)CBe,\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1}ϵ\left(1\frac{\gamma }{2}\right)\omega .$$
(30)
With these definitions, we find that
$`u={\displaystyle \frac{1}{r}}=B(1`$ $`+`$ $`e\mathrm{cos}\omega \varphi )+{\displaystyle \frac{1}{2}}\epsilon B^2[e^2\mathrm{cos}2\omega \varphi `$ (31)
$`+`$ $`6(1\gamma +{\displaystyle \frac{1}{2}}e^2)]+O(\epsilon ^2),`$
where $`\epsilon M`$ is a small parameter as before, and
$$\omega =13\epsilon B\left(1+\frac{f}{6d}\right)+O(\epsilon ^2).$$
(32)
The first term on the right-hand side of equation (31) is of order $`\epsilon ^0`$, and shows explicitly the elliptical shape of the orbit. This is then modified by the second term, of order $`\epsilon ^1`$. Note that $`e`$ is just the eccentricity of the ellipse. The angular shift between two successive perihelia is given by
$$\delta \varphi =\varphi 2\pi =\frac{6\pi M^2d}{h^2}\left(1+\frac{f}{6d}\right)+O(\epsilon ^2),$$
(33)
in agreement with the final result (38.1) of KWE. It should be emphasized that the angular momentum $`h`$ is not in general the same quantity in 5D as it is in 4D. In particular, putting equations (2) and (3) into the second of equations (8), we find
$`h`$ $`=`$ $`A^{1ab}r^2{\displaystyle \frac{d\varphi }{dS}}=A^{1ab}r^2{\displaystyle \frac{d\varphi }{ds}}\sqrt{1+A^bk^2}`$ (34)
$`=`$ $`h_{\left(4\mathrm{D}\right)}\sqrt{1+A^bk^2}.`$
If $`k0`$, therefore, it follows that $`hh_{\left(4\mathrm{D}\right)}`$.
To eliminate $`h`$ from equation (33), let us consider the points along the orbit where $`r`$ takes its minimum value $`r_{}`$ and maximum value $`r_+`$ respectively. From inspection of equation (31) we see that $`r_{}=B^1(1+e)^1+O(\epsilon )`$ at $`\omega \varphi =0`$ and $`r_+=B^1(1e)^1+O(\epsilon )`$ at $`\omega \varphi =\pi `$. The semimajor axis $`a_\mathrm{o}`$ of the ellipse is then
$$a_\mathrm{o}\frac{1}{2}(r_{}+r_+)=\frac{1}{B(1e^2)}+O(\epsilon ),$$
(35)
so that
$$B\frac{Md}{h^2}=\frac{1}{a_\mathrm{o}(1e^2)}+O(\epsilon ),$$
(36)
or
$$h^2=\epsilon (1e^2)a_\mathrm{o}d+O(\epsilon ^2).$$
(37)
Substituting equation (36) into equation (33), we find
$$\delta \varphi =\frac{6\pi M}{a_\mathrm{o}(1e^2)}\left(1+\frac{f}{6d}\right)+O(\epsilon ^2).$$
(38)
Only one term in this result remains physically obscure, and that is the ratio $`f/d`$. This is given in terms of $`\mathrm{}`$ and $`k`$ by the definitions (28). The latter two constants are related by equation (26) as follows
$$\mathrm{}^{\mathrm{\hspace{0.17em}2}}=h^2\left[\left(\frac{du}{d\varphi }\right)^2+Au^2\right]A^{2a+b2}+k^2A^{ab}+A^a.$$
(39)
Since $`h^2`$ is of order $`\epsilon ^1`$ by equation (37), while $`u`$ and $`u^{}`$ are of order $`\epsilon ^0`$ by equation (31), it follows from equation (39) that $`\mathrm{}^{\mathrm{\hspace{0.17em}2}}=1+k^2+O(\epsilon )`$. Using the definitions (28), we therefore obtain
$$\frac{f}{6d}=1+a+\frac{2b}{3}+\frac{k^2b\left(ab\right)/3}{a+k^2\left(ab\right)}+O(\epsilon ),$$
(40)
so that the final perihelion precession angle (38) becomes
$$\delta \varphi =\frac{6\pi M}{a_\mathrm{o}(1e^2)}\left[a+\frac{2}{3}b+\frac{k^2(ab)b/3}{a+k^2(ab)}\right]+O(\epsilon ^2).$$
(41)
This represents the generalization of KWE equation (38.2) to cases in which $`k0`$ (and eccentricity $`e0`$). In the special case $`b=0`$ (and $`a=+1`$), for which the metric (3) reduces to Schwarzschild form on hypersurfaces $`y=`$ const, it is interesting to note that one recovers the standard 4D general relativity result, regardless of the value of $`k`$. In this limit, therefore, the perihelion shift test is insensitive to the momentum of the test body along the extra coordinate. And in general, one must choose a soliton with $`b0`$ in order to distinguish experimentally between test particles with different values of $`k`$.
As usual, let us parametrize our result in terms of the departure from 4D general relativity so that
$$\delta \varphi =\delta \varphi _{\mathrm{GR}}(1+\mathrm{\Delta }_{\mathrm{PP}}),$$
(42)
where (to first order in $`\epsilon `$):
$`\delta \varphi _{\mathrm{GR}}`$ $``$ $`{\displaystyle \frac{6\pi M}{a_\mathrm{o}(1e^2)}},`$
$`\mathrm{\Delta }_{\mathrm{PP}}`$ $``$ $`a+{\displaystyle \frac{2}{3}}b+{\displaystyle \frac{k^2(ab)b/3}{a+k^2(ab)}}1.`$ (43)
Theoretical work indicates that $`k`$, which is a measure of momentum along the fifth dimension, is related to the charge-to-mass ratio of the test body (Wesson and Liu (1997)). For a planet such as Mercury, we may take $`k=0`$. Putting equation (19) into equation (43), we therefore have
$$\mathrm{\Delta }_{\mathrm{PP}}=b/63b^2/8+O(b^4).$$
(44)
Perihelion precession is thus a potentially more sensitive probe of higher-dimensional gravity than either light deflection or time delay, in that it depends on the first, as well second order in $`b`$.
Unfortunately, however, this increased sensitivity is offset in the case of Mercury’s orbit about the Sun by uncertainty in the solar oblateness. The latter introduces a new term $`\xi J_2`$ inside the brackets on the right-hand side of equation (42), where $`\xi R_{}^2/2M_{}a_\mathrm{o}(1e^2)`$ and $`J_2`$ is the solar quadrupole moment (Campbell et al. (1983)). Dividing through by the orbital period $`T`$, we may therefore write for the rate of perihelion advance (to order $`b^{\mathrm{\hspace{0.17em}3}}`$)
$$\mathrm{\Delta }\omega \frac{\delta \varphi }{T}=\mathrm{\Delta }\omega _{\mathrm{GR}}(1+\xi J_2+b/63b^2/8),$$
(45)
where $`\mathrm{\Delta }\omega _{\mathrm{GR}}\delta \varphi _{\mathrm{GR}}/T=42.98`$ arcsec/century. The observed value of Mercury’s perihelion precession rate is quite close to this value, $`\mathrm{\Delta }\omega =43.11\pm 0.21`$ arcsec per century (Shapiro, Counselman and King 1976). Experimental data on $`J_2`$ is a good deal more controversial and has ranged over two orders of magnitude, from a maximum value of $`(23.7\pm 2.3)\times 10^6`$ (Dicke and Goldenberg (1967)) to a minimum of $`(0.17\pm 0.02)\times 10^6`$ (Duvall et al. (1984)). One straightforward least-sqares fit to a number of published measurements leads to intermediate value of $`J_2=5.0\times 10^6`$, which however implies a general relativistic precession rate more than two standard deviations away from that observed (Campbell et al. (1983)). Such a discrepancy could be explained in the context of higher-dimensional gravity by modelling the Sun as a soliton with $`b=0.062`$. This is just consistent with the constraint $`|b|0.07`$ from light deflection (§4) and time delay (§5), which is intriguing since these tests probe somewhat independent aspects of relativistic gravity. Improved experimental data relating to any of the three tests would be of great interest.
Conservative limits on $`b`$ from perihelion precession may be obtained by quoting the results of a recent review in which all available data (to 1997) have been combined to give a weighted mean value for the solar oblateness of $`J_2=(3.64\pm 2.84)\times 10^6`$ (Rozelot and Rösch (1997)). Using this uncertainty range, together with that in the observed value of $`\mathrm{\Delta }\omega `$ for Mercury’s orbit, we find that
$$b=0.03\pm 0.07,$$
(46)
for the Sun. This is consistent with the bounds obtained from light deflection and time delay. Sensitivity of the perihelion precession test to the value of $`b`$ could be improved by an order of magnitude if better data on $`J_2`$ were to become available; the proposed ASTROD mission, for example, might measure this parameter to an accuracy of $`5\times 10^8`$ (Ni (1998)).
## 7 Geodetic effect
We now move on to consider spinning massive test particles with velocity 5-vectors $`u^Cdx^C/dS`$ and spin 5-vectors $`S^C`$. The motion of these objects is governed by three central equations; namely, the geodesic equation
$$\frac{d^2x^C}{dS^{\mathrm{\hspace{0.17em}2}}}+\widehat{\mathrm{\Gamma }}_{AB}^Cu^Au^B=0,$$
(47)
the parallel transport equation
$$\frac{dS^C}{dS}+\widehat{\mathrm{\Gamma }}_{AB}^CS^Au^B=0,$$
(48)
and the orthogonality condition
$$u^CS_C=0.$$
(49)
Here $`\widehat{\mathrm{\Gamma }}_{AB}^C`$ refers to the 5D Christoffel symbol for the metric (3). This is defined in exactly the same manner as the usual 4D Christoffel symbol, with indices running over five values instead of four (see KWE, Appendix A1 for details<sup>2</sup><sup>2</sup>2There are some minor typographical errors in this appendix, which we note briefly here. The factors of $`(12M)/r`$ in equations (A2.2), (A2.6) and (A2.7) should read $`12M/r`$. The same thing applies to equations (57) and (58) in the main body of KWE. Also, the exponents $`(1/2)`$ and $`1/2`$ in equations (57) and (58) should be switched, in agreement with equations (A2.7) and (A2.2) respectively. These discrepancies do not affect any of the other equations or conclusions reported in KWE, and do not appear in the new reference book on Kaluza-Klein gravity by Wesson (1999).).
In order to simplify the problem, we follow KWE in assuming that the test particle moves in a circular orbit with $`\theta =\pi /2`$, $`r=r_\mathrm{o}`$ and $`\dot{\theta }=\dot{r}=0`$. Its velocity $`u^C`$ may then be expressed as follows in terms of the constants of motion $`\mathrm{},h`$ and $`k`$, as given by equations (8)
$$u^C\frac{dx^C}{dS}=(\mathrm{}A^a,\mathrm{\hspace{0.33em}0},\mathrm{\hspace{0.33em}0},hr_\mathrm{o}^2A^{a+b1},kA^b).$$
(50)
From the metric (3), we have
$$1=A^a\left(u^0\right)^2A^{1ab}r_\mathrm{o}^2\left(u^3\right)^2A^b\left(u^4\right)^2,$$
(51)
which, with equation (50), implies
$$\mathrm{}^{\mathrm{\hspace{0.17em}2}}h^2r_\mathrm{o}^2A^{2a+b1}k^2A^{ab}A^a=0.$$
(52)
It may be shown that the motion of the test body as given by equations (50) and (52) is geodesic in the sense of equation (47).
We now propose to generalize the treatment of KWE by leaving the extra component $`S^{\mathrm{\hspace{0.17em}4}}`$ of spin unrestricted, rather than setting it to zero. In fact, writing explicitly $`S^C(S^{\mathrm{\hspace{0.17em}0}},S^{\mathrm{\hspace{0.17em}1}},S^{\mathrm{\hspace{0.17em}2}},S^{\mathrm{\hspace{0.17em}3}},S^{\mathrm{\hspace{0.17em}4}})`$, we find that the orthogonality condition (49) imposes the following restriction on the spin components
$$\mathrm{}S^{\mathrm{\hspace{0.17em}0}}hS^{\mathrm{\hspace{0.17em}3}}kS^{\mathrm{\hspace{0.17em}4}}=0,$$
(53)
so that $`S^{\mathrm{\hspace{0.17em}4}}`$ will not vanish in general, if the parameter $`k`$ is well-defined.
We now proceed to solve the parallel transport equation (48), taking one value of the index $`C`$ at a time. To begin with, the $`C=2`$ component gives
$$S^{\mathrm{\hspace{0.17em}2}}=\frac{H_2}{r_\mathrm{o}}=\text{ const},H_2=\text{ const}.$$
(54)
(Note that, due to our choice of coordinates, $`S^{\mathrm{\hspace{0.17em}0}},S^{\mathrm{\hspace{0.17em}1}}`$ and $`S^{\mathrm{\hspace{0.17em}4}}`$ are dimensionless while $`S^{\mathrm{\hspace{0.17em}2}}`$ and $`S^{\mathrm{\hspace{0.17em}3}}`$ have units of inverse length.) Defining a new function $`g=g(S)`$ of the 5D proper time, we may write without loss of generality
$$S^{\mathrm{\hspace{0.17em}0}}H_0g,H_0=\text{ const}.$$
(55)
The $`C=0`$ component of equation (48) then reads
$$S^{\mathrm{\hspace{0.17em}1}}=\frac{H_0}{a\mathrm{}M}r_\mathrm{o}^2A^{a+1}\frac{dg}{dS}.$$
(56)
The $`C=4`$ component, meanwhile, takes the form
$`S^{\mathrm{\hspace{0.17em}4}}`$ $`=`$ $`H_4g+K_4,`$
$`H_4`$ $`=`$ $`{\displaystyle \frac{bk}{a\mathrm{}}}H_0A^{ab}=\text{ const},`$
$`K_4`$ $`=`$ $`\text{ const},`$ (57)
where we have used equation (56). In a similar way, the $`C=3`$ component of equation (48) gives
$`S^{\mathrm{\hspace{0.17em}3}}`$ $`=`$ $`{\displaystyle \frac{H_3}{r_\mathrm{o}}}g+K_3,`$
$`H_3`$ $`=`$ $`{\displaystyle \frac{hH_0}{a\mathrm{}M}}A^{2a+b1}\left[1(1+a+b){\displaystyle \frac{M}{r_\mathrm{o}}}\right]=\text{ const},`$
$`K_3`$ $`=`$ $`\text{ const}.`$ (58)
We now solve the $`C=1`$ component of equation (48), assuming for simplicity that $`K_3=K_4=0`$. Using equations (50), (55), (57) and (58), we find
$`{\displaystyle \frac{dS^{\mathrm{\hspace{0.17em}1}}}{dS}}`$ $`=`$ $`{\displaystyle \frac{H_0M}{a\mathrm{}r_\mathrm{o}^2}}A^{a+b1}\{a^2\mathrm{}^{\mathrm{\hspace{0.17em}2}}b^2k^2A^{ab}{\displaystyle \frac{}{}}`$ (59)
$``$ $`{\displaystyle \frac{h^2}{M^2}}[1(1+a+b){\displaystyle \frac{M}{r_\mathrm{o}}}]^2A^{2a+b1}\}g.`$
Differentiating equation (56) with respect to $`S`$, meanwhile, gives
$$\frac{dS^{\mathrm{\hspace{0.17em}1}}}{dS}=\frac{H_0}{a\mathrm{}M}r_\mathrm{o}^2A^{a+1}\frac{d^2g}{dS^{\mathrm{\hspace{0.17em}2}}}.$$
(60)
Equating these two expressions, we obtain
$$\frac{d^2g}{dS^{\mathrm{\hspace{0.17em}2}}}=\mathrm{\Omega }^2g,$$
(61)
where
$`\mathrm{\Omega }^2`$ $``$ $`{\displaystyle \frac{h^2}{r_\mathrm{o}^4}}A^{b2}\{[1(1+a+b){\displaystyle \frac{M}{r_\mathrm{o}}}]^2A^{2a+b1}`$ (62)
$`{\displaystyle \frac{M^2}{h^2}}(a^2\mathrm{}^{\mathrm{\hspace{0.17em}2}}b^2k^2A^{ab})\}.`$
The general solution of equation (61) is $`g(S)=K_1\mathrm{sin}(\mathrm{\Omega }S)+K_2\mathrm{cos}(\mathrm{\Omega }S)`$. We choose $`K_1=1`$ and $`K_2=0`$ for simplicity. The spin components are then given by
$`S^{\mathrm{\hspace{0.17em}0}}`$ $`=`$ $`H_0\mathrm{sin}(\mathrm{\Omega }S),S^{\mathrm{\hspace{0.17em}1}}=H_1\mathrm{cos}(\mathrm{\Omega }S),`$
$`S^{\mathrm{\hspace{0.17em}2}}`$ $`=`$ $`{\displaystyle \frac{H_2}{r_\mathrm{o}}},S^{\mathrm{\hspace{0.17em}3}}={\displaystyle \frac{H_3}{r_\mathrm{o}}}\mathrm{sin}(\mathrm{\Omega }S),`$
$`S^{\mathrm{\hspace{0.17em}4}}`$ $`=`$ $`H_4\mathrm{sin}(\mathrm{\Omega }S),`$ (63)
where $`H_1`$ and $`H_2`$ are arbitrary constants and
$`H_0`$ $`=`$ $`{\displaystyle \frac{a\mathrm{}M}{r_\mathrm{o}^2\mathrm{\Omega }}}A^{a1}H_1,`$
$`H_3`$ $`=`$ $`{\displaystyle \frac{h}{r_\mathrm{o}^2\mathrm{\Omega }}}A^{a+b2}\left[1(1+a+b){\displaystyle \frac{M}{r_\mathrm{o}}}\right]H_1,`$
$`H_4`$ $`=`$ $`{\displaystyle \frac{bkM}{r_\mathrm{o}^2\mathrm{\Omega }}}A^{b1}H_1.`$ (64)
The spatial part of $`S^C`$ is thus seen to rotate in the plane of the orbit with angular speed $`\mathrm{\Omega }`$. Substituting these results into equation (53) yields
$`a\mathrm{}^{\mathrm{\hspace{0.17em}2}}`$ $``$ $`{\displaystyle \frac{h^2}{Mr_\mathrm{o}}}A^{2a+b1}\left[1(1+a+b){\displaystyle \frac{M}{r_\mathrm{o}}}\right]`$ (65)
$``$ $`bk^2A^{ab}=0.`$
Solving simultaneously with equation (52), we obtain for the constants of motion
$`\mathrm{}^{\mathrm{\hspace{0.17em}2}}`$ $`=`$ $`A^a\{1+k^2A^b+{\displaystyle \frac{M}{r_\mathrm{o}}}[{\displaystyle \frac{a+(ab)k^2A^b}{1(1+2a+b)M/r_\mathrm{o}}}]\},`$
$`h^2`$ $`=`$ $`Mr_\mathrm{o}A^{1ab}\left[{\displaystyle \frac{a+(ab)k^2A^b}{1\left(1+2a+b\right)M/r_\mathrm{o}}}\right].`$ (66)
These expressions can be written in terms of a small parameter $`\epsilon M`$ as usual
$`\mathrm{}^{\mathrm{\hspace{0.17em}2}}`$ $`=`$ $`(1+k^2)\left[1\left(a{\displaystyle \frac{bk^2}{1+k^2}}\right){\displaystyle \frac{M}{r_\mathrm{o}}}\right]+O(\epsilon ^2),`$
$`h^2`$ $`=`$ $`Mr_\mathrm{o}[a+(ab)k^2]\{1+[(4a+3b1)`$ (67)
$`+{\displaystyle \frac{2b(ab)k^2}{a+(ab)k^2}}]{\displaystyle \frac{M}{r_\mathrm{o}}}\}+O(\epsilon ^3).`$
With the aid of equation (62), we then find for the angular speed of the spin vector
$`\mathrm{\Omega }`$ $`=`$ $`\sqrt{{\displaystyle \frac{[a+(ab)k^2]M}{r_\mathrm{o}^3}}}\{1+{\displaystyle \frac{M}{2r_\mathrm{o}}}[3(1ab)`$ (68)
$`+{\displaystyle \frac{b(ab)k^2}{a+(ab)k^2}}]+O(\epsilon ^2)\}.`$
This quantity is not the same as the test body’s orbital angular speed, which is given in terms of the 5D proper time $`dS`$ as
$`\omega `$ $``$ $`{\displaystyle \frac{d\varphi }{dS}}=hr_\mathrm{o}^2A^{a+b1}`$
$`=`$ $`\sqrt{{\displaystyle \frac{M}{r_\mathrm{o}^3}}}A^{(a+b1)/2}\sqrt{{\displaystyle \frac{a+(ab)k^2A^b}{1(1+2a+b)M/r_\mathrm{o}}}},`$
where we have used equations (8) and (66). In terms of $`\epsilon `$
$`\omega `$ $`=`$ $`\sqrt{{\displaystyle \frac{\left[a+(ab)k^2\right]M}{r_\mathrm{o}^3}}}\{[1+{\displaystyle \frac{M}{r_\mathrm{o}}}({\displaystyle \frac{3b}{2}}`$ (70)
$`+{\displaystyle \frac{b(ab)k^2}{a+(ab)k^2}})]+O(\epsilon ^2)\}.`$
It is precisely the excess of $`\mathrm{\Omega }`$ over $`\omega `$ that gives rise to the geodetic effect.
Suppose the spin vector $`S^C`$ is initially oriented in the radial direction; ie, $`H_2=0`$ at $`S=0`$. During one orbit, the test body’s angular displacement $`\varphi `$ goes from $`0`$ to $`2\pi `$, so that $`\delta S=2\pi /\omega `$. In the same period, $`S^{\mathrm{\hspace{0.17em}3}}`$ goes from its initial value of zero at $`S=0`$ to its final value at proper time $`S`$. To first order in $`\epsilon `$, the spin vector has advanced through an angle
$`\delta \varphi `$ $`=`$ $`{\displaystyle \frac{r_\mathrm{o}[S^{\mathrm{\hspace{0.17em}3}}(S)S^{\mathrm{\hspace{0.17em}3}}(0)]}{S^{\mathrm{\hspace{0.17em}1}}(0)}}+O(\epsilon ^2),`$ (71)
$`=`$ $`2\pi {\displaystyle \frac{H_3}{H_1}}\left({\displaystyle \frac{\mathrm{\Omega }}{\omega }}1\right)+O(\epsilon ^2),`$
$`=`$ $`2\pi \left({\displaystyle \frac{\mathrm{\Omega }}{\omega }}1\right)+O(\epsilon ^2),`$
where we have used equations (64), (67) and (68). Combining equations (68) and (70), we find that
$$\frac{\mathrm{\Omega }}{\omega }1=\frac{3M}{2r_\mathrm{o}}\left[a+\frac{2}{3}b+\frac{k^2b(ab)/3}{a+(ab)k^2}\right]+O(\epsilon ^2),$$
(72)
so that the geodetic precession angle can finally be expressed as follows in terms of its deviation from the prediction of 4D general relativity
$$\delta \varphi =\delta \varphi _{\mathrm{GR}}(1+\mathrm{\Delta }_{\mathrm{GP}}),$$
(73)
where (to first order in $`\epsilon `$)
$`\delta \varphi _{\mathrm{GR}}`$ $``$ $`3\pi M/r_\mathrm{o},`$
$`\mathrm{\Delta }_{\mathrm{GP}}`$ $``$ $`a+{\displaystyle \frac{2}{3}}b+{\displaystyle \frac{k^2(ab)b/3}{a+k^2(ab)}}1.`$ (74)
This represents the generalization of KWE equation (66) to cases in which $`S^{\mathrm{\hspace{0.17em}4}}0`$. Deviations from 4D general relativity have exactly the same form for geodetic precession as they do for perihelion precession. Taking $`k=0`$ and using equation (19), as in §6, we find that
$$\mathrm{\Delta }_{\mathrm{GP}}=b/63b^2/8+O(b^4).$$
(75)
Like the perihelion shift, geodetic precession depends on $`b`$ to first as well as second order, and is thus a potentially more sensitive probe of the theory than either light deflection or time delay.
The Gravity Probe B satellite, currently scheduled for launch in early 2001, has been designed to measure deviations from 4D general relativity with a precision of better than $`|\mathrm{\Delta }_{\mathrm{GP}}|2.5\times 10^5`$ (Buchman et al. (1996)). Using equation (75), we find that this corresponds to a sensitivity to values as small as
$$|b|1\times 10^4,$$
(76)
or better for the Earth — a constraint some five hundred times stronger than any other solar system bound obtained to date, and five thousand times stronger than the only other Earth-based test (light deflection using Hipparcos; §4).
We conclude this section by noting that a complementary analysis of geodetic precession has been carried out for a static, spherically-symmetric 5D metric different from that given by equation (3), one in which the fifth dimension is flat (Mashhoon, Liu and Wesson 1994; Mashhoon, Wesson and Liu 1998). The inclusion of spin is of special importance in this case since the classical tests (based on the equations of motion) alone cannot discriminate between 4D and 5D effects. The geodetic precession rate has been computed, and differs from the 4D Einstein value in the weak-field, low velocity limit (Liu and Wesson (1996)). A preliminary interpretation of the discrepancy indicates, however, that it is likely to be somewhat below the threshold of detection by Gravity Probe B (Overduin and Wesson (1998)).
## 8 Discussion
Having obtained upper limits on $`|b|`$ of order $`0.07`$ (and possibly $`10^4`$) from experiment, we consider here the range of values that might be expected for this parameter on theoretical grounds. These turn out to be small (perhaps of order $`10^8`$ to $`10^2`$) in the solar system, but could be larger (of order $`0.1`$) in larger systems such as clusters of galaxies.
These estimates are based on the fact that the soliton’s effective 4D mass is not concentrated at a point, like that of a black hole, but has instead a finite (though sharply peaked) density profile whose steepness depends on the metric parameters (Liu and Wesson (1992); Wesson and Ponce de Leon (1994)). Quoting the latter authors, but replacing their metric parameters $`\stackrel{~}{a},ϵ,k`$ (due to Davidson and Owen 1985) with our $`M,a,b`$ via $`M2/\stackrel{~}{a}`$, $`aϵk`$ and $`bϵ`$, we find for the density of the soliton
$$8\pi \rho (r)=\frac{abM^2/r^4}{\left[1(M/2r)^2\right]^4}\left(\frac{1M/2r}{1+M/2r}\right)^{2(a+b)}.$$
(77)
Pressure is given by $`p=\rho /3`$, so that the matter described by equation (77) could be radiationlike, or composed of ultrarelativistic particles such as neutrinos. Total gravitational mass (as deduced from the asymptotic form of the metric) is $`M_g=aM`$, so it is clear that $`b`$ must be negative for positive density. Numerical analysis further reveals that the mass of the soliton is increasingly concentrated at small $`r`$ as $`|b|`$ approaches zero, and that the 4D Schwarzschild limit ($`b=0`$) can in fact be viewed as a maximally compressed soliton (Wesson and Ponce de Leon (1994)). Physically, this means that solar system bodies, which (viewed as solitons) are essentially point masses, are likely to be associated with very small values of $`|b|`$.
To attach some numbers to these qualitative remarks, we make use of equation (19) and consider the weak-field ($`rM/2`$), small-$`b`$ limit, in which
$$\rho (r)bGM_g^2/(8\pi c^2r^4),$$
(78)
where we have reverted to physical units. Equation (78) allows us to associate ranges of $`b`$-values with solitons of mass $`M_g`$, if the density $`\rho `$ can be estimated at some radius $`r`$. It has, for instance, been suggested (eg, Freese (1986); Gould (1992)) that relativistic hot dark matter in the form of massive neutrinos could be trapped inside the Earth. Krauss et al. (1986) have derived one possible density profile for such particles, assuming that equilibrium is established between those undergoing capture, annhilation, and escape from the Earth’s gravitational potential. We do not attempt to fit our equation (78) to this profile at all radii, but merely take the predicted neutrino density at the Earth’s surface as illustrative. From Fig. 2 of Krauss et al. (1986), the expected escape rate for 10 GeV neutrinos is $`2\times 10^{16}`$ s<sup>-1</sup>, which translates into a density at the Earth’s surface of $`\rho (R_{})=3\times 10^{20}`$ kg m<sup>-3</sup> (about 50 times the canonical local halo dark matter density of $`5\times 10^{22}`$ kg m<sup>-3</sup>). If we suppose that this is rather associated with solitonic matter making up some fraction $`\zeta `$ of the Earth’s total mass ($`M_g=\zeta M_{}`$), then equation (78) gives $`b=4\times 10^{14}\zeta ^2`$. For dark matter of this kind to be significant, $`b`$ must be small for solar system bodies. With $`\zeta 10^3`$, for example, we have $`b4\times 10^8`$, while $`\zeta 10^6`$ would correspond to $`b0.04`$. These numbers are consistent with the experimental limits obtained in §§ 4 \- 7 above. It may be possible to constrain the theory more tightly by looking at violations of the weak equivalence principle by solar system bodies (Overduin (2000)).
On larger scales, systems such as galaxies and clusters of galaxies are suspected by many to harbor significant amounts of relativistic hot dark matter. We take here as an example a recent numerical simulation (Kofman et al. (1996)) in which light (2.3 eV) neutrinos make up 20% (by mass) of a cluster whose total mass $`M_T=6\times 10^{14}M_{}`$. Fig. 3 of this paper shows a typical neutrino density of $`\rho (r)200\rho _c`$ at $`r=0.03`$ Mpc, where $`\rho _c=2\times 10^{26}h_0^2`$ kg m<sup>-3</sup> is the critical density. If this were instead attributed to solitonic dark matter of total mass $`M_g=\zeta M_T`$, then the latter would have $`b=0.01\zeta ^2`$ by equation (78), where we have taken $`h_0=0.65`$. If all the hot dark matter were solitonic ($`\zeta =0.2`$), then $`|b|`$ could be as large as 0.3. These values are illustrative only, since density profiles of hot dark matter in clusters are likely somewhat shallower than that indicated by equation (78)<sup>1</sup><sup>1</sup>1 Density profiles with $`\rho r^4`$ at large $`r`$ have however been discussed in other contexts, such as elliptical galaxies (Jaffe (1983); de Zeeuw (1985); Hernquist (1990)).. Nevertheless they establish that values of $`|b|`$ in galaxy clusters might in principle be significantly larger than those in the solar system, and this encourages us to speculate that stronger tests of higher-dimensional gravity might be carried out using the excellent observational data now available on gravitational lensing by these objects.
## 9 Conclusions
We have re-examined the classical tests of general relativity, as well as the geodetic precession test, when Einstein’s theory is extended from four to five dimensions. The physical meaning of previous calculations for light deflection and time delay have been clarified physically, and the restriction of zero momentum and/or spin along the extra coordinate that characterized the earlier calculations of perihelion shift and geodetic precession has been lifted.
Our results show that Kaluza-Klein gravity remains consistent with experiment. The free parameter of the theory, however, is increasingly constrained to small values. Thus, data on light deflection, radar ranging to Mars and the perihelion precession of Mercury all imply a value of $`|b|0.07`$ for the Sun. Improved data on solar oblateness should improve the sensitivity of the perihelion precession bound by as much as an order of magnitude. And the upcoming launch of Gravity Probe B will allow us to measure values of $`|b|`$ for the Earth with an accuracy of one part in $`10^4`$ or better.
We thank P. S. Wesson for comments. H. L. acknowledges the support of the National Natural Science Foundation of China (grant no. 19975007). J. O. thanks R. I. Bush and W.-T. Ni for comments on solar oblateness and radar ranging, and acknowledges the support of the National Science and Engineering Research Council of Canada. He also expresses his gratitude for the hospitality of C. W. F. Everitt and the theory group at Gravity Probe B, Stanford University, where part of this work was carried out. |
warning/0003/math0003104.html | ar5iv | text | # Untitled Document
Relations among divisors on the moduli space of curves with marked points
Adam Logan
1. Introduction
Let $`\overline{}_g`$ be the coarse moduli space of stable curves of genus $`g`$. Eisenbud, Harris and Mumford proved a relation between certain divisors on $`\overline{}_g`$ (the Brill-Noether divisors, to be described below). Calculating their classes in $`Pic\overline{}_g𝐐`$, they succeeded in proving that $`\overline{}_g`$ is of general type for $`g>23,g+1`$ composite. In subsequent work, the restriction that $`g+1`$ be composite was removed.
In this paper, I will generalize their relation to $`\overline{}_{g,n}`$, the moduli space of stable curves of genus $`g`$ with $`n`$ marked points. This does not yield new results on the Kodaira dimension of $`\overline{}_{g,n}`$, as the “divisors of Brill-Noether type”, which I introduce below, are less effective for this purpose than certain other divisors which I studied in my dissertation \[L\]. (I expect to publish the other results of \[L\] shortly.)
The remainder of this introduction will be devoted to stating the results. For basic facts on $`\overline{}_g`$ and $`\overline{}_{g,n}`$, the reader is referred to \[HM\] and \[K\] respectively.
Definition. Fix a nonnegative integer $`g`$, and let $`r>1,d>1`$ be integers such that $`g(r+1)(gd+r)=1`$. (Here, the left side is the expected dimension of the space of $`g_d^r`$’s on a curve of genus $`g`$.) A divisor of Brill-Noether type on $`\overline{}_g`$ is a codimension-1 component of the locus of curves which have an admissible $`g_d^r`$.
(All $`g_d^r`$’s on nonsingular irreducible curves are admissible; in general, an admissible $`g_d^r`$ is a $`g_d^r`$ on each component with certain ramification conditions at the singular points of the curve. Again, see \[HM\] for details.)
The Brill-Noether Ray Theorem of Eisenbud, Harris, and Mumford then asserts that, for fixed $`g`$, all of these divisors on $`\overline{}_g`$ are linearly dependent, and calculates their class. We generalize their definition and theorem as follows:
Definition. Fix $`g`$ and $`n`$, and let integers $`r,d`$ and sequences $`\{Z_1\},\mathrm{},\{Z_n\}`$ of length $`r+1`$ be such that
$$g(r+1)(gd+r)Z_{i,j}=1.$$
(The left-hand side is the expected dimension of the set of $`g_d^r`$’s on a curve of genus $`g`$ with ramification sequences $`\{Z_1\},\mathrm{},\{Z_n\}`$ at points $`p_1,\mathrm{},p_n`$.) A divisor of Brill-Noether type on $`\overline{}_{g,n}`$ is a divisorial component of the locus of curves and sets of points such that there is a $`g_d^r`$ on the curve whose ramification sequences at the $`p_i`$ are at least $`\{Z_i\}`$.
Theorem 1.1 For any $`n>2`$, every divisor of Brill-Noether type on $`\overline{}_{g,n}`$ is a linear combination of pullbacks of divisors of Brill-Noether type from $`\overline{}_{g,2}`$. Also, the space of divisors of Brill-Noether type on $`\overline{}_{g,1}`$ has dimension $`2`$, for any $`g>2`$.
Corollary 1.2 The dimension of the subspace of $`Pic\overline{}_{g,n}𝐐`$ spanned by the divisors of Brill-Noether type is $`1+n+\left(\genfrac{}{}{0pt}{}{n}{2}\right)`$.
Acknowledgments. This paper is a lightly revised version of parts of my 1999 doctoral dissertation \[L\], written under the direction of Joe Harris, to whom I am very grateful. I would also like to thank Mira Bernstein for teaching me how to work with moduli spaces. Support for this work was provided by a Sloan Dissertation Fellowship.
2. The First Case
Throughout the paper, we work over an algebraically closed field of characteristic 0 (though this is surely unnecessary), and we deal only with genus $`3`$. We start by introducing some notation. Arbarello and Cornalba have given a basis for $`Pic\overline{}_{g,n}𝐐`$ \[AC, Thm. 2\]:
Theorem 2.1 $`Pic_{\text{fun}}\overline{}_{g,n}`$, the Picard group of the moduli stack $`\overline{}_{g,n}`$, is free on the following generators: $`\lambda ,\delta _0,\psi _i(1in),`$ and $`\delta _{i;S}(0ig/2),S\{1,2,\mathrm{},n\}),cardS>1\text{ if }i=0`$.
Here $`\lambda `$ is the pullback of $`\lambda `$ on $`\overline{}_g`$, $`\delta _0`$ the divisor corresponding to the locus of curves with a nondisconnecting node, $`\delta _{i,S}`$ the divisor corresponding to the locus of curves with a node whose removal leaves one component of genus $`i`$ with precisely the marked points indexed by $`S`$, and $`\psi _i`$ is the divisor class on the moduli stack which takes the value $`\pi _{}(\sigma _i^2)`$ on the family $`𝒳\stackrel{\pi }{}B`$ with the $`\sigma _i`$ as sections.
Usually we will prefer to work with $`\omega _i`$, the relative dualizing sheaf of the $`i`$th projection map from $`\overline{}_{g,n}`$ to $`\overline{}_{g,n1}`$. If we replace $`\psi _i`$ by $`\omega _i`$ in the above, the theorem remains true, because $`\omega _i=\psi _i_{iS}\delta _{0;S}`$, as will be seen below.
We will be pulling back divisors, so it also seems appropriate to state the results concerning this. Again, the answer is given in \[AC, p. 161\]:
Theorem 2.2 The pullback map is given as follows:
$$\begin{array}{cc}\hfill \pi _n^{}\lambda & =\lambda ,\hfill \\ \hfill \pi _n^{}\delta _0& =\delta _0,\hfill \\ \hfill \pi _n^{}\omega _i& =\omega _i,\hfill \\ \hfill \pi _n^{}\psi _i& =\psi _i\delta _{0;\{i,n\}}\hfill \\ \hfill \pi _n^{}\delta _{i;S}& =\delta _{i;S}+\delta _{i;S\{n\}},\hfill \end{array}$$
except that
$$\pi _1^{}\delta _{g/2;\mathrm{}}=\delta _{g/2;\mathrm{}}$$
for $`n=1`$.
For notational simplicity, we have only stated this for $`\pi _n`$, but the action of the symmetric group $`𝒮_n`$ on $`\overline{}_{g,n}`$ makes it easy to describe the effects of the other $`\pi `$. Also observe that, by an easy induction starting from $`\psi =\omega `$ on $`\overline{}_{g,1}`$, we get
$$\psi _i=\omega _i+\underset{iS\{1,2,\mathrm{},n\},S\{i\}}{}\delta _{0;S}$$
on $`\overline{}_{g,n}`$, as asserted above.
We now start by dealing with the case $`n=1`$.
Definition. Let $`BN`$ be the divisor class
$$(g+3)\lambda \frac{(g+1)}{6}\delta _0\underset{i=1}{\overset{g/2}{}}i(gi)\delta _i$$
on $`\overline{}_g`$.
If $`g+1`$ is not prime, then $`BN`$ is a positive multiple of the class of an effective divisor, namely that of the union of the codimension-$`1`$ components of the locus of curves admitting a $`g_d^r`$, where $`(r+1)(gd+r)=g+1`$ \[EH, Thm. 1\] (the condition on $`r,d`$, and $`g`$ ensures that there are such components).
This will be one of the generators of the Brill-Noether space on $`\overline{}_{g,1}`$ when $`g+1`$ is not prime; the other will be the Weierstrass divisor, that is, the locus of Weierstrass points. Its class was computed by Cukierman \[Cuk, Thm. 2.0.12 and following remark\]. His result is as follows:
Theorem 2.3 Let $`g2`$. The class on $`\overline{}_{g,1}`$ of the Weierstrass divisor $`𝒲`$, which is the closure of the divisor on $`_{g,1}`$ given by Weierstrass points of smooth curves, is
$$\frac{g(g+1)}{2}\omega \lambda \underset{i=1}{\overset{g1}{}}\frac{(gi)(gi+1)}{2}\delta _i.$$
Proof. See \[Cuk\]. Alternatively, this can be done by the method of test curves.
We recall the Plücker formula, which counts the ramification points of a $`g_d^r`$ on a smooth curve. It asserts that if $`C`$ is a smooth curve of genus $`g`$ and $`V`$ a $`g_d^r`$, then $`_{pC}\beta (V,p)=(r+1)(d+r(g1))`$. Here $`\beta (V,p)=_0^ra_i(V,p)i`$, where the $`a_i`$ give the sequence of orders of vanishing of $`V`$ at $`p`$.
Proposition 2.4 The classical Plücker formula remains valid for reducible curves with no nondisconnecting nodes. (Of course, “$`g_d^r`$” is understood to mean “limit linear series”. Also, the ramification conditions at the node imposed by the definition of “limit linear series” are not considered as contributing to ramification.)
Proof. By induction, it suffices to consider curves with exactly two components. Applying the Plücker formula to each component separately, we get a total of $`(r+1)(2d+r(g2))`$ there. However, the definition of limit linear series imposes exactly $`(r+1)(dr)`$ conditions at the node, leaving $`(r+1)(d+r(g1))`$ in total, as claimed.
Definition. Let $`r`$ and $`d`$ be positive integers such that $`a=g(r+1)(gd+r)1`$, and $`Z`$ a possible ramification sequence for a $`g_d^r`$ which sums to $`a`$. We define the divisor $`D_{g,r,d,Z}`$ on $`\overline{}_{g,1}`$ to be the divisor of curves and points $`(C,p)`$ such that $`C`$ admits a $`g_d^r`$, $`V`$, whose ramification sequence at the marked point equals or exceeds $`Z`$.
The $`D_{g,r,d}`$ are the divisors of Brill-Noether type on $`\overline{}_{g,1}`$. In order to study the $`D_{g,r,d}`$, we consider the following two maps to $`\overline{}_{g,1}`$:
1. Map $`\overline{}_{0,g+1}`$ to $`\overline{}_{g,1}`$ by attaching a fixed elliptic curve at each of the first $`g`$ points.
2. Map $`\overline{}_{2,1}𝒲`$ to $`\overline{}_{g,1}`$ by attaching a fixed general curve of genus $`g2`$ with two marked points.
We claim that the images of these maps are disjoint from $`D`$, and that this implies the linear dependence claimed. For the first of these maps, it is enough to count ramification points; there are not enough to spare on the component of genus $`0`$. In particular, at each point of attachment we must have $`\beta r`$ on this component, making at least $`gr`$. Even if the remaining $`(r+1)(dr)gr=g(r+1)(gd+r)=a`$ units of ramification are all concentrated at one point, that isn’t quite enough. For the second map, this follows from the extended Brill-Noether theorem, \[EH, Thm. 1.1\].
In order to deduce the dependence from this, it is necessary to determine the pullback map on divisors induced by the two maps given above. In both cases, most of the work takes care of itself, since these maps fit into commutative diagrams in which one arrow is the map in question, one arrow is a map for which the pullback is computed in \[EH, proof of Thms. 2.1, 3.1\], and the other arrow(s) are easily understood. Specifically, to study the first map, insert it into a commutative diagram as follows, in which the vertical projections are simply the forgetful maps as shown:
$$\begin{array}{ccc}\overline{}_{0,g+1}& \stackrel{f^{}}{}& \overline{}_{g,1}\\ \pi _{g+1}& & \pi _1\\ \overline{}_{0,g}& \stackrel{f}{}& \overline{}_g\end{array}$$
This diagram being commutative, the two pullback maps on $`Pic𝐐`$ must be equal. Since the pullback maps on the vertical arrows are as described in Theorem 2.2, and the map on the bottom arrow is given by Eisenbud and Harris, we can easily determine $`f^{}`$ of any class on $`\overline{}_{g,1}`$ pulled back from $`\overline{}_g`$. For example, since $`\pi _{g+1}^{}f^{}\lambda =0`$ (indeed, $`f^{}\lambda =0`$), as in \[EH\], and since $`\pi _1^{}\lambda =\lambda `$, it follows that $`f^{}\lambda =0`$, and likewise for $`\delta _0`$ in place of $`\lambda `$.
Definition. Let $`\theta _iPic\overline{}_{0,g+1}=_{ST}\delta _{0,S}`$, where $`T`$ runs over subsets of $`\{1,\mathrm{},g+1\}`$ of cardinality $`i+1`$ that contain $`g+1`$. Also, let $`ϵ_i=_{cardS=i}\delta _{0;S}.`$
It is easy to see that $`\theta _i+\theta _{gi}=\pi _{g+1}^{}ϵ_i.`$ (As usual, the case $`2i=g`$ is an exception: then the pullback is just $`\theta _i`$.) Since $`ϵ_i=f^{}\delta _i`$ for $`i>1`$, and since for $`i<g1`$ we have that $`f^{}\delta _i`$ is supported on $`\theta _i`$, it follows immediately that $`f^{}\delta _i=\theta _i`$ in this range.
On the other hand,
$$\pi _{g+1}^{}f^{}\delta _1=\underset{i=2}{\overset{g2}{}}\frac{i(gi)}{(g1)}\theta _i,$$
as follows from the calculation of $`f^{}\delta _1`$ in \[EH, Thm. 3.1\]. This, therefore, is $`f^{}(\delta _1+\delta _{g1})`$, and so
$$f^{}(\delta _{g1})=\underset{i=1}{\overset{g2}{}}\frac{i(gi)}{(g1)}\theta _i.$$
Finally it is necessary to compute $`f^{}\omega `$. The easiest way to do this is simply to use the fact that $`f^{}𝒲=0`$. Using Cukierman’s formula (Theorem 2.3), one translates this into the statement that
$$f^{}\omega =\underset{i=1}{\overset{g2}{}}\frac{(gi)(gi1)}{g(g1)}\theta _i.$$
Next we consider the second map, to which a similar procedure applies, complete with a similar commutative diagram:
$$\begin{array}{ccc}\overline{}_{2,1}& \stackrel{g^{}}{}& \overline{}_{g,1}\\ =& & \pi _1\\ \overline{}_{2,1}& \stackrel{g}{}& \overline{}_g\end{array}$$
This time, according to \[EH, Sect. 2\], we have
$$\begin{array}{cc}\hfill g^{}\delta _0& =\delta _0,\hfill \\ \hfill g^{}\delta _1& =\delta _1,\hfill \\ \hfill g^{}\delta _2& =\omega ,\hfill \\ \hfill g^{}\lambda & =\lambda =\delta _0/10+\delta _1/5,\hfill \\ \hfill g^{}\delta _i& =0\mathrm{for}i>2.\hfill \end{array}$$
Therefore we will get
$$\begin{array}{cc}\hfill g^{}\delta _0& =\delta _0,\hfill \\ \hfill g^{}\lambda & =\lambda =\delta _0/10+\delta _1/5,\hfill \\ \hfill g^{}\delta _{g1}& =\delta _1,\hfill \\ \hfill g^{}\delta _{g2}& =\omega ,\hfill \end{array}$$
with all $`\delta `$ not yet mentioned going to $`0`$. After all, $`\pi _1{}_{}{}^{}\delta _{i}^{}=\delta _i+\delta _{gi}`$, and as $`g^{}\delta _i`$ reflects the nodes of the curve of genus $`2`$, the marked points are not on the genus-$`i`$ side. Note also that the relation $`\lambda =\delta _0/10+\delta _1/5`$ on $`\overline{}_2`$ pulls back to $`\overline{}_{2,1}`$ without any change in its appearance.
To complete the description, we show that $`g^{}\omega =0`$. Using the push-pull formula, we see that it is enough to show that $`\omega g_{}^{}C=0`$ for any curve $`C`$ contained in $`\overline{}_{2,1}`$. But this is easy, as $`\omega `$ is the self-intersection of a constant section on the component of genus $`g2`$. The variation of the other component, which is attached at some other point, has no effect on this.
To prove the theorem, we show that the intersection of the kernels has dimension $`2`$. This is most easily seen as follows: if we know the coefficients of $`\delta _0`$ and $`\lambda `$ in a divisor contained in $`\mathrm{ker}g^{}`$, that determines its coefficients of $`\delta _{g1}`$ and $`\delta _{g2}`$. However, a knowledge of these two coefficients determines the coefficient of $`\omega `$ (in order that the coefficient of $`\theta _{g2}`$ in the pullback by $`f^{}`$ be $`0`$), and that forces the coefficients of all of the rest. Assuming that the $`\theta _i`$ are independent, we conclude that the dimension is at most 2; it is at least 2 because the Weierstrass and Brill-Noether classes are contained in the kernel. (Of course, we have only shown that the Brill-Noether class is in the kernel when it is the class of an effective divisor, but one can check with no difficulty that it never survives the map $`f^{}`$.)
Now we prove that the $`\theta _i`$ are actually independent in $`Pic(\overline{}_{0,g+1})𝐐`$. This is not quite trivial: there are relations between the different $`\delta _{0;S}`$ on $`\overline{}_{0,n}`$. Since $`\overline{}_{0,4}𝐏^1`$, for example, and the Picard group of $`𝐏^1`$ has rank $`1`$, we must have $`\delta _{0;\{1,2\}}=\delta _{0;\{1,3\}}=\delta _{0;\{1,4\}}`$. But it is not difficult either.
To start with, only $`\theta _1`$ has nonzero degree on a fiber of $`\pi _{g+1}`$, so its coefficient must be $`0`$ in any relation. Then $`\theta _{g1}`$ is the only other $`\theta `$ with nonzero degree on a fiber of $`\pi _1`$, so its coefficient must be $`0`$ as well. For the rest, put a fixed number $`i`$, from $`1`$ to $`g3`$, of the first $`g`$ points on a $`𝐏^1`$ together with $`p_{g+1}`$, attach another $`𝐏^1`$ at another point with the remaining marked points, and consider a curve in which one of these points moves along its component. This curve will intersect $`\theta _{g1}`$ (where the moving point meets another marked point on its component), $`\theta _i`$ (generically), and $`\theta _{i+1}`$ (where the moving point reaches the point of attachment), and the coefficient of $`\theta _{i+1}`$ will be $`1`$. Therefore, by easy induction on $`i`$, all the coefficients of $`\theta _i`$ from $`2`$ to $`g2`$ in our putative relation must be $`0`$, and we are done.
Moreover, it is easy to see that the space spanned by these divisors is actually of dimension $`2`$ (not $`1`$). If $`g+1`$ is not prime, this is clear, because $`𝒲`$ has a nonzero coefficient of $`\omega `$ while the pullback of the Brill-Noether class from $`\overline{}_g`$ does not. If $`g+1`$ is prime, it is not much harder. (Details will appear in a paper presenting the results of \[L\].)
3. The General Case
Recall the definition of the divisor class $`BN`$ on $`\overline{}_g`$ as
$$(g+3)\lambda \frac{g+1}{6}\delta _0i(gi)\delta _i.$$
If $`g+1`$ is composite, $`BN`$ is effective, but otherwise it may not be. We may define divisors on $`\overline{}_{g,n}`$ for any $`g,n`$ by taking any $`r`$ and $`d`$ for which the expected dimension of the space of $`g_d^r`$’s on a curve of genus $`g`$ is $`k(1)`$, and imposing exactly enough ramification conditions separately at the $`n`$ distinct points to reduce the expected dimension to $`1`$. This produces a locus in $`\overline{}_{g,n}`$ which may in general have components of various dimensions. We will refer to any codimension-1 component of any of these loci as a Brill-Noether type divisor, and to the subspace of $`Pic\overline{}_{g,n}𝐐`$ generated by all Brill-Noether type divisors as the Brill-Noether subspace. Observe that if we pull back a divisor of Brill-Noether type from $`\overline{}_{g,n1}`$ to $`\overline{}_{g,n}`$, we get another divisor of Brill-Noether type: the same conditions are imposed at $`p_1,\mathrm{},p_{n1}`$, and none at $`p_n`$. The main purpose of this paper is to study the classes of these divisors: this is what we will do in this section.
The proofs of the earlier results on linear dependence suggest a way to proceed, and we will follow it. As before, we consider maps
$$\overline{}_{0,g+n}\overline{}_{g,n}$$
given by attaching a fixed elliptic curve at each of the first $`g`$ points, and
$$\overline{}_{2,1}\overline{}_{g,n}$$
given by attaching a fixed $`n+1`$-pointed curve of genus $`g`$ at a marked point. And as before, the Plücker formula shows that the first map misses all of these divisors, while the extended Brill-Noether theorem proves that the image of $`\overline{}_{2,1}𝒲`$ by second map has trivial intersection with them. Therefore, their classes must lie in the intersections of the kernels of these two maps. These facts are not quite enough to characterize the Brill-Noether subspace, though. We will prove the following theorem instead:
Theorem 3.1 For any $`g>2`$, the Brill-Noether subspace of $`Pic(\overline{}_{g,n})𝐐`$ has dimension $`1+n+\left(\genfrac{}{}{0pt}{}{n}{2}\right)`$, unless $`g+1`$ is prime and $`n=0`$. In particular, its projection with respect to the standard basis onto the subspace spanned by $`\lambda `$, the $`\omega `$, and the $`\delta _{0;\{i,j\}}`$ is an isomorphism.
Proof. To start the argument, observe that $`BN`$ has a nonzero coefficient of $`\lambda `$, and the pullback of $`𝒲`$ by the map which forgets all but the $`i`$th point has nonzero coefficient of $`\omega _i`$ and $`\lambda `$, while all of their other coefficients in the space we have projected to are $`0`$. This, together with the results of the previous section, takes care of the cases $`n=0,1`$.
We start by proving that the projection is surjective. The case $`n=2`$ is disposed of as soon as we find a divisor of Brill-Noether type with nonzero coefficient of $`\delta _{0;\{1,2\}}`$. To do this, we need a bit of notation.
Definition. Let $`A(g,m,n)=A^{}(g,d,n)`$ be the number of $`g_d^1`$’s on a general curve of genus $`g`$ ramified to order $`m1`$ at a specified point and to order $`n1`$ at an unspecified point, where $`2d=g+m+n1`$.
Theorem 3.2 (\[L\], thm. 3.2)
$$A^{}(g,d,n)=g!(n^21)\underset{j=\mathrm{max}(0,m+nd1)}{\overset{\mathrm{min}(m1,n1,d)}{}}\frac{(m+n2j1)}{((dmn+j+1)!(dj)!)}.$$
(It seems almost certain that this was known long before \[L\], but I know of no reference.)
In odd genus $`g`$, we consider the divisor $`D`$ on $`\overline{}_{g,2}`$ of curves and points such that there is a $`g_{(g+3)/2}^1`$ ramified at both of the points. On the one hand, there are $`c_{(g+1)/2}g_{(g+3)/2}^1`$’s ramified at a given point, and each of them is ramified at exactly $`3g`$ other points, so the degree on a fiber is $`3gc_{(g+1)/2}`$. (Here $`c_n`$ is the $`n`$th Catalan number $`(2n)!/n!(n+1)!`$.) On the other hand, consider $`\pi _1{}_{}{}^{}D\delta _{0;\{1,2\}}`$. An easy argument with linear series shows this to consist of the sum of the pullback of a Brill-Noether divisor from $`\overline{}_g`$ and a divisor of Brill-Noether type on $`\overline{}_{g,1}`$, to wit, that of curves and points where the curve has a $`g_{(g+3)/2}^1`$ doubly ramified at the marked point. The degree of this on a fiber is $`A(g,1,3)=24\left(\genfrac{}{}{0pt}{}{g}{(g+3)/2}\right)`$, and the coefficient of $`\delta _{0;\{1,2\}}`$ will be $`0`$ iff this is equal to $`6gc_{(g+1)/2}`$. Expanding out both sides and multiplying through by
$$\frac{((g+1)/2)!((g+3)/2)!}{6g!},$$
we get $`g(g+1)=(g1)(g+1)`$, which holds for no positive integer, so the coefficient is nonzero as desired.
In even genus, consider the divisor $`D`$ on $`\overline{}_{g,2}`$ of curves and points such that there is a $`g_{(g+4)/2}^1`$ ramified doubly at the first point and singly at the second. When we cut and push forward, we get the divisor $`D_{g;2}+D_{g;4}`$ on $`\overline{}_{g,1}`$. Again, to show that $`D`$ has nonzero coefficient of $`\delta _{0;\{1,2\}}`$, we must prove that the sum of the degrees on the fibers of $`D`$ is not equal to the degree on the fibers of its pushforward; in other words, that $`A(g,2,3)+A(g,3,2)A(g,1,2)+A(g,1,4)`$. Writing these out in terms of the formula for $`A`$ given in Theorem 3.2 and dividing through by $`g!`$, we get
$$11\underset{j=0}{\overset{1}{}}\frac{(42j)}{(g/2+j2)!(g/2j+2)!}\frac{6}{(g/21)!(g/2+1)!}+\frac{60}{(g/22)!(g/2+2)!},$$
which on multiplying through by $`(g/21)!(g/2+2)!`$ becomes
$$44(g/21)+22(g/2+2)6(g/2+2)+60(g/21),$$
a statement that is always true.
Then, on $`\overline{}_{g,n}`$, we can prescribe the coefficients of the $`\delta _{0;\{i,j\}}`$ by pulling back these divisors from $`\overline{}_{g,2}`$ in appropriate ways. The $`\omega `$’s are dealt with by pulling back $`𝒲`$ from $`\overline{}_{g,1}`$, and $`\lambda `$ by pulling back $`BN`$ from $`\overline{}_g`$. We must show, now, that a divisor in the Brill-Noether subspace whose coefficients of $`\lambda `$, the $`\omega `$, and the $`\delta _{0;\{i,j\}}`$ are all $`0`$ is $`0`$.
Lemma 3.3 Let $`S=\{i,j\}\{1,\mathrm{},n+1\}`$. The map
$$Pic\overline{}_{g,n+1}𝐐Pic\overline{}_{g,n}𝐐$$
given by $`D\pi _j{}_{}{}^{}(D\delta _{0;S})`$ maps the Brill-Noether subspace to the Brill-Noether subspace. (Note that this is the pullback map on Picard groups induced by the map $`\overline{}_{g,n}\overline{}_{g,n+1}`$ which takes the point representing $`(C,p_1,\mathrm{},p_i,\mathrm{},p_n)`$ to $`(C^{},p_1,\mathrm{},P_1,\mathrm{},P_2,\mathrm{},p_n)`$, where $`C^{}`$ is $`C`$ with a copy of $`𝐏^1`$ attached at $`p_i`$ and the $`P_i`$ are points on this copy of $`𝐏^1`$.
Proof. It suffices to prove this on a set of generators for the Brill-Noether subspace. So let $`D`$ be a divisor of Brill-Noether type. Then for a point of $`\delta _{0;S}`$ to be contained in $`D`$ means that there is a limit linear series on the reducible curve the point corresponds to that satisfies the necessary ramification conditions. On the $`𝐏^1`$ containing $`p_i`$ and $`p_j`$, there is no choice in the matter: we know what the total order of vanishing must be at the point of attachment. This forces the total order of vanishing on the genus-$`g`$ component, and it is easily checked that this results in a codimension-$`1`$ condition for each way to distribute the vanishing between base points and ramification. Thus $`D`$ maps to a sum of divisors of Brill-Noether type, and the lemma is proved.
Lemma 3.4 Let $`DBNS`$ be a divisor whose coefficients of $`\lambda `$, the $`\omega `$, and the $`\delta _{0;\{i,j\}}`$ are all $`0`$, with respect to the standard basis. Then for any $`S`$, the coefficient of $`\delta _{0;S}`$ in $`D`$ is $`0`$ as well. (Note that we do not yet assert that its coefficient of $`\delta _0`$ must be $`0`$).
Proof. Let $`S\{1,\mathrm{},n\}(cardS>2)`$, and fix a $`cardS`$-pointed $`𝐏^1`$ and a general $`(n+1cardS)`$-pointed curve of genus $`g`$. Consider a family of curves whose base is isomorphic to $`𝐏^1`$, and whose fiber at a point $`P`$ has the $`cardS`$-pointed $`𝐏^1`$ attached at $`P`$ to the curve of genus $`g`$ at the first point. For each element $`xS`$, this family meets $`\delta _{0;S\{x\}}`$ once, and it meets $`\delta _{0;S}`$ with multiplicity $`2cardS`$—the section on the curve of genus $`g`$ is constant, so has self-intersection $`0`$, while the section on the $`𝐏^1`$ is a diagonal on $`𝐏^1\times 𝐏^1`$, blown up at $`cardS`$ points. On the other hand, it is plain that this curve does not meet any of the other boundary components or $`\lambda `$. In addition, I claim that the intersection with each $`\omega _i`$ is $`0`$.
For $`iS`$, this is obvious. For $`iS`$, the self-intersection of the $`i`$th section is $`1`$, so this contributes $`1`$. For every element $`x`$ of $`S`$ other than $`i`$, we get a contribution to $`\omega _i`$ of $`1`$ from $`\delta _{0;S\{x\}}`$, thus $`cardS1`$. Finally, we add the intersection with $`\delta _{0;S}`$, which is $`2cardS`$; total $`0`$.
By the extended Brill-Noether theorem \[EH, Thm. 1.1\], this curve misses all divisors of Brill-Noether type entirely, and so for every $`S`$ of cardinality greater than $`2`$ we get a relation
$$(2cardS)d_{0;S}+\underset{xS}{}d_{0;S\{x\}}=0.$$
(Roman letters are the coefficients of divisors named by similar Greek letters.) Therefore, if $`D`$ is a divisor in the Brill-Noether subspace with coefficients of $`\delta _{0;\{i,j\}}`$ equal to $`0`$, an induction on $`cardS`$ proves the lemma.
Now we are ready to settle the case $`n=2`$ in Theorem 3.1. Suppose that $`DBNS`$ has all coefficients of $`\lambda ,\omega ,`$ and $`\delta _{0;\{i,j\}}`$ equal to 0, but suppose that the coefficient of $`\delta _{i,S}`$ is nonzero. If $`S`$ is empty, then consider $`\pi _{}\delta _{0;\{1,2\}}D`$—it is in the Brill-Noether subspace and has a nonzero coefficient of $`\delta _{i,\mathrm{}}`$, contradiction. We may thus assume that all of the $`\delta _{i;\mathrm{}}`$ are $`0`$.
We now know that $`D`$ pulls back to $`0`$ on $`Pic\overline{}_{0,g+2}`$. Define $`\theta _{i;S}`$ on $`\overline{}_{0,g+2}`$ to be the sum
$$\underset{\genfrac{}{}{0pt}{}{T\{1,2,\mathrm{},g\}}{cardT=i}}{}\delta _{0;T(S+n)},$$
where $`S+n`$ is the set obtained by adding $`n`$ to each element of $`S`$, so that $`\theta _{i;S}`$ is the pullback of $`\delta _{i;S}`$. Let
$$D=a_i\delta _{i;\{1\}}+d\delta _0;$$
then $`a_i\theta _{i;\{1\}}=0`$. We show that the $`\theta _{i;\{1\}}`$ are linearly independent.
To start, only $`\theta _{1;\{1\}}`$ has nonzero degree on a fiber of $`\pi _{g+1}`$, so it cannot appear in a relation. Next, only $`\theta _{g1;\{1\}}`$ of the remaining $`\theta `$ has nonzero degree on a fiber of any of the other projection maps. For the rest, fix $`1<i<g2`$, put $`i`$ of the first $`g`$ points and $`p_{g+2}`$ on a $`𝐏^1`$ and attach this $`𝐏^1`$ to another $`𝐏^1`$ which has the other $`gi`$ of the first $`g`$ marked points and $`p_{g+1}`$. Then let one of the points vary on the first $`𝐏^1`$. This meets only $`\theta _{gi;\{1\}}`$ and $`\theta _{gi+1;\{1\}}`$ of the $`\theta `$, the latter with multiplicity $`1`$. Again it is immediate that the $`\theta `$ are in fact linearly independent.
This proves that all of the $`a_i`$ are $`0`$, so $`D`$ is just a multiple of $`\delta _0`$. That means that $`D=0`$, though, because the map $`\overline{}_{2,1}\overline{}_{g,2}`$ does not pull any nonzero multiple of $`\delta _0`$ back to a multiple of $`𝒲`$. This completes the proof in the case $`n=2`$.
Finally, for general $`n`$, given $`D`$ we start by removing its coefficients of $`\lambda `$, the $`\omega `$, and the $`\delta _{0;\{i,j\}}`$. Its coefficients of $`\delta _{0;S}`$ are automatically $`0`$; suppose that its coefficient of $`\delta _{i;S}`$ is nonzero. Choose $`j,k`$ so that either $`j,kS`$ or $`j,kS`$, cut with $`\delta _{0;\{j,k\}}`$, and forget the $`k`$th point. This produces a smaller $`n`$ and $`D`$ with such a nonzero coefficient, and proceeding in this way we get down to the case $`n=2`$, contradiction. So all the $`\delta _{i;S}`$ have a coefficient of $`0`$. We finish the proof of the theorem by concluding as before that $`\delta _0`$ cannot appear.
References
\[AC\] E. Arbarello, M. Cornalba, The Picard groups of the moduli spaces of curves. Topology 26, 153–171, 1987.
\[Cuk\] F. Cukierman, Families of Weierstrass points. Duke Math. J. 58, 317–346, 1989.
\[EH\] D. Eisenbud, J. Harris, The Kodaira dimension of the moduli space of curves of genus $`23`$. Inv. Math. 90, 359–387, 1987.
\[HM\] J. Harris, I. Morrison, Moduli of Curves. GTM 167. Springer-Verlag, 1998.
\[K\] F. Knudsen, The projectivity of the moduli space of stable curves, II. The stacks $`_{g,n}`$, Math. Scand. 52 no. 2, 161–199, 1983.
\[L\] A. Logan, Moduli spaces of curves with marked points. Harvard University doctoral dissertation, June 1999. |
warning/0003/hep-th0003030.html | ar5iv | text | # SU(2) GAUGE THEORY IN COVARIANT (MAXIMAL) ABELIAN GAUGES
## 1 Introduction
An $`SU(2)`$ Lattice Gauge Theory (LGT) on a finite lattice is invariant under the compact group $`𝒢`$
$$𝒢=_{\mathrm{sites}}SU(2).$$
(1)
$`𝒢`$ does not have a smooth continuum limit and it is not entirely clear whether some of the effects observed on the lattice, such as absolute confinement, are due to the compact nature of $`𝒢`$ and are absent in a local Quantum Field Theory (QFT) with non-compact $`SU(2)`$ gauge invariance. The fact that compact (lattice) QED does differ markedly from the continuum theory due to lattice “artefacts” makes this question all the more relevant. These “artefacts” are Abelian monopoles specific to lattice QED – their existence is intimately related to the compactness of the (Abelian) lattice gauge group and they have no continuum analogues. In the Abelian case, one can remove these lattice artefacts by imposing the constraint
$$A_\mu =^1_\mu F_{\mu \nu },$$
(2)
where $`P_{\mu \nu }(x)=e^{iF_{\mu \nu }(x)}`$ is the plaquette variable and $`U_\mu (x)=e^{iA_\mu (x)}`$ is the $`U(1)`$ link variable. On the lattice, Eq. (2) is not just a gauge fixing condition, but in addition eliminates the monopoles associated with harmonic one forms. The continuum limit of this (projected) Abelian LGT is free QED in Landau gauge. If Eq. (2) was not imposed, the gauge invariant transverse photon correlation function of lattice QED was found to differ markedly from what one expects in the continuum.
In view of this example, the question whether the critical limit of a non-Abelian LGT can be described by “QCD” with a non-compact gauge group is legitimate. It has been conjectured that a non-Abelian LGT may not be asymptotically free and exhibit a Kosterlitz-Thouless phase transition at a finite value of the coupling constant.
In resolving the issue of the critical limit of a non-Abelian LGT it may be useful to have a physically equivalent local LGT with an equivariant BRST-symmetry whose structure group has been reduced to the maximal Abelian subgroup of the original LGT. The two LGT’s in question are physically equivalent because the expectation values of gauge-invariant operators (i.e. Wilson loops and their linked generalizations) are the same. Although this equivalent Abelian LGT has only been constructed for an $`SU(2)`$-LGT, the method can be generalized to any $`SU(n)`$-LGT. The lattice group $`𝒢`$ of an $`SU(2)`$-LGT is reduced to the maximal Abelian subgroup $``$ by a local Topological Lattice Theory (TLT) that computes the Euler number of the coset $`𝒢/`$,
$$\chi (𝒢/)=_{\mathrm{sites}}\chi (SU(2)/U(1)S_2)=2^{\mathrm{\#}\mathrm{sites}},$$
(3)
on each orbit of a lattice configuration using Morse Theory. I should stress that this construction of a TLT is conceptually quite different from the usual Faddeev-Popov procedure and does not require the uniqueness of the solution to a “gauge condition” – the Euler number of the manifold would in fact have to be 1 for this to be the case. There are at least $`2^{\mathrm{\#}\mathrm{sites}}`$ gauge equivalent Gribov copies that contribute to $`\chi (𝒢/)`$ on any orbit of the (finite) lattice. The construction of the TLT is mathematically rigorous, because the coset manifold $`𝒢/`$ is compact and finite-dimensional on a finite lattice (albeit of rather large dimension) and the orbit-space of the original LGT is connected. One cannot reduce the full lattice group $`𝒢`$ in this manner because $`\chi (𝒢)=0`$. The best one can do is to reduce the lattice gauge group to the smallest subgroup $``$ for which the Euler number of the coset manifold $`𝒢/`$ does not vanish. In the case of compact $`SU(n)`$ the smallest subgroup which satisfies this requirement is the maximal Abelian one.
Of interest here will be the continuum theory that describes the critical limit of this “partially gauge fixed” LGT. The equivariant BRST-symmetry and $`U(1)`$-invariance together with locality and power-counting renormalizability determine the continuum limit up to lattice artefacts associated with the compactness of the residual $`U(1)`$-structure group. Assuming these Abelian artefacts have been removed in a manner similar to the one prescribed above, the critical limit of this LGT is unique because the BRST-invariance of the LGT is a global one. The continuum model is then described by the local action given below.
Because physical correlation functions are the same in the reduced Abelian LGT and can be shown to satisfy reflection positivity in the original $`SU(2)`$-LGT, the physical states of the “partially gauge fixed” Abelian LGT also have positive norm. By proving the equivalence of the two LGT’s for gauge-invariant correlation functions one thus also verifies the unitarity of the partially gauge-fixed Abelian LGT. The continuum theory describing the critical limit of this LGT should therefore be unitary as well. Note that this proof of the unitarity of the continuum theory is valid non-perturbatively and not just to all orders in perturbation theory. Instead of investigating the critical limit of the original $`SU(2)`$-LGT, we thus will consider the critical limit of the (physically) equivalent LGT with an Abelian structure group and an equivariant BRST-symmetry.
## 2 Continuum SU(2) Gauge Theory in Abelian Gauges
Up to Abelian lattice artifacts mentioned above, the continuum theory describing the critical limit of the “partially gauge fixed” Abelian LGT is completely specified by the global symmetries and power counting. It is described by the Lagrangian
$$=_{\mathrm{inv}.}+_{\mathrm{AG}}+_{\mathrm{aGF}}.$$
(4)
Here $`_{\mathrm{inv}.}`$ is the usual $`SU(2)`$-invariant Lagrangian with the SU(2)-connection $`\stackrel{}{V}_\mu =(W_\mu ^1,W_\mu ^2,A_\mu )`$ written<sup>1</sup><sup>1</sup>1Latin indices take values in $`\{1,2\}`$ only, Einstein’s summation convention applies and $`\epsilon ^{12}=\epsilon ^{21}=1`$, vanishing otherwise. All results are given in the $`\overline{MS}`$ renormalization scheme. in terms of two (real) vector bosons $`W`$, and an Abelian “photon” $`A`$,
$$_{\mathrm{inv}.}=_{\mathrm{matter}}+\frac{1}{4}(G_{\mu \nu }G_{\mu \nu }+G_{\mu \nu }^aG_{\mu \nu }^a),$$
(5)
with
$`G_{\mu \nu }`$ $`=`$ $`_\mu A_\nu _\nu A_\mu g\epsilon ^{ab}W_\mu ^aW_\nu ^b`$
$`G_{\mu \nu }^a`$ $`=`$ $`D_\mu ^{ab}W_\nu ^bD_\nu ^{ab}W_\mu ^b=_\mu W_\nu ^a_\nu W_\mu ^a+g\epsilon ^{ab}(A_\mu W_\nu ^bA_\nu W_\mu ^b).`$ (6)
$`_{\mathrm{AG}}`$ reduces the invariance to the maximal Abelian subgroup $`U(1)`$ of $`SU(2)`$ in a covariant manner,
$$_{AG}=\frac{F^aF^a}{2\alpha }\overline{c}^aM^{ab}c^bg^2\frac{\alpha }{2}(\overline{c}^a\epsilon ^{ab}c^b)^2,$$
(7)
with
$$F^a=D_\mu ^{ab}W_\mu ^b=_\mu W_\mu ^a+gA_\mu \epsilon ^{ab}W_\mu ^b\mathrm{and}M^{ab}=D_\mu ^{ac}D_\mu ^{cb}+g^2\epsilon ^{ac}\epsilon ^{bd}W_\mu ^cW_\mu ^d.$$
(8)
Like the corresponding Abelian LGT, $`L_{U(1)}=L_{\mathrm{inv}.}+_{\mathrm{AG}}`$ is invariant under $`U(1)`$-gauge transformations and an on-shell BRST symmetry $`s`$ and anti-BRST symmetry $`\overline{s}`$, whose action on the fields is
$$\begin{array}{ccccccc}\hfill sA_\mu & =& g\epsilon ^{ab}c^aW_\mu ^b\hfill & & \hfill \overline{s}A_\mu & =& g\epsilon ^{ab}\overline{c}^aW_\mu ^b\hfill \\ \hfill sW_\mu ^a& =& D_\mu ^{ab}c^b\hfill & & \hfill \overline{s}W_\mu ^a& =& D_\mu ^{ab}\overline{c}^b\hfill \\ \hfill sc^a& =& 0\hfill & & \hfill \overline{s}\overline{c}^a& =& 0\hfill \\ \hfill s\overline{c}^a& =& F^a/\alpha \hfill & & \hfill \overline{s}c^a& =& F^a/\alpha ,\hfill \end{array}$$
(9)
with an obvious extension to include matter fields. On the connections $`A_\mu `$ and $`W_\mu ^a`$ this BRST-variation effects an infinitesimal transformation in the coset $`𝒢/`$ parameterized by the two ghosts $`c^a(x)`$. Note that $`sc^a=0`$ here, because the coset is not a group manifold.
The BRST algebra Eq. (9) closes on-shell on the set of $`U(1)`$-invariant functionals: on functionals that depend only on $`W,A,c`$ and the matter fields, $`s^2`$ for instance effects an infinitesimal U(1)-transformation with the parameter $`\frac{g}{2}\epsilon ^{ab}c^ac^b`$. The algebra Eq. (9) thus defines an equivariant cohomology. It was derived from a more extensive nilpotent (off-shell) BRST-algebra on the lattice by integrating out some of the additional fields. As mentioned in the introduction, the renormalizability and unitarity of this continuum theory is guaranteed because it describes the critical limit of an Abelian LGT that was proven to have the same gauge invariant correlation functions as the original $`SU(2)`$-LGT. Note that the physical sector comprises states created by composite operators of $`A,W`$ and the matter fields in the equivariant cohomology of $`s`$ (or $`\overline{s}`$). They are BRST closed, $`U(1)`$-invariant and do not depend on the ghosts.
For $`\alpha >0`$, Eq. (7) could be considered a “soft” gauge fixing to the Maximal Abelian Gauge (MAG). It differs from what one naively obtains using a Faddeev-Popov procedure by a quartic ghost interaction proportional to $`\alpha `$. Eq. (7) also does not implement the non-linear constraint $`F^a=0`$ exactly. However, setting $`\alpha =0`$ and perturbatively solving the constraint $`F^a=0`$ is not consistent and not the same as taking the limit $`\alpha 0`$. One could have inferred the highly singular nature of this limit from the fact that the 4-ghost interaction diverges at one loop even when the photon- and vector boson propagators are transverse. A 4-ghost counterterm thus is required even in the (formal) limit $`\alpha 0`$ and the leading term in the anomalous dimension of the gauge parameter in fact is $`3g^2/(8\pi ^2\alpha )`$ in this limit. The physical reason for the singular behavior of the limit $`\alpha 0`$ is inherently non-perturbative and nicely exhibited by the lattice calculation: without the quartic ghost interaction, Gribov copies of a configuration conspire to give vanishing expectation values for all physical observables. No matter how small, the quartic ghost interaction is required to have a normalizable partition function and expectation values of physical observables that are identical with those of the original SU(2)-LGT. From a perturbative point of view, $`\alpha 0`$ at finite coupling $`g^2`$ corresponds to a strong coupling limit that is not accessible perturbatively.
$`_{\mathrm{aGF}}`$ in Eq. (4) fixes the remaining $`U(1)`$ gauge invariance and thus defines the perturbative series unambiguously. I will consider a conventional covariant gauge-fixing of the form,
$$L_{\mathrm{aGF}}=\delta [\overline{\omega }(_\mu A_\mu \frac{\xi }{2}b)]=b_\mu A_\mu \frac{\xi }{2}b^2+\overline{\omega }\omega \frac{1}{2\xi }(_\mu A_\mu )^2,$$
(10)
where the last equivalence is obtained by decoupling the Abelian ghosts $`\omega `$ and $`\overline{\omega }`$ and the Nakanishi-Lautrup field $`b`$. $`\delta `$ is a BRST-symmetry defined on the fields as
$$\begin{array}{ccccccc}\hfill \delta A_\mu & =& _\mu \omega \hfill & & \hfill \delta W_\mu ^a& =& g\omega \epsilon ^{ab}W_\mu ^b\hfill \\ \hfill \delta c^a& =& g\omega \epsilon ^{ab}c^b\hfill & & \hfill \delta \overline{c}^a& =& g\omega \epsilon ^{ab}\overline{c}^b\hfill \\ \hfill \delta \omega & =& 0\hfill & & & & \\ \hfill \delta \overline{\omega }& =& b\hfill & & \hfill \delta b& =& 0\hfill \end{array}$$
(11)
Trivially extending $`s`$ and $`\overline{s}`$ to the additional fields
$$s\omega =s\overline{\omega }=sb=\overline{s}\omega =\overline{s}\overline{\omega }=\overline{s}b=0,$$
(12)
one can verify that $`\delta `$ is nil-potent and anticommutes with $`s`$ and $`\overline{s}`$
$$\delta ^2=s\delta +\delta s=\overline{s}\delta +\delta \overline{s}=0$$
(13)
As far as algebraic renormalization is concerned, the action Eq. (4) thus is composed of a term in the cohomology of $`\delta `$ that is invariant under the global symmetries $`s`$ and $`\overline{s}`$ given in Eq. (9) and a $`\delta `$-exact term. The latter is not invariant under $`s`$ nor $`\overline{s}`$. Since the global symmetries commute with $`\delta `$, the situation is the same as in any gauge-fixing that breaks some of the global symmetries (or supersymmetries) of the theory. There is a well-defined procedure to handle this case in algebraic renormalization. From a more heuristic point of view, we already know that the $`s`$ and $`\overline{s}`$ symmetries of the theory are not anomalous from the lattice regularization of this model. I will therefore not give the algebraic proof here.
What has been gained compared to conventional covariant gauge fixing? Since the present gauge fixing in a sense is “hierarchical”, we are able to single out the maximal Abelian subgroup. As emphasized before, the $`s`$ and $`\overline{s}`$ symmetries can be implemented on the lattice and the resulting Abelian LGT shown to be physically equivalent to one with an SU(2) structure group. It may eventually be possible to construct the dual of this Abelian LGT. In addition, the theory described by Eq. (4) does not suffer from a generic Gribov problem due to zero modes of the ghosts. Since the global $`s`$ and $`\overline{s}`$ symmetries are preserved by the lattice regularization, Eq. (4) probably descibes the critical limit of a LGT better than any other set of covariant gauges. Finally, because the Abelian $`\omega `$-ghost and $`\overline{\omega }`$-antighost as well as the Nakanishi-Lautrup field $`b`$ decouple, we effectively end up with a local and covariant gauge-fixed theory with fewer ghosts. This reduction in the number of fields is at the expense of a quartic ghost interaction that gives the ghosts some interesting dynamics of their own.
## 3 The Dynamically Broken SL(2,R) Symmetry
The Lagrangian Eq. (4) also exhibits a global bosonic SL(2,R) symmetry that is generated by
$$\mathrm{\Pi }^+=c^a(x)\frac{\delta }{\delta \overline{c}^a(x)},\mathrm{\Pi }^{}=\overline{c}^a(x)\frac{\delta }{\delta c^a(x)},$$
(14)
and the ghost number $`\mathrm{\Pi }=[\mathrm{\Pi }^+,\mathrm{\Pi }^{}]`$. This SL(2,R) symmetry is preserved by the regularization (for instance dimensional) and thus is not anomalous. The conserved currents corresponding to $`\mathrm{\Pi }^\pm `$ are U(1)-invariant and BRST, respectively anti-BRST exact,
$$j_\mu ^+=c^aD_\mu ^{ab}c^b=sc^aW_\mu ^a,j_\mu ^{}=\overline{c}^aD_\mu ^{ab}\overline{c}^b=\overline{s}\overline{c}^aW_\mu ^a.$$
(15)
I will argue that the global $`SL(2,R)`$ symmetry of the theory is spontaneously broken to the Abelian subgroup generated by the ghost number $`\mathrm{\Pi }`$. An order parameter for the spontaneous breakdown of the SL(2,R) symmetry thus is
$$\overline{c}^a\epsilon ^{ab}c^b=\frac{1}{2}\mathrm{\Pi }^{}(c^a\epsilon ^{ab}c^b)=\frac{1}{2}\mathrm{\Pi }^+(\overline{c}^a\epsilon ^{ab}\overline{c}^b).$$
(16)
Because the currents Eq. (15) are (anti)-BRST exact, a spontaneously broken SL(2,R) symmetry is accompanied by a BRST-quartet of massless Goldstone states with ghost numbers $`2,1,1`$ and $`2`$. They are U(1)-invariant $`cc`$, $`cW`$, $`\overline{c}W`$ and $`\overline{c}\overline{c}`$ bound states. It is important to note in this context that BRST quartets do not contribute to physical quantities<sup>2</sup><sup>2</sup>2This is analogous to the decoupling of the Goldstone quartets of the weak interaction in renormalizable $`R_\xi `$ gauges.. The spontaneous breaking of the $`SL(2,R)`$ symmetry in a sense is similar to a dynamical Higgs mechanism in the adjoint. The vector bosons $`W`$ aquire a mass (see below) but in contrast to a conventional Higgs mechanism in the adjoint, this mass is not a free parameter of the theory, but can be determined in terms of $`\mathrm{\Lambda }_{\overline{MS}}`$.
To see that ghost condensation will almost invariably occur at weak coupling in the model described by Eq. (4) it is illustrative to compare with the BCS-theory of superconductivity. In BCS-theory, an at low momentum transfers attractive 4-fermion interaction leads to the condensation of certain fermion pairs and the formation of a gap in the quasi-particle spectrum near the Fermi surface. An analogous phenomenon occurs here for ghost and anti-ghost modes corresponding to small eigenvalues of the FP-Operator $`M^{ab}`$ defined by Eq. (8): for zero modes, the bilinear term in Eq. (7) vanishes and the quartic ghost interaction selects the channel in which condensation occurs. The quartic ghost interaction in Eq. (7) is attractive when the color of the ghost and anti-ghost are opposite: it thus leads to the formation of a $`\overline{c}^a\epsilon ^{ab}c^b`$ condensate at arbitrarily weak coupling $`\alpha g^2`$ by the BCS-mechanism. \[The two spin states of a fermion here has an analog in the two color orientations of the ghosts. We choose ghost number to be conserved and observe $`\overline{c}c`$, rather than $`cc`$ or $`\overline{c}\overline{c}`$ condensation, as would be the case if we chose $`\mathrm{\Pi }^{}`$ or $`\mathrm{\Pi }^+`$ as unbroken generators.\] The analogy with BCS-theory is particularly appealing because the operator $`M^{ab}`$ has small eigenvalues whenever the gauge field configuration is in the vicinity of a Gribov horizon. That the ground state may be dominated by such configurations was previously suggested in an attempt to restrict the functional integral to the fundamental modular region of Landau gauge. In the present context, gauge field configurations with non-Abelian monopoles are on the Gribov horizon, since the failure of the gauge fixing condition to an Abelian subgroup is necessary for the presence of monopoles.
Thus, if monopoles are relevant in describing the ground state of the theory, it is not inconceivable that the ghosts will condense. The converse is not necessarily true because $`M^{ab}`$ can have arbitrarily small eigenvalues at field configurations with vanishing monopole number. We will see below that the ghosts already condense in the vicinity of the trivial gauge field configuration. The analogy with BCS-theory suggests that they condense for any value of the quartic coupling $`\alpha g^2`$, with a gap that depends exponentially on $`1/(\alpha g^2)`$.
To perturbatively investigate the consequences of $`\overline{c}^a\epsilon ^{ab}c^b0`$, the quartic ghost interaction in Eq. (7) is linearized by introducing an auxiliary scalar field $`\rho (x)`$ of canonical dimension two. Adding the quadratic term
$$_{\mathrm{aux}}=\frac{1}{2g^2}(\rho g^2\lambda \overline{c}^a\epsilon ^{ab}c^b)^2$$
(17)
to the Lagrangian of Eq. (4), the tree level quartic ghost interaction vanishes at $`\lambda ^2=\alpha `$ and is then formally of $`O(g^4)`$, proportional to the difference $`Z_\lambda ^2Z_\alpha `$ of the renormalization constants of the two couplings<sup>3</sup><sup>3</sup>3The discrete symmetry $`c^a\overline{c}^a,\overline{c}^ac^a,\rho \rho `$ relating $`s`$ and $`\overline{s}`$ also ensures that $`\rho `$ only mixes with $`\overline{c}^a\epsilon ^{ab}c^b`$..
We will see that the perturbative expansion about a non-trivial solution $`\rho =v0`$ to the gap equation
$$\frac{v}{g^2}=\sqrt{\alpha }c^a(x)\epsilon ^{ab}\overline{c}^b(x)|_{<\rho >=v},$$
(18)
is stable and much better behaved in the infrared.
Defining the quantum part $`\sigma (x)`$ of the auxiliary scalar $`\rho `$ by
$$\rho (x)=v+\sigma (x)\mathrm{with}\sigma =0,$$
(19)
the momentum representation of the Euclidean ghost propagator at tree level becomes
$$c^a\overline{c}^b_p=\frac{p^2\delta ^{ab}+\sqrt{\alpha }v\epsilon ^{ab}}{p^4+\alpha v^2}=_0^{\mathrm{}}𝑑\omega e^{\omega p^2}[\delta ^{ab}\mathrm{cos}(\omega v\sqrt{\alpha })+\epsilon ^{ab}\mathrm{sin}(\omega v\sqrt{\alpha })].$$
(20)
Feynman’s parameterization of this propagator leads to an evaluation of loop integrals using dimensional regularization that is only slightly more complicated than usual.
Using Eq. (20) in Eq. (18) the gap equation to one-loop in $`d=42\epsilon `$ dimensions is,
$`{\displaystyle \frac{v}{\widehat{g}^2}}`$ $`=`$ $`2\mu ^4\sqrt{\alpha }{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\omega }{(4\pi \mu ^2\omega )^{2\epsilon }}}\mathrm{sin}(\omega \sqrt{v^2\alpha })`$ (21)
$`=`$ $`{\displaystyle \frac{\alpha v}{8\pi ^2}}\left[{\displaystyle \frac{1}{\epsilon }}\mathrm{ln}{\displaystyle \frac{\pi y^2T^2}{\mu ^2e^{1\gamma _E}}}+O(\epsilon )\right].`$
Here $`\gamma _E`$ is Euler’s constant. Including the counterterm
$$\frac{v}{\widehat{g}^2}\frac{v}{\widehat{g}^2}Z_v^2Z_g^2=\frac{v}{\widehat{g}^2}+\frac{\alpha v}{8\pi ^2\epsilon }+O(\widehat{g}^2)$$
(22)
on the left-hand side of Eq. (21) cancels the $`1/\epsilon `$ divergence of the right-hand side of Eq. (21). The renormalized (non-trivial) solution to the gap equation in four dimensions thus is,
$$\alpha v^2=e^2\mathrm{\Lambda }^4,\mathrm{with}\mathrm{\Lambda }^2(\alpha ,g,\mu )=4\pi \mu ^2e^{\gamma _E\frac{8\pi ^2}{\alpha g^2}}.$$
(23)
Note the exponential dependence of the gap $`v`$ on the quartic coupling $`\alpha g^2`$ expected from BCS-theory. One can show that the solution Eq. (23) corresponds to the global minimum of the effective potential by either directly computing the (one-loop renormalized) effective potential,
$$V(v,\mu ,g,\alpha )=\frac{\alpha v^2}{32\pi ^2}\mathrm{ln}\frac{\alpha v^2}{e^3\mathrm{\Lambda }^4}+O(g^2),$$
(24)
or by integrating Eq. (21) and noting that $`V(\mathrm{\Lambda },0)=0`$. Consistency requires that the 1PI two-point function of the scalar $`\sigma `$ is positive definite for all Euclidean momenta when Eq. (18) is satisfied. From Eq. (22) one obtains for the anomalous dimension of $`v`$ to one loop
$$\gamma _v=\frac{d\mathrm{ln}Z_v}{d\mathrm{ln}\mu }=\frac{g^2}{16\pi ^2}(2\alpha \beta _0)+O(g^4)$$
(25)
where $`\beta _0`$ is the lowest order coefficient of the $`\beta `$-function. At $`\alpha =\beta _0/2`$, the anomalous dimension of $`v`$ is of order $`g^4`$ and corrections to the asymptotic $`v(g0)`$ solution in this particular critical gauge therefore are analytic in $`g^2`$ and may be computed order by order in perturbation theory. In this critical gauge we thus have that the scale $`\mathrm{\Lambda }`$ describing the minimum of the effective potential $`V(v,\mu ,g,\alpha )`$ is analytically related to $`\mathrm{\Lambda }_{\overline{MS}}`$:
$$\mathrm{\Lambda }(\alpha =\beta _0/2,g,\mu )=\mathrm{\Lambda }_{\overline{MS}}(1+O(g^2))$$
(26)
In other gauges $`v0`$ at weak coupling is either much larger than $`\mathrm{\Lambda }_{\overline{MS}}`$ (for $`\alpha \beta _0/2`$), or much smaller (for $`\alpha \beta _0/2`$). To leading order in the loop expansion, the anomalous dimension Eq. (25) does not vanish in these cases and higher order loop corrections to $`v0`$ are of comparable magnitude. \[As noted before, the extreme limit $`\alpha 0`$ in which the non-trivial solution Eq. (23) becomes degenerate with the trivial one, in particular corresponds to a strong coupling problem.\] In the critical gauge $`\alpha =\beta _0/2`$, the perturbative 1-loop calculation of $`v0`$ is self-consistent in the sense that all higher order corrections to the expectation value are of order $`g^2`$ because the anomalous dimension Eq. (25) is of order $`g^4`$ at this point. I wish to emphasize that this does not imply that physical effects associated with ghost condensation in this gauge are themselves gauge dependent. It only implies that a non-trivial solution to the gap equation is perturbatively consistent at $`\alpha =\beta _0/2`$. \[That some gauges are better suited than others for a non-perturbative evaluation of gauge invariant quantities is well known from QED: the gauge invariant hydrogen spectrum is qualitatively best obtained in Coulomb gauge. Below we relate $`\alpha v^2`$ to the vacuum expectation value of the trace of the energy-momentum tensor.\]
## 4 The Vector Boson Mass
Recent lattice simulations indicate that the $`W`$-bosons are massive in maximal Abelian gauges. At least qualitatively, this may be explained by the mechanism discussed here. Although the tree level contribution to $`m_W^2`$ vanishes by Bose symmetry, ghost condensation induces a finite mass $`m_W^2=g^2\sqrt{\alpha v^2}/(16\pi )`$ at one loop as shown in Fig. 1.
$$=\frac{g^2\sqrt{v^2\alpha }}{16\pi }\delta _{\mu \nu }\delta ^{ab}.$$
(27)
Fig. 1. The finite one-loop contribution to the $`W`$ mass.
Technically, the one-loop contribution is finite because the integral in Eq. (27) involves only the $`\delta ^{ab}`$-part of the ghost propagator Eq. (20). Since $`p^2/(p^4+\alpha v^2)=\alpha v^2/(p^2(p^4+\alpha v^2))+1/p^2`$, the $`v`$-dependence of the loop integral is IR- and UV-finite. The quadratic UV-divergence of the $`1/p^2`$ subtraction at $`v=0`$ is canceled by the other, $`v`$-independent, quadratically divergent one-loop contributions – (in dimensional regularization this scale-invariant integral vanishes by itself). $`m_W^2`$ furthermore is positive due to the overall minus sign of the ghost loop. The sign of $`m_W^2`$ is crucial. It is a further indication that the model is stable and (as far as the loop expansion is concerned) does not develop tachyonic poles at Euclidean $`p^2`$ for $`v0`$. Conceptually, the local mass term proportional to $`\delta _{\mu \nu }\delta ^{ab}`$ is finite due to the BRST symmetry Eq. (9), which excludes a mass counter-term. The latter argument implies that contributions to $`m_W^2`$ are finite to all orders of the loop expansion.
Although this “mass”-term regulates the IR-behavior of the $`W`$-propagator perturbatively, the 1-loop calculation above should quantitatively describe the behavior of the $`W`$-propagator at high momenta, where $`g^2(p^2)`$ is a small parameter and this calculation is consistent. Ghost condensation thus should lead to a leading power correction $`v/p^2`$ to the $`W`$-propagator at high momenta. The consistency of this behavior with the Operator Product Expansion (OPE) is examined below.
## 5 Discussion of Physical Consequences
We have seen that ghost condensation in covariant Abelian gauges is associated with the spontaneous breaking of a global SL(2,R) symmetry whose diagonal generator is the ghost number. The currents of the broken symmetries are (anti-)BRST exact and the Goldstone states form a BRST-quartet that decouples from physical observables. What are observable consequences? We gain some insight by computing the contribution to the expectation value of the trace of the energy momentum tensor. From the effective potential Eq. (24) one obtains:
$$\theta _\mu ^\mu =\frac{\alpha v^2}{8\pi ^2}=\frac{e^2}{8\pi ^2}\mathrm{\Lambda }^4,$$
(28)
At the minimum of the effective potential. $`v0`$ thus lowers the vacuum energy density and ghost condensation may be interpreted as a low-energy manifestation of the trace anomaly.
Let me also comment on the consistency of the approach from the point of view of the Operator Product Expansion (OPE). In generic covariant gauges, the OPE implies that power corrections to physical correlators (but also to Green functions) are at least suppressed by an order $`M^4/p^4`$ relative to the perturbative behavior at high momenta (in the absence of quarks). This is simply because the operator of lowest dimension in the BRST-cohomology of these gauges has canonical dimension four. Consistency of the present approach requires that one explain
* why the ground state in the present case can support a vacuum expectation value $`v\overline{c}^a\epsilon ^{ab}c^b0`$ of canonical dimension two and simultaneously,
* why power corrections of order $`v/p^2`$ are absent in physical correlation functions
To address the first question, we construct a $`U(1)`$-invariant local operator of dimension two whose vacuum expectation value manifestly is invariant under the BRST-algebra of Eq. (9). Note that
$`s[\alpha \overline{c}^ac^a+{\displaystyle \frac{1}{2}}W_\mu ^aW_\mu ^a]`$ $`=`$ $`_\mu (W_\mu ^ac^a)`$
$`\overline{s}[\alpha \overline{c}^ac^a+{\displaystyle \frac{1}{2}}W_\mu ^aW_\mu ^a]`$ $`=`$ $`_\mu (W_\mu ^a\overline{c}^a).`$ (29)
$`O_2:=[\overline{c}^ac^a+\frac{1}{2\alpha }W_\mu ^aW_\mu ^a]`$ is not in the equivariant cohomology of the on-shell BRST-algebra Eq. (9), because $`s^20`$ only on U(1)-invariant operators that do not depend on $`\overline{c}`$ whereas $`O_2`$ does. Nevertheless, the zero-momentum component of $`O_2`$ is invariant under $`s`$, $`\overline{s}`$ and $`\delta `$. $`O_20`$ therefore is an invariant statement as far as the algebra of symmetries is concerned. We have seen that $`O_2=\frac{\sqrt{\alpha v^2}}{16\pi }0`$ in fact appears to be the dynamically favored possibility. The existence of an operator of dimension two whose zero-momentum component is invariant explains the leading power corrections we find for the $`W`$\- $`A`$\- and ghost- propagators from the point of view of the OPE.
Fig. 2. Schematic representation of leading contributions that give rise to power corrections at large momenta in the gauge invariant correlators Eq. (30). For reasons given in the main text, the power leading correction $`g^2v`$ of the propagators cancels in this gauge invariant combination.
On the other hand, since $`O_2`$ is not invariant under global $`SU(2)`$ transformations, we do not expect $`O_2`$ to appear in the OPE of gauge invariant correlators. The leading power correction $`v`$ in the propagators therefore should cancel in gauge invariant correlation functions such as
$$(G_{\mu \nu }G_{\rho \sigma }+G_{\mu \nu }^aG_{\rho \sigma }^a)(G_{\alpha \beta }G_{\gamma \delta }+G_{\alpha \beta }^bG_{\gamma \delta }^b)_p\mathrm{}$$
(30)
To leading order in the loop expansion this can be verified explicitly. To this order, the three diagrams of Fig. 2 with two (transverse) photons and two (transverse) vector bosons as intermediate states lead to power corrections. In the limit $`p^2\mathrm{}`$ at least one of the photons, respectively vector bosons in the loop integrals has momentum much larger than $`v`$. The leading power correction $`v`$ in Eq. (30) from the vector bosons and the photon thus cancel if and only if the photon polarization,
$$\mathrm{\Gamma }_{\mu \nu }^{AA}(p,v)=(\delta _{\mu \nu }p^2p_\mu p_\nu )\mathrm{\Pi }^{AA}(p^2,v)$$
(31)
for $`p^2v`$ has the asymptotic power expansion
$$\mathrm{\Pi }^{AA}(p^2v,v)12\frac{m_W^2}{p^2}+O(g^2\mathrm{ln}p^2,v^2/p^4).$$
(32)
Here $`m_W^2=g^2\frac{\sqrt{\alpha v^2}}{16\pi }`$ is the $`W`$-boson mass Eq. (27). Evaluating the one ghost loop contributions to the photon polarization at high external momentum one indeed verifies Eq. (32). Note that the factor $`2`$ in Eq. (32) is essential for the cancellation of the leading power correction in Eq. (30), since twice as many $`W`$’s as photons contribute. Power corrections of order $`\alpha v^2/p^4`$ in the asymptotic expansion of Eq. (30) do not cancel and are related to $`\theta _\mu ^\mu `$ by Eq. (28).
We have here proposed a mechanism by which the $`W`$-bosons of an SU(2) gauge theory in Abelian gauges essentially become massive while the leading power corrections to gauge-invariant correlation functions nevertheless are of order $`\mathrm{\Lambda }_{\overline{MS}}^4/p^4`$ only. Although numerical lattice simulations show similar effects for the off-diagonal gluons, the numerical gauge fixing to MAG is not described by a local effective action and the results therefore cannot be directly compared with the ones presented here. One unfortunately cannot even extract the anomalous dimension of $`m_W^2`$ from the present numerical studies due to their rather narrow range of couplings. Although lattice simulations presently do not unambiguously confirm the mechanism of mass generation by ghost condensation I discussed, the simplicity and inherent consistency of this approach may warrant further study.
## References |
warning/0003/quant-ph0003046.html | ar5iv | text | # Strict Holism in a Quantum Superposition of Macroscopic States
## I Introduction
In his famous elementary textbook, Richard Feynman claims that the only mystery of quantum mechanics is exemplified by the electron self-interference in the two-slit experiment Feynmann . Interference is a consequence of the superposition principle, and indeed most of the puzzling aspects of quantum mechanics are related to the superposition of two or more states, as is the case in the Einstein-Podolsky-Rosen (EPR) paradox EPR or the Greenberger-Horne-Zeilinger (GHZ) theorem GHZ . Both EPR and GHZ show a striking characteristic of quantum mechanics: the nonseparability of systems situated far apart from each other. In quantum mechanics, systems that interacted with each other in the past may become entangled, and, even if they are separated by a great distance later on, their properties can be correlated in a way that would evade any attempt to give a classical explanation Bell . This nonseparability has as a consequence the nonexistence of a joint probability distribution, and hence of a local hidden-variable theory, that explains the outcome of the experiments suppeszannoti . More recently, Mermin Mermin showed that if we allow states with a large number $`N`$ of particles to be superposed in a way similar to the superposition of particles in the GHZ theorem, then quantum mechanics deviates exponentially with $`N`$ from the classical case (i.e., one that could be understood by a local hidden-variable).
The nonexistence of local-hidden variables that can account for all the experimental outcomes suggests that quantum mechanics has some holistic characteristic. Holism is the idea that the whole cannot be considered as the sum of its individual parts. The fact that systems far apart are nonseparable has led some authors to suggest that quantum mechanics has in its core a holistic characteristic Ghirardi ; Primas . Nonseparability, in the sense used in EPR or GHZ, means that a local hidden-variable theory that predicts the outcome of the experiments is impossible. Of course, nonseparability implies holism, but that the converse is not true is what we show in this paper. To do this, we will first show that a GHZ $`N`$-particle quantum mechanical system behaves in a deterministic way, when considered as a whole, but that every proper subsystem of this system behaves in a completely random way. This is done by first showing that any subsystem has maximal entropy, whereas the whole system has entropy zero. Then, we analyze, from a probabilistic point of view, the $`N`$-particle GHZ example. We show that quantum mechanics is more restrictive on the subsystems than pure probability considerations, even though, for the particular observables in question, a joint probability distribution exists. Then, we propose a definition of holism that is distinct from the concept of separability, and discuss this definition by means of simple examples. Our definition of holism is satisfied by the GHZ quantum mechanical system presented earlier.
## II Quantum Mechanical Holism
Let us start with the entangled GHZ-like $`N`$-particle state
$$|\psi =\frac{1}{\sqrt{2}}\left[\underset{k=1}{\overset{N}{}}|+_k+\underset{k=1}{\overset{N}{}}|_k\right],$$
(1)
where $`\widehat{\sigma }_{iz}|+_i=|+_i`$, $`\widehat{\sigma }_{iz}|_i=|_i`$, with $`\widehat{\sigma }_{iz}`$ being the spin operator in the $`z`$ direction acting on the $`i`$-th particle. It is easy to show that this state is an eigenstate, with eigenvalue $`1`$, of the observable operator
$$\widehat{\mathrm{\Sigma }}=\widehat{\sigma }_{1x}\widehat{\sigma }_{2x}\mathrm{}\widehat{\sigma }_{Nx}.$$
(2)
In other words, the observable $`\widehat{\mathrm{\Sigma }}`$, made out of the product of all $`N`$ spin observables, is deterministic, as a measurement of it always results in the value $`1.`$ In a similar way, this determinism is also true for the observables
$$\underset{i}{}\widehat{\sigma }_{iy}\underset{j}{}\widehat{\sigma }_{jx},$$
(3)
where the index $`i`$ is any subset with even cardinality of $`2^{\{1,2,\mathrm{},N\}}`$, and $`j`$ is the complement of $`i.`$
The state (1) has been the focus of several interesting papers, all of them related to the deterministic aspects of the above observables GHZ ; Mermin ; GHZ1 ; GHZ2 ; GHZ3 ; GHZ4 ; GHZ5 ; GHZ6 . However, in this paper we will be interested in observables acting only on a subset of the set of all particles in (1). We start with the following.
Given the ket
$$|\psi =\frac{1}{\sqrt{2}}(|++\mathrm{}++|\mathrm{}),$$
(4)
and the spin operators $`\widehat{\sigma }_{id}`$, where $`i=1\mathrm{}N`$ and $`d=x,y,z,`$ then any product of $`n<N`$ distinct spin operators has expectation zero.
*Proof.* Let us start with a hermitian operator $`\widehat{\mathrm{\Sigma }}^{^{}}`$ that is the product of $`n<N`$ distinct spin operators, such that we can write $`\widehat{\mathrm{\Sigma }}^{^{}}`$ as
$$\widehat{\mathrm{\Sigma }}^{^{}}=\underset{k=1}{\overset{a}{}}\widehat{\sigma }_{k,x}\underset{k=a+1}{\overset{b}{}}\widehat{\sigma }_{k,y}\underset{k=b+1}{\overset{c}{}}\widehat{\sigma }_{k,z}\underset{k=c+1}{\overset{N}{}}\widehat{1}_k,$$
(5)
with $`0<a<b<c<n`$, and $`a+b+c=n`$. We want to compute $`\psi |\widehat{\mathrm{\Sigma }}^{^{}}|\psi ,`$ the expected value of this operator, so
$`\psi |\widehat{\mathrm{\Sigma }}^{^{}}|\psi `$ $`=`$ $`{\displaystyle \frac{1}{2}}i^{ba1}[{\displaystyle \underset{k=1}{\overset{N}{}}}+|_k+{\displaystyle \underset{k=1}{\overset{N}{}}}|_k]\times `$ (6)
$`[{\displaystyle \underset{k=1}{\overset{b}{}}}|_k{\displaystyle \underset{k=b+1}{\overset{N}{}}}|+_k`$
$`(1)^{c+a}i^a{\displaystyle \underset{k=1}{\overset{b}{}}}|+_k{\displaystyle \underset{k=b+1}{\overset{N}{}}}|_k].`$
From the equation above, it is immediate that the inner product is zero if $`b<N,`$ as we wanted to prove.
Proposition 1 shows that the correlations for the $`N`$-particle system are quite strange. We have a set of $`N`$ particles that has always the same observable associated to its totality, but when we look at any of its parts, then the parts are completely uncorrelated. In this system the presence of a nonzero correlation appears only when we look at the system as a whole, and not at its parts. In the next section we will analyze in details the probabilistic properties of the probability distribution associated to, say, the operator $`\widehat{\mathrm{\Sigma }}`$.
## III Probabilistic Properties
It is interesting to note the consequences of the previous result. Say we are measuring the spin in the $`x`$ direction for $`n<N`$ particles. In this case all the particles are independent, and also behave in a completely random way, as the probability of measuring $`1`$ is the same as the probability of measuring $`1`$. However, if we measure the spin of *all* $`N`$ particles, the whole system is deterministic in a sense that will be made clear later. First, let us start with the following Proposition.
Let
$$|\psi =\frac{1}{\sqrt{2}}(|++\mathrm{}++|\mathrm{}),$$
(7)
$`\widehat{\mathrm{\Sigma }}=_{k=1}^N\widehat{\sigma }_{k,x},`$ and to each particle $`i`$, $`1iN`$, we associate the random variable $`𝐒_i`$, representing the value of its spin measurements, taking values $`\pm 1`$. If $`t=n\mathrm{\Delta }t`$, $`n=0,1,2,\mathrm{}`$ and we measure $`|\psi `$ using $`\widehat{\mathrm{\Sigma }}`$ at each $`t`$. We define the random variables $`𝐗_t^{\{k\}}=_{\{k\}}𝐒_k`$, where $`\{k\}`$ is any proper subset of $`\{1,\mathrm{},N\}`$ and $`𝐗_t=_{k=1}^N𝐒_k`$. Then each $`𝐗_t^{\{k\}}`$, and $`𝐗_t`$ define Bernouilli processes.
*Proof.* First we should note that $`|\psi `$ is an eigenstate of $`\widehat{\mathrm{\Sigma }}`$, such that we can measure $`\widehat{\mathrm{\Sigma }}`$ as many times as we want without affecting $`|\psi `$. If we keep measuring spin in the $`x`$ direction for all particles in equal intervals of time $`\mathrm{\Delta }t`$, we can make a data table for the experimental result that would look like Table 1, where we associate to each of the spin measurements for particle $`i`$ the random variable $`𝐒_i`$ taking values $`\pm 1`$.
Each column of this table would be completely uncorrelated to the any other column or combinations of columns with less than $`N`$ columns involved. Similar independence and randomness hold for any row of length at most $`N1`$, i.e., at least one entry is deleted. However, if we multiply $`𝐒_1`$, $`𝐒_2`$,*$`\mathrm{}`$*, $`𝐒_N`$, we always obtain the same value $`_{i=1}^N𝐒_i=1`$. Furthermore, since the wave function $`|\psi `$ is unchanged, the equal probabilities of obtaining a $`1`$ or $`1`$ for each of the columns or shortened rows are also unchanged. As a consequence, the temporal sequence of product random variables $`𝐗_t^{\{k\}}=_{\{k\}}𝐒_k`$, where $`\{k\}`$ is any proper subset of $`\{1,\mathrm{},N\}`$, form a Bernouilli process, i.e. at each time $`t`$ the random variables $`𝐗_t^{\{k\}}`$ are independently and identically distributed, as we wanted to show. It is straightfoward to extend the same argument to $`𝐗_t.`$
We are now in a position to make explicit the statement that the system as a whole is deterministic and its subsystems are random.
The random variables $`𝐗_t^{\{k\}}=_{\{k\}}𝐒_k`$, where $`\{k\}`$ is any proper subset of $`\{1,\mathrm{},N\}`$, defined in a way similar to Proposition 2, have maximal entropy for such process, whereas the random variable $`𝐗_t=_{k=1}^N𝐒_k`$ has zero entropy.
*Proof.* Since both $`𝐗_t^{\{k\}}`$ and $`𝐗_t`$ define a Bernouilli process, their entropy is $`H=p_i\mathrm{log}p_i`$, where $`p_i`$ is the probability of each possible outcome, in this case $`\pm 1`$. $`𝐗_t=_{i=1}^N𝐒_i`$, representing the system as a whole, has entropy zero, since for all $`t`$ $`P(𝐗_t=1)=1`$ and $`P(𝐗_t=1)=0`$. Yet, any proper subset $`\{k\}`$ of $`\{1,\mathrm{},N\}`$ will define a random variable $`𝐗_t^{\{k\}}=_{\{k\}}𝐒_k`$ whose entropy is maximal for such a process, as $`P(𝐗^{\{k\}}=1)=1/2`$ and $`P(𝐗^{\{k\}}=1)=1/2`$, i.e. the entropy $`H=p_i\mathrm{log}p_i=1`$, where $`\mathrm{log}`$ is to base 2, as we wanted to prove.
The results just obtained show that the system in question is strongly holistic, in the sense that a measurement of $`\widehat{\mathrm{\Sigma }}`$ containing *all* particles in the system yields a deterministic result, whereas any spin measurement made on a subsystem has a perfectly random outcome. However, since we can measure all the $`N`$ spin values simultaneously, we can also write a data table for the experimental outcomes, and a joint probability distribution exists. In this sense, the system is holistic but is separable, as we can factor the joint probability distribution.
Even though a joint probability distribution exists, we stress that such a strange distribution, where only when we consider all particles is the system deterministic, is rarely if ever found in any empirical domain. In fact, quantum mechanics provides, as far as we know, the only example in nature of a case where we have perfect correlation for a triple and zero correlation for pairs. This is the case if we take a three-particle GHZ system, as it yields $`𝐗_i`$ $`\pm 1`$ random variables, with $`E(𝐗_1𝐗_2𝐗_3)=1`$, $`E(𝐗_i)=0`$, $`i=1,\mathrm{},3`$. It is also interesting to stress that, in the three-particle GHZ case, the pair correlations are zero as a consequence of the triple correlation and the individual expectations. This can be verified by direct computation. Say we have $`E(𝐗_1𝐗_2𝐗_3)=1`$. Then, all terms with 0 or 2 negative components sum to 1, i.e.,
$$x_1x_2x_3+\overline{x}_1\overline{x}_2x_3+\overline{x}_1x_2\overline{x}_3+x_1\overline{x}_2\overline{x}_3=1,$$
(8)
where we use the notation $`x_1`$to represent $`P(𝐗_1=1)`$, $`\overline{x}_1`$to represent $`P(𝐗_1=1)`$, $`\overline{x}_1x_2`$ to represent $`P(𝐗_1=1,𝐗_2=1)`$, and so on. We also have that
$`x_1x_2=x_1x_2x_3=x_1x_3=x_2x_3`$ $`=`$ $`a,`$ (9)
$`\overline{x}_1\overline{x}_2=\overline{x}_1\overline{x}_2x_3=\overline{x}_1x_3=\overline{x}_2x_3`$ $`=`$ $`b,`$ (10)
$`\overline{x}_1x_2=\overline{x}_1x_2\overline{x}_3=\overline{x}_1\overline{x}_3=x_2\overline{x}_3`$ $`=`$ $`c,`$ (11)
$`x_1\overline{x}_2=x_1\overline{x}_2\overline{x}_3=x_1\overline{x}_3=\overline{x}_2\overline{x}_3`$ $`=`$ $`d,`$ (12)
with $`a+b+c+d=1.`$ Next, from (9)–(12), $`x_1=a+d,`$ $`\overline{x}_1=b+c,`$ $`x_2=a+c,`$ $`\overline{x}_2=b+d,`$ $`x_3=a+b,`$ $`\overline{x}_3=c+d`$, and from $`E(𝐗_i)=0`$, $`x_1=x_2=\overline{x}_1=\overline{x}_2=\frac{1}{2}`$. From (9)–(12) and the following equations, we obtain at once $`a=b=c=d`$ and
$$E(𝐗_1𝐗_2)=E(𝐗_2𝐗_3)=E(𝐗_1𝐗_3)=0.$$
(13)
However, contrary to the three-particle case, if we increase the number of particles to four, the correlations are not dictated by $`E(𝐗_1𝐗_2𝐗_3𝐗_4)=1`$, $`E(𝐗_i)=0`$, $`i=1,\mathrm{},4`$ anymore. For the four-particle case, we can compute, in a manner similar to the three-particle one, that $`E(𝐗_i𝐗_j𝐗_k)=0`$, $`i<j<k.`$ However, the correlations $`E(𝐗_i𝐗_j)`$ can individually, but not independently, take any value in the closed interval $`[1,1]`$. On the other hand, if all the correlations are zero, then the positive atoms have a uniform distribution, by an argument similar to the one given above. In fact, we can show the following.
Given $`E(𝐗_1\mathrm{}𝐗_n)=0`$ and the product of any nonempty subset of the random variables $`𝐗_1\mathrm{}𝐗_n`$ also has expectation zero, including $`E(𝐗_i)=0`$, $`1in`$. Then the $`2^n`$ atoms of the probability space supporting $`𝐗_1\mathrm{}𝐗_n`$ has a uniform probability distribution, i.e., each atom has probability $`1/2^n.`$
*Proof.* We show this by induction. For $`n=1`$, we have by hypothesis that $`E(𝐗_i)=0`$, so, as required, $`P(𝐗_i=1)=x_1=1/2.`$ Next, our inductive hypothesis is that for every subsystem having $`m<n`$, the $`2^m`$ atoms have a uniform distribution, and we need to show this holds for $`n`$. Using the induction hypothesis for $`n1`$, we have at once the following pair of equations:
$`x_1x_2\mathrm{}x_{n1}`$ $`=`$ $`x_1x_2\mathrm{}x_{n1}x_n+x_1x_2\mathrm{}x_{n1}\overline{x}_n=2^{1n},`$
$`x_1x_2\mathrm{}x_{n2}x_n`$ $`=`$ $`x_1x_2\mathrm{}x_{n1}x_n+x_1x_2\mathrm{}\overline{x}_{n1}x_n=2^{1n}.`$
Subtracting one equation from the other we have at once $`x_1x_2\mathrm{}\overline{x}_{n1}x_n=x_1x_2\mathrm{}x_{n1}\overline{x}_n.`$ By similar arguments, we show that all atoms that have exactly one negative value of $`\overline{x}_i`$ for the $`n`$-particle case are equal in probability. Moreover, without any new complication this argument extends to equal probability for any atom having exactly $`k`$ negative values, $`2kn`$.
Next, we can easily show that those atoms differing by 2, and therefore by an even number of, negative values have equal probability. We give the argument for $`k=0`$ and $`k=2`$:
$`x_1x_2\mathrm{}x_{n1}`$ $`=`$ $`x_1x_2\mathrm{}x_{n1}x_n+x_1x_2\mathrm{}x_{n1}\overline{x}_n=2^{1n},`$
$`\overline{x}_1x_2\mathrm{}x_{n2}x_n`$ $`=`$ $`\overline{x}_1x_2\mathrm{}x_{n1}x_n+\overline{x}_1x_2\mathrm{}\overline{x}_{n1}x_n=2^{1n}.`$
Using the previous result and subtracting we get $`x_1x_2\mathrm{}x_{n1}x_n=\overline{x}_1x_2\mathrm{}x_{n1}\overline{x}_n.`$ Finally, we use the hypothesis that $`E(𝐗_1\mathrm{}𝐗_n)=0`$. This zero expectation requires that the sum of all the terms with 0 or an even number of negative values have the same sum as all the terms with an odd number of negative values. This implies at once that all atoms have equal probability, and so each has probability $`1/2^n`$, proving Proposition 4.
We also prove a more restricted result, but a sifnificant one, by purely probabilistic means, i.e., no quantum mechanical concepts or assumptions are needed in the proof.
Given $`E(𝐗_1\mathrm{}𝐗_N)=\pm 1`$ and $`E(𝐗_i)=0`$, $`i=1,\mathrm{},N`$, then any correlation of $`N1`$ particles is zero, e.g., $`E(𝐗_1\mathrm{}𝐗_{N1})=0`$, $`E(𝐗_1\mathrm{}𝐗_{N2}𝐗_N)=0`$, etc.
*Proof.* We give the proof for $`E(𝐗_1\mathrm{}𝐗_N)=1`$. Then there are $`2^N`$ atoms in the probability space. Given the expectation equal to 1, half ot the atoms must have probability 0, namely all those representing negative spin products. Now, we consider all the terms expressing $`E(𝐗_1\mathrm{}𝐗_{N1})`$. On the positive side, we have all those with even or zero negative values:
$$x_1x_2\mathrm{}x_{N1}+\overline{x}_1\overline{x}_2\mathrm{}x_{N1}+\mathrm{}+\overline{x}_1\overline{x}_2\mathrm{}\overline{x}_{N1}$$
(14)
if $`N1`$ is even and as the last term if $`N1`$ is odd $`x_1\overline{x}_2\mathrm{}\overline{x}_{N1}`$. To be extended to atoms, a positive $`x_N`$ must be added. So, in probability
$$x_1x_2\mathrm{}x_{N1}=x_1x_2\mathrm{}x_{N1}x_N,$$
because, given $`E(𝐗_1\mathrm{}𝐗_N)=1`$
$$x_1x_2\mathrm{}x_{N1}\overline{x}_N=0,$$
and similar for the other terms in (14).
The same thing applies in similar fashion to the negative side, e.g.,
$$\overline{x}_1x_2\mathrm{}x_{N1}=\overline{x}_1x_2\mathrm{}x_{N1}\overline{x}_N,$$
since the atom on the right must have zero or an even number of negative values.
But we observe that, by hypothesis, $`E(𝐗_1\mathrm{}𝐗_{N2}𝐗_N)=0`$, but the probability $`x_N`$ is just equal to the sum of the probabilities of the positive terms of $`E(𝐗_1\mathrm{}𝐗_{N2}𝐗_N)`$ and $`\overline{x}_N`$ is just equal to the sum of the probabilities of the negative terms above. Since, $`x_N\overline{x}_N=0`$, we conclude $`E(𝐗_1\mathrm{}𝐗_{N1})=0`$. The same argument can be extended to the other $`N1`$ combinations of $`𝐗_i`$, and this completes the proof.
## IV $`\mathrm{\Pi }`$-Holism
The remarkable property that a quantum system has a perfect correlation for its whole but a totally random behavior for *any* of its part seems to us to represent a holistic characteristic of quantum mechanics. This holism is, however, quite distinct from what is known in the literature as separability. For that reason, we propose the following definition for strict holism.
Let $`\mathrm{\Omega }=(\mathrm{\Omega },,𝒫)`$ be a finite probability space and let $`𝐅=\left\{𝐗_i,1iN\right\}`$ be a family of $`\pm 1`$ random variables defined on $`\mathrm{\Omega }`$. Let $`\mathrm{\Pi }`$ be a property defined for finite families of random variables. Then $`𝐅`$ is strictly $`\mathrm{\Pi }`$-holistic iff
(i) $`𝐅`$ has $`\mathrm{\Pi }`$;
(ii) No subfamily of $`𝐅`$ has $`\mathrm{\Pi }`$.
Moreover, if $`\mathrm{\Pi }`$ is a numerical property,
(iii) No subfamily of $`𝐅`$ approximates $`\mathrm{\Pi }`$.
To understand this definition, let us give some examples from classical mechanics. It is well know in classical gravitation theory that a two-particle system has a well defined solution. However, if we add to this system an extra particle, no closed solutions to this system exist in some cases, and in fact its behavior can be completely random Alekseev . One may be tempted to think that this chaotic behavior is a holistic property, but according to the definition above, it is not. For instance, let us take the restricted three-body problem analyzed by Alekseev, where two particles with large mass orbit around their Center of Mass (CM), while a third small particle oscillates in a line passing through the CM and perpendicular to the plane of orbit of the two large masses. The whole system behaves randomly, as well at least one subsystem, the one defined by the small particle. Hence, this system is not $`\mathrm{\Pi }`$-holistic, if we choose $`\mathrm{\Pi }`$ to be the property of being random.
As yet another example, let us consider a glass of water. The water is a large system that does not behave like a water molecule, but in a coordinated way dictated by hydrodynamics. Is then this system holistic? If we take, say, half the glass of water, the properties of this half of water are the same as the whole glass, except its mass, hence the system is not $`\mathrm{\Pi }`$-holistic for the other macroscopic properties of the water. What about properties like, say, mass? Say we take the full glass and remove only a water molecule from it. The new subsystem approximates the mass of the original one, violating hypothesis (iii) from the Definition, and so if we choose $`\mathrm{\Pi }`$ to be the property mass, the system is not $`\mathrm{\Pi }`$-holistic.
Let $`𝐅=\{𝐒_i,i=1,\mathrm{},N\}`$ be the set of random variables of all the spin measurements of the state
$$|\psi =\frac{1}{\sqrt{2}}(|++\mathrm{}++|\mathrm{}),$$
and let $`𝐗_t`$ be the product random variable of Proposition 3, and let $`𝐗_t^{\{k\}}`$ be the product random variable of any subfamily $`\{k\}`$ as defined earlier. Let the entropy be the $`\mathrm{\Pi }`$ property of these product random variables. Then $`𝐅`$ is $`\mathrm{\Pi }`$-holistic.
*Proof.* Immediate, from Proposition 3, since the entropy of $`𝐗_𝐭`$ is $`0`$ and, for any *$`\{k\}`$,* the entropy of $`𝐗_𝐭^{\{k\}}`$ is $`1`$.
## V Final Remarks
To summarize, we found that an $`N`$-particle GHZ state has a strong holistic property. However, it may be difficult to detect experimentally a quantum mechanical holistic characteristic with a large number of particles, as decoherence may play an important role, given that the decoherence time decreases rapidly if we increase the number of particles decoherence1 ; decoherence2 ; decoherence3 . A promising setup where this holism could be verified for a reasonably large number of particles is the one proposed by Cirac and Zoller Garg ; Garg2 . We found that for $`N4`$, the measurements of $`E(𝐗_1𝐗_2𝐗_3\mathrm{}𝐗_N)`$ and of $`E(𝐗_i)`$ do not fix a probability distribution, and extra measurements are necessary for the pairs, triples, and so on, for the probability distribution to be fixed. We believe that these measurements, which should yield many zero correlations, could be used to put additional constraints on some local-hidden variable models that exploit the detection loophole detectionloophole ; loophole2 ; loophole3 ; loophole4 . |
warning/0003/cond-mat0003299.html | ar5iv | text | # Defect-unbinding and the Bose-glass transition in layered superconductors
\[
## Abstract
The low-field Bose–glass transition temperature in heavy-ion irradiated Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+δ</sub> increases progressively with increasing density of irradiation–induced columnar defects, but saturates for densities in excess of $`1.5\times 10^9`$ cm<sup>-2</sup>. The maximum Bose-glass temperature corresponds to that above which diffusion of two–dimensional pancake vortices between different vortex lines becomes possible, and above which the “line–like” character of vortices is lost. We develop a description of the Bose–glass line that is in excellent quantitative agreement with the experimental line obtained for widely different values of track density and material parameters.
\]
Heavy-ion irradiated (HII) layered superconductors have recently been at the focus of attention, because the irradiation–induced amorphous columnar tracks help overcome the detrimental effects of the high material anisotropy , and partially re-establish long-range superconducting phase order . In the layered superconductor Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+δ</sub>, heavy–ion irradiation increases the irreversibility field $`B_{irr}(T)`$ below which the $`I(V)`$-curve is nonlinear due to vortex pinning on the tracks to values well above the field $`B_{FOT}(T)`$ at which the first order vortex lattice–to liquid transition field takes place in the pristine material . At inductions $`B_{FOT}(T)BB_{irr}(T)`$, the irradiated superconductor displays the phenomenology of the Bose–glass phase of localized vortices; moreover, the transport properties show a distinct anisotropy related to the presence of the tracks that is absent in the pristine material , and which suggests that the vortices behave as well-defined separate lines, i.e. vortex lines in the Bose-glass phase are disentangled .
The position of $`B_{irr}(T)`$ and the occurence of flux–line entanglement were shown to be intimately related in moderately anisotropic HII superconductors such as YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> . There, $`B_{irr}(T)B_{BG}(T)`$ corresponds to the second order phase transition line between the Bose-glass and the vortex liquid ; $`B_{BG}(T)`$ progressively increases with increasing columnar defect density $`n_d`$, to an upper limit attained when $`n_d1\times 10^{11}\mathrm{cm}^2`$ (corresponding to an ion dose-equivalent “matching” field $`B_\varphi \mathrm{\Phi }_0n_d=2`$ T). Departing from the correspondence of vortex lines with the world lines of interacting bosons in two dimensions (2D) , it was argued that the upper limit of $`B_{BG}(T,n_d)`$ corresponds to the field beyond which the (entangled) vortex liquid becomes stable with respect to the introduction of linear defects. In this temperature and field regime, these produce a rather weak pinning due to intervortex repulsion and the averaging effect of vortex thermal excursions . The situation in layered superconductors such as Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+δ</sub> is apriori different. First, the weak coupling between adjacent CuO<sub>2</sub> bilayers (with separation $`s`$ 1.5 nm) implies that vortex lines are extremely soft. Over the larger part of the first Brillouin zone of the vortex lattice, the contribution of the dipole interaction between “pancake” vortices in adjacent bilayers to the tilt modulus $`c_{44}`$ is expected to exceed that of the line tension $`\epsilon _1`$, which is determined by the interlayer Josephson effect . Flux lines are then better described as stacks of pancakes and the analogy with the 2D boson system is no longer valid . Another consequence is that pinning near $`B_{irr}(T)`$ in HII Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+δ</sub> is not weak: at fields $`B\frac{1}{6}B_\varphi `$ pancake vortices gain maximum free energy by remaining localized on the columnar defects, even in the vortex liquid phase . Nevertheless, a well–defined transition from very slow nonlinear vortex dynamics in the Bose–glass , to Ohmic response in the vortex liquid does exist near the irreversibility line in Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+δ</sub>.
In order to cast light on the mechanism of this transition, we have measured $`B_{irr}(T)`$ in HII Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+δ</sub> for widely varying track density and oxygen content. The latter determines the values of the anisotropy parameter $`\gamma `$ and the penetration depth $`\lambda _{ab}(T)`$ : both increase as one decreases the oxygen content towards optimal doping. It turns out that the phenomenological behavior of $`B_{irr}(T)`$ at low fields is rather similar to that found in YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub>. $`B_{irr}(T)`$ increases progressively towards higher values with increasing defect density, but saturates for $`n_d1.5\times 10^9`$ cm<sup>-2</sup>, or $`B_\varphi 30`$ mT. For higher matching fields, the low–field portion ($`B\frac{1}{6}B_\varphi `$) of $`B_{irr}(T)`$ adopts a strictly exponential temperature dependence; although $`T_{irr}(B)/T_c`$ strongly increases when the value of $`\lambda _{ab}`$ decreases, we will see that this increase is not due to the increase of the pinning energy.
Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8</sub> crystals were grown at the University of Tokyo using the travelling-solvent floating-zone method, and then postannealed at either 800C or 500 C. This produces $`T_c`$’s of 89 K (optimal doping) and 82 K (overdoped) respectively . The crystals were cut to rectangles of size 500 (l) $`\times `$ 400 (w) $`\times `$ 20 (t) $`\mu `$m<sup>3</sup>, and irradiated with varying fluences of 5.8 GeV Pb ions at the Grand Accélérateur National d’Ions Lourds (GANIL) at Caen, France. The ion beam was directed parallel to the sample $`c`$-axis. Each ion impact created an amorphous columnar track of radius 3.5 nm traversing the sample along its entire thickness. Samples were prepared with track densities $`1\times 10^9`$ cm$`{}_{}{}^{2}<n_d<2\times 10^{11}`$ cm<sup>-2</sup>, corresponding to 20 mT $`B_\varphi `$ 4 T. The irradiation caused $`T_c`$ to decrease: $`T_c/n_d=3.6\times 10^{11}`$ Kcm<sup>2</sup> or $`T_c/B_\varphi =1.8`$ KT<sup>-1</sup>.
Subsequent measurements were performed using the Local Hall Probe Magnetometer in AC mode . A small ac field of amplitude $`h_{ac}=1`$ G and frequency $`f=7.75`$ Hz is applied parallel to the sample $`c`$–axis, colinearly with the DC field used to create the vortices. The ac field leads to a periodic electric field gradient of magnitude $`2\pi \mu _0h_{ac}f`$ across the sample. Using a miniature Hall probe one measures the RMS induction $`B_{ac}(f,T)`$ at the center of the sample top surface, which is simply related to the sample screening current .
Figure 1 shows the fundamental and third harmonic transmittivities, defined as $`T_H[B_{ac}(f,T)B_{ac}(f,TT_c)]/B`$ and $`T_{H3}B_{ac}(3f,T)/B`$ respectively, with
$`BB_{ac}(f,TT_c)B_{ac}(f,TT_c)`$ , measured for the optimally doped crystal with $`B_\varphi =2`$ T. The presence of a third harmonic response $`|T_{H3}|`$ implies the non–linearity of the sample’s $`I(V)`$–characteristic. The fields $`B_{irr}(T)`$, or temperatures $`T_{irr}(B)`$, below which the third harmonic signal can be first observed upon cooling are plotted in Fig. 2. At $`B_{irr}`$, the working point enters the nonlinear part of the sample $`I(V)`$ curve at an electric field of the order of $`10^7`$ Vm<sup>-1</sup>; this corresponds to a voltage drop of $``$ 50 pV across the sample. At such low voltages, the measurement of the ac screening current may be sensitivity–limited . In practice, the use of $`h_{ac}=1`$ Oe means that the minimum measurable current density $`j_{min}5\times 10^2`$ Am<sup>-2</sup>, comparable to a transport current of 0.1 mA. The coincidence of our $`B_{irr}(T)`$ with the $`B_{BG}`$–data determined by Seow et al. indicates that for all practical purposes $`B_{irr}`$ is a good approximation of the Bose-glass transition field (Fig. 2).
The evolution of $`B_{irr}(T)`$ with $`B_\varphi `$ is plotted in Fig. 3. For $`B_\varphi B_\varphi ^{min}=30`$ mT, $`B_{irr}(T)`$ increases monotonically with increasing $`B_\varphi `$ . As long as $`B_{irr}(T)<B_\varphi `$, it depends exponentially on temperature, while for $`B_{irr}>B_\varphi `$, $`B_{irr}(1T/T_c)`$. For large $`B_\varphi B_\varphi ^{min}`$, one can distinguish three distinct sections of $`B_{irr}(T)`$. At all but the very lowest fields, $`B_{irr}`$ again depends exponentially on $`T`$, but with the specificity that it is independent of $`B_\varphi `$. In this regime \[(I) in Fig. 2\], there thus exists a an upper limit of $`B_{BG}(T)`$ in Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8</sub>, as is the case in YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> . At $`B_{irr}B_{int}\frac{1}{6}B_\varphi `$, the exponential increase abruptly changes into a nearly vertical rise (regime II); $`B_{int}`$ is the field at which intervortex repulsion start to determine the pinned vortex configuration. This transition is also manifest in the vortex liquid phase as that at which a “recoupling” transition was measured using Josephson Plasma Resonance (JPR) . As $`B_{irr}`$ increases to a sizeable (but not constant) fraction of $`B_\varphi `$, a weak temperature dependence is once
again adopted, $`B_{irr}(1T/T_c)^\alpha `$ with $`\alpha 1`$ (regime III). The behavior of the irreversibility line for different oxygen content is shown in Fig. 4. The exponential decrease at high $`T`$ is steeper for the overdoped crystal, which has the smaller $`\lambda _{ab}`$ and $`\gamma `$, and therefore the larger condensation energy $`\epsilon _0/4\pi \xi ^2`$, and the stronger intervortex– and vortex-defect interaction (both are proportional to the typical vortex energy scale $`\epsilon _0=\mathrm{\Phi }_0^2/4\pi \mu _0\lambda _{ab}^2`$; $`\xi `$ is the coherence length).
In order to describe the low–field exponential temperature dependence of $`B_{irr}(T)`$, we exploit the fact that in this regime (I) vortex interactions are irrelevant to the column occupation . In other words, each vortex line can become localized on an appropriate defect site. If a sufficient number of columns is available to every line (i.e. at large $`B_\varphi `$), extra free energy can be gained by the redistribution of pancakes constituting a given line over different columns. At low $`B_\varphi 0.5`$ T this entropy gain is insufficient to balance the loss in vortex interaction energy, and pancakes belonging to the same line remain aligned on the same columnar defect. In either case, all pancakes are localized on a column. The superposition of the exponential portions of $`B_{irr}(T)`$ for all matching fields 30 mT $`<B_\varphi <4`$ T implies that at $`B_{irr}`$ the redistribution over different column sites is irrelevant for the pancake delocalization mechanism — i.e. at and below $`B_{irr}(T)`$, pancakes belonging to the same vortex line are well aligned on the same site, belonging to a set of allowed “columnar-defect” sites. Then, we no longer need to consider other positions in the “intercolumn space” in our model description, which becomes that of a “discrete” superconductor. Since the allowed sites are more or less equivalent, the circumstance that they are in fact columnar defect sites becomes immaterial: namely, the free energy of all vortices is lowered by approximately the same amount. The low–field vortex state therefore does not differ fundamentally from that in the unirradiated crystal. Only, the vortex confinement in the defect potential, which plays the role of the “substrate potential” of Ref. , inhibits thermal line wandering and the elastic relaxation of the vortex lattice.
The main thermal excitations in this situation are expected to be small defects in the vortex lattice. In a Josephson–coupled layered superconductor, these amount to bound pancake vacancy–interstitial pairs within the same layer, i.e. the “quartets” of Ref. . In the HII layered superconductor, such a pair corresponds to the “exchange” of one or more pancakes between two sites. The energy of the quartets is
$$\epsilon _q4c_{66}a_0^2s\left(R/\mathrm{\Lambda }\right)^2\epsilon _0s\left(R/\mathrm{\Lambda }\right)^2,(R\mathrm{\Lambda })$$
(1)
with $`c_{66}`$ the vortex lattice shear modulus, $`a_0=(\mathrm{\Phi }_0/B)^{1/2}`$ the vortex spacing, $`R`$ the distance between a bound vacancy and interstitial, $`\mathrm{\Lambda }=[\lambda _{ab}^1+(\gamma s)^1]^1`$ the generalized penetration depth taking into account both magnetic and Josephson coupling, and $`\gamma s`$ the Josephson length . It was shown in Ref. that the glass transition in a layered superconductor corresponds to the pair–unbinding transition. Correspondingly, the Bose glass transition is the unbinding temperature of the dislocation pairs in the “discrete” superconductor, i.e.the temperature above which pancakes can diffuse from line to line; it can be estimated as $`k_BT_{BG}=\epsilon _q(R_l)`$, where $`R_ln_l^{1/2}`$ and $`n_l`$ is the equilibrium density of free dislocation pairs in the vortex liquid. Taking into account that only small pairs of size $`a_0`$ ( i.e. vacancies / interstitials) matter, one has $`na_0^2\mathrm{exp}(\epsilon _0s/k_BT)`$; the activation energy $`\epsilon _0s`$ is larger than that in the unirradiated superconductor because of the lack of lattice relaxation around the pair. Gathering terms, $`k_BT_{BG}=\epsilon _0s(a_0/\mathrm{\Lambda })^2\mathrm{exp}(\epsilon _0s/k_BT_{BG})`$ or
$$B_{BG}=B_\mathrm{\Lambda }\left(\frac{\epsilon _0s}{k_BT}\right)\mathrm{exp}\left(\frac{\epsilon _0s}{k_BT}\right),(B_\mathrm{\Lambda }BB_\varphi )$$
(2)
with $`B_\mathrm{\Lambda }=\mathrm{\Phi }_0/\mathrm{\Lambda }^2`$, $`\epsilon _0`$ and $`\mathrm{\Lambda }`$ to be evaluated at $`T_{BG}`$. The Bose–glass transition line does not depend on the details of the columnar defect potential such as pinning energy, column radius, or matching field. In Fig.4, we compare Eq. (2) to the experimental results for Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8</sub> with different oxygen content. We obtain excellent quantitative agreement using $`\lambda _{ab}`$–values from the literature and the single free parameter $`\gamma `$ (contained only in $`B_\mathrm{\Lambda }`$). The values $`\gamma =360`$ and 550 found for the overdoped and the optimally doped material respectively are well within accepted experimental limits. Thus, Eq. (2) reproduces the correct field, temperature, $`B_\varphi `$, and doping–dependence of the Bose-glass delocalization line.
The Bose-glass transition separates the low–$`T`$ phase in which pancake vortices can wander between columnar defects, but always remain bound to the same site, or vortex line, from the high $`T`$–phase in which diffusion of pancakes between sites (lines) is possible. Below $`T_{BG}`$, individual pancakes (“defect pairs”) cannot provide flux transport at low currents; vortex lines can only move as a whole, leading to the line–like behavior observed in Ref. . Oppositely, the pancake diffusion above $`T_{BG}`$ not only implies an Ohmic resistivity , but also that the linear nature of the vortex lines should no longer be apparent in the angular dependence of the transport properties. The intersite diffusion of pancakes (or “pancake–exchange”) means that at $`T>T_{BG}`$ vortex lines are effectively entangled, on a scale $`l=s\mathrm{exp}(\epsilon _0s/k_BT)`$. The upper limit of the Bose-glass transition in layered superconductors is thus analogous to that in moderately anisotropic compounds in that it represents the boundary at which vortices become delocalized into an entangled flux liquid and where the angular dependence of the columnar defects’ contribution to the resistivity vanishes.
The meaning of $`B_\mathrm{\Lambda }`$ becomes apparent from Eq. (1), where $`\mathrm{\Lambda }`$ appears as the typical interaction distance between small dislocation pairs. For separations $`\mathrm{\Lambda }`$, i.e. for matching fields $`B_\varphi B_\mathrm{\Lambda }`$, the formation of quartets demands a (shear) energy cost that is much greater than either $`\epsilon _0s`$ or $`k_BT`$. The excitation of a pancake onto another site is then unlikely even in the vortex liquid. Delocalization will occur as it does in a moderately anisotropic superconductor: thermal wandering of the vortices into the intercolumn space lowers their binding energy until they can be liberated from the columns by thermal activation . Similarly, in regime II ($`B>\frac{1}{6}B_\varphi `$), where the number of available tracks is insufficient not because of the intercolumn distance but because of track occupation by other vortices, delocalization is initiated by pancake wandering into the intercolumn space. Then, $`B_{irr}(T)`$ increases with decreasing $`B_\varphi `$ because the number of available sites decreases. The experimentally found condition $`B<\frac{1}{6}B_\varphi `$ on the validity of Eq. (2) also describes its validity range at low $`B_\varphi `$: $`B_{irr}`$ becomes ion–dose dependent when $`B_\varphi <B_\varphi ^{min}=6B_\mathrm{\Lambda }`$, in excellent agreement with the experimental values $`B_\varphi ^{min}30`$ mT and $`B_\mathrm{\Lambda }=7.0`$ mT (4.5 mT) found for overdoped (optimally doped) Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8</sub> respectively.
Summarizing, in the regime of individual vortex line pinning by columnar defects in a layered superconductor, the upper limit of the Bose–glass transition corresponds to the onset of pancake vortex diffusion between different lines and flux-line entanglement. However, in both the glass and the vortex liquid phases, only columnar defect sites are available to the pancakes; hence, a modelisation in terms of a “discrete superconductor” is appropriate. The onset of diffusion in this “discrete superconductor” presents a realization of the glass transition by unbinding of small defect pairs proposed in Ref. ; in the original continuous problem of an unirradiated layered superconductor, the defect–unbinding mechanism is masked by strong Gaussian thermal fluctuations .
We thank V.B. Geshkenbein, P.H. Kes, A.E. Koshelev, P. LeDoussal, and V.M. Vinokur for stimulating discussions. The work of M.V.F. was supported by DGA grant No. 94-1189. |
warning/0003/hep-ph0003019.html | ar5iv | text | # Chiral symmetry breaking, color superconductivity and quark matter phase diagram: a variational approach
## I Introduction
The structure of vacuum in Quantum Chromodynamics (QCD) is one of the most interesting questions in strong interaction physics . The evidence for quark and gluon condensates in vacuum is a reflection of its complex nature , whereas chiral symmetry breaking is an essential feature in the description of the low mass hadron properties. Due to the nonperturbative nature of QCD in this regime different effective models have been used to understand the nature of chiral symmetry breaking. These have been constructed, for the most part, in the framework of a Nambu-Jona- Lasinio (NJL) model with a four fermion interaction. There also have been attempts to generalise this in the case of Coulomb gauge QCD with an effective propagator simulating the effects of confining potentials . Studies at finite temperatures using imaginary time formulation of finite temperature field theory have also been carried out.
Recently there has been a lot of interest in strongly interacting matter at high densities. In particular, a color superconducting phase for it involving diquark condensates has been considered with a gap of about 100 MeV . The studies have been done with an effective four fermion interaction between quarks , direct instanton approach or a perturbative QCD calculation at finite density . There has also been a study of this phase in NJL model . In the NJL model however,the aspect of chiral symmetry breaking in presence of diquark condensates has not been considered. Such a question has been considered in Ref. in an instanton induced four fermion interaction model within a mean field approximation. In this note, we use a different method to study the problem. We consider a variational approach with an explicit assumption for the ground state having both quark antiquark and diquark condensates. The actual calculations are carried out for the NJL model such that the minimisation of the free energy density determines which condensate will exist at what density.
We organise this paper as follows. In section II we construct the ground state as a vacuum realignment with quark antiquark as well as diquark condensates. The finite temperature and density effects are included within the framework of thermofield dynamics through a thermal Bogoliubov transformation involving doubling of the Hilbert space. In section III we shall consider an effective model as is considered for chiral symmetry breaking in Coulomb gauge QCD . To solve the gap equation exactly we shall further take a four fermion point interaction limit of the same similar to NJL model. In section IV we solve the gap equations for the mass and the superconducting gap,determine the phase diagram and discuss the results. Finally, we give a summary in section V.
## II An ansatz for the ground state
As noted earlier we shall include here the effects of both chiral symmetry breaking as well as diquark pairing. For the consideration of chiral symmetry breaking, we shall take the perturbative vacuum state with chiral symmetry as $`|0`$. In this basis we shall take the quarks as massless. We shall then assume a specific vacuum realignment which breaks chiral symmetry because of interaction.
We have seen earlier that chiral symmetry breaking takes place with the formation of quark antiquark condensates in perturbative vacuum . We consider the quark field operator expansion in momentum space given as
$`\psi (𝐱)`$ $`{\displaystyle \frac{1}{(2\pi )^{3/2}}}{\displaystyle \stackrel{~}{\psi }(𝐤)e^{i𝐤𝐱}𝑑𝐤}`$ (1)
$`=`$ $`{\displaystyle \frac{1}{(2\pi )^{3/2}}}{\displaystyle \left[U_0(𝐤)q_I^0(𝐤)+V_0(𝐤)\stackrel{~}{q}_I^0(𝐤)\right]e^{i𝐤𝐱}𝑑𝐤},`$ (2)
where
$`U_0(𝐤)=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{c}1\\ \sigma \widehat{k}\end{array}\right),`$ (5)
$`V_0(𝐤)=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{c}\sigma \widehat{k}\\ 1\end{array}\right).`$ (8)
The subscript $`I`$ and the superscript $`0`$ indicate that the operators $`q_I^0`$ and $`\stackrel{~}{q}_I^0`$ are two component ones which annihilate or create quanta acting upon the perturbative or the chiral vacuum.
We now consider vacuum destabilisation leading to chiral symmetry breaking described by,
$$|vac=𝒰_Q|0,$$
(9)
where
$$𝒰_Q=\mathrm{exp}(q_I^0(𝐤)^{}(𝝈𝐤)h(𝐤)\stackrel{~}{q}_I^0(𝐤)d𝐤h.c.).$$
(10)
In the above, we have suppresed the flavor and color indices on the two component quark and antiquark operators. Further $`h(𝐤)`$ is a real function of $`|𝐤|`$ which describes vacuum realignment for quarks of a given flavor. We consider here two flavors and take the condensate function $`h(𝐤)`$ to be the same for u and d quarks. Clearly, a nontrivial $`h(𝐤)`$ shall break chiral $`SU(2)_L\times SU(2)_R`$ symmetry to the custodial symmetry $`SU(2)_V`$ for the light quark doublet . In what follows we shall exploit this result.
Having defined the state as in Eq.(9) for chiral symmetry breaking, we shall next define the state involving diquarks. We note that as per BCS result such a state will be dynamically favored if there is an attractive interaction between the quarks . Such an interaction exists in QCD in the qq color antitriplet, Lorentz scalar and isospin singlet channel. In the flavor antisymmetric channel the interaction can be scalar, pseudoscalar or vector whereas in flavor symmetric channel only the axial vector channel is attractive. In the present work, we shall consider the ansatz state involving diquarks as
$$|\mathrm{\Omega }=U_d|vac=\mathrm{exp}(B_d^{}B_d)|vac,$$
(11)
where
$$_d^{}=\frac{1}{2}\left[q_r^{ia}(𝐤)^{}f(𝐤)q_r^{jb}(𝐤)^{}ϵ_{ij}ϵ_{3ab}+\stackrel{~}{q}_r^{ia}(𝐤)f_1(𝐤)\stackrel{~}{q}_r^{jb}(𝐤)ϵ_{ij}ϵ_{3ab}\right]𝑑𝐤.$$
(12)
In the above, $`i,j`$ are flavor indices, $`a,b`$ are the color indices and $`r(=\pm 1/2)`$ is the spin index. As noted earlier we shall consider systems with two flavors and three colors. We have also introduced here two trial functions $`f(𝐤)`$ and $`f_1(𝐤)`$ respectively for the diquark and diantiquark channel. As may be noted the state constructed in eq.(11) is spin singlet and is antisymmetric in color and flavor. The corresponding Bogoliubov transformation for the operators is given by
$`\left(\begin{array}{c}q_{Ir}^{ia}(𝐤)\\ q_{Ir}^{kc}(𝐤)^{}\end{array}\right)`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{cos}f(𝐤)& r\frac{f(𝐤)}{|f(zbfk)|}ϵ_{ik}ϵ_{3ac}\mathrm{sin}f(𝐤)\\ r\frac{f^{}(𝐤)}{|f(𝐤)|}ϵ_{ki}ϵ_{3ca}\mathrm{sin}f(𝐤)& \mathrm{cos}f(𝐤)\end{array}\right)\left(\begin{array}{c}q_{Ir}^{ia}(𝐤)\\ \stackrel{~}{q}_{Ir}^{kc}(𝐤)\end{array}\right).`$ (19)
In a similar manner one can write down the Bogoliubov transformation for the antiquark operators corresponding to $`|\mathrm{\Omega }`$ basis.
Finally, to include the effect of temperature and density we obtain the state at finite temperature and density $`|\mathrm{\Omega }(\beta ,\mu )`$ by a thermal Bogoliubov transformation over the state $`|\mathrm{\Omega }`$ using thermofield dynamics (TFD) as described in ref.s . We have,
$$|\mathrm{\Omega }(\beta ,\mu )=𝒰_{\beta ,\mu }|\mathrm{\Omega }$$
(20)
where $`𝒰_{\beta ,\mu }`$ is
$$𝒰_{\beta ,\mu }=e^{^{}(\beta ,\mu )(\beta ,\mu )},$$
(21)
with,
$$^{}(\beta ,\mu )=q_I^{}(𝐤)^{}\theta _{}(𝐤,\beta ,\mu )\underset{¯}{q}_I^{}(𝐤)^{}+\stackrel{~}{q}_I^{}(𝐤)\theta _+(𝐤,\beta ,\mu )\underset{¯}{\overset{~}{q}}_I^{}(𝐤)d𝐤.$$
(22)
In Eq.(22) the ansatz functions $`\theta _\pm (𝐤,\beta ,\mu )`$ will be related to quark and antiquark distributions and the underlined operators are the operators in the extended Hilbert space associated with thermal doubling in TFD method. Thus the ansatz functions at finite temperature and density given in Eq.(20) invloves five functions - $`h(𝐤)`$, for the quark anti quark condensates, $`f(𝐤)`$ and $`f_1(𝐤)`$ describing respectively the diquark and diantiquark condensates and $`\theta _\pm (𝐤,\beta ,\mu )`$ to include the temperature and density effects. All these functions are to be obtained by minimising the thermodynamic potential. This will involve an assumption about the effective hamiltonian . We shall carry out this minimisation in the next section.
## III estimation of the pressure and the model
Since the low and medium energy behaviour of QCD is not well understood one has to make assumptions about the effective interaction between the quarks. To consider chiral symmetry breaking in the Coulomb gauge pairing model one had the effective action
$$𝒮_{eff}=d^4x\overline{\psi }(x)(i\gamma ^\mu _\mu )\psi (x)+\frac{1}{2}d^4xd^4yg^2(xy)J_\mu ^aD_{ab}^{\mu \nu }(xy)J_\nu ^b(y),$$
(23)
where,
$$J_\mu ^a(x)=\overline{\psi }^k(x)\gamma _\mu (\frac{\lambda ^a}{2})\psi ^k(x).$$
A sum over color indices $`a`$ and $`b`$ is implied, $`g`$ is the strong coupling constant and, $`D_{ab}^{\mu \nu }(xy)`$ is the full gluon propagator. This effective action may be thought of as resulting from integrating out the gluonic degrees of freedom from the full QCD Lagrangian keeping only the bilinear terms of the quark currents. While dealing with the gluon sector, to make things computable, one uses different approximations for the full propagator which is written down in Coulomb gauge . The most obvious one is to neglect transeverse gluons completely maintaining only the instantaneous Coulomb interactions. This was done in Ref. for chiral symmetry breaking. Another possibilty is to neglect retardation effects for the transverse gluons and have e.g. $`D^{ij}(p)=D^{ij}(𝐩,p_0)`$ \- the so called Breit interaction. In such instantaneous approximations one can interpret the product of the coupling constant and the gluon propagator as the Fourier tranform of the effective quark antiquark potential. Different effective potentials such as confining, Richardson, and a screened potential for the transverse part have been considered for chiral symmetry breaking in Coulomb gauge pairing model.
To deal with the case of finite temperature and density one might take one loop resummed perturbative QCD results at finite temperature and densities for the propagator. However, we shall here assume a point like interaction approximation. The reason we shall choose this is its simplicity regarding its solvability. Further we can compare the results with those of NJL model calculations . The delta function interaction produces short distance singularities and so to regulate the integrals we shall restrict the phase space inside the sphere $`|𝐩|<\mathrm{\Lambda }`$. Thus the effective Hamiltonian we shall be considering is given by
$$=\psi ^{}(i𝜶\mathbf{})\psi +\frac{g^2}{2}J_\mu ^aJ^{\mu a}.$$
(24)
We next write down the expectation values of various operators in the thermal vacuum given in Eq.(20). These expressions would be used to calculate thermal expectation value of the Hamiltonian to compute the thermodynamic potential. We define
$$\mathrm{\Omega }(\beta ,\mu )|\stackrel{~}{\psi }_\alpha (𝐤)\stackrel{~}{\psi }_\beta (𝐤^{})^{}|\mathrm{\Omega }(\beta ,\mu )=\mathrm{\Lambda }_{+\alpha \beta }(𝐤,\beta ,\mu )\delta (𝐤𝐤^{}),$$
(25)
and,
$$\mathrm{\Omega }(\beta ,\mu )|\stackrel{~}{\psi }_\beta (𝐤)^{}\stackrel{~}{\psi }_\alpha (𝐤^{})|\mathrm{\Omega }(\beta ,\mu )=\mathrm{\Lambda }_{\alpha \beta }(𝐤,\beta ,\mu )\delta (𝐤𝐤^{}),$$
(26)
where,
$$\mathrm{\Lambda }_+(𝐤,\beta ,\mu )=\frac{1}{2}\left(1+F_1(𝐤)F(𝐤)+\left(\gamma ^0\mathrm{sin}2h(𝐤)+𝜶\widehat{𝐤}\mathrm{cos}2h(𝐤)\right)\left(1F(𝐤)F_1(𝐤)\right)\right),$$
(27)
and,
$$\mathrm{\Lambda }_{}(𝐤,\beta ,\mu )=\frac{1}{2}\left(1+F(𝐤)F_1(𝐤)\left(\gamma ^0\mathrm{sin}2h(𝐤)+𝜶\widehat{𝐤}\mathrm{cos}2h(𝐤)\right)\left(1F(𝐤)F_1(𝐤)\right)\right).$$
(28)
Here, the effect of diquark condensate and the temperature and/or density dependence is encoded in the function $`F(𝐤)`$ and $`F_1(𝐤)`$ given as
$$F(𝐤)=\mathrm{sin}^2f(𝐤)+\mathrm{cos}2f(𝐤)\mathrm{sin}^2\theta _{}(𝐤,\beta ,\mu ),$$
(29)
and,
$$F_1(𝐤)=\mathrm{sin}^2f_1(𝐤)+\mathrm{cos}2f_1(𝐤)\mathrm{sin}^2\theta _+(𝐤,\beta ,\mu ).$$
(30)
Clearly, at zero temperature and zero density the functions $`F`$ and $`F_1`$ vanish and the projection operators reduce to the forms considered earlier,
$$\mathrm{\Lambda }_\pm (\stackrel{}{k},\beta )=\frac{1}{2}\left[1\pm (\gamma ^0\mathrm{sin}2h(k)+\stackrel{}{\alpha }\widehat{k}\mathrm{cos}2h(k))\right].$$
(31)
Further, at zero density but finite temperature $`\theta _{}=\theta _+`$ and projection operators reduce to ,
$$\mathrm{\Lambda }_\pm (\stackrel{}{k},\beta )=\frac{1}{2}\left[1\pm \mathrm{cos}2\theta (\gamma ^0\mathrm{sin}2h(k)+\stackrel{}{\alpha }\widehat{k}\mathrm{cos}2h(k))\right].$$
(32)
We also have
$`\mathrm{\Omega }(\beta ,\mu )|\psi _\alpha ^{ia}(\stackrel{}{x})\psi _\gamma ^{jb}|\mathrm{\Omega }(\beta ,\mu )`$ $`=`$ $`{\displaystyle \frac{ϵ^{ij}ϵ^{3ab}}{(2\pi )^3}}{\displaystyle e^{i𝐤𝐱}𝒫_{+\gamma \alpha }(𝐤,\beta ,\mu )𝑑𝐤},`$ (33)
$`\mathrm{\Omega }(\beta ,\mu )|\psi _\alpha ^{ia}(\stackrel{}{x})\psi _\gamma ^{jb}|\mathrm{\Omega }(\beta ,\mu )`$ $`=`$ $`{\displaystyle \frac{ϵ^{ij}ϵ^{3ab}}{(2\pi )^3}}{\displaystyle e^{i𝐤𝐱}𝒫_{\alpha \gamma }(𝐤,\beta ,\mu )𝑑𝐤},`$ (34)
where,
$$𝒫_+(𝐤,\beta ,\mu )=\frac{1}{4}\left[S(𝐤)+\left(\gamma ^0\mathrm{sin}2h(𝐤)𝜶\widehat{𝐤}\mathrm{cos}2h(𝐤)\right)A(𝐤)\right]\gamma _5C,$$
(35)
and,
$$𝒫_{}(𝐤,\beta ,\mu )=\frac{C\gamma _5}{4}[(S(𝐤)+(\gamma ^0\mathrm{sin}2h(𝐤)𝜶\widehat{𝐤}\mathrm{cos}2h(𝐤))A(𝐤)].$$
(36)
Here, $`C=i\gamma ^2\gamma ^0`$ is the charge conjugation matrix (we use the notation of Bjorken and Drell) and the functions $`S(𝐤)`$ and $`A(𝐤)`$ are given as,
$$S(𝐤)=\mathrm{sin}2f(𝐤)\mathrm{cos}2\theta _{}(𝐤,\beta ,\mu )+\mathrm{sin}2f_1(𝐤)\mathrm{cos}2\theta _+(𝐤,\beta ,\mu ),$$
(38)
and,
$$A(𝐤)=\mathrm{sin}2f(𝐤)\mathrm{cos}2\theta _{}(𝐤,\beta ,\mu )\mathrm{sin}2f_1(𝐤)\mathrm{cos}2\theta _+(𝐤,\beta ,\mu ),$$
(39)
Using Eq. (26) we have for the kinetic energy part
$$T\mathrm{\Omega },\beta ,\mu |\psi ^{}(i𝜶\mathbf{})\psi |\mathrm{\Omega }\beta \mu =\frac{\gamma }{(2\pi )^3}𝑑𝐤\mathrm{cos}2h(𝐤)(1FF_1),$$
(40)
where, the degenracy factor $`\gamma =12`$ corresponding to spin(2), color(3) and flavor(2) degrees of freedom. $`T`$ in Eq.(40) above, however, includes the contribution of the perturbative zero point energy corresponding to $`h(𝐤)=0=f(𝐤)=f_1(𝐤)`$. Subtracting this out we have the contribution from the kinetic energy
$$𝒯=T\widehat{T}|_{h=f=f_1=\theta _\pm =0}=\frac{\gamma }{(2\pi )^3}|𝐤|\left[2\mathrm{sin}^2h(𝐤)+\mathrm{cos}2h(𝐤)(F+F_1)\right].$$
(41)
Similarly the contribution from the interaction term in Eq.(24) after subtracting out the zero point pertuturbative energy turns out to be
$$𝒱\mathrm{\Omega },\beta ,\mu |\frac{g^2}{2}J_\mu ^aJ^{\mu a}|\mathrm{\Omega }\beta \mu =V_1+V_2$$
(42)
where, the contribution $`V_1`$ arises from using Eq.s (25), (26) and is given as
$$V_1=\frac{g^2}{2}N_f\frac{(N_c^21)}{2}2(I_1^22I_2^2)=8g^2(I_1^22I_2^2)$$
(43)
with,
$$I_1=\frac{1}{(2\pi )^3}𝑑𝐤(FF_1)$$
(44)
and
$$I_2=\frac{1}{(2\pi )^3}𝑑𝐤(1FF_1)\mathrm{sin}2h(𝐤).$$
(45)
The term $`V_2`$ arises from using Eq.(35)and (36) and we have
$$V_2=\frac{2g^2}{3}(2I_3^2+I_4^2),$$
(46)
where
$$I_3=\frac{1}{(2\pi )^3}𝑑𝐤S(𝐤)$$
(47)
and
$$I_4=\frac{1}{(2\pi )^3}𝑑𝐤A(𝐤)\mathrm{sin}2h(𝐤).$$
(48)
Then the free energy density is
$$=ϵ\mu N=𝒯+𝒱\mu N$$
(49)
where, $`\mu `$ is the chemical potential and $`N`$ is quark number density. Further,
$$N=\psi ^{}\psi =\frac{\gamma }{(2\pi )^3}𝑑𝐤(FF_1)=\gamma I_1.$$
(50)
Finally, for the entropy density we have
$$s=\frac{\gamma }{(2\pi )^3}𝑑𝐤\left(\mathrm{sin}^2\theta _{}ln\mathrm{sin}^2\theta _{}+\mathrm{cos}^2\theta _{}ln\mathrm{cos}^2\theta _{}+sin^2\theta _+lnsin^2\theta _++\mathrm{cos}^2\theta _+ln\mathrm{cos}^2\theta _+\right)$$
(51)
and the thermodynamic potential which is negative of pressure is given by
$$\mathrm{\Omega }=𝒫=ϵ\mu N\frac{1}{\beta }s$$
(52)
Now if we minimise the thermodynamic potential $`\mathrm{\Omega }`$ with respect to $`h(𝐤)`$,we get
$$tan2h(𝐤)=\frac{8g^2I_2}{3|𝐤|}\frac{M}{|𝐤|}.$$
(53)
where, $`M=8g^2I_2/3`$. Substituting this back in Eq.(45) we have the mass gap equation
$$M=\frac{8g^2I_2}{3}=\frac{8g^2}{3}\frac{1}{(2\pi )^3}\frac{M}{\sqrt{𝐤^2+M^2}}(1FF_1)𝑑𝐤.$$
(54)
Clearly, the above includes the effect of diquark condensates as well as temperature and density through the functions $`F`$ and $`F_1`$ given in Eq.s (29) and (30) respectively. The zero temperature and density limit is given by setting $`F=0=F_1`$ which has the same structure as in NJL model .
Next, minimising the thermodynamic potential $`\mathrm{\Omega }`$ with respect to $`f`$ and $`f_1`$ yields
$$tan2f(𝐤)=\frac{2g^2}{9}\frac{\left(2I_3I_4\mathrm{sin}2h\right)}{\sqrt{𝐤^2+M^2}\left(\mu \frac{4g^2}{3}I_1\right)}$$
(55)
and,
$$tan2f_1(𝐤)=\frac{2g^2}{9}\frac{\left(2I_3+I_4\mathrm{sin}2h\right)}{\sqrt{𝐤^2+M^2}+\left(\mu \frac{4g^2}{3}I_1\right)}.$$
(56)
We might further simplify equations (55) and (56) by noting from Eq.s (47), (48), (39) and (38) that the integral $`I_4`$ is small compared to $`I_3`$, as the integrand in the former is a difference of two postive quantities whereas the latter is a sum of these two quantities. In addition, the integrand in $`I_4`$ is suppresed by the quark antiquark condensate function $`\mathrm{sin}2h(𝐤)`$ which is not the case in $`I_3`$. Finally, this apart, in Eq.s (55) and (56) the numerical coefficient in the second term of the numerator is small compared to that in the first term. Thus, we can approximate
$$tan2f(𝐤)=\frac{\mathrm{\Delta }}{E\nu }$$
(57)
$$tan2f_1(𝐤)=\frac{\mathrm{\Delta }}{E+\nu },$$
(58)
where, the superconducting gap $`\mathrm{\Delta }=(4g^2/9)I_3`$, $`E=\sqrt{𝐤^2+M^2}`$ and $`\nu =\mu 4g^2/3I_1`$ is the chemical potential in presence of interaction . From the definition of superconducting gap and Eq.(47) we have the superconducting gap equation given by
$$\mathrm{\Delta }=\frac{4g^2}{9}I_3=\frac{4g^2}{9(2\pi )^3}𝑑𝐤\left(\frac{\mathrm{\Delta }}{\sqrt{\mathrm{\Delta }^2+(E\nu )^2}}\mathrm{cos}2\theta _{}(𝐤,\beta ,\mu )+\frac{\mathrm{\Delta }}{\sqrt{\mathrm{\Delta }^2+(E+\nu )^2}}\mathrm{cos}2\theta _+(𝐤,\beta ,\mu )\right)$$
(59)
Finally, minimisation of the thermodynamic potential with respect to the thermal functions $`\theta _\pm (𝐤,\beta ,\mu )`$ gives
$$\mathrm{sin}^2\theta _{}=\frac{1}{\mathrm{exp}(\beta \omega _{})+1}$$
(60)
and
$$\mathrm{sin}^2\theta _+=\frac{1}{\mathrm{exp}(\beta \omega _+)+1}$$
(61)
where, $`\omega _\pm =\sqrt{\mathrm{\Delta }^2+\xi _\pm ^2}`$ and $`\xi _\pm =(E\pm \nu )`$. Therefore, with all the ansatz functions determined from the minimisation of the thermodynamic potential $`\mathrm{\Omega }(\beta ,\mu )`$, the mass gap equation (54) and superconducting gap equations are given respectively as
$$\frac{4g^2}{3}\frac{1}{(2\pi )^3}𝑑𝐤\frac{1}{\sqrt{𝐤^2+M^2}}\left(\frac{\xi _{}}{\omega _{}}tanh(\frac{\beta \omega _{}}{2})+\frac{\xi _+}{\omega _+}tanh(\frac{\beta \omega _+}{2})\right)=1$$
(62)
$$\frac{4g^2}{9}\frac{1}{(2\pi )^3}𝑑𝐤\left(\frac{tanh(\frac{\beta \omega _{}}{2})}{\omega _{}}+\frac{tanh(\frac{\beta \omega _+}{2})}{\omega _+}\right)=1$$
(63)
Equation (62) is the generalisation of the mass gap equation considered at zero temperature and density of Ref. to include the effect of finite temperature and density alongwith the effect of diquark condensates. Similarly, Eq.(63) is the relativistic generalisation of the BCS gap equation in the presence of a dynamically generated mass gap through a quark antiquark condensate structure for the vacuum. These are the main new features of the present work.
## IV Solution of the gap equations and results
Before trying to solve the coupled gap equations (62) and (63) we discuss the solutions in the following two different limiting cases, so that we make connection with earlier studies .
* Case–I :Only quark antiquark condensate i.e. $`h(𝐤)0;f(𝐤)=0=f_1(𝐤)`$
In this limit $`\mathrm{\Delta }=0`$ and the gap equation (62) reduces to
$$M=\frac{8g^2}{3}\frac{1}{(2\pi )^3}\frac{M}{\sqrt{𝐤^2+M^2}}(1n_+n_{})𝑑𝐤$$
(64)
where,
$$n_\pm =\frac{1}{\mathrm{exp}(\beta (E(𝐤)\pm \nu ))+1}$$
This is the same equation as obtained in Ref. except for numerical factors before the integrand. This is due to the difference in the Lorentz structure of the four fermion interaction term. In Ref. the interaction term is $`(\overline{\psi }\psi )^2+((\overline{\psi }\gamma _5\tau \psi )^2`$ that can originate from a Fierz transformation of the $`J_\mu ^aJ^{\mu a}`$ interaction as considered here. In Ref. only $`J_0^aJ^{0a}`$ term was included. We might however mention here that keeping $`J_0J_0`$ term only will make the pressure negative at high densities even when the condensate vanishes. At zero temperature the mass gap is given by the solution of the equation
$$\frac{4g^2}{3\pi ^2}_{k_f}^\mathrm{\Lambda }\frac{k^2dk}{\sqrt{k^2+M^2}}=1$$
(65)
where the fermi mometum $`k_f`$is defined through the (interacting) chemical potential $`\nu `$ as $`k_f^2+M^2=\nu ^2`$. We also note that this gap equation is derived here by minimising the thermodynamic potential over a variational ansatz and not through a mean field approximation or a Hartree Fock calculation. The critical fermi momentum is given by the density where $`M=0`$ and is given at zero temperature by,
$$k_f^c=\mathrm{\Lambda }(1\frac{3\pi ^2}{2g^2\mathrm{\Lambda }^2})^{1/2}$$
(66)
and the corresponding quark number density is given by $`n_0=\gamma k_f^{c3}/6\pi ^2`$. The variation of mass gap as a function of fermi momentum is shown in Fig. (1). We have chosen the value of the coupling $`g^2=55.98GeV^2`$ and $`\mathrm{\Lambda }=0.67GeV`$ as typical values. In fact, they have been chosen so that the diquark condensate function vanishes at the same critical temperature as in Ref..
With this choice of parameters, the dynamically generated quark mass at zero temperature and density becomes about 490 GeV which is rather high compared to the standard value of around 300 MeV. With a lower value of $`g^2\mathrm{\Lambda }^2`$ the same could be made smaller. The critical density in this model turns out to be about $`1.7/fm^3`$ corresponding to a value of the fermi momentum $`k_f^c`$ of $`0.4GeV`$. In Fig.2 we have plotted the pressure as a function of fermi momentum. Here we have added the bag constant $`ϵ_0`$ which is the energy density at zero quark number density at zero temperature . This turns out to be about $`(0.2GeV)^4`$. As observed earlier the pressure has a cusp like structure and becomes negative at finite density. The portion of the curve that goes down with fermi momentum corresponds to a nontrivial mass gap solution and the portion that increases with density correponds to zero mass solution of the mass gap equation (62). As in Ref, and unlike Ref. the effects of interactions do not disappear beyond the chiral symmetry restoration phase. This is due to the fact that there is no running coupling involved here and in fact, the pressure grows as $`n^2`$ or $`k_f^6`$ at high densities as may be expected from the interaction term $`V_1`$ involving $`I_1^2=N^2/\gamma ^2`$ given in Eq.(43). However, one has to keep in mind that to be consistent with the philosophy of high momentum cutoff $`\mathrm{\Lambda }`$, the fermi momentum should be less than the cut off momentum. The negative pressure at intermediate densities can be understood in terms of mechanical instability and can have the interpretation that uniform nonzero density quark matter will break up into droplets of finite density in which chiral symmetry is restored surrounded by empty space with zero pressure and density. It is tempting to identify these droplets with nucleons within which the density is nonzero and $`\overline{q}q`$ =0— a fact reminiscent of bag models. Nothing within the model however implies that the droplets have quark number three .
* Case– II:No chiral condensate but only diquark condensate i.e. $`h(𝐤)=0`$ $`f(𝐤)0`$, $`f_1(𝐤)0`$
In this case, the quark mass arising from chiral condensate is zero and the superconducting gap equation (63) reduces to
$$1=\frac{2g^2}{9\pi ^2}_0^\mathrm{\Lambda }k^2𝑑k\left[\frac{1}{\omega _{}}tanh(\frac{\beta \omega _{}}{2})+\frac{1}{\omega _+}tanh(\frac{\beta \omega _+}{2})\right]$$
(67)
where, as before $`\omega _\pm =\sqrt{\mathrm{\Delta }^2+\xi _\pm ^2}`$ but, with $`\xi _\pm =|𝐤|\pm \nu `$. This is the relativistic generalisation of BCS gap equation in superconductivity . It has the same structure as in Ref in the limit the form factor is replaced by a sharp cutoff as in Ref. and apart from the numerical factor preceeding the integrand. Further, this can be rearranged to interpret separately contributions arising from particles, antiparticles and holes . We have chosen the values of the coupling and the cutoff as in Ref. so as to get the same critical temperature as in Ref.. The gap equation at zero temperature is given as
$$1=\frac{2g^2}{9\pi ^2}_0^\mathrm{\Lambda }k^2𝑑k\left[\frac{1}{\omega _{}}+\frac{1}{\omega _+}\right]$$
(68)
The superconducting gap $`\mathrm{\Delta }`$ is plotted in Fig 3. While obtaining the solution of the gap equation (68) we also insist that a nontrivial solution of the gap equation is acceptable only if the corresponding free energy is smaller than that with no gap. The gap increases with fermi momentum, has a maximum of about 90 MeV around fermi momentum 550 MeV, beyond which the effect of cutoff is felt and it vanishes at around $`k_f=600`$ MeV.
The resulting equation of state i.e. pressure as a function of fermi momentum is plotted in Fig.4. There is no unstable phase as discussed for chiral condensate case with negative pressure. Because of the interaction term, the equation of state does not go over to the free massless equation of state beyond the critical density. We have plotted the pressure in Fig.4 for fermi momentum upto 600 MeV noting that the cutoff here is about 650 MeV.
* Case– III: $`h(𝐤)0;`$$`f(𝐤)0`$and $`f_1(𝐤)0`$.
Here we solve the coupled gap equations (62) and (63) which at zero temperature but finite density reduce to
$$\frac{2g^2}{3\pi ^2}𝑑k\frac{k^2}{E}\left(\frac{\xi _{}}{\omega _{}}+\frac{\xi _+}{\omega _+}\right)=1$$
(69)
$$\frac{2g^2}{9\pi ^2}𝑑kk^2\left(\frac{1}{\omega _{}}+\frac{1}{\omega _+}\right)=1$$
(70)
where, $`\omega _\pm =\sqrt{\mathrm{\Delta }^2+\xi _\pm ^2}`$; $`\xi _\pm =E\pm \nu `$ and $`E=\sqrt{𝐤^2+M^2}`$. To solve the above equations numerically, we take $`\nu `$ as an input and obtain the values of mass gap $`M`$ and the superconducting gap $`\mathrm{\Delta }`$ so that both the equations are satisfied simultaneously. The chemical potential $`\nu `$ is related to the fermi momentum as $`k_f=\sqrt{\nu ^2M^2}`$. The resulting mass gap and the superconducting gap are plotted in Fig. 5 and Fig 6 respectively. For the sake of comparison we have also plotted the mass gap without diquark condensate (case–I) in Fig 5,and, superconducting gap without chiral gap (case–II) in Fig. 6. As may be noted, the mass gap does not change very much through the inclusion of diquark condensates. The superconducting gap however starts at a lower threshold and varies slowly compared to the case of no mass gap. Both the curves however merge at the chiral restoration point.
We have also plotted the equation of state , pressure as a function of fermi momentum in Fig.7. The equation of state does not change very much compared to the case with only quark antiquark condensates.
This should be expected as the effect of diquark condensate is small in the region where chiral condensate is nonvanishing. As may be evident from Fig. there is a region from $`k_f0.1`$ fm till $`k_f0.4`$ fm where simultaneous existence of both types of condensates is possible. However, in this region the pressure becomes negative. As in case I, this will correspond to mechanical instability. The finite density matter will break into droplets of finite density in which chiral symmetry is restored and the matter is in a superconducting phase surrounded by empty space with zero pressure and density. Below $`k_f=0.1`$ fm the phase with only quark antiquark condensates and for $`k_f>0.4`$ fm the phase with diquark condensates is thermodynamically feasible. However, one should not extrapolate to fermi momenta higher than the cutoff momentum of the model to be consistent with the philosophy of high momentum cutoff.
### A The phase diagram
To discuss the phase diagram we have solved the two gap equations Eq.(62) and Eq. (63) at finite temperature and density to calculate different thermodynamic quantities. The result of such a calculation for the mass gap is shown in Fig.(8). Similarly the superconducting gap at different temperatures is shown in Fig.(9).
The superconducting gap decreases with temperature and vanishes at around 55 MeV. Similarly the mass gap also decreases with temperature and vanishes at about 270 MeV at zero density. The superconducting transition is a second order phase transition. The chiral phase transition is first order at zero density but is second order at high temperature. The pressure at different temperatures is shown in Fig.10 as a function of $`(n/n_0)^{1/3}`$, where n is number density and $`n_0=(2/\pi ^2)\nu _c^3`$ is the critical number density at zero temperature. The cusp in the pressure density curve, indicating a first order phase transition, vanishes at about 84 MeV. At temperaures beyond that the pressure increases monotonically with density.
To discuss the critical line for the chiral transition in the number density and temperature plane we need to solve the gap equation (62) with mass gap equated to zero i.e. the critical line satisfies the equation
$$\frac{4g^2}{3\pi ^2}_0^\mathrm{\Lambda }𝑑kk\left(1\frac{1}{\mathrm{exp}(\beta (k\nu ))+1}\frac{1}{\mathrm{exp}(\beta (k+\nu ))+1}\right)=1$$
(71)
As the distribution functions are rapidly decreasing functions of momentum we may approximate the upper limit $`\mathrm{\Lambda }\mathrm{}`$. In that case Eq.(71) reduces to
$$\nu ^2+\tau ^2=\nu _c^2$$
(72)
where,$`\tau =\pi T/\sqrt{3}`$ and $`\nu _ck_f^c=\mathrm{\Lambda }(13\pi ^2/(2g^2\mathrm{\Lambda }^2)`$, the critical chemical potential or fermi momentum at zero temperature calculated earlierin Eq.(66). It may be useful to write down this relation in terms of number densities rather than the chemical potential using the relation
$$n=\frac{6}{\pi ^2}_0^\mathrm{\Lambda }𝑑kk^2\left(\frac{1}{\mathrm{exp}(\beta (k\nu ))+1}\frac{1}{\mathrm{exp}(\beta (k+\nu ))+1}\right)\frac{2}{\pi ^2}\nu (\nu ^2+\pi ^2T^2)$$
(73)
where, we again take infinity as the upper limit in the momentum integration as was done in Eq.(71). This leads to the equation for the critical line in the number density plane as
$$n/n_0=\left(1(\frac{\tau }{\nu _c})^2\right)^{1/2}\left(1+2(\frac{\tau }{\nu _c})^2\right)$$
(74)
where, $`n_0=(2/\pi ^2)\nu _c^3`$ is the critical density at zero temperature. This approximate solution is shown by the dotted line in Fig11. The numerical solution for the critical line with the finite cutoff for the momentum is shown by the solid curve in the same figure. As the chiral phase transition is a first order phase transition at zero temperature and is second order at high temperature there will be a tricritical point for the chiral phase transition below which there will be a mixed phase. This mixed phase here will correspond to the phase with droplets of finite density quark matter in the superconducting but chiral symmetric phase surrounded by empty space with $`n=0,P=0,\mathrm{\Delta }=0`$ and chiral symmetry broken phase. This can be obtained by a Maxwell construction to exclude the part of the phase diagram where pressure decreases with density. This is shown as the dashed curve in Fig.11. The tricritical point in this model turns out to be 84 MeV.
We have also shown in Fig.11 the critical line for the color superconducting phase by the dot-dashed line. The decreasing part of this line is due to the fact that we have taken a cut off in the momentum and at high enough density the superconducting gap vanishes as shown in Fig. 9.
A detailed quantitative comparison of our results with those in Ref. is not very meaningful since our model is a simple, schematic one. However, it is very reassuring that the structure of the phase diagram is very similar to the one obtained by Berges and Rajagopal .
## V summary
We have analysed in a current current point interaction model the structure of vacuum with quark antiquark as well as diquark pairs. The methodology used is a variational one with an explicit construct of the trial state. The present work is not based on a mean field approximation . Because of the point interaction structure we could solve for the gap functions explicitly. If we had taken the interaction term with a potential the gap equation would become an integral equation. It will be interesting to include a realistic effective potential and solve for the gap equation. There appears to be a region depending upon the coupling where both chiral condensates as well as diquark condensates are thermodynamically feasible. Presence of diquark condensate does not modify the dynamical mass of the quarks. The dynamical mass however affects the threshold for superconducting gap. The equation of state does not differ very much with inclusion of diquark condensates. We also obtain the complete phase diagram of the system in overall agreement with that of Ref.. Finally, our results suggests that our simple model contains the essential physics of the system.
## ACKNOWLEDGMENTS
HM would like to thank Amruta Mishra for numerous discussions. |
warning/0003/cond-mat0003503.html | ar5iv | text | # Correlations due to localization in quantum eigenfunctions of disordered microwave cavities
\[
## Abstract
Non-universal correlations due to localization are observed in statistical properties of experimental eigenfunctions of quantum chaotic and disordered microwave cavities. Varying energy$`E`$ and mean free path $`l`$ enable us to experimentally tune from localized to delocalized states. Large level-to-level Inverse Participation Ratio (IPR $`I_2`$) fluctuations are observed for the disordered billiards, whose distribution is strongly asymmetric about $`I_2`$. The density auto-correlations of eigenfunctions are shown to decay exponentially and the decay lengths are experimentally determined. All the results are quantitatively consistent with calculations based upon nonlinear $`\sigma `$models.
\]
The universal properties of quantum chaotic systems have been extensively studied in terms of their eigenvalue and eigenfunction statistics . In Random Matrix Theory (RMT), the Gaussian distribution of eigenfunction amplitudes $`\psi (\stackrel{}{r})`$ leads to the universal Porter-Thomas (PT) distribution for the densities $`|\psi |^2`$, which has been observed in microwave cavities as well as other systems. However, non-universality has important manifestations, for instance due to periodic orbits which lead to scars in eigenfunctions, and localization arising from quantum interference in diffusion. While there have been many theoretical treatments, from semiclassical periodic orbit theories to nonlinear sigma models , there have been few experimental studies of eigenfunctions because of their lack of accessibility.
In this paper, we present several striking manifestations of localization in experimental eigenfunctions of disordered microwave billiard cavities. Localization due to boundary or impurity scattering results in correlations that affect statistics and spatial correlations of the eigenfunctions in several measures, leading to deviations from their universal values for quantum chaotic systems. The moments of the density distribution, $`I_n=|\mathrm{\Psi }(\overline{r})|^{2n}d^3r`$, in particular the Inverse Participation Ratio $`I_2`$ $`(IPR)`$, and their distributions $`P_{I_n}(I_n)`$, are important measures of the properties of the chaotic and disordered eigenfunctions. In chaotic billiards, $`I_2`$ has a mean value $`I_2`$ close to that of the universal $`2`$-dimensional $`(2D)`$ limiting value of $`3.0`$, with small level-to-level fluctuations $`\delta I_2I_2`$, resulting in a nearly symmetric distribution about $`I_2`$. In disordered billiards not only is the mean value $`I_23.0`$, but the fluctuations are also much greater $`\delta I_2I_2`$. The IPR distribution for the disordered billiards is strongly asymmetric about $`I_2`$, and is quantitatively consistent with the calculations based upon the nonlinear $`\sigma `$models of supersymmetry, parametrized by a conductivity $`g`$ . Spatial correlations are studied in terms of the density auto-correlation $`|\mathrm{\Psi }(r)|^2|\mathrm{\Psi }(r^{^{}})|^2`$ and are shown to die out more rapidly in the disordered billiards compared with the chaotic geometries with a characteristic decay length given by the mean free path $`l`$. Here again the data are in quantitative agreement with the $`\sigma `$model calculations. Our results represent the first quantitative comparison of experiments and theory.
The experiments were carried out using thin $`(`$ height $`d`$ $`<6`$ mm$`)`$ cavities, whose cross-sections can be shaped in essentially arbitrary geometries. For these two dimensional cavities, the operational wave equation is $`(^2+k^2)\mathrm{\Psi }=0`$, where $`\psi =E_z`$ is the microwave electric field. Similar microwave experiments, which exploit this QM-EM mapping, have been used to study quantum chaos in closed and open systems . Eigenfunctions $`|\psi (r)|^2`$ were directly measured using a cavity perturbation technique . Localization effects were observed by fabricating billiards in which $`1cm`$ circular tiles were placed in a $`44cm`$ x $`21.8cm`$ rectangle at random locations (Fig.1) to act as hard scatterers. (The locations were generated using a random number generator and the tiles placed manually). Several realizations of the disordered geometry were experimentally studied, by varying the density of scatterers from $`12`$ to $`71`$, thus varying mean free path $`l8.9cm`$ to $`3.6cm`$. Earlier experiments had shown that these disordered geometries are an excellent experimental realization of the classic problem of an electron in a 2-D disordered potential. Subsequent to this work, there have been important theoretical developments , some motivated by our earlier experiments.
The experimental eigenfunctions directly demonstrate the trend toward decreasing localization from disordered-chaotic-integrable as seen in Fig.1, which shows representative experimental eigenfunctions, along with their corresponding IPR $`I_2`$, for several billiards. The strongly localized state Fig.1(a) at $`f=3.04GHz`$ of the disordered billiard with $`N=36`$ scatterers has very large $`I_2=13.42`$. In contrast the more delocalized state at higher frequency $`f=7.33GHz`$ (Fig.1(b)) explores almost all the available coordinate space similar to chaotic cavities and has a smaller $`I_2=4.06`$. For the chaotic Sinai stadium, typical values of the IPR are around $`3.0`$ ($`I_2=3.01`$ for this eigenfunction Fig.1(c)) while for the integrable rectangle billiard Fig.1(d) $`IPR`$ $`I_2=2.25`$ for all eigenfunctions. Fig.1 demonstrates a key advantage of our experiments, which is that by varying $`l`$ and wavevector $`k`$ we are able to access a wide range of the disorder strength $`kl`$, from strongly localized states for $`kl<1`$ to delocalized states with $`kl1`$.
Eigenfunctions such as in Fig.1 were then analyzed in terms of $`I_2`$ and $`P_{I_2}(I_2)`$. In the following, for convenience, we use the notation $`I_2=|\mathrm{\Psi }(r)|^4𝑑v=u^2𝑑v,w=(I_23)/6`$ , $`u=|\mathrm{\Psi }(r)|^2`$, and $`I_1=|\mathrm{\Psi }(r)|^2𝑑v=1`$, and the integral is over the volume $`v`$ (area in $`2D`$ ). Nearly 250 eigenfunctions were analyzed each containing more than 3200 eigenfunction values.
Fig.2 shows level-to-level variations $`I_2(E)`$ for the Sinai stadium. Here $`I_2`$ is mainly clustered around a mean value of $`I_2=3.0`$, with relatively small level-to-level fluctuations. The IPR distribution $`P_{I_2}(I_2)`$ of the chaotic Sinai-stadium billiard is shown in Fig.3(top). $`P_{I_2}(I_2)`$ is seen to be a nearly symmetric distribution with a small width $`I_2I_2I_2`$.
Ref. demonstrated that the Sinai-stadium billiard data obey the universal Porter-Thomas (P-T) distribution $`P_u(u)=(2\pi u)^{1/2}\mathrm{exp}(u/2)`$, with $`u=|\mathrm{\Psi }|^2`$ to remarkable degree, while deviations from P-T were demonstrated due to localization in the disordered billiards. The IPR for P-T distribution can be immediately obtained $`I_2=_0^{\mathrm{}}u^2P_u(u)𝑑u=3.0`$, which is a universal value. Note that there are no fluctuations expected about this universal value in RMT, i.e., $`P_{I_2}(I_2)`$ is a $`\delta `$-function at $`I_2=3`$ . Clearly boundary scattering on the system length scale $`R`$ leads to non-universal correlations (e.g. from periodic orbits leading to scars in the wavefunctions). This breaks the assumption in RMT of Gaussian fluctuations of the eigenfunction amplitude, and in turn leads to fluctuations in the distribution $`P_{I_2}(I_2)`$ (although of narrow width) observed in Fig.3(top).
Even more strikingly, the IPRs of the disordered billiards shown in Fig.2 display a remarkably large spectrum of level-to-level fluctuations (Fig.2). For small $`f`$, when $`\lambda >l`$, strong localization leads to large values of $`I_2`$ in the disordered cavity, which can be as high as $`20`$. It is worth noting that the density distributions $`P_u(u)`$of the eigenfunctions deviate strongly from the P-T distribution and are consistent with the large IPR values.(In this paper we have focussed on the billiard with $`l=5.1cm.`$, but similar results were obtained for other billiards and will be presented in a larger publication.)
As $`f`$ is increased (or $`\lambda `$ is decreased), the IPR progressively decreases, approaching the universal limiting value of $`3`$, as shown in Fig.2. The corresponding distribution $`P_{I_2}(I_2)`$ shown in Fig.3(bottom) is strikingly different from the Sinai-stadium case, in that, it is strongly asymmetric, reflecting the very large $`I_2I_2`$ values that are present, and is strongly non-Gaussian. The IPR fluctuations in Fig.2 are closely similar to the famous universal conductance fluctuations of a mesoscopic system.
Recently several theoretical studies have calculated the IPR distribution based on the supersymmetry method . For a mesoscopic system the IPR distribution depends on the conductivity $`g`$ of the system, defined as $`g=\mathrm{ln}(R/l)/w`$ , where $`R`$ is the system size and $`l`$ is the mean free path and $`..`$ is the realization average for a fixed “disorder strength” $`2kl`$. The resulting probability distribution $`P(I_2)`$ for $`I_2<I_2=3`$ is :
$$P_{I_2}(I_2)=C_1\frac{g}{2}\mathrm{exp}[\frac{g}{6}(I_2I_2)\frac{\pi }{2}e^{\frac{g}{3}(I_2I_2)}],$$
(1)
and the corresponding distribution for $`I_2I_2`$ :
$$P_{I_2}(I_2)=C_2\sqrt{\frac{g}{I_2}}\mathrm{exp}(\frac{\pi }{6}gI_2)$$
(2)
where $`C_{1\text{ }}`$ and $`C_2`$ are normalization constants.
The solid line in Fig.3(bottom) represents Eq.(2) for $`I_2I_2`$ using a small conductivity $`g=1.0`$ and is seen to describe the data very well. Another way of presenting the comparison with Eq.(2) is by using a scaled variable $`q=gI_2`$ whence the distribution changes to a Porter-Thomas distribution in $`q`$. The system is within the weak disorder limit, i.e. the random phase approximation still valid, and $`g`$ depends weakly on the disorder factor $`2kl`$ . Now rescaling the $`I_2`$ with $`g`$ we obtain the distribution $`P_q(q)`$ (unnormalized) and plot $`\mathrm{ln}[P_q(q)\sqrt{q}]`$ vs $`q`$ in Fig.4. This shows a straight line which implies a good agreement of theory and experiment for IPR distribution in a moderate disordered region where $`I_2I_2`$. This is the first experimental observation of the asymmetric distribution predicted by Prigodin and Altshuler .
We now return to the case of the Sinai-stadium. Although a formulation in terms of a ballistic sigma model has been presented for chaotic cavities, it is not simply amenable to experimental comparison. Instead we use Eq.1 with the assumption that since $`lR,`$ a suitably large conductivity $`(g>>1)`$ can be used. In our experimental case, $`g7.8`$ matches Eq.1 not only for $`I_2<I_2`$, but also for $`I_2>I_2,`$ as shown in Fig.3(top). The nearly symmetric distribution can be understood since the fluctuations arise from correlations at the scale of the system length $`R`$, which is the only length scale in the problem. The values of $`g`$ corresponding to the disordered billiards and the Sinai-stadium billiard are consistent with the $`w`$ values and the mean free paths.
We now show that the experimental data also directly demonstrate the decay of spatial correlations $`\mathrm{\Psi }^2(r)\mathrm{\Psi }^2(r^{})`$ of the eigenfunctions and obtain the decay length, which corresponds to the classical mean free path. To calculate the correlation with arbitrary disorder strength $`2kl`$, defining $`K(r)=|ImG(r^{})|^2/(\pi \nu )^2`$ , where $`G(|rr^{^{}}|)=<r|(EH)^1|r^{}>`$ is the Green function of the disordered system Hamiltonian, then it can be shown that $`K(r)=|\frac{1}{\pi }_{\mathrm{}}^{\mathrm{}}\frac{1}{1+y^2}J_0[kr(1+\frac{1}{2kl}y)]𝑑y|^2`$.
For the chaotic or ballistic limit $`(2kl>>1),`$ the result for Gaussian fluctuations is $`\mathrm{\Psi }^2(r)\mathrm{\Psi }^2(r^{})1+(I_21)J_0^2(rr^{}),`$ where $`J_0`$ is the first order Bessel function . The correlation for a moderate disordered case with correct limits can be derived repeating the calculations of Ref:
$`\mathrm{\Psi }^2(r)\mathrm{\Psi }^2(r^{})`$ $``$ $`1+(I_21)K(k|rr^{}|)`$ (3)
$``$ $`1+(I_21)J_0^2(k|rr^{}|)e^{\frac{k|rr^{^{}}|}{kl}}`$ (4)
These are valid when $`rr^{}l`$. We have solved $`K(k|rr^{^{}}|)`$ numerically. In Fig.5, we plot the average correlation derived from experimental data, numerical calculations of Eq.(3), and analytical expression of Eq.(4) for different disorder strengths $`2kl.`$ For the Sinai billiards, the average correlation starts at $`3`$ and tend to $`1`$ via Friedel oscillations, consistent with Eq.(4) with $`2kl=37`$, which is very large and hence the result is close to that of the universal dependence given by $`kl=\mathrm{}`$. For disordered billiards, the auto-correlation is very large ($`I_2`$) for short lengths, i.e. around $`|rr^{}|0`$ due to localization, but the auto-correlation decays with a decay length scale, responsible for fast fall, and it should go to zero at large $`k|rr^{}|`$. The experimental data shows good agreement with the numerical and analytical calculations for moderate disorder, as shown in Fig.5 for values $`2kl=12`$ and $`7`$. These values are in excellent agreement with the corresponding mean free paths $`(l=3.6`$, $`5.1`$, $`5.9`$ and $`8.9)`$ obtained by directly considering the number of scatterers $`N`$, so that $`l\sqrt{ab/N}`$ and $`a`$ and $`b`$ are the sides of the enclosing rectangle. Thus our analysis directly demonstrates localization and yields a quantitative measure in terms of the correlation decay length.
Returning to the level-to-level $`I_2`$ data in Fig.2 we note that they can be viewed as a form of localization-delocalization transition, tuned by the frequency $`f`$ ! Representing the IPR as $`I_2(E)=I_{2,sm}(E)+\delta I_2(E)`$, i.e. a smooth part $`I_{2,sm}(E)`$ plus fluctuations $`\delta I_2(E)`$, we discuss the trends in $`I_2(E)`$. Calculations in Ref. based on infinite dimensional tight binding model shows that $`I_2(E)`$ diverges exponentially near the critical point $`E_c`$: $`I_2(E)=I_2(E=\mathrm{})\mathrm{exp}(A/|EE_c|^\mu )`$, with $`\mu =\frac{1}{2}`$ , due to very strong correlations of the wave function near $`E_c`$. $`I_2(E=\mathrm{})`$ will be obviously the asymptotic universal value $`3,`$ as predicted by RMT in $`2D`$. Our experimental $`IPR`$ data is as large as $`IPR`$ $``$ $`20`$ (strongly localized), decaying to $`4`$ (weakly localized) at high frequencies upto $`10GHz.`$ In the present case, the smallest scale of the system $`5cm`$ sets a natural lower cut-off frequency $`f_c=c/2l=3GHz`$, so that there are no eigenstates for $`E<`$ $`E_c=f_c^2.`$ For $`EE_c`$ expanding the above equation in first order, we obtain $`I_{2,sm}(E)3+B/|EE_c|^\mu `$. In Fig.2 we have plotted this expression with $`\mu =0.5`$ and $`B=9.0`$ (obtained by a best fit) and compared with the experimental data. The above expression captures the trend of the data. While an exact comparison with any expression is difficult since the fluctuations $`\delta I_2(E)>I_{2,sm}(E),`$ the analysis illustrates that we are observing a frequency driven localized to de-localized transition in a disordered medium in terms of the $`IPR.`$
We have shown that the IPR is an extremely valuable parameter to study real-space localization in quantum eigenfunctions. We have demonstrated for the first time a quantitative analysis of features well beyond universality due to localization in experimental eigenfunctions. The observed IPR distribution is strongly non-Gaussian due to the correlations induced by scattering. The nonlinear sigma-model is in quantitative agreement with the experimental data for moderate localization. Our work thus provide experimental support for the approach based upon quantum diffusion in the localization regime.
This work was supported by NSF-PHY-9722681. S.S. thanks the hospitality of the Quantum Chaos Workshop at the Australian National University, Canberra, where part of this work was carried out. We thank A. Kudrolli for useful discussions. |
warning/0003/hep-th0003075.html | ar5iv | text | # 1 Introduction
## 1 Introduction
In conventional (20th century) physics, high energy or high momentum came to be associated with small distances. The physics of the 21st century is likely to be dominated by a very different perspective. According to the Infrared/Ultraviolet connection which underlies much of our new understanding of string theory and its connection to gravity, physics of increasing energy or momentum is governed by increasingly large distances. Examples include the growth of particle size with momentum , the origin of the long-range gravitational force in Matrix Theory from high energy quantum corrections and the IR/UV connection in AdS spaces. Another important manifestation is the spacetime uncertainty principle of string theory
$$\mathrm{\Delta }x\mathrm{\Delta }t\alpha ^{}.$$
(1.1)
Similar uncertainty principles occur in non-commutative geometry where the coordinates of space do not commute. An important consequence of the non-commutativity is the fact that the particles described by non-commutative field theories have a spatial extension which is proportional to their momentum . This in turn leads to unfamiliar violations of the conventional decoupling of IR and UV degrees of freedom in these theories . In this paper we will describe another example of IR/UV non-decoupling that occurs in AdS/CFT theories. The relevant space-times have the form $`AdS_n\times S^m`$. We are interested in the motion of the graviton and other massless bulk particles on the $`S^m`$. The motion is characterized by an angular momentum $`L`$ or more exactly a representation of the rotation group $`O(m+1)`$. In 20th century physics such particles are regarded as point or almost point particles regardless of $`L`$. In fact we will see that as $`L`$ increases the particles blow up in size very much like the quanta of non–commutative field theories. When the size reaches the radius of the $`S^m`$, the growth can no longer continue and the tower of Kaluza–Klein states terminates. This is the origin of the stringy exclusion principle .
In section 2 we will review the theory of electric dipoles moving in a magnetic field. This system is the basic object of non-commutative field theory. When the theory is defined on a 2-sphere there is a bound on the angular momentum when the ends of the dipole separate to the antipodes of the sphere. In sections 3, 4 and 5 we consider the cases of $`AdS_7\times S^4`$, $`AdS_4\times S^7`$ and $`AdS_5\times S^5`$. In each case we find that the spectrum of angular momentum is bounded and that the bound agrees with expectations from the stringy exclusion principle.
## 2 Dipoles in Magnetic Fields
In this section we briefly review the dipole analogy for non-commutative field theory . We begin with a pair of unit charges of opposite sign moving on a plane with a constant magnetic field $`B`$. The Lagrangian is
$$=\frac{m}{2}\left(\dot{x}_1^2+\dot{x}_2^2\right)+\frac{B}{2}ϵ_{ij}\left(\dot{x}_1^ix_1^j\dot{x}_2^ix_2^j\right)\frac{K}{2}(x_1x_2)^2.$$
(2.1)
Let us suppose that the mass is so small so that the first term in Eq.(2.1) can be ignored. Let us also introduce center of mass and relative coordinates
$`X`$ $`=`$ $`(x_1+x_2)/2`$ (2.2)
$`\mathrm{\Delta }`$ $`=`$ $`(x_1x_2)/2.`$ (2.3)
The Lagrangian becomes
$$=Bϵ_{ij}\dot{X}^i\mathrm{\Delta }^j2K\mathrm{\Delta }^2.$$
(2.4)
From Eq.(2.3) we first of all see that $`X`$ and $`\mathrm{\Delta }`$ are non-commuting variables satisfying
$$[X^i,\mathrm{\Delta }^j]=i\frac{ϵ^{ij}}{B}.$$
(2.5)
Furthermore the center of mass momentum conjugate to $`X`$ is
$$P_i=Bϵ_{ij}\mathrm{\Delta }^j.$$
(2.6)
Thus when the dipole is moving with momentum $`P`$ in some direction it is stretched to a size
$$|\mathrm{\Delta }|=|P|/B.$$
(2.7)
in the perpendicular direction. This is the basis for the peculiar non-local effects in non–commutative field theory.
Now suppose the dipole is moving on the surface of a sphere of radius $`R`$. Assume also that the sphere has a magnetic flux $`N`$. In other words there is a magnetic monopole of strength
$$2\pi N=\mathrm{\Omega }_2BR^2.$$
(2.8)
at the center of the sphere. We can get a rough idea of what happens by just saying that when the momentum of the dipole is about $`2BR`$ the dipole will be as big as the sphere. At this point the angular momentum is the maximum value
$$L=PRBR^2.$$
(2.9)
This is of order the total magnetic flux $`N`$.
We will now do a more precise analysis and see that the maximum angular momentum is exactly $`N`$. Parameterize the sphere by two angles $`\theta ,\varphi `$. The angle $`\varphi `$ measures angular distance from the equator. It is $`\pm \pi /2`$ at the poles. The azimuthal angle $`\theta `$ goes from $`0`$ to $`2\pi `$. We work in a gauge in which the $`\theta `$ component of the vector potential is non-zero. It is given by
$$A_\theta =N\frac{1\mathrm{sin}\varphi }{2R\mathrm{cos}\varphi }$$
(2.10)
For a unit charged point particle moving on the sphere the term coupling the velocity to the vector potential is
$$_A=A_\theta R\mathrm{cos}\varphi \dot{\theta }=NR\frac{1\mathrm{sin}\varphi }{2R}\dot{\theta }.$$
(2.11)
Now consider a dipole with its center of mass moving on the equator. The positive charge is at position $`(\theta ,\varphi )`$ and the negative charge is at $`(\theta ,\varphi )`$. For the motion we consider $`\varphi `$ is time independent. Eq.(2.10) becomes
$$_A=N(\frac{1\mathrm{sin}\varphi }{2})\dot{\theta }N(\frac{1+\mathrm{sin}\varphi }{2})\dot{\theta }$$
(2.12)
or
$$_A=N\mathrm{sin}\varphi \dot{\theta }.$$
(2.13)
Again we want to consider a slow-moving dipole whose mass is so small that its kinetic term may be ignored compared to the coupling to the magnetic field, i.e. $`mR<<N`$. Let us also add a spring coupling
$$_S=\frac{k}{2}R^2\mathrm{sin}^2\varphi ;$$
(2.14)
for simplicity, we have used the chordal distance in this potential. The total Lagrangian is
$$=\frac{k}{2}R^2\mathrm{sin}^2\varphi N\mathrm{sin}\varphi \dot{\theta }$$
(2.15)
and the angular momentum is
$$L=N\mathrm{sin}\varphi .$$
(2.16)
The angular momentum will reach its maximum when $`\varphi =\pi /2`$ at which point
$$|L_{max}|=N.$$
(2.17)
The fact that the angular momentum of a single field quantum in non-commutative field theory is bounded by $`N`$ is well known in the context of non–commutative field theory on a sphere . Here we see that it is a large distance effect.
## 3 $`AdS\times S`$
We now study BPS particles moving on the sphere of maximally supersymmetric $`AdS`$ vacua of string and M theory. For simplicity, we present all details of the argument in the case of the M5-brane geometry. Results for the other cases are given in the latter subsections. Note that in all of these cases, the energy of our objects is well below the energy of a stable AdS black hole.
### 3.1 $`AdS_7\times S^4`$
We are interested in the motion of a BPS particle on the 4-sphere of $`AdS_7\times S^4`$. We will assume that the radius of curvature $`R`$ is much larger than the 11 dimensional Planck length $`l_p`$. The analogy with the previous example is very close. The role of the magnetic field is played by the 4-form field strength on the sphere. We call the flux density $`B`$. Quantization of flux requires
$$\mathrm{\Omega }_4BR^4=2\pi N.$$
(3.1)
From the supergravity equations of motion it can be seen that $`R`$ is given by
$$R=l_p(\pi N)^{\frac{1}{3}}.$$
(3.2)
The assumption of large $`R`$ means $`N>>1`$.
We want to know what happens to a graviton or any other massless 11 dimensional particle when it moves on the 4-sphere in the presence of the 4-form field strength. As long as the angular momentum is small ($`L<<N`$) the graviton is expected to be much smaller than $`R`$, and we can make the approximation that space is locally flat. Furthermore from Eq.(3.1) we see that the field strength $`B`$ is small
$$BN^{\frac{1}{3}}l_p^4.$$
(3.3)
Since the space is almost flat we can locally introduce flat space coordinates $`x^0,\mathrm{}.,x^{10}`$. Let us take the $`B`$ field to lie in the $`(7,8,9,10)`$ directions. The graviton is moving along the $`x^{10}`$ direction. Its momentum is $`P_{10}=L/R`$. This setup is very close to one that was studied by Myers using Matrix Theory .
In matrix theory the 11 dimensional graviton is viewed as a threshold bound state of $`n=P_{10}R_{10}`$ D0-branes. Myers shows that in a background 4-form fieldstrength the D0-brane configuration is described as a spherical membrane with a radius $`r`$ that grows with $`P_{10}`$ according to
$$rBP_{10}l_p^6.$$
(3.4)
Let us assume that this formula is approximately valid until $`r`$ becomes of order $`R`$. In that case when the graviton size becomes $`R`$ it will have momentum
$$P_{10}(max)R/Bl_p^6$$
(3.5)
and angular momentum
$$L_{max}RP_{10}(max)R^2/Bl_p^6N.$$
(3.6)
Thus as in the previous section we find that the maximum single particle angular momentum is $`N`$.
We will now give a more precise calculation which parallels that of section 1. We are interested in the dynamics of a relativistic spherical membrane moving in $`S^4`$. The membrane has zero net charge but it couples to the background field strength. It behaves like the dipole of section 1.
Let us parametrize $`S^4`$ using cartesian coordinates $`X_1,\mathrm{}..,X_5`$ so that
$`X_1`$ $`=`$ $`R\mathrm{cos}\theta _1`$ (3.7)
$`X_2`$ $`=`$ $`R\mathrm{sin}\theta _1\mathrm{cos}\theta _2`$ (3.8)
$`X_3`$ $`=`$ $`R\mathrm{sin}\theta _1\mathrm{sin}\theta _2\mathrm{cos}\theta _3`$ (3.9)
$`X_4`$ $`=`$ $`R\mathrm{sin}\theta _1\mathrm{sin}\theta _2\mathrm{sin}\theta _3\mathrm{cos}\theta _4`$ (3.10)
$`X_5`$ $`=`$ $`R\mathrm{sin}\theta _1\mathrm{sin}\theta _2\mathrm{sin}\theta _3\mathrm{sin}\theta _4.`$ (3.11)
The angles $`\theta _1,\mathrm{}.,\theta _3`$ go from $`0`$ to $`\pi `$. The angle $`\theta _4`$ is the azimuthal angle and goes from $`0`$ to $`2\pi `$. Then
$$X_1^2+X_2^2+X_3^2+X_4^2+X_5^2=R^2.$$
(3.12)
Next we embed a spherical membrane in $`S^4`$. We choose to parametrize the surface of the membrane by $`\theta _3,\theta _4`$. The brane is allowed to move in the $`X_1,X_2`$ plane. Its size depends on its location in the $`X_1,X_2`$ plane according to
$$r=R\mathrm{sin}\theta _1\mathrm{sin}\theta _2.$$
(3.13)
We see that when the size is at its maximum value, $`r=R`$, the membrane is at the origin $`X_1=X_2=0`$ like the two charges at the ends of the dipole in section 1. Since
$$X_1^2+X_2^2=R^2r^2,$$
(3.14)
the membrane can move around a circle in the plane and have constant size. We also set
$`X_1`$ $`=`$ $`\sqrt{R^2r^2}\mathrm{cos}\varphi `$ (3.15)
$`X_2`$ $`=`$ $`\sqrt{R^2r^2}\mathrm{sin}\varphi .`$ (3.16)
In terms of the coordinates $`r,\varphi ,\theta _3,\theta _4`$, the metric on the 4–sphere becomes
$$ds^2=\frac{R^2}{(R^2r^2)}dr^2+(R^2r^2)d\varphi ^2+r^2d\mathrm{\Omega }_2^2,$$
(3.17)
where $`d\mathrm{\Omega }_2^2`$ is the metric of a unit 2–sphere parametrized by $`\theta _3`$ and $`\theta _4`$. From the metric we see that the volume element is just
$$Rr^2drd\varphi d\mathrm{\Omega }_2.$$
(3.18)
The kinetic energy of the membrane is given by the Dirac–Born–Infeld Lagrangian. We are mostly interested in the case when the size of the sphere $`r`$ is constant and close to its maximum value $`R`$. In this case the membrane moves around a circle in the $`X_1,X_2`$ plane of radius $`(R^2r^2)^{1/2}`$ with angular velocity $`\dot{\varphi }`$. Dropping time derivatives of $`r`$, we have
$$_K=T\mathrm{\Omega }_2r^2\sqrt{1(R^2r^2)\dot{\varphi }^2}.$$
(3.19)
Here, $`T`$ is the membrane tension which is given by
$$T=\frac{1}{4\pi ^2l_p^3}$$
(3.20)
in 11-dimensional Planck units.
Next we add the Chern–Simons coupling involving the background field. The contribution of the four-form field strength to the action of the brane per orbit around the $`S^4`$ is
$$S_B=_{wv}C=_\mathrm{\Sigma }F.$$
(3.21)
The first integral is over the world-volume of the brane. $`F=dC`$ is the four-form flux, and $`\mathrm{\Sigma }`$ is a four-manifold in the $`S^4`$ whose boundary is the 3-dimensional surface swept out by the brane during one orbit. Since the background flux is just $`F=Bd\mathrm{vol}`$, where $`B`$ is the constant flux density and $`d\mathrm{vol}`$ is the volume form on $`S^4`$, we have
$$S_B=B\mathrm{vol}(\mathrm{\Sigma }).$$
(3.22)
Therefore the Chern-Simons term in the Lagrangian is
$$_B=\frac{S_B}{T}=B\mathrm{vol}(\mathrm{\Sigma })\frac{\dot{\varphi }}{2\pi }$$
(3.23)
where $`\dot{\varphi }`$ is the (constant) angular velocity of the brane. Parametrizing the motion as above, the volume of $`\mathrm{\Sigma }`$ is
$$\mathrm{vol}(\mathrm{\Sigma })=R\mathrm{\Omega }_2_0^{2\pi }𝑑\varphi _0^rr^2𝑑r^{}=\frac{8\pi ^2}{3}Rr^3.$$
(3.24)
So the Chern-Simons term is
$$_B=\frac{\dot{\varphi }}{2\pi }B\mathrm{\Omega }_4Rr^3=\dot{\varphi }N\frac{r^3}{R^3}$$
(3.25)
where we used the flux quantization condition, Eq.(3.1).
Therefore, the full bosonic Lagrangian is
$$=m\sqrt{1\dot{\varphi }^2(R^2r^2)}+N\frac{r^3}{R^3}\dot{\varphi }$$
(3.26)
with $`m=\mathrm{\Omega }_2Tr^2`$. Using Eq.(3.2) and Eq.(3.15), we also see that
$$\frac{N}{R^3}=T\mathrm{\Omega }_2.$$
(3.27)
From the Lagrangian we find that the angular momentum is given by
$$L=\frac{m\dot{\varphi }(R^2r^2)}{\sqrt{1\dot{\varphi }^2(R^2r^2)}}+N\frac{r^3}{R^3}.$$
(3.28)
The maximum size a membrane can have is $`R`$. Also, the velocity of its center of mass, $`\dot{\varphi }R`$, cannot exceed the speed of light. This implies that the angular momentum has a maximum value given by $`N`$ <sup>1</sup><sup>1</sup>1There is an exception to this statement at the pathological value $`r=0`$ which we discuss at the end of this section.:
$$L_{max}=N.$$
(3.29)
When the membrane has maximal size $`R`$, the angular momentum is the maximum value $`N`$. We see that the Kaluza–Klein graviton has a maximum angular momentum in agreement with the stringy exclusion principle. For the energy, we find
$$E=\dot{\varphi }L=\sqrt{\left(\frac{Nr^2}{R^3}\right)^2+\frac{(LNr^3/R^3)^2}{R^2r^2}}.$$
(3.30)
Varying the energy with respect to $`r`$ at fixed $`L`$, we find
$$\frac{dE}{dr}=\frac{r}{E(R^2r^2)^2}\left(LN\frac{r}{R}\right)\left(L2N\frac{r}{R}+N\frac{r^3}{R^3}\right).$$
(3.31)
We see that for $`L<N`$ there exists a stable minimum at
$$r=\frac{L}{N}R.$$
(3.32)
Therefore, the membrane grows as we increase the angular momentum. This is in agreement with Eq.(3.4) <sup>2</sup><sup>2</sup>2The cubic factor in $`dE/dr`$ has a positive root at $`0<r<(L/N)R`$ at which the energy has a local maximum and another root at $`r>R`$.. When $`r=R`$ and $`L=N`$, a more careful analysis for the stability of the solution is needed and we do so at the end of this section. The value of the energy at the minimum is
$$E=\frac{L}{R},$$
(3.33)
which is the energy of a Kaluza–Klein graviton with angular momentum $`L`$. From Eq.(3.23), we also find that the velocity of the center of mass equals the speed of light, $`\dot{\varphi }R=1`$.
We now show that there is a stable solution at $`r=R`$ and $`L=N`$. Setting $`\stackrel{~}{r}=Rr`$ and expanding the Lagrangian up to quadratic powers in $`\stackrel{~}{r}`$, we obtain
$$_K=T\mathrm{\Omega }_2R^2+T\mathrm{\Omega }_2R\stackrel{~}{r}(2+\dot{\varphi }^2R^2)T\mathrm{\Omega }_2\stackrel{~}{r}^2(\dot{\varphi }^4R^4+\frac{5}{2}\dot{\varphi }^2R^2+1),$$
(3.34)
and
$$_B=\mathrm{\Omega }_4BR^4\dot{\varphi }(13\frac{\stackrel{~}{r}}{R}+3\frac{\stackrel{~}{r}^2}{R^2})=N\dot{\varphi }R\stackrel{~}{r}\frac{3N}{R^3}\dot{\varphi }R+\stackrel{~}{r}^2\frac{3N}{R^3}\dot{\varphi }R.$$
(3.35)
Using $`N/R^3=T\mathrm{\Omega }_2`$, the total Lagrangian becomes
$$=T\mathrm{\Omega }_2R^2+N\dot{\varphi }+T\mathrm{\Omega }_2R\stackrel{~}{r}(23\dot{\varphi }R+\dot{\varphi }^2R^2)T\mathrm{\Omega }_2\stackrel{~}{r}^2(\dot{\varphi }^4R^4+\frac{5}{2}\dot{\varphi }^2R^23\dot{\varphi }R+1).$$
(3.36)
There is an extremum at $`\stackrel{~}{r}=0`$ provided that
$$23\dot{\varphi }R+\dot{\varphi }^2R^2=0.$$
(3.37)
This can be achieved if the velocity $`\dot{\varphi }R=1`$. Thus, when the size of the membrane is $`R`$ its center of mass moves with the speed of light. Furthermore, the extremum is stable since
$$\frac{d^2V(r)}{d^2r}|_{r=R}>0.$$
(3.38)
Thus when the size of the membrane is $`R`$, the angular momentum has its maximum value $`N`$ and the energy is given by
$$E=T\mathrm{\Omega }_2R^2=\frac{N}{R}.$$
(3.39)
At the classical level, this is in exact agreement with the energy of a Kaluza–Klein graviton having angular momentum $`N`$ about the sphere. When $`N>>1`$, the maximal angular momentum is large and the classical formula for the energy agrees with the BPS bound for the energy given the angular momentum. Quantum corrections are suppressed by $`1/N`$. The Kaluza–Klein graviton is a BPS state and its energy should not change under the process of blowing up into a membrane. Again, the size of the brane is determined by the angular momentum. Since the maximum size a brane can have is $`R`$, there is a maximum angular momentum as predicted by the dual conformal field theory. The fact that the energy of the brane agrees with the BPS formula for given angular momentum is a non–trivial test of our model.
We also note that there is a minimum of the energy at $`r=0`$ as well. Classically, it corresponds to a massless particle moving around the equator with angular momentum $`L`$.<sup>3</sup><sup>3</sup>3We thank Sunny Itzhaki for discussions of this point. Such a solution is singular from the perspective of the gravitational field equations since for angular momenta of order $`N`$, it represents a huge energy (of order $`N^{2/3}`$) concentrated at a point. Therefore it is subject to uncontrolled quantum corrections. In particular, there are quantum corrections proportional to powers of the momentum times the flux density, which are large at angular momenta of order $`N`$.
We have shown in this paper that such a singular solution can be resolved by blowing into a smooth macroscopic membrane of size $`(L/N)R`$. Our classical analysis is expected to be valid for the large membrane. The smooth membrane solution certainly has much more nearby phase space and as a result the true quantum ground state will be overwhelmingly supported at the membrane solution. We believe that this is another example of how string/$`M`$–theory resolves singularities of the type studied rececently in ref..
### 3.2 $`AdS_5\times S^5`$
The extension of our analysis to the other two maximally supersymmetric cases is straightforward and we will be much less explicit. Consider the case of $`AdS_5\times S^5`$ first. The radius of the five sphere is given by
$$R=(4\pi g_sN)^{\frac{1}{4}}l_s,$$
(3.40)
where $`g_s`$, $`l_s`$ are the string coupling constant and string length scale and $`N`$ is the number of units of five–form flux on the sphere. We take $`N`$ large keeping $`g_sN`$ fixed and large.
In type IIB string theory on $`AdS_5\times S^5`$ with $`N>>1`$, the maximum angular momentum of a BPS particle on the $`S^5`$ is $`N`$ . From the gauge theory perspective, this can be seen from the fact that one builds up such states by a single trace of the $`N\times N`$ scalars in the $`\mathrm{𝟔}`$ of $`SO(6)`$. The largest representation of $`SO(6)`$ one can build in this way is the spin-$`N`$ representation, $`\mathrm{Sym}^N\mathrm{𝟔}`$.
From our perspective this is because the particle moving on the sphere expands into a spherical D3-brane. In this case we present only the exact classical analysis of the D3-brane wrapping an $`S^3`$ that moves in $`S^5`$. The bosonic Lagrangian is
$$=_{DBI}+_{CS}=T_{D3}\mathrm{\Omega }_3r^3\sqrt{1(R^2r^2)\dot{\varphi }^2}+\dot{\varphi }N\frac{r^4}{R^4}.$$
(3.41)
The tension of the D3-brane is
$$T_{D3}=\frac{1}{(2\pi )^3l_s^4g_s}.$$
(3.42)
We will use the relation
$$T_{D3}\mathrm{\Omega }_3=\frac{N}{R^4}.$$
(3.43)
The angular momentum in terms of $`\dot{\varphi }`$ is
$$L=\frac{m\dot{\varphi }(R^2r^2)}{\sqrt{1\dot{\varphi }^2(R^2r^2)}}+N\frac{r^4}{R^4}$$
(3.44)
where $`m=T_{D3}\mathrm{\Omega }_3r^3=(N/R^4)r^3`$. Again we see that the angular momentum is bounded by $`N`$ since $`0rR`$ and $`0\dot{\varphi }R1`$. The energy is
$$E=\sqrt{m^2+\frac{(LNr^4/R^4)^2}{R^2r^2}}.$$
(3.45)
Varying the energy with respect to $`r`$ at fixed $`L`$, we find in this case a stable minimum when
$$r^2=\frac{L}{N}R^2.$$
(3.46)
The value of the energy at this minimum again matches the BPS bound when $`L`$ is large, for $`N>>1`$:
$$E=\frac{L}{R}.$$
(3.47)
This is strong evidence that at any appreciable momentum, at least at the (semi-) classical level, the good description of Kaluza–Klein gravitons is in terms of branes, rather than fundamental strings. From the dual CFT, we know that there is a unique BPS state with these quantum numbers; consistency with the exclusion principle implies that it is the one described by the spherical brane.
### 3.3 $`AdS_4\times S^7`$
In this case, we expect the graviton to expand into an M5-brane which is an $`S^5S^7`$. The radius of the sphere is given by
$$R=(2^5\pi ^2N)^{\frac{1}{6}}l_p.$$
(3.48)
The tension of the 5–brane is given by
$$T=\frac{1}{(2\pi )^5l_p^6}$$
(3.49)
and we have the relation
$$m=T\mathrm{\Omega }_5r^5=\frac{N}{R^6}r^5.$$
(3.50)
The Lagrangian is
$$=T\mathrm{\Omega }_5r^5\sqrt{1(R^2r^2)\dot{\varphi }^2}+\dot{\varphi }N\frac{r^6}{R^6}.$$
(3.51)
The angular momentum in terms of $`\dot{\varphi }`$ is
$$L=\frac{m\dot{\varphi }(R^2r^2)}{\sqrt{1\dot{\varphi }^2(R^2r^2)}}+N\frac{r^6}{R^6}.$$
(3.52)
The energy is
$$E=\sqrt{m^2+\frac{(LNr^6/R^6)^2}{R^2r^2}}.$$
(3.53)
Varying the energy with respect to $`r`$ at fixed $`L`$, we find in this case a stable minimum when
$$r^4=\frac{L}{N}R^4.$$
(3.54)
The value of the energy at this minimum again matches the BPS bound when $`L`$ is large:
$$E=\frac{L}{R}.$$
(3.55)
### 3.4 Remarks about $`AdS_3`$
We will conclude this section with some comments about $`AdS_3\times S^3\times M_4`$. This case is distinguished in several ways.
First, it is not clear into what the graviton should expand. Consider the geometry built from the D1-D5 system with $`Q_1`$ D-strings and $`Q_5`$ five-branes. The stringy exclusion bound on the angular momentum is $`LQ_1Q_5`$. The graviton is expected to blow up into a circular string moving on the $`S^3`$, but should it be a D5-brane wrapped on the four-manifold, or a D-string?
Secondly, the energetic considerations degenerate in this case. If we assume for argument’s sake that the graviton blows up into either a D-string or wrapped fivebrane which is a circle on the $`S^3`$, we find that the energy at fixed $`L`$ has no nontrivial minimum. Considering some incarnation of the “fractional strings” of does not help.
It may help to consider the S-dual situation of the F1-NS5 system. In this case, the dynamics of fundamental strings on the relevant sphere are described by a level-$`Q_5`$ $`SU(2)`$ WZW model. One expects the exclusion principle to be related to the affine cutoff on $`SU(2)`$ representations. Finally, a clarification of this case should match the result of that the exclusion bound occurs at the critical value of the energy for black hole formation. A better understanding remains for future work.
## 4 Conclusions
Physics in non-commutative spaces is characterized by a simple signature – the increase of size on systems with increasing momentum. In this paper we have seen that the motion of massless quanta on the $`S`$ factor of $`AdS_n\times S^m`$ has exactly this behavior. The massless particle blows up into a spherical brane of dimensionality $`m2`$ whose radius increases with increasing momentum. Eventually the radius of the blown up brane becomes equal to the radius of the sphere that contains it. It can no longer grow and the spectrum is terminated. This is the origin of the stringy exclusion principle. Thus we see one more piece of evidence for non-commutativity of space in quantum gravity.
## 5 Acknowledgments
We would like to thank Allan Adams, Micha Berkooz, Raphael Bousso, Sumit Das, Antal Jevicki, Jason Prince Kumar, Albion Lawrence, Alec Matusis, and Joe Polchinski for conversations. We are especially grateful to Joe Polchinski for explaining to us the flat-space energetics. We thank Hirosi Ooguri and Sunny Itzhaki for comments about the first version. We would also like to thank Steve Shenker for help with the title. The work of JM is supported in part by the Department of Defense NDSEG fellowship program. The work of LS and NT is supported in part by NSF grant 9870115. |
warning/0003/hep-th0003241.html | ar5iv | text | # 1 Introduction
## 1 Introduction
It has been well-known that there could be various types of extended objects in string theory and 11-dimensional M-theory. D-branes in superstring theories made us possible to study the duality relations among different theories. In particular, duality relation between gauge theory and supergravity or superstring theory, particularly AdS/CFT correspondence , has been investigated by using the D-branes as intermediates. Moreover, it has been discussed so much that moduli space structure of gauge theory can be read from the various brane configurations in background of superstring theory or M-theory . Thus, it is important to study the properties of branes in various backgrounds in order to deepen the understanding of gauge theory and string theory in the above viewpoints.
Branes in 10-dimensional type IIA string theory originate from 11-dimensional M-theory at least in the sense of low energy effective supergravity theory . For example, D4-branes and solitonic 5-branes (NS5-branes) in IIA supergravity theory are obtained by dimensional reduction of solitonic M5-branes, and fundamental strings and D2-branes are from membranes. As for D6-branes which carry Kaluza-Klein magnetic charges, original 11-dimensional counterpart is Kaluza-Klein monopole solution . It is described as the 11-dimensional background which is direct product of 7-dimensional Minkowski space $`^{1,6}`$ and the Euclidean Taub-NUT space $`M_{TN}`$.
We can further consider stable M-theory backgrounds where several different types of branes live together. A family of configurations representing process of brane creation from another brane belongs to this class . However, since little is known about the explicit supergravity solutions corresponding to such complicated configurations, the process of brane creation in terms of supergravity is still not understood well enough.
Now, consider a supersymmetric configuration of M-theory which has $`N`$ coincident KK-monopoles and an M5-brane of world-volume $`R^{1,3}\times \mathrm{\Sigma }`$ with $`\mathrm{\Sigma }M_{TN}`$. Upon compactification to 10-dimensional IIA theory, this becomes a configuration of an NS5-brane, a D4-brane or some bound state of these two in the background of $`N`$ D6-branes. Of such a class of systems, there is a family of configurations representing brane creation phenomenon. In a flat spacetime background, it is represented as the following process: If an NS5-brane of world-volume $`\{x_0,\mathrm{},x_3,x_7,x_8\}`$ crosses a D6-brane of world-volume $`\{x_0,\mathrm{},x_3,x_4,x_5,x_6\}`$, a new finite D4-brane between these two branes is created. Here $`x_i`$ $`(i=0,\mathrm{},9)`$ denote the spacetime coordinates. This process is U-dual to the original Hanany-Witten configuration . An attempt of investigating such a phenomenon in 11-dimensional viewpoint was done in e.g., refs.. In particular, in ref., one parameter family of M5-branes in the supergravity background of a KK-monopole and its compactified counterpart is investigated.
On the other hand, in another context, explicit supergravity solutions describing M2- or M5-branes localized near core of KK-monopoles were constructed by using the fact that the metric becomes flat in the vicinity of KK-monopoles . The solution of an M5-brane with world-volume $`R^{1,3}\times \mathrm{\Sigma }_0`$ ($`\mathrm{\Sigma }_0M_{TN}`$) near core of KK-monopoles was found in ref.. It is naturally expected that the near core version of brane creation phenomenon explained in the last paragraph can be described by dimensional reduction of this solution. In practice, in refs., some discussion about brane creation in 10-dimensional viewpoint was done.
It was pointed out in ref. that in the near core region of KK-monopoles, there is some ambiguity in the definition of Ramond-Ramond four-form field strength and NS-NS three-form field strength. This means that the identification of D4-branes and NS5-branes cannot be done exactly. In ref., this problem was further studied and a resolution by introducing a certain non-conserved charge was suggested. Using this procedure, it was argued there that an NS5-brane is transmuted into a D4-brane in a certain limit and that this phenomenon is a near core version of the brane creation. There still remains a problem in the sense that the argument is restricted in the near core region of KK-monopoles and that the role of non-conserved charge is unclear.
Our aim of the present paper is to confirm the definition of brane current (i.e., the identification of branes) and to investigate the brane creation phenomena in 10-dimensions for systems of an M5-brane in the background of $`N`$ KK-monopoles. As for the definition of current associated with branes, we adopt the natural definition of current which assigns conserved charge to branes in the background of D6-branes. The ambiguity presented above can be resolved by noticing that the location of Dirac string type singularity coming up from the D6-branes is relevant to the identification of D4-branes. This interpretation is essentially the same as the charge assignment of the systems defined by T-dual of $`(p,q)`$5-branes or $`(p,q)`$-strings in the 7-brane background .
We also give an explicit relation between the complex structures of $`M_{TN}`$ and some specific complex coordinates of the near core flat space. As a result, it becomes possible to discuss the brane creation process of ref. in terms of exact supergravity solution in the near core region of KK-monopoles. By applying our definition of brane charge and the way of identifying branes, we will argue that D4-branes seem to come up from the Dirac string type singularity and go along the NS5-brane. If we take a limit such that net NS5-brane charge disappears, it can be seen that only D4-branes come up from D6-branes.
The organization of this paper is as follows. We begin in section 2 to review the Kaluza-Klein monopole solutions in 11-dimensional background. We explain the properties of Euclidean Taub-NUT space $`M_{TN}`$ and explicitly give three independent complex structures by defining holomorphic coordinates corresponding to them. In section 3 we see that in the near core region of KK-monopoles, the background $`M_{TN}`$ reduces to the flat space. Moreover, we specify the complex coordinates of this flat space so that they are connected to the complex structures of $`M_{TN}`$. In section 4 we introduce a family of complex one-dimensional curves $`\{\mathrm{\Sigma }\}`$ in $`M_{TN}`$ which respectively correspond to configurations of an M5-brane of world-volume $`R^{1,3}\times \mathrm{\Sigma }`$ in the KK-monopole backgrounds. Each of the curves are taken to be holomorphic with respect to one of the complex structures. In section 5, we further study such systems by restricting in the near core region of KK-monopoles and give their supergravity solutions. In section 6, we perform compactification of the systems and explain the ambiguity concerning the identification of D4-branes. In section 7 we digress from our systems for a moment and argue the configuration of M-branes in the stringy cosmic string background and its compactification. In particular, we see that the Dirac string type singularity, which appears as IIA counterpart of the branch cut related to SL(2,Z) invariance of type IIB theory, has an important role in identifying the type of branes. Then we return to the KK-monopole backgrounds in senction 8 and by using the analogy with the way of identifying branes in the stringy cosmic string background, we give an unambiguous identification of D4-branes. In section 9, using our definition of brane currents, we investigate the phenomena of brane creation. In the final section 10, we conclude with some discussion.
## 2 Kaluza-Klein monopoles and the Taub-NUT space
We now review the 11-dimensional supergravity solution of $`N`$ coincident Kaluza-Klein monopoles which represents $`N`$ D6-branes after compactifying to 10-dimensions. The solution is given by taking the 11-dimensional spacetime as a product of 7-dimensional Minkowski space $`^{1,6}`$ and the four dimensional Euclidean multi-centered Taub-NUT space $`M_{TN}`$ . The Taub-Nut space is a Hyper-Kähler manifold which admits three independent complex structures.
For the Taub-NUT space with $`A_{N1}`$ singularity at $`r=0`$, the metric is given by
$$ds_{TN}^2=Vd\stackrel{}{r}d\stackrel{}{r}+V^1(dx_{11}+\stackrel{}{A}d\stackrel{}{r})^2$$
(1)
where $`x_{11}`$ is a compact direction whose radius is $`R`$ : $`x_{11}x_{11}+2\pi R`$, and
$$\stackrel{}{r}=(r_1,r_2,r_3),V=1+\frac{NR}{2r},\stackrel{}{}\times \stackrel{}{A}=\stackrel{}{}V.$$
(2)
We explicitly choose $`\stackrel{}{A}`$ to be
$$\stackrel{}{A}d\stackrel{}{r}=\frac{NR}{2}(\mathrm{cos}\theta 1)d\psi $$
(3)
where
$$r_1=r\mathrm{cos}\theta ,r_2=r\mathrm{sin}\theta \mathrm{cos}\psi ,r_3=r\mathrm{sin}\theta \mathrm{sin}\psi $$
(4)
and $`0\psi 2\pi `$. Since the space is hyper-Kähler and has Ricci-flat metric, it solves Einstein equations and admits half of supersymmetry. We will only deal with the bosonic part of the theory.
The metric eq.(1) with eq.(3) has a coordinate singularity at $`\theta =\pi `$, i.e., the negative $`r_1`$-axis. After dimensional reduction, seen in terms of ten dimensions, this singularity is identified as the Dirac string singularity coming up from the D6-branes. This singularity can be moved to another place by carrying out a coordinate transformation. The simplest example is the transformation $`(\stackrel{}{r},x_{11})(\stackrel{}{r},y_{11})`$ where $`y_{11}x_{11}NR\psi `$. Then, by the relation $`dx_{11}+\stackrel{}{A}d\stackrel{}{r}=dy_{11}+\stackrel{}{A^{}}d\stackrel{}{r}`$ with
$$\stackrel{}{A^{}}d\stackrel{}{r}=\frac{NR}{2}(\mathrm{cos}\theta +1)d\psi ,$$
(5)
the coordinate singularity is moved to the positive $`r_1`$-axis.
Now we give explicitly three independent complex structures of the space $`M_{TN}`$. First we present a natural choice of complex structure for the coordinate system $`(\stackrel{}{r},x_{11})`$. It is specified by the holomorphic complex variables $`(v_0,w_0)`$ as
$`v_0`$ $`=`$ $`{\displaystyle \frac{1}{NR}}(r_2+ir_3),`$ (6)
$`w_0`$ $`=`$ $`\sqrt{{\displaystyle \frac{r+r_1}{NR}}}\mathrm{exp}\left({\displaystyle \frac{1}{NR}}(r_1+ix_{11})\right).`$ (7)
Here we take $`v_0`$ and $`w_0`$ to be dimensionless. Remembering $`x_{11}x_{11}+2\pi R`$, $`w_0^N`$ may be more useful than $`w_0`$ as a coordinate variable. Other two complex structures that are orthogonal to the above one are represented by the complex variables $`(v_1,w_1)`$ and $`(v_2,w_2)`$ which are holomorphic with respect to the complex structures respectively:
$$\{\begin{array}{cc}\hfill v_1& =\frac{1}{NR}(r_1+ir_2),\hfill \\ \hfill w_1& =\sqrt{\frac{r+r_3}{NR}}\mathrm{exp}\left(\frac{1}{NR}(r_3+ix_{11}^{(1)})\right)\hfill \end{array}$$
(8)
and
$$\{\begin{array}{cc}\hfill v_2& =\frac{1}{NR}(r_3+ir_1),\hfill \\ \hfill w_2& =\sqrt{\frac{r+r_2}{NR}}\mathrm{exp}\left(\frac{1}{NR}(r_2+ix_{11}^{(2)})\right).\hfill \end{array}$$
(9)
Here $`x_{11}^{(1)}`$ and $`x_{11}^{(2)}`$ are compact coordinates with period $`2\pi R`$ determined by the differential equations
$$dx_{11}^{(1)}=dx_{11}\frac{NR}{2r}\left(\frac{r_2dr_3r_3dr_2}{r+r_1}\frac{r_1dr_2r_2dr_1}{r+r_3}\right)(=dx_{11}NRd\chi )$$
(10)
and
$$dx_{11}^{(2)}=dx_{11}\frac{NR}{2r}\left(\frac{r_2dr_3r_3dr_2}{r+r_1}\frac{r_3dr_1r_1dr_3}{r+r_2}\right).$$
(11)
These equations can be integrated explicitly except on singularity: negative $`r_1`$ and $`r_3`$ axes for $`x_{11}^{(1)}`$, and negative $`r_1`$ and $`r_2`$ axes for $`x_{11}^{(2)}`$. These singularities are due to the coordinate transformation from $`(\stackrel{}{r},x_{11})`$ (which is singular on negative $`r_1`$ axis) to $`(\stackrel{}{r},x_{11}^{(1)})`$ (which is singular on negative $`r_3`$ axis) or to $`(\stackrel{}{r},x_{11}^{(2)})`$.
Note that the metric in the coordinate system $`(\stackrel{}{r},x_{11}^{(1)})`$ is
$$ds_{TN}^2=Vd\stackrel{}{r}d\stackrel{}{r}+V^1(dx_{11}^{(1)}+\stackrel{}{A^{\prime \prime }}d\stackrel{}{r})^2$$
(12)
where
$$\stackrel{}{A^{\prime \prime }}d\stackrel{}{r}=\frac{NR}{2r}\frac{r_1dr_2r_2dr_1}{r+r_3}.$$
(13)
If we compactify along the $`x_{11}^{(1)}`$ direction instead of $`x_{11}`$ direction, then the resulting 10-dimensional D6-brane solution has Dirac string singularity on negative $`r_3`$ axis. The movement of the singularity corresponds to gauge transformation: $`\stackrel{}{A}\stackrel{}{A^{\prime \prime }}=\stackrel{}{A}+d\chi `$. In the following discussion, we fix the compactification direction as along $`x_{11}`$.
## 3 Near core region of KK-monopoles
In this section, we see that the metric $`ds_{TN}^2`$ reduces to the flat metric in the near core region of KK-monopoles . We also show that a flat complex coordinate system in this region is exactly related to one of the complex structures described in section 2. Note that by the near core region, we mean the region $`rNR`$ in terms of the coordinate system $`(\stackrel{}{r},x_{11})`$.
In the limit $`r/NR0`$, the metric (1) becomes
$$ds_{r0}^2=\frac{NR}{2r}d\stackrel{}{r}d\stackrel{}{r}+\frac{2r}{NR}\left(dx_{11}+\frac{NR}{2}(\mathrm{cos}\theta 1)d\psi \right)^2$$
(14)
since $`V(=1+\frac{NR}{2r})\frac{NR}{2r}`$. By changing variables from $`(r,\theta ,\psi ,x_{11})`$ to $`(\rho ,\stackrel{~}{\theta },\stackrel{~}{\psi },\stackrel{~}{\varphi })`$ such as
$$\rho =\sqrt{2NRr},\stackrel{~}{\theta }=\frac{\theta }{2},\stackrel{~}{\psi }=\psi +\frac{x_{11}}{NR},\stackrel{~}{\varphi }=\frac{x_{11}}{NR},$$
(15)
the metric (14) can be rewritten:
$$ds_{r0}^2=d\rho ^2+\rho ^2d\stackrel{~}{\theta }^2+\rho ^2(\mathrm{sin}^2\stackrel{~}{\theta }d\stackrel{~}{\psi }^2+\mathrm{cos}^2\stackrel{~}{\theta }d\stackrel{~}{\varphi }^2).$$
(16)
Here the range of the variables is
$$\rho 0,0\stackrel{~}{\theta }\frac{\pi }{2},0\stackrel{~}{\varphi },\stackrel{~}{\psi }2\pi $$
(17)
with the $`𝐙_N`$ identification $`(\stackrel{~}{\varphi },\stackrel{~}{\psi })(\stackrel{~}{\varphi },\stackrel{~}{\psi })+(2\pi /N,2\pi /N)`$.
Furthermore, defining the complex variables as<sup>1</sup><sup>1</sup>1Here we choose the complex variables $`V`$ and $`W`$ such that they are connected to holomorphic variables of the complex structure $`(v_0,w_0)`$.
$$W=\rho e^{i\stackrel{~}{\varphi }}\mathrm{cos}\stackrel{~}{\theta },V=\rho e^{i\stackrel{~}{\psi }}\mathrm{sin}\stackrel{~}{\theta },$$
(18)
eq.(16) becomes
$$ds_{r0}^2=dWd\overline{W}+dVd\overline{V}$$
(19)
where $`(W,V)(We^{2\pi i/N},Ve^{2\pi i/N})`$. Thus we see that the space is a flat orbifold $`𝐂^2/𝐙_N`$ which is considered as the ALE space with $`A_{N1}`$ singularity at the origin. For later convenience we represent $`(W,V)`$ by the original variables:
$`W`$ $`=`$ $`\sqrt{NR}e^{i\frac{x_{11}}{NR}}\sqrt{r+r_1},`$ (20)
$`V`$ $`=`$ $`\sqrt{NR}e^{i\frac{x_{11}}{NR}}{\displaystyle \frac{r_2+ir_3}{\sqrt{r+r_1}}}`$ (21)
$`=`$ $`NRW^1(r_2+ir_3).`$
Note that there is no coordinate singularity in $`(W,V)`$ space except at the orbifold point. The coordinate singularity $`\theta =\pi `$ in the $`(\stackrel{}{r},x_{11})`$ coordinate system corresponds to $`W=0`$, which is completely smooth in the new coordinate system. Similarly, $`\theta =0`$, which corresponds to coordinate singularity in the $`(\stackrel{}{r},y_{11})`$ system, appears as $`V=0`$ which is regular in $`(W,V)`$ space.
Next, we study the relation between $`(W,V)`$ and the complex structure of $`M_{TN}`$ represented by $`(w_0,v_0)`$. The behavior of $`w_0`$ and $`v_0`$ in the near core limit is obtained by extracting the lowest order terms of $`r_i/NR`$ in the eqs.(6) and (7):
$`w_0`$ $``$ $`\sqrt{{\displaystyle \frac{r+r_1}{NR}}}e^{i\frac{x_{11}}{NR}}(w_{0(r0)}),`$ (22)
$`v_0`$ $`=`$ $`{\displaystyle \frac{1}{NR}}(r_2+ir_3)(v_{0(r0)}).`$ (23)
Thus, we see that the complex variables $`(W,V)`$ are related to the $`(w_0,v_0)`$ of $`M_{TN}`$ as
$`w_{0(r0)}`$ $`=`$ $`{\displaystyle \frac{1}{NR}}W,`$ (24)
$`v_{0(r0)}`$ $`=`$ $`{\displaystyle \frac{1}{(NR)^2}}VW.`$ (25)
This means that in the near core limit, we can take $`(W,V)`$ as variables representing the complex structure.
We can also relate the variables $`(W,V)`$ with other complex structures of $`M_{TN}`$. In the near core flat space, a complex structure orthogonal to $`(W,V)`$ is given by a linear combination of two complex structures :
$$(z_7+iz_8,z_9+iz_{10})\mathrm{and}(z_9+iz_7,z_8+iz_{10})$$
(26)
where we introduced real variables $`z_i`$ as
$$W=z_7+iz_{10},V=z_8+iz_9.$$
(27)
In practice, we are able to show as the same way as the above discussion that two complex structures presented in eq.(26) are considered as the limiting representations of the complex structures $`(w_1,v_1)`$ and $`(w_2,v_2)`$ of $`M_{TN}`$ respectively. For example, in the $`r/NR0`$ limit, the variables $`(w_1,v_1)`$ become
$`w_{1(r0)}`$ $`=`$ $`{\displaystyle \frac{1}{NR}}W^{(1)},`$ (28)
$`v_1`$ $`=`$ $`{\displaystyle \frac{1}{NR}}W^{(1)}V^{(1)}`$ (29)
where
$`W^{(1)}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}((z_7+iz_8)+(z_9+iz_{10})),`$ (30)
$`V^{(1)}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}((z_7+iz_8)(z_9+iz_{10})).`$ (31)
We again have a relation $`(W^{(1)},V^{(1)})(W^{(1)}e^{2\pi i/N},V^{(1)}e^{2\pi i/N})`$. Note that as for a choice of origin of the two periodic coordinates $`x_{11}`$ and $`x_{11}^{(1)}`$ in the flat $`(W,V)`$ space, we have taken $`x_{11}NR\mathrm{arg}(z_7+iz_{10})`$ and $`x_{11}^{(1)}NR\mathrm{arg}(W^{(1)}+V^{(1)})`$.
## 4 M-branes in the KK-monopole background
We consider to put an M5-brane or an M2-brane in the KK-monopole background $`^{1,6}\times M_{TN}`$ such that two spatial dimensions of the M-brane are included in $`M_{TN}`$, i.e., the world volume of the M5-brane (or M2-brane) is $`R^{1,3}\times \mathrm{\Sigma }`$ (or $`R^{1,0}\times \mathrm{\Sigma }`$) where $`\mathrm{\Sigma }`$ is a two-dimensional surface contained in $`M_{TN}`$. By the discussion in refs. , the curve $`\mathrm{\Sigma }M_{TN}`$ must be holomorphic with respect to a complex structure in order to preserve a certain number, in this case $`1/4`$ of 32, of supersymmetry. From now on, we concentrate on the case of M5-branes.
Now we consider two particular curves in $`M_{TN}`$ discussed in refs. which are holomorphic with respect to $`w_0`$ and $`v_0`$:
$`(A)`$ $`(w_0)^N=e^{\frac{b+i\alpha }{R}},`$ (32)
$`(B)`$ $`(w_0)^N=e^{\frac{b+i\alpha }{R}}(v_0)^N`$ (33)
where $`b`$ and $`\alpha `$ are some real parameters.
Remember that $`M_{TN}`$ is spanned by the coordinate system $`(\stackrel{}{r},x_{11})`$ except on the negative $`r_1`$ axis. By representing above two curves in this coordinate system, we can investigate the shape of curves as embedding objects in $`(r_1,r_2,r_3)`$ and the behavior of the parameter $`x_{11}`$ on these curves. It can easily be seen that the two curves with the same $`b`$ have exactly the same shape in $`(r_1,r_2,r_3)`$. The curve $`(A)`$ is
$$\sqrt{\frac{r+r_1}{NR}}=e^{\frac{r_1+b}{NR}}$$
(34)
and the curve $`(B)`$ is represented as the same equation but with $`r_1r_1`$. The behavior of the parameter $`x_{11}=x_{11}(\stackrel{}{r})`$ is as follows: The curve $`(A)`$ is at $`x_{11}=\alpha `$, while $`(B)`$ is winding around $`x_{11}`$ along the $`\psi `$ \[$`=\mathrm{arg}(r_2+ir_3)`$\] direction as $`x_{11}=NR\psi +\alpha `$. The important point is that the locations of coordinate singularity, which will be interpreted as the Dirac string singularity in 10-dimendions, relative to the curves $`(A)`$ and $`(B)`$ are different from each other. (See Fig.1.)
In the case of $`N=1`$, it was argued in refs. that each of these two curves corresponds to a single NS5-brane or a configuration consisting of a finite D4-brane between an NS5-brane and a D6-brane in the compactified 10-dimensions. We postpone the analysis of compactification until section 6. We notice that there are many other possible curves which leave the same number of supersymmetry and have the same shape as the above two curves: We can take the curves holomorphic to any linear combination of all three complex structures of $`M_{TN}`$. They all have the same shape in $`(r_1,r_2,r_3)`$ space and are represented by the equations rotating eq.(34) around the $`r=0`$ point suitably. The location of the singularity is still on the negative $`r_1`$ axis. The value of $`x_{11}`$ on each of the curves is represented as $`x_{11}=f(\stackrel{}{r})\alpha `$. The form of the function $`f=f(\stackrel{}{r})`$ is generically not a simple form except for the curve based on the complex structure $`(w_0,v_0)`$. For example, the curves holomorphic with respect to $`(v_1,w_1)`$ are given as
$`(A)^{}`$ $`(w_1)^N=e^{\frac{b+i\alpha }{R}},`$ (35)
$`(B)^{}`$ $`(w_1)^N=e^{\frac{b+i\alpha }{R}}(v_1)^N.`$ (36)
For the curve $`(A)^{}`$, $`f(\stackrel{}{r})=NR\chi `$ where $`\chi `$ is given by eq.(10).
In the following discussion, we only deal with a class of M5-branes specified by the curves we have given above. In particular, we often proceed with discussion by choosing the curves $`(A)`$, $`(B)`$ and $`(A)^{}`$ as examples.
## 5 M5-branes in the near core region of KK-monopoles and the supergravity solutions
In the last section, we considered a family of curves which correspond to configurations of an M5-brane in the background of KK-monopoles. Supergravity solutions representing such systems have not been constructed. However, if we restrict ourselves to the near core region of KK-monopoles, solutions can be obtained by using the method of . We know that the near core region of KK-monopole solution is represented by the flat complex coordinates $`(W,V)`$. Thus, by considering some two-dimensional flat plane $`\mathrm{\Sigma }_0`$ in this region, supergravity solution of an M5-brane with world-volume $`R^{1,3}\times \mathrm{\Sigma }_0`$ is obtained. By transforming back the coordinates to $`(\stackrel{}{r},x_{11})`$, we know the behavior of the M5-brane in the near core region of KK-monopoles .
Now, from the analysis in section 3, we have the relation between the near core coordinates $`(W,V)`$ and the complex structures of the whole Taub-NUT space. Using this knowledge, we study the behavior of the class of holomorphic curves considered in the last section in the near core region. In order to do this, parts of the curves must be in the near core region, i.e., the distance between the curve and the core of the monopoles $`r=0`$ has to be much smaller than $`NR`$. This is equivalent to the condition $`e^{b/NR}1`$. Assuming this, by the relations (24) and (25), the two curves $`(A)`$ and $`(B)`$ given in eqs. (32) and (33) can be approximated in the $`r/NR0`$ limit as
$`\stackrel{~}{(A)}`$ $`W^N=c^N,`$ (37)
$`\stackrel{~}{(B)}`$ $`V^N=c^N`$ (38)
where
$$c=NRe^{\frac{b+i\alpha }{NR}}.$$
(39)
Also, the curves $`(A)^{}`$ and $`(B)^{}`$ in (35) and (36) can be rewritten in this region as
$`\stackrel{~}{(A)^{}}`$ $`(W^{(1)})^N=c^N,`$ (40)
$`\stackrel{~}{(B)^{}}`$ $`(V^{(1)})^N=c^N.`$ (41)
These are $`N`$ copies of flat planes in the flat $`(W,V)`$ space as is expected. Similarly, all the other curves in the class considered in the previous section are reduced to flat curves in this region. Note that each of these curves is a paraboloid if seen in $`(r_1,r_2,r_3)`$. For example, the curve $`\stackrel{~}{(A)}`$ is
$$r_1=\frac{r_2^2+r_3^2}{2a}+\frac{a}{2}$$
(42)
where $`a=|c|^2/NR`$.
The supergravity solution realizing one of these curves $`\mathrm{\Sigma }_0`$ as an M5-brane of world-volume $`R^{1,3}\times \mathrm{\Sigma }_0`$ can be obtained as the same way as in ref. . Also, the explicit forms of four-form field strength $`F_{[4]}`$ and the current $`J_{[5]}^{(11)}`$ of M5-brane can be given .
Note that although for any flat plane in the space $`(W,V)`$ we can obtain the corresponding supergravity solution, we limit ourselves to the class of flat curves such that they are represented as the $`r/NR0`$ limit of the curves described in the last section. For example, curves like $`(aW+bV)^N=c^N`$ or $`(aW+b\overline{V})^N=c^N`$ are excluded unless $`a=0`$ or $`b=0`$.
Now, we give the solution of M5-brane corresponding to the curve $`\stackrel{~}{(A)}`$, $`W^N=c^N`$, explicitly. The metric is given by
$$ds_{11}^2=f_5^{\frac{1}{3}}(\eta _{\mu \nu }dx^\mu dx^\nu +dVd\overline{V})+f_5^{\frac{2}{3}}(dy^mdy^m+dWd\overline{W})$$
(43)
where $`\mu ,\nu =0,\mathrm{}3`$, $`m=4,5,6`$, and
$$f_5=1+\underset{l=1}{\overset{N}{}}\frac{k}{(s^2+|Wce^{\frac{2\pi l}{N}i}|^2)^{3/2}},s^2=y_4^2+y_5^2+y_6^2.$$
(44)
Four-form field strength $`F_{[4]}`$ is
$$F_{p_1p_2p_3p_4}=3kϵ_{p_1p_2p_3p_4q}\underset{l=1}{\overset{N}{}}\frac{\stackrel{~}{y^q}}{(s^2+|Wce^{\frac{2\pi l}{N}i}|^2)^{5/2}}$$
(45)
where $`\stackrel{~}{y}^p=(x_4,x_5,x_6,z_7,z_{10})`$. The Hodge dual of current associated with the M5-brane is given by
$`dF_{[4]}`$ $``$ $`\stackrel{~}{}J_{[5]}^{(11)}`$ (46)
$`=`$ $`{\displaystyle \underset{l=1}{\overset{N}{}}}3k\mathrm{\Omega }_4\delta (y_4)\delta (y_5)\delta (y_6)\delta (Wce^{\frac{2\pi l}{N}i})dy_4dy_5dy_6dz_7dz_{10}.`$
where $`\stackrel{~}{}`$ denotes the 11-dimensional Hodge dual. Similarly, we can give the same information as above for other M5-brane systems.
At this point, we comment on the definition of current with an M5-brane in the background of KK-monopoles besides the near core region. In the generic $`r`$:finite region, we do not have the supergravity solution and the definite form of the current cannot be obtained. However, only from the knowledge of the location of the M5-brane in a certain coordinate system of $`M_{TN}`$, e.g., $`(\stackrel{}{r},x_{11})`$, an approximate form of the Hodge dual of current associated with the M5-brane is determined. That is, for a two dimensional curve described by $`g_i(\stackrel{}{r},x_{11})=0`$ ($`i=1,2`$), Hodge dual of the current of the M5-brane at $`g_i=0`$, $`y_4=y_5=y_6=0`$ is denoted by the 5-form which has the nonzero value only on the M5-brane:
$$\stackrel{~}{}J_{[5]}^{(11)}\delta (y_4)\delta (y_5)\delta (y_6)\delta (g_1)\delta (g_2)dy_4dy_5dy_6dg_1dg_2.$$
(47)
We will analyze the brane identification in 10 dimensions approximately from the knowledge of the form of the wedge product $`dg_1dg_2`$, in particular, the behavior of the term including $`dx_{11}`$.
## 6 Ten-dimensional analysis
We investigate the compactification of the systems of an M5-brane in the KK-monopole background. We fix the compactification direction along $`x_{11}`$. It is known that an M5-brane becomes an NS5-brane, a D4-brane or a bound state of these two types of branes upon compactification. However it is argued in refs. that there is a sort of subtlety in identification of branes when the original 11-dimensional system has M5-branes and KK-monopoles together. We clarify the problem by using the exact form of metric and four-form field strength in the $`r0`$ region and by performing the compactification explicitly.
Assuming the isometry along $`x_{11}`$ direction, 11-dimensional theory is related to 10-dimensions as follows:
$`ds_{11}^2`$ $`=`$ $`e^{\frac{2}{3}\mathrm{\Phi }}ds_{10}^2+e^{\frac{4}{3}\mathrm{\Phi }}(dx_{11}+A_{[1]})^2`$ (48)
$`F_{[4]}^{(11)}`$ $`=`$ $`G_{[4]}+G_{[3]}dx^{11}.`$ (49)
Here $`ds_{10}^2`$ is 10-dimensional string metric, $`\mathrm{\Phi }`$ is 10-dimensional dilaton and $`A_{[1]}`$ is Ramond-Ramond 1-form potential. As for the four-form field strength, we used the notation $`G_{[4]}dB_{[3]}`$ and $`G_{[3]}dB_{[2]}`$ where $`B_{[3]}`$ and $`B_{[2]}`$ are Ramond-Ramond 3-form potential and NS-NS 2-form potential respectively. There is another description of 10-dimensional four-form field strength $`\stackrel{~}{G}_{[4]}`$ given by
$$F_{[4]}^{11}=\stackrel{~}{G}_{[4]}+G_{[3]}(dx^{11}+A_{[1]})$$
(50)
where $`\stackrel{~}{G}_{[4]}=G_{[4]}G_{[3]}A_{[1]}`$. This definition may be more convenient since $`A_{[1]}`$ and $`G_{[4]}`$ couple with each other and the related term in 10-dimensional supergravity action is collected by the form $`\stackrel{~}{G}_{[4]}^2`$ after all. There is another advantage in using the latter definition which is related to the ‘gauge invariance.’ In terms of 11-dimensions, the gauge symmetry is represented as a coordinate transformation $`(\stackrel{}{r},x_{11})(\stackrel{}{r},x_{11}^{})`$ where $`x_{11}^{}=x_{11}+\gamma (\stackrel{}{r})`$. As a metric of the form eq.(48), it is written as gauge transformation of $`A_{[1]}`$ :
$$\{\begin{array}{cc}\hfill dx_{11}& dx_{11}^{}=dx_{11}+d\gamma \hfill \\ \hfill A_{[1]}& A_{[1]}^{}=A_{[1]}d\gamma .\hfill \end{array}$$
(51)
It is easily seen that under this gauge transformation, $`\stackrel{~}{G}_{[4]}`$ is invariant but $`G_{[4]}`$ is not. We will argue this point later again.
We will apply the general formula of compactification to our configurations given in the last section. In order to compactify the system along $`x_{11}`$, there should be an isometry along the direction. Thus for any of our M5-brane systems with parameter $`c0`$, we have to put infinite number of M5-branes uniformly along the $`x_{11}`$ direction, i.e., we consider a family of curves with $`\alpha =2\pi Rk/n`$ ($`k=1,2,\mathrm{}n`$), and take the limit $`n\mathrm{}`$. Note that if $`c=0`$, there is already an isometry along $`x_{11}`$. For example, for the M5-brane system given by the curve $`W^N=c^N`$ in eq.(43), we should put infinite M5-branes on the circle of radius $`|W|=|c|`$, which corresponds to take
$$f_5=1+k^{}_0^{2\pi }(s^2+(W|c|e^{i\xi })^2)^{\frac{3}{2}}𝑑\xi $$
(52)
where $`k^{}`$ is some regularized parameter. Then, the compactification can be performed, and $`ds_{10}^2`$, $`\mathrm{\Phi }`$ and $`A_{[1]}`$ are calculated according to eq.(48).
Now we define the currents associated with D4-branes and with NS-NS 5-branes in 10-dimensions. The most straightforward definition is
$$dG_{[3]}=j_{[5]},dG_{[4]}=j_{[4]}$$
(53)
where $`j_{[5]}`$ and $`j_{[4]}`$ are the 10-dimensional Hodge dual of the currents of NS-NS 5-branes and D4-branes respectively. The relation of these currents with 11-dimensional current $`J_{[5]}^{(11)}`$ is obtained by differentiating eq.(49) as
$$\stackrel{~}{}J_{[5]}^{(11)}=j_{[4]}+j_{[5]}dx_{11}.$$
(54)
In ref., another definition of D4-brane current $`\stackrel{~}{j}_{[4]}`$ is proposed as
$$\stackrel{~}{j}_{[4]}=d\stackrel{~}{G}_{[4]}F_{[2]}G_{[3]}$$
(55)
where $`F_{[2]}=dA_{[1]}`$. Two different definitions of D4-brane current coincide with each other when there is no R-R 1-form potential. The main difference is that the latter definition, $`\stackrel{~}{j}_{[4]}`$, respects the gauge invariance eq.(51) of $`A_{[1]}`$ but the former does not, and that the former assures the conservation of charge but the latter does not. This means that the D4-brane charge defined by the current $`\stackrel{~}{j}_{[4]}`$ remains unchanged if we change the compactification direction as $`x_{11}x_{11}^{}`$. In ref., by taking the current $`\stackrel{~}{j}_{[4]}`$ as physically relevant definition of D4-branes, charge non-conservation and brane transmutation are argued.
On the contrary, here we choose a standpoint of interpreting the naive definition of D4-brane current $`j_{[4]}`$ as physically relevant one. This viewpoint gives the physical meaning to the Dirac string type singularity associated with $`A_{[1]}`$. It seems strange at first sight, however, we can offer a reasonable explanation. First of all, we remember that the gauge invariance of $`A_{[1]}`$, whose determination specifies the location of the singularity, is not a symmetry of the compactified 10-dimensional theory in itself, but a symmetry in 11-dimensions as in eqs.(51). In compactifying to 10-dimensions, we have to specify $`x_{11}`$ or $`x_{11}^{}`$ definitely as an eleventh direction along which the reduction is executed. It is nothing but fixing of a gauge in terms of (51). The Dirac string singularity in 10-dimensions cannot be moved to another place without changing the compactification direction. Thus, looking in 10-dimensions, this singularity is not necessary to be unphysical at least concerning the identification of branes. In practice, in our generic configurations, we will see that D4-branes come up from the singularity and go to infinity along an NS5-brane.
Before demonstrating the consequence of the above interpretation of 10-dimensional currents, in the next section we deal with similar configurations in which there exists Dirac string type singularity. Using the models, we examine the relation between the singularity and the brane identification.
## 7 M-branes in the stringy cosmic string backgrounds
We take the stringy cosmic string solution as 11-dimensional supergravity background. The solution is known to represent a configuration with some parallel D7-branes ($`[1,0]`$-branes) or their SL$`(2,𝐙)`$ dual $`[p,q]`$7-branes if it is taken as a IIB background . It is also interpreted as the compactification of 12-dimensional F-theory on K3, or non-compact K3, which admits elliptic fibration. We will only deal with non-compact K3 manifolds, especially with at most one point-like singularity of $`A_N`$ type as a total space. We first explain the corresponding IIB background, and then, by using the conjectured duality between F/(K3$`\times S^1`$) and M/K3, or IIB/$`S^1`$ and M/$`T^2`$, we construct the corresponding background in terms of M-theory.
It is important that the complex modulus $`\tau `$ of the fiber torus of the non-compact K3 is only determined up to SL$`(2,𝐙)`$ as a function on the base manifold which is isomorphic to some orbifold of $`𝐂`$. This SL$`(2,𝐙)`$ is the IIB S-duality group and parametrized as $`\tau =\chi +ie^\varphi `$ where $`\chi `$ is axion and $`\varphi `$ is dilaton field. If we try to assign some definite value of $`\tau `$ at each point in the base manifold, there should exist branch cut coming up from each 7-brane to infinity . The behavior of the $`(p,q)`$-string, or some extended configuration including three string junctions, in the 7-brane background has been investigated in various situations. Here a $`(p,q)`$-string is a bound state of $`p`$ fundamental strings and $`q`$ D-strings and can end on a $`[p,q]`$7-brane.
Here we focus on the role of branch cut and take some duality transformation to obtain the M-theory counterpart. Moreover, since we want M5-branes in the dual 11-dimensional theory, we apply the argument to the case of $`(p,q)`$5-branes of world-volume $`R^{1,4}\times l`$ where $`l`$ is a one dimensional object in the base space. Here a bound state of $`p`$ D5-branes and $`q`$ NS5-branes is denoted by a $`(p,q)`$5-brane. If a $`(p,q)`$5-brane goes around $`N`$ D7-branes counterclockwise, the brane is converted into the $`(p+Nq,q)`$5-brane with the shift $`\tau \tau +N`$ as in Fig.2. This means that if the 5-brane crosses the branch cut singularity, then the charge assignment of the brane changes as in Fig.2, and if it crosses $`N`$ 7-branes, $`Nq`$ new D5-branes are created from the 7-branes to form a three 5-brane junction . This is regarded as U-dual of original Hanany-Witten effect where a D3-brane is created if D5-brane crosses an NS5-brane .
Using the duality conjecture between F-theory and M-theory, we can obtain the M-theory background from this 7-brane background. Consequently, the fiber torus turns up as a geometrical object in 11-dimensions, i.e., the background becomes elliptic (non-compact) K3 manifold times 7-dimensional Minkowski space. There is another method to reach the M-theory background; beginning with the explicit background of 7-branes, we perform T-duality along the appropriate world component of 7-branes as well as $`(p,q)`$5-branes, and then go up to 11-dimensions . The $`(p,q)`$5-brane turns into a bound state of $`p`$ D4-branes and $`q`$ NS5-branes in IIA theory, and then becomes an M5-brane winding around $`(p,q)`$-cycle of the torus. In particular, if we begin with $`N`$ D7-branes, we obtain the IIA background with $`N`$ D6-branes and thus we see that this background resembles the compactification of KK-monopole background in the sense that both deal with D6-branes. However, unlike the KK-monopole case, the number $`N`$ must be less than 24 for geometrical reason.
Since the supergravity solutions including both 7-branes and $`(p,q)`$5-branes are not known, we study the configuration by putting 5-branes as probes in the background of 7-branes. Since this configuration reserves some supersymmetry, the M5-brane transformed from IIB $`(p,q)`$5-brane must be holomorphically embedded in the corresponding M-theory background . We use the same embedding as in the case of M2-branes given in ref..
For the background with $`N`$ D7-branes at the origin of the base space spanned by $`(z,\overline{z})`$, the metric in terms of dual 11-dimensional theory is
$$ds_{scs}^2=\eta _{\mu \nu }dx^\mu dx^\nu +e^{\mathrm{\Phi }(z,\overline{z})}dzd\overline{z}+\tau _2^1d\zeta d\overline{\zeta }$$
(56)
where
$$e^{\mathrm{\Phi }(z,\overline{z})}=\tau _2\eta ^2\overline{\eta }^2|z|^{\frac{N}{6}},$$
(57)
$`\tau =\tau _1+i\tau _2`$ and $`\zeta =\stackrel{~}{u}+\tau \stackrel{~}{v}`$. The torus is spanned by the periodic coordinates $`(\stackrel{~}{u},\stackrel{~}{v})(\stackrel{~}{u}+2\pi R,\stackrel{~}{v}+2\pi R)`$. And the curve wound along $`(p,q)`$-cycle of the torus is represented as $`q\stackrel{~}{u}p\stackrel{~}{v}=const.`$. The modulus $`\tau `$ is some holomorphic function of $`z`$, and the behavior in the $`z0`$ limit is
$$\tau (z)\frac{N}{2\pi i}\mathrm{log}z,$$
(58)
which requires a branch cut from the origin. As we mentioned earlier, this branch cut is coordinate singularity in 11-dimensions. Other notations are the same as ref.. By compactifying this metric along $`\stackrel{~}{u}`$ direction, we obtain the IIA background of D6-branes with
$`e^{\frac{4}{3}\mathrm{\Phi }}`$ $`=`$ $`\tau _2^1,`$ (59)
$`A_{[1]}`$ $`=`$ $`\tau _1d\stackrel{~}{v}\stackrel{z0}{}{\displaystyle \frac{N}{2\pi }}\mathrm{arg}zd\stackrel{~}{v}.`$ (60)
We see that the coordinate singularity is now interpreted as the Dirac string singularity with respect to one-form potential $`A_{[1]}`$ in IIA theory. In the three-dimensional space spanned by $`(z,\overline{z},\stackrel{~}{v})`$, the singularity is two-dimensional, which is different from the case of KK-monopole solution.
Now in this background we put an M5-brane of world-volume $`R^{1,3}\times \mathrm{\Sigma }`$ where $`\mathrm{\Sigma }`$ is a holomorphic curve winding around the $`(p,q)`$-cycle of the torus. We choose the complex structure such that it is orthogonal to that defined by $`(z,\zeta )`$ . If we use the discussion in the end of section 5, the Hodge dual of the current associated with this M5-brane is given as
$$\stackrel{~}{}J_{[5]}^{(11)}\delta (y_4)\delta (y_5)\delta (y_6)\delta (q\stackrel{~}{u}p\stackrel{~}{v})\delta (g)dy_4dy_5dy_6d(q\stackrel{~}{u}p\stackrel{~}{v})dg.$$
(61)
Here $`g=g(z,\overline{z})`$ is defined as a real one-dimensional line in the $`z`$-plane representing the location of branes in the plane . If the brane is located over the both sides of the singularity as in Fig.2(a), coordinates $`(\stackrel{~}{u},\stackrel{~}{v})`$ are discontinuous at the singularity, and the definition of $`(p,q)`$-cycle changes at the singularity: If we cross the singularity counterclockwise, $`(p,q)`$-cycle changes to $`(p+Nq,q)`$-cycle.
In the flat spacetime $`^{1,8}\times T^2`$, an M5-brane winding along $`(p,q)`$-cycle of the torus leads to the bound state of $`p`$ D4-branes and $`q`$ NS5-branes upon compactification along the $`p`$-cycle ($`\stackrel{~}{u}`$ direction). This formula is also applicable to our case and the current associated with D4-brane charge is represented by the NS5-brane current if $`q0`$ as
$$j_{[4]}=j_{[5]}\frac{p}{q}d\stackrel{~}{v}$$
(62)
where we used eq.(54).
If the M5-brane is put such that it crosses the singularity, the interpretation of D4-brane charge after compactification depends on the place of the brane, left or right of the singularity. Note that the current associated with NS5-brane is determined independently of the location of the singularity. In the case of the brane depicted in Fig.2(a), D4-brane current is given as
$$\{\begin{array}{cc}\hfill j_{[4]}^l& =j_{[5]}\frac{p}{q}d\stackrel{~}{v},\hfill \\ \hfill j_{[4]}^r& =j_{[5]}(\frac{p}{q}+N)d\stackrel{~}{v}\hfill \end{array}$$
(63)
where the superscript $`l`$ and $`r`$ denote left and right hand side of the singularity respectively as in Fig.2(a).
From this assignment of currents, it can be seen as if new D4-brane charges are created from the string-like singularity, though there is no discontinuity in the original 11-dimensional viewpoint. This is not so peculiar since if we perform T-duality transformation along $`q`$-cycle ($`\stackrel{~}{v}`$ direction), the determination of current is consistent with that of $`(p,q)`$5-branes in terms of the SL$`(2,𝐙)`$ symmetry of IIB background with 7-branes.
## 8 Definition of D4-brane charges : Identification of D4-branes
We know from the discussion of M5-branes in the stringy cosmic string background where the location of the Dirac string type singularity plays an important role in identifying the brane as an NS5-brane, a D4-brane, or a bound state of them in the compactified 10-dimensions. Now we return to the case of KK-monopole background and consider the identification of 10-dimensional brane dimensional reduced from an M5-brane.
We explicitly argue three types of M5-branes specified by the following curves in $`M_{TN}`$ among others :
$`(i)`$ $`f_1|w_0|e^{\frac{b}{NR}}=0[(A)]`$ (64)
$`(ii)`$ $`f_2\left|{\displaystyle \frac{v_0}{w_0}}\right|e^{\frac{b}{NR}}=0[(B)]`$ (65)
$`(iii)`$ $`f_3|w_1|e^{\frac{b}{NR}}=0[(A)^{}]`$ (66)
They are the typical curves concerning the location of the singularity relative to the curve seen in $`(r_1,r_2,r_3)`$ space (see Fig.1 or 3).
The curve $`(i)`$ does not intersect the singularity for any value of $`b`$, $`(ii)`$ intersects the singularity at the vertex of the curve for any $`b`$ and $`(iii)`$ intersects the singularity asymmetrically. Note that all other curves belonging to our family are classified into the same class as $`(iii)`$ in the sense that they intersect the singularity asymmetrically. Also notice that this class is further divided into two classes: One class \[$`(iii)`$A\] consists of curves which cross the singularity only one time for all values of $`b`$, and the other \[$`(iii)`$-2\] is the set of curves that do not cross the singularity if we take $`b`$ small enough. The curve $`(iii)`$ in (66) itself belongs to the former class $`(iii)`$-1. A curve in the class $`(iii)`$-2 can be realized by choosing a complex structure, e.g., given by the holomorphic coordinates $`(w_0+w_1,v_0+v_1)`$. In the limit $`b/NR\mathrm{}`$, the point of intersection becomes closer to the D6-branes for all curves. Note that the curves in the class $`(iii)`$-2 has another intersection point at large $`r`$ in the limit.
We analyze the configuration of D4-branes and NS5-branes by the currents $`j_{[4]}`$ and $`j_{[5]}`$ in a similar way as in the previous section. First we consider the 11-dimensional current $`J_{[5]}^{(11)}`$ for each of three curves. They are obtained by applying the general formula eq.(47) to the above cases :
$`(i)`$ $`\stackrel{~}{}J_{[5]}^{(11)}\delta (y_4)\delta (y_5)\delta (y_6)\delta (f_1(\stackrel{}{r}))dy_4dy_5dy_6df_1dx_{11},`$ (67)
$`(ii)`$ $`\stackrel{~}{}J_{[5]}^{(11)}\delta (y_4)\delta (y_5)\delta (y_6)\delta (f_2(\stackrel{}{r}))dy_4dy_5dy_6df_2(dx_{11}NRd\psi ),`$ (68)
$`(iii)`$ $`\stackrel{~}{}J_{[5]}^{(11)}\delta (y_4)\delta (y_5)\delta (y_6)\delta (f_2(\stackrel{}{r}))dy_4dy_5dy_6df_2(dx_{11}NRd\chi )`$ (69)
where $`f_i`$ is given in eqs.(64)$``$ (66) and $`\chi `$ is in eq.(10). Note that in the near core limit, we can show the definite form of the 11-dimensional current as in eq.(46). In the reduced 10-dimensional theory, the relation between D4-brane current $`j_{[4]}`$ and NS5-brane current $`j_{[5]}`$ for each brane is determined independent of the scale of $`r`$ as
$`(i)`$ $`j_{[4]}=0,`$ (70)
$`(ii)`$ $`j_{[4]}=j_{[5]}(NRd\psi ),`$ (71)
$`(iii)`$ $`j_{[4]}=j_{[5]}(NRd\chi ).`$ (72)
This means that there is no D4-brane current in the case $`(i)`$. In other cases, D4-branes can be present only on an NS5-brane. By investigating the form of $`j_{[4]}`$ on the NS5-brane, we see how D4-branes are stretched on the NS5-brane. If the NS5-brane intersects with the singularity only once, then it is interpreted that the smeared D4-branes come up from the Dirac singularity, i.e., the D4-branes are created from the singularity, and go to infinity. For the curves in the class $`(iii)`$-2, there are two points on the corresponding NS5-brane where the D4-branes can be created or absorbed. Note that since $`j_{[4]}`$ is proportional to the number $`N`$, the number of D4-branes smeared on the NS5-brane is also proportional to $`N`$. The results are depicted in Fig.3.
## 9 On the brane creation phenomena
Brane creation was first discussed in the flat IIB background as the Hanany-Witten effect representing a phenomenon that a new brane is created by crossing of certain two types of branes. This effect is confirmed by charge conservation, although the process of brane creation has not yet been clarified in the framework of supergravity theory. One reason for this is that it is difficult to describe the situation of a brane ending on another brane as a solution of supergravity. Now, our discussion given in the preceding sections enables us to study the process of brane creation in the vicinity of D6-branes in terms of the corresponding exact solution of supergravity.
Using the argument, we now discuss how the brane creation is explained in the near core region of KK-monopoles and also consider the extension to the $`r`$:finite region of $`M_{TN}`$. In particular, we again deal with three types of M5-branes $`(i)(iii)`$ and their compactification.
First, consider the curves $`(ii)`$ and $`(iii)`$. They intersect the Dirac string singularity if the parameter $`b`$ is taken to be large enough. As was indicated in the last section, it is interpreted in 10-dimensions that D4-branes come up from the point of intersection with the singularity and go to infinity or another intersection point along the NS5-brane. In the near core limit $`r/NR0`$, each curve is a paraboloid and the shape of the paraboloid changes with the parameter $`|c|`$ as in eq.(42). In the limit $`c/NR0`$ (i.e., $`b/NR\mathrm{}`$), the paraboloid becomes like a thin tube and degenerates to an object like a half-string. As $`c`$ approaches $`0`$, the NS5-brane bends so as to wrap the D6-branes and gives multipole moment if measured far away from the brane . Finally in the limit $`c0`$, the total NS5-brane charge vanishes. On the other hand, in this limit, all the D4-branes spread over the NS5-brane gather to form a bundle of coincident D4-branes coming up from D6-branes. This process of disappearing the NS5-brane and the assembling of D4-branes into one place may be regarded as the brane creation in the near core limit.
Extending this process of brane creation to the outer region, $`r`$:finite region, we see a transition of branes with respect to the value $`b`$. For the curves $`(ii)`$ and $`(iii)`$-1, the transition process is essentially the same as in the near core region, since these curves already intersect with the singularity in the $`b/NR\mathrm{}`$. A curve in the class $`(iii)`$-2 begins with a pure NS5-brane if $`b`$ is small enough. Then as $`b/NR`$ becomes larger, the NS5-brane begins to bend and at some value of $`b`$ the curve comes in contact with the singularity, and after that we have two points of intersection and the D4-brane charges turn out. Note that in the limit $`b/NR\mathrm{}`$, NS5-brane charge disappears and only D4-brane charge remains roughly in the region $`r`$ smaller than $`b`$ for these cases.
On the other hand, consider the curve $`(i)`$ for which there is no D4-brane charge on the NS5-brane. In the limit $`b/NR\mathrm{}`$, the brane is bent as the same way as other cases, however, we have to take another explanation as others. In particular, in the region near the D6-branes, this curve reduces to an object which has an appearance of a D4-brane ending on D6-branes, but has neither NS5-brane charges nor D4-brane charges. We have a problem how to interpret the object. The resolution of this can be done by noticing the singularity again: The singularity is placed as overlapped on the D4-brane-like object in the limit $`b/NR\mathrm{}`$. Thus the situation is exceptional and the identification of the brane-like object cannot be performed.
The analogous structure is found in the IIB stringy cosmic string background where some 5-branes ending on $`N`$ coincident D7-branes are exactly on the branch cut. In terms of IIA theory, the branes must be identified as D4-branes, while it cannot be seen from charge conservation. The identification can be done by shifting the branes away from the branch cut singularity as in Fig.5(b).
In our case of KK-monopole background, shifting the D4-brane away from the singularity corresponds to taking another curves holomorphic with respect to different complex structure, or shifting the compactification direction as in eq.(51). By using either method, the identification of the branes can be done. This resolves the puzzle we stated.
## 10 Summary and Discussion
We have proposed a method of identifying the branes in the IIA background especially in the presence of D6-branes obtained by compactification of 11-dimensional Kaluza-Klein monopole solution. It is essentially the same as the identification of branes in the IIB background with 7-branes. We have also discussed the brane creation from D6-branes. In particular, since we know exact supergravity solutions of M5-branes in the near core region of KK-monopoles in 11-dimensions, we have clarified the mechanism of brane creation in this region as a compactified 10-dimensional theory.
Now we explain the relation between the brane creation based on our supergravity argument and the original Hanany-Witten argument based on flat branes in the flat background. In the flat space argument, created branes between two branes have finite length. On the other hand, in our argument the created branes are interpreted as half-infinite branes: The D4-branes coming up from the D6-branes are not cut on the NS5-brane, but continue along the NS5-brane. The difference may be related to the fact that there exist Dirac string singularities in the background with D6-branes. If we could obtain supergravity solutions of the original Hanany-Witten configurations, we would be able to clarify whether a finite brane can exist.
We comment on the case of M2-branes instead of M5-branes. In this case, we have to take care of the definition of currents after compactification: In order to proceed with the argument on identification of F-strings or D2-branes as the same method as in the case of M5-branes, it seems to be necessary to use the definition of 10-dimensional field strength as $`\stackrel{~}{}F_{[4]}^{(11)}=\widehat{G}_{[3]}+\widehat{G}_{[4]}dx^{11}`$ instead of $`G_{[3]}`$ and $`G_{[4]}`$ in eq.(49). Although the two definitions coincide with each other if there are no Kaluza-Klein charges, in general $`G_{[n]}\widehat{G}_{[n]}`$ in the presence of KK-charges. Nevertheless we do not know direct reason to change the definition of 10-dimensional field strengths as above depending on the objects we deal with. Even if we decide to use such a definition, the problem still remains in the situations where both M5-branes and M2-branes exist simultaneously. (In such a situation, Chern-Simons term $`FFA`$ in 11-dimensional supergravity action may play an important role.) Note that a same kind of confusion in defining the field strengths in 10-dimensions was pointed out in the previous argument . We do not have definite explanation of this puzzle at the present point.
Moreover, note that our 10-dimensional configuration with Dirac string singularity is based on singular compactification from the 11-dimensional background. We do not know if such dimensional reduction is consistently described as a compactification of string theory.
### Acknowledgments
This work was supported in part by JSPS Research Fellowships for Young Scientists. |
warning/0003/nucl-th0003036.html | ar5iv | text | # Statistical analysis of light fragment production from medium energy proton-induced reactions
## I Introduction
Fragment and residual nuclei production has attracted many people’s interests, not only nuclear physicists but also astrophysicists and nuclear engineers. The production of intermediate mass fragments from high energy proton-nucleus reactions or nucleus-nucleus reactions has been a hot topic in nuclear physics for a decade . Astrophysics and cosmic ray physics have been interested in residual nuclei production in order to calculate the production of cosmogenic nuclides in extraterrestrial matter by solar and galactic cosmic rays. Recently, nuclear engineering has needed the particle production cross sections for the development of accelerator-based systems for transmutation of radioactive nuclear waste. From the radiation safety aspect, it has also become more important to estimate the amount of radioactivity produced from various targets, as new applications of high energy proton accelerators, such as spallation neutron sources and the production of beams of unstable nuclei are being developed.
In 1997, the Organization for Economic Cooperation and Development (OECD) Nuclear Energy Agency (NEA) conducted benchmark calculations on activation yields to determine the predictive power of current nuclear reaction models and codes . The results calculated using many different codes which are based on a combination of different models, such as the intranuclear cascade model (INC), the exciton model, the evaporation-fission model, the quantum molecular dynamics model and the statistical multifragmentation model, were compared with experimental data. It became clear that most of the computer codes did not reproduce light fragment production reactions, such as Fe(p,X)<sup>7</sup>Be, especially at low proton-incident energy. They considered that an adequate description of the Fermi break-up model and a fragmentation model was urgently needed.
The evaporation model has been very successful in describing residual nuclei production from hot nuclei. In many codes, not only those codes that were used in the benchmark calculation by OECD/NEA, but also the codes that are widely used for shielding calculations, such as the LAHET code , the evaporation model is used to describe the de-excitation of thermalized nuclei. Despite its success, the model has not been used to describe light fragment emission, except for a few studies concerning break-up of highly excited nuclei .
In this study, we propose a generalized evaporation model (GEVAP) for Monte Carlo simulation, based on the Weisskopf-Ewing model . Nucleons and helium nuclei are the dominant particles emitted from an excited nucleus. Therefore, only these particles are treated as ejectiles in the Dostrovsky’s evaporation models implemented in the LAHET code . On the other hand, some studies consider light nuclei heavier than $`\alpha `$ particles as ejectiles since there is no reason that those particles can not be emitted from excited nuclei via evaporation process. In our generalized evaporation model, 66 nuclides up to Mg are included as ejectiles, not only in their ground states but also in their excited states. Besides, we use the accurate level density function for the total decay width calculation instead of an approximate form of level density function which is used in the Dostrovsky’s evaporation models .
Light fragments produced from proton-induced reactions are analyzed by the combination of the INC model implemented in the LAHET code and the generalized evaporation model (GEVAP). In order to estimate light particle production from a nucleon-nucleus reaction by the generalized evaporation model, we have to assume the ensemble of hot nuclei which are produced after the initial non-equilibrium stage. Since the excitation energy, the mass, and the charge of the hot residual nuclei produced from high energy reaction are widely distributed, we can not use the simple assumption that a single excited nucleus represents the ensemble of hot thermalized nuclei. However, the INC model can provide an ensemble of residual nuclei with broad distribution in excitation energy, nuclear mass and charge. The LAHET code employs the Bertini intranuclear cascade model for a non-equilibrium stage of nuclear reaction, and the Fermi break-up model and the evaporation model proposed by Dostrovsky et al. for a thermalized stage. Mass $`A_i`$, charge $`Z_i`$, excitation energy $`E`$, recoil energy, and the direction of recoil motion are extracted from the INC calculation done by the LAHET code. Then the de-excitation process of the hot nucleus with these quantities are calculated by GEVAP, instead of by the Fermi break-up model and the evaporation model employed in the LAHET code. In the following, we call this calculation procedure ‘INC/GEVAP’.
We focus mainly on <sup>7</sup>Be produced by proton-induced reactions in the energy range from 10 MeV to 3 GeV, because <sup>7</sup>Be is the most intensively measured light fragment produced from various targets, and many experimental data are available for comparison. We compare the INC/GEVAP results with experimental data as well as the results calculated by using LAHET, to make the effect of using different de-excitation models clear.
## II The generalized evaporation model
Let us consider that a parent nucleus $`i`$ with an excited energy $`E`$\[MeV\], a mass number $`A_i`$, and a charge number $`Z_i`$ emits a particle $`j`$ in its ground state with $`A_j`$ and $`Z_j`$, and becomes a daughter nucleus $`d`$ with $`A_d`$ and $`Z_d`$. According to the Weisskopf’s formulation , the decay probability $`P_j`$ with total kinetic energy in the center-of-mass system between $`ϵ`$ and $`ϵ`$ \+ d$`ϵ`$ is expressed as
$$P_j(ϵ)dϵ=g_j\sigma _{inv}(ϵ)\frac{\rho _d(EQϵ)}{\rho _i(E)}ϵdϵ,$$
(1)
where $`\sigma _{inv}`$ is the cross section for the inverse reaction, $`\rho _i`$ and $`\rho _d`$ are level densities \[MeV<sup>-1</sup>\] of the parent and the daughter nucleus, respectively. With the spin $`S_j`$ and the mass $`m_j`$ of the emitted particle $`j`$, $`g_j`$ is expressed as $`g_j=(2S_j+1)m_j/\pi ^2\mathrm{}^2`$. In this study we use the Audi-Wapstra mass table to calculate the Q-values $`Q`$ for emission of particle $`j`$.
The cross section for the inverse reaction $`\sigma _{inv}`$ is expressed as
$$\sigma _{inv}(ϵ)=\{\begin{array}{cc}\sigma _gc_n\left(1+b/ϵ\right)\hfill & \mathrm{for}\mathrm{neutrons}\hfill \\ \sigma _gc_j\left(1V/ϵ\right)\hfill & \mathrm{for}\mathrm{charged}\mathrm{particles}\hfill \end{array}\sigma _g\alpha \left(1+\frac{\beta }{ϵ}\right),$$
(2)
where $`\sigma _g=\pi R_{b}^{}{}_{}{}^{2}`$ \[fm<sup>2</sup>\] is the geometric cross section, and $`V=Z_jZ_de^2/R_c`$ is the Coulomb barrier.
In this study, we use the parameter set determined by Dostrovsky et al. and Matsuse et al . Dostrovsky et al. determined $`c_n`$, $`c_j`$, $`b`$, $`R_b`$, and $`R_c`$ for n, p, d, t, <sup>3</sup>He, and $`\alpha `$ emission by fitting the expression to the theoretical calculation done by Shapiro and Blatt and Weisskopf , so that the effect of overlapping wave functions was taken into account. These parameters are used in the Dostrovsky’s evaporation model implemented in LAHET. Meanwhile, Matsuse et al. determined the critical distance ($`R_b`$ and $`R_c`$, with $`c_j=1`$) by fitting Eq. (2) to experimental fusion cross sections for heavy ion reactions. We use the Dostrovsky’s parameters for n, p, d, t, <sup>3</sup>He, and $`\alpha `$ emission and the Matsuse’s parameters for other particles. In the following we call these parameters “the precise parameter set”. Besides the calculation with the precise parameter set, we use the simple parameter set, given by $`c_n=c_j=1`$, $`b=0`$ and $`R_b=R_c=r_0(A_j^{1/3}+A_d^{1/3})`$ \[fm\] for the inverse cross section. In the calculation with the simple parameter set, values of $`r_0=1.2`$, $`1.5`$, and $`2.0`$ are tried to test the stability of our model.
The total decay width $`\mathrm{\Gamma }_j`$ can be calculated by integrating Eq. (1) with respect to the total kinetic energy $`ϵ`$ from the Coulomb barrier $`V`$ up to the maximum possible value $`(EQ)`$. By using Eq. (2) for $`\sigma _{inv}`$, the total decay width for the particle emission is expressed as
$$\mathrm{\Gamma }_j=\frac{g_j\sigma _g\alpha }{\rho _i(E)}_V^{EQ}ϵ\left(1+\frac{\beta }{ϵ}\right)\rho _d(EQϵ)𝑑ϵ.$$
(3)
According to the Fermi-gas model, the total level density $`\rho (E)`$ of a nucleus summed over all the possible states with the angular momenta is given by the expression
$$\rho (E)=\frac{\pi }{12}\frac{e^{2\sqrt{a(E\delta )}}}{a^{1/4}(E\delta )^{5/4}}\text{ for }EE_x,$$
(4)
where $`a=A_d/8`$ \[MeV<sup>-1</sup>\] is the level density parameter, and $`\delta `$\[MeV\] is the pairing energy of the daughter nucleus evaluated by Cook et al. . For those values not evaluated by Cook et al., $`\delta `$ obtained by Gilbert and Cameron are used. $`E_x`$ is determined by Gilbert and Cameron as $`E_x=U_x+\delta `$ where $`U_x=2.5+150/A_d`$. In the calculation with the precise parameter set, we use the Gilbert-Cameron-Cook-Ignatyuk (GCCI) level density parameter , in which the pairing corrections and the energy dependence of the level density parameter are taken into account, instead of the simple expression $`a=A_d/8`$. The GCCI level density parameter is employed in the LAHET code.
When $`E`$ is below $`E_x`$, instead of Eq. (4) the following formula gives a good fit to the experimental level densities :
$$\rho (E)=\frac{1}{T}e^{(EE_0)/T}\text{ for }E<E_x,$$
(5)
where $`T`$ is the nuclear temperature given by $`1/T=\sqrt{a/U_x}1.5/U_x`$. To connect Eq. (4) and Eq. (5) smoothly, $`E_0`$ is defined as $`E_0=E_xT(\mathrm{log}T0.25\mathrm{log}a1.25\mathrm{log}U_x+2\sqrt{aU_x})`$.
We use the expressions Eq. (4) and Eq. (5) to calculate the total decay width. The simple form $`\rho \mathrm{exp}(2\sqrt{a(E\delta )})`$, which is used in the Dostrovsky’s evaporation models , is a good approximation when the residual excitation energy is high, however, it is not applicable for residual nuclei with small mass and low excitation energy.
When $`EQV`$ is below $`E_x`$, Eq. (3) can be solved analytically, by substituting Eq. (5) into Eq. (3).
$$\mathrm{\Gamma }_j=\frac{\pi g_j\sigma _g\alpha }{12\rho _i(E)}\{I_1(t,t)+(\beta +V)I_0(t)\}\text{ for }EQV<E_x,$$
(6)
where $`I_0(t)`$ and $`I_1(t,t_x)`$ are expressed as:
$`I_0(t)`$ $`=`$ $`e^{E_0/T}(e^t1),`$ (7)
$`I_1(t,t_x)`$ $`=`$ $`e^{E_0/T}T\{(tt_x+1)e^{t_x}t1)\},`$ (8)
where $`t=(EQV)/T`$ and $`t_x=E_x/T`$. When $`EQV`$ is greater than $`E_x`$, the integral of Eq. (3) can not be solved analytically because of the denominator in Eq. (4). However, it is expressed approximately as
$$\mathrm{\Gamma }_j=\frac{\pi g_j\sigma _g\alpha }{12\rho _i(E)}\left[I_1(t,t_x)+I_3(s,s_x)e^s+(\beta +V)\left\{I_0(t_x)+I_2(s,s_x)e^s\right\}\right]\text{ for }EQVE_x.$$
(9)
where $`I_2(s,s_x)`$ and $`I_3(s,s_x)`$ are given by:
$`I_2(s,s_x)`$ $`=`$ $`2\sqrt{2}\left\{s^{3/2}+1.5s^{5/2}+3.75s^{7/2}(s_x^{3/2}+1.5s_x^{5/2}+3.75s_x^{7/2})e^{s_xs}\right\},`$
$`I_3(s,s_x)`$ $`=`$ $`(\sqrt{2}a)^1[2s^{1/2}+4s^{3/2}+13.5s^{5/2}+60.0s^{7/2}+325.125s^{9/2}\{(s^2s_x^2)s_x^{3/2}`$
$`+`$ $`(1.5s^2+0.5s_x^2)s_x^{5/2}+(3.75s^2+0.25s_x^2)s_x^{7/2}+(12.875s^2+0.625s_x^2)s_x^{9/2}`$
$`+`$ $`(59.0625s^2+0.9375s_x^2)s_x^{11/2}+(324.8s^2+3.28s_x^2)s_x^{13/2}\}e^{s_xs}],`$
with $`s=2\sqrt{a(EQV\delta )}`$ and $`s_x=2\sqrt{a(E_x\delta )}`$.
In the present Monte Carlo simulation, ejectile $`j`$ is selected according to the probability distribution calculated as $`p_j=\mathrm{\Gamma }_j/_j\mathrm{\Gamma }_j`$, where $`\mathrm{\Gamma }_j`$ is given by Eqs. (6) or (9). The total kinetic energy $`ϵ`$ of the emitted particle $`j`$ and the daughter nucleus is chosen according to the probability distribution given by Eq. (1). The angular distribution of the motion is randomly selected to be isotropic in the center-of-mass system. The excitation energy of the daughter nucleus $`E_d`$ is calculated as $`E_d=EQϵ`$.
In this study, we consider 66 nuclides as ejectiles, not only in their ground states but also in their excited states. It is important to include excited states in the particles emitted via the evaporation process, because it greatly enhances the yield of heavy particle emission . The selected ejectiles satisfy the following criteria: (1) isotopes with $`Z_j12`$; (2) naturally existing isotopes or isotopes near the stability line; (3) isotopes with half-life longer than 1 ms. The selected ejectiles are listed in Table I.
If the mean lifetime of a resonance is longer than the decay width of the resonance emission, such a resonance can survive during the evaporation process. The excited state is included if its half lifetime $`T_{1/2}`$ \[sec\] satisfies the following condition:
$$\frac{T_{1/2}}{\mathrm{ln}2}>\frac{\mathrm{}}{\mathrm{\Gamma }_j^{}},$$
(10)
where $`\mathrm{\Gamma }_j^{}`$ is the decay width of the resonance emission. $`\mathrm{\Gamma }_j^{}`$ can be calculated in the same manner as for a ground state particle emission. The Q-value for the resonance emission is expressed as $`Q^{}=Q+E_j^{}`$, where $`E_j^{}`$ is the excitation energy of the resonance. The spin state of the resonance $`S_j^{}`$ is used in the calculation of $`g_j`$, instead of the spin of the ground state $`S_j`$. We use the ground state mass $`m_j`$ for excited states because the difference between the masses is negligible.
Instead of treating a resonance as an independent particle, we simply enhance the decay width of the ground state particle emission. We redefine the decay width $`\mathrm{\Gamma }_j`$ as
$$\mathrm{\Gamma }_j=\mathrm{\Gamma }_j^0+\underset{n}{}\mathrm{\Gamma }_j^n,$$
(11)
where $`\mathrm{\Gamma }_j^0`$ is the decay width of the ground state particle $`j`$ emission, and $`\mathrm{\Gamma }_j^n`$ is that of the $`n`$th excited state of the particle $`j`$ emission which satisfies Eq. (10).
The total kinetic energy distribution of the excited particle emission is assumed to be the same as that of the ground state particle emission. $`S_j^{}`$, $`E_j^{}`$, and $`T_{1/2}`$ used in this study are extracted from the Evaluated Nuclear Structure Data File (ENSDF) database maintained by the National Nuclear Data Center.
## III Comparison with experimental data
The excitation functions of <sup>7</sup>Be produced by proton reactions on <sup>16</sup>O, <sup>27</sup>Al, <sup>nat</sup>Fe, and <sup>93</sup>Nb are shown in Fig. 1. The results calculated by INC/GEVAP with the precise parameter set, which consists of the parameters for inverse reactions determined by Dostrovsky et al. and Matsuse et al and the GCCI level density parameter, are shown by the solid lines. The estimates by INC/GEVAP with $`r_0=1.5`$ and those by LAHET are also shown in the figures as well as the experimental data collected in Ref.. The results by INC/GEVAP with $`r_0=1.2`$ and $`2.0`$ are represented only for the Nb target, because there is less difference in the results for other targets. For the O target, the differences in the estimates between by INC/GEVAP with $`r_0=1.5`$ and by that with $`r_0=1.2`$ or $`1.5`$ are 20 %, except at the threshold energy. For Al target, the differences are within a factor of two in the whole energy region. INC/GEVAP with $`r_0=2.0`$ produces almost the same cross sections for the Fe(p,X)<sup>7</sup>Be reaction as those with precise parameter set, and the differences are within 20 %. In the whole energy region, for all the targets, INC/GEVAP produces more <sup>7</sup>Be as $`r_0`$ increases. The estimates by INC/GEVAP for Al with $`r_0=1.5`$ give the best agreement with the experimental data, as seen the dashed line lying underneath the measurement points in the whole energy region in Fig. 1 (b). Whereas for Fe and Nb INC/GEVAP with the precise parameter set reproduce the excitation functions better than that with the simple parameter set, and the estimates agree with most of the experimental data within 50 %.
INC/GEVAP reproduces the excitation functions for all the targets, whereas LAHET fails to reproduce the shape of the excitation functions except for the O target. The shapes of the excitation functions estimated by INC/GEVAP do not change with the choice of the parameter sets. Since LAHET severely underestimates the <sup>7</sup>Be productions from Al below 300 MeV, Fe and Nb below 3 GeV, it is obvious that the Fermi break-up is not the dominant process for the <sup>7</sup>Be productions in these reactions.
The isotopic distributions of H, He, Li, and Be nuclei produced from 2100 MeV proton incident on <sup>16</sup>O and from 480 MeV proton incident on <sup>nat</sup>Ag are shown in Fig 2. The results estimated by INC/GEVAP with the precise parameter set (the open squares) are shown as well as the experimental data for the O target measured by Olsen et al. , and the data for the Ag target by Green et al. (the closed circles). INC/GEVAP reproduces the isotopic distributions for both these reactions, and the estimates agree with most of the measurements with 50 % accuracy.
## IV Summary and Conclusion
We have formulated a generalized evaporation model (GEVAP) based on the Weisskopf-Ewing model . The features of the model are: (1) the accurate level density function is used for deriving the decay width of particle emission; (2) sixty-six nuclides up to Mg, not only in their ground state but also in their excited states, are taken into account in this study.
The combination of the intranuclear cascade model (INC) and GEVAP successfully reproduces the excitation functions of <sup>7</sup>Be produced by protons incident on <sup>16</sup>O, <sup>27</sup>Al, <sup>nat</sup>Fe, and <sup>93</sup>Nb. The choice of the parameter set in GEVAP does not affect the resulting shapes of the excitation functions. INC/GEVAP also predicts the isotopic distributions of H, He, Li, and Be produced from O and Ag with 50 % accuracy.
From the results, it is concluded that the evaporation process is the main process via which particles lighter than or equal to Be are produced by protons incident on targets heavier than O. INC/GEVAP can predict the production cross sections of particles lighter or equal to Be between 50 % to a factor of two in accuracy depending on the choice of the parameter set used in GEVAP. In this study, the precise parameter set gives the best results for overall reactions, and the accuracy is 50 % on the average.
###### Acknowledgements.
The author would like to thank Dr. T. Numao for useful discussion and encouragements, and also Dr. L. Moritz, Dr. G. Greeniaus, Dr. P. Jackson and Dr. R. Korteling for valuable comments on this paper, and Dr. R. E. Prael for supplying the LAHET code. The author also appreciates the TRIUMF computer group for providing a computer for the calculations. |
warning/0003/math-ph0003013.html | ar5iv | text | # Coherent and Squeezed States in Shape Invariant Potentials Obtained from the Master Function
## 1 Introduction
Coherent states, known as the closest states to classical ones, play an important role in many different contexts of the theoretical and experimental physics, specially quantum optics and multiparticle dynamics. Schrodinger first discovered the coherent states of the harmonic oscillator potential in $`1926`$ and much work has been done since then on their properties and applications. The coherent states have also been found in systems with the Lie group symmety . Recently, coherent states have been found in special Hamiltonians . These coherent states are called minimum uncertainty coherent states(MUCS). In coherent states the standard deviation of $`x`$ and $`p`$ are equal and their product is minimum over these states. There are also quantum states where, though we have minimum uncertainty for the standard deviation of coordinate and momentum, they are not equal any more; these states are called squeezed states. These quantum states are as important as coherent ones and their generation play an important role in many different branch of physics and communication engineering. Nieto and his colleagues have developed an interesting algorithm for the generation of the coherent and squeezed states for special potential, where the product of the standard deviation generalized Harmonic phase variable $`X_c`$ and $`P_c`$ are minimum over these states. Here in this article following the algorithm of reference we obtain these coherent and squeezed states for all the shape invariant potential obtained from the master function of reference . The Hamiltonians of the reference are the special cases of the most general shape invariant Hamiltonian to be treated in this paper and the result thus obtained in this work, are in good agreement with those of thier reference in these few special cases.
This paper is organized as follows. In section II we explain very briefly the shape invariant potential obtained from the master functions. In section III first we show that the generalized Harmonic variable $`X_c`$ of reference is linear function of the $`x`$ coordinate of orthogonal function and $`P_c`$ the generalized Harmonic momentum is proportional to master function. In section IV following the reference , we obtain the raising and lowering operator of these potentials. In section V we obtain the most general minimum uncertainty states for the general potential obtained from the master function. For particular choice of even master and weight function together with the symmetric interval we can obtain even or odd minimum uncertainty coherent (squeezed)states or well known Schrodinger-Cat state. It is also shown that the ground state of the Hamiltonian is one of the MUCS of the system. Eigenstates of Annihilation operators namlely the generalization of annihilation operators coherent states(AOCS) is derived in section VI. Here in this section it is shown that in general the minimum uncertainty states are different from AOCS ones. Section VII devoted to investigate the time evolution of the Minimum Uncertainty states. Here in this section it is shown that, the time evolution of the generlized quantum phase coordinates is almost similar to the time evolution of the phase coordinates of the qunantum oscillator except for the appearance of the constant phase $`\omega _0`$ and also the Hamiltonian depenence of the frequency $`\omega _H`$. The paper ended with a brief conclusion.
## 2 The Shape InvariantPotentials Obtained From the Master function
According to references by introducing the master function $`A(x)`$ as a polynomial of at most second order one can define a non-negative weight function $`W(x)`$ in the interval $`[a,b]`$ such that the expression $`\frac{1}{W(x)}\frac{d}{dx}(A(x)W(x))`$ be a polynomial of at most first order and the function $`A(x)W(x)`$ to vanish at the ends of the interval. Now we can define second order differential operator $`L=\frac{1}{W(x)}\frac{d}{dx}A(x)W(x)\frac{d}{dx}`$ with the following properties:
1) $`L`$ is a self-adjoint linear operator.
2) $`L`$ transforms a given polynomial of order $`m`$ to another polynomial of order $`m`$ at most.
3) The expression $`\frac{1}{W(x)}(\frac{d}{dx})^^n(A^^n(x)W(x))`$ is a polynomial of order at most $`n`$, which is indeed Rodrigues formula for the classical orthogonal polynomials.
4) The polynomials
$$\varphi _n(x)=\frac{a_n}{W(x)}(\frac{d}{dx})^^n(A^^n(x)W(x))$$
are orthogonal with respect to the weight function $`W(x)`$ in the interval $`[a,b]`$ as defined above, and one can find $`a_n`$ simply by comparing the coefficient of highest power of $`\varphi _n(x)`$ with those of the traditionally defined special orthogonal polynomials.
5) The polynomials $`\varphi _n(x)`$ are eigenfunctions of operator $`L`$, and therefore satisfy the following second order linear differential equation
$$\frac{1}{W(x)}\frac{d}{dx}(A(x)W(x)\frac{d}{dx}\varphi _n(x))=\gamma _n\varphi _n(x).$$
(2-1)
In order the differential $`Eq.(21)`$ have polynomial solution of degree $`n`$, $`\gamma _n`$ must be given by
$$\gamma _n=n(\frac{(A(x)W(x))^{}}{W(x)})^{}\frac{n(n1)}{2}A^{\prime \prime }(x).$$
Thus, the general form of the differential equation is as follows
$$A(x)\varphi _n^{\prime \prime }(x)+\frac{(A(x)W(x))^{}}{W(x)}\varphi _n^{}(x)[n(\frac{(A(x)W(x))^{}}{W(x)})^{}+\frac{n(n1)}{2}A^{\prime \prime }(x)]\varphi _n(x)=0.$$
(2-2)
By differentiating the differential Eq.$`(22)`$ $`m`$ times and then multiplying it by $`(1)^^mA^{\frac{m}{2}}(x)`$ we get the following associated differential equation
$$A(x)\varphi _{_{n,m}}^{\prime \prime }(x)+\frac{(A(x)W(x))^{}}{W(x)}\varphi _{_{n,m}}^{}(x)+[\frac{1}{2}(n^2+nm^2)A^{\prime \prime }(x)+(mn)(\frac{A(x)W^{}(x)}{W(x)})^{}$$
$$\frac{m^2}{4}\frac{(A^{}(x))^^2}{A(x)}\frac{m}{2}\frac{A^{}(x)W^{}(x)}{W(x)}]\varphi _{_{n,m}}(x)=0$$
(2-3)
where
$$\varphi _{_{n,m}}(x)=(1)^^mA^{\frac{m}{2}}(x)(\frac{d}{dx})^^m\varphi __n(x).$$
Now, changing the variable $`\frac{dx}{d\xi }=\sqrt{A(x)}`$ in associated differential equations of Eq.$`(23)`$ and defining the new function $`\psi _n^m(\xi )=A^{\frac{1}{4}}(x)W^{\frac{1}{2}}(x)\varphi _{n,m}(x)`$ we obtain the Schrödinger equation
$$\frac{d^2}{d\xi ^2}\psi _n^m(\xi )+V_m(x(\xi ))\psi _n^m(\xi )=E(n,m)\psi _n^m(\xi ),m=0,1,2,\mathrm{},n,$$
(2-4)
where the most general shape invariant potential is
$$V_m(x)=W^2(x)+\frac{d}{d\xi }W(x)=\frac{1}{2}(\frac{A(x)W^{}(x)}{W(x)})^{}\frac{2m1}{4}A^{\prime \prime }(x)+\frac{1}{4A(x)}(\frac{A(x)W^{}(x)}{W(x)})^2+$$
$$\frac{m}{2}\frac{A^{}(x)W^{}(x)}{W(x)}+\frac{4m^21}{16}\frac{A^2(x)}{A(x)},$$
(2-5)
and prime stands for derivative with respect to $`x`$. The spectrum $`E(n,m)`$ is
$$E(n,m)=(nm+1)[(\frac{A(x)W^{}(x)}{W(x)})^{}+\frac{1}{2}(n+m)A^{\prime \prime }(x)].$$
(2-6)
## 3 Generalized Harmonic Phase space variables $`X_c`$ and $`P_c`$
Following the prescription of references in a one dimensional hamiltonian:
$$H=\frac{1}{2}(\frac{d\xi }{dt})^2+V_m(x(\xi )),$$
(3-1)
with $`V_m(x(\xi ))`$ given in Eq.$`(25)`$, the classical paths of constant energy around the minimum points of the potential form closed paths in the phase space $`(\xi ,P_\xi )`$. Therefore, there is an injective canonical map from the phase space $`(\xi ,P_\xi )`$ into the new phase space $`(X_c,P_c)`$ , such that the closed constant energy paths turn into elliptic constant energy paths. Hence, in the phase space $`\xi P_\xi `$ the time dependence of these closed paths can be written as
$$X_c=A(E)\mathrm{sin}\omega _c(E)t,$$
(3-2)
$$P_c=mA(E)\omega _c(E)\mathrm{cos}\omega _c(E)t.$$
From the constancy of the hamiltonian, Eq.$`(31)`$, along the paths we have
$$t+t_0=\frac{d\xi }{\sqrt{\frac{2}{m}(EV_m(x(\xi ))}}.$$
By changing the variable $`\frac{dx}{\sqrt{A(x)}}=d\xi `$, we get
$$t+t_0=\sqrt{\frac{m}{2}}\frac{dx}{\sqrt{A(x)(EV_m(x))}}.$$
(3-3)
Inserting the expression $`(25)`$ for $`V_m(x)`$ and considering the fact that $`A(x)`$ is at most second order and $`A(x)\frac{d\mathrm{log}W(x)}{dx}`$ is at most first order, one can show that the expression $`A(x)(EV_m(x))`$ is quadratic. Hence, integrating Eq.$`(35)`$ we obtain
$$x+\frac{\eta _2}{2\eta _1}=\sqrt{(\frac{\eta _2}{2\eta _1})^2\frac{\eta _3}{\eta _1}}\mathrm{sin}\sqrt{\frac{2\eta _1}{m}}(t+t_0)$$
(3-4)
with
$$\eta _1=\frac{1}{2}A^{\prime \prime }(E\gamma +\frac{2m1}{4}A^{\prime \prime })+\frac{12m}{4}A^{\prime \prime }(\frac{AW^{}}{W})^{}\frac{1}{4}(\frac{AW^{}}{W})^{}\frac{4m^21}{16}(\frac{AW^{}}{W})^{}$$
$$\eta _2=A^{}(0)(E\gamma +\frac{2m1}{4}A^{\prime \prime })+\frac{1}{2}A^{}(0)(\frac{AW^{}}{W})^{}\frac{1}{2}(\frac{AW^{}}{W})^{}(\frac{AW^{}}{W})(0)\frac{m}{2}(A^{\prime \prime }(0)+A^{}(0))\frac{AW^{}}{W})(0)$$
$$\eta _3=A(0)(E\gamma +\frac{2m1}{4}A^{\prime \prime })+\frac{1}{2}A(0)(\frac{AW^{}}{W})^{}\frac{1}{4}(\frac{AW^{}}{W}(0))^2\frac{m}{2}A^{}(0)(\frac{AW^{}}{W})(0)\frac{4m^21}{16}(A^{}(0))^2.$$
(3-5)
Comparing the relations $`(32)`$ and $`(36)`$ we have
$$X_c=x_0(x+\frac{\eta _2}{2\eta _1})$$
$$\omega _c=\sqrt{\frac{2\eta _1}{m}}$$
$$A(E)=x_0\sqrt{(\frac{\eta _2}{2\eta _1})^2\frac{\eta _3}{\eta _1}},$$
(3-6)
where $`x_0`$ and $`t_0`$ are arbitrary constants of integration. Also $`\gamma `$ is a constant which is added to the potential $`V_m(x(\xi ))`$ for convenience. Similarly, for the momentum $`P_c=m\frac{dX_c}{dt}`$ we have
$$P_c=mx_0\frac{dx}{dt}=x_0m\frac{d\xi }{dt}\frac{dx}{d\xi }=x_0mP_\xi \sqrt{A(x)}.$$
(3-7)
The quantum operators corresponding to $`X_c`$ and $`P_c`$, denoted by $`\widehat{X}`$ and $`\widehat{P}`$, are defined as following according to :
$$\widehat{X}=X_c(x)=x_0(x+\frac{\eta _2}{2\eta _1})$$
$$\widehat{P}=\frac{1}{2i}(X_c^{}p+pX_c^{}).$$
Making a change variable from $`\xi `$ to $`x`$, we get
$$\widehat{P}=\frac{p_0}{2i}(A(x)\frac{d}{dx}+\frac{d}{dx}A(x))$$
where $`x_0`$ and $`p_0`$ are arbitrary constants.
## 4 Raising And Lowering Operators
In the algorithm of generation of the minimum uncertainty states of the hamiltonian $`\frac{d^2}{d\xi ^2}+V_m(x(\xi ))`$, the raising and lowering operators of its discrete eigenstates $`\psi _n^m`$, play an important role. These operators are denoted by $`\stackrel{~}{B}_{n,m}`$ and $`\stackrel{~}{A}_{n,m}`$ respectively. To obtain them, following the refernce , first we factorize the differential equation $`(23)`$ in a shape invariant form as:
$$\{\begin{array}{cc}B(n,m)A(n,m)\varphi _{_{n,m}}(x)=E(n,m)\varphi _{_{n,m}}(x)& \\ A(n,m)B(n,m)\varphi _{_{n1,m}}(x)=E(n,m)\varphi _{_{n1,m}}(x).& \end{array}$$
(4-1)
From the equations $`(41)`$ it is straightforward to derive the following recursion relations
$$\{\begin{array}{cc}B(n,m)\varphi _{_{n1,m}}(x)=\mu _{m,n}\varphi _{_{n,m}}(x)& \\ A(n,m)\varphi _{_{n,m}}(x)=\frac{E(n,m)}{\mu _{m,n}}\varphi _{_{n1,m}}(x),& \end{array}$$
(4-2)
where $`E(n,m)`$, $`B(n,m)`$ and $`A(n,m)`$ are given by
$$E(n,m)=$$
$$\frac{[(\frac{A(x)W^{}(x)}{W(x)})^{}\frac{AW^{}}{W}(0)+n^^2A^{\prime \prime }(x)A^{}(0)+(2nm)(\frac{A(x)W^{}(x)}{W(x)})^{}A^{}(0)+mA^{\prime \prime }(x)\frac{AW^{}}{W}(0)]^^2}{4[(\frac{A(x)W^{}(x)}{W(x)})^{}+nA^{\prime \prime }(x)]^^2}$$
$$\frac{1}{4}(\frac{AW^{}}{W}(0))^^2(nm)(\frac{A(x)W^{}(x)}{W(x)})^{}A(0)\frac{1}{4}m^^2(A^{}(0))^^2\frac{1}{2}mA^{}(0)(\frac{AW^{}}{W})(0)$$
$$\frac{1}{2}(n^^2m^^2)A^{\prime \prime }(x)A(0)$$
$$B(n,m)=A(x)\frac{d}{dx}((\frac{A(x)W^{}(x)}{W(x)})^{}+\frac{1}{2}nA^{\prime \prime }(x))x$$
$$\frac{2(\frac{A(x)W^{}(x)}{W(x)})^{}\frac{AW^{}}{W}(0)+n^^2A^{\prime \prime }(x)A^{}(0)+(2nm)(\frac{A(x)W^{}(x)}{W(x)})^{}A^{}(0)+(m+n)A^{\prime \prime }(x)\frac{AW^{}}{W}(0)}{2[(\frac{A(x)W^{}(x)}{W(x)})^{}+nA^{\prime \prime }(x)]}$$
$`A(n,m)=A(x)\frac{d}{dx}\frac{1}{2}nA^{\prime \prime }(x)x`$
$$\frac{n^^2A^{\prime \prime }(x)A^{}(0)+(2nm)(\frac{A(x)W^{}(x)}{W(x)})^{}A^{}(0)+(mn)A^{\prime \prime }(x)\frac{AW^{}}{W}(0)}{2[(\frac{A(x)W^{}(x)}{W(x)})^{}+nA^{\prime \prime }(x)]}.$$
(4-3)
In order to evaluate $`\mu _{n,m}`$, it is sufficient to divide both sides of Eq.$`(42)`$ by $`(A(x))^{\frac{m}{2}}`$ and compare the coefficients of the highest degree terms of boths sides, resulting in
$$\mu _{n,m}=[\frac{1}{2}A^{\prime \prime }(x)(n1)((\frac{AW^{}}{W})^{}+\frac{1}{2}nA^{\prime \prime })].$$
(4-4)
From $`\psi _n^m(\xi )=A^{\frac{1}{4}}W^{\frac{1}{2}}\varphi _{n,m}`$, it follows that $`\widehat{A}_{n,m}`$ and $`\widehat{B}_{n,m}`$, that is
$$\{\begin{array}{cc}\stackrel{~}{A}_{n,m}=A^{\frac{1}{4}}W^{\frac{1}{2}}A_{n,m}A^{\frac{1}{4}}W^{\frac{1}{2}}& \\ \stackrel{~}{B}_{n,m}=A^{\frac{1}{4}}W^{\frac{1}{2}}B_{n,m}A^{\frac{1}{4}}W^{\frac{1}{2}}& \end{array}$$
(4-5)
are the required raising and lowering operators of the wave functions, that is, we have
$$\{\begin{array}{cc}\stackrel{~}{B}_{n,m}\psi _{n1}^m(\xi )=\mu _{n,m}\psi _n^m(\xi )& \\ \stackrel{~}{A}_{n,m}\psi _n^m(\xi )=\frac{E(n,m)}{\mu _{n,m}}\psi _{n1}^m(\xi ).& \end{array}$$
(4-6)
where $`\stackrel{~}{A}_{n,m}`$ and $`\stackrel{~}{B}_{n,m}`$ are
$$\stackrel{~}{A}=A(x)\frac{d}{dx}(\frac{1}{2}nA^{\prime \prime }(x)+\frac{1}{4}A^{\prime \prime }(x)+\frac{1}{2}(\frac{A(x)W^{}(x)}{W(x)})^{})x$$
$$\frac{1}{2}A^{}(0)\frac{A(x)W^{}(x)}{W(x)}(0)$$
$$\frac{n^2A^{\prime \prime }(x)A^{}(0)+(2nm)A1A^{}(0)+(mn)A^{\prime \prime }(x)C}{2(A1+nA^{\prime \prime }(x))}$$
$$\stackrel{~}{B}=A(x)\frac{d}{dx}(\frac{1}{2}nA^{\prime \prime }(x)+(\frac{AW^{}}{W(x)})^{}(\frac{1}{2}\frac{A(x)W^{}(x)}{W(x)})^{}\frac{1}{4}A^{\prime \prime }(x))x$$
$$+\frac{1}{2}A^{}(0)+\frac{A(x)W^{}(x)}{W(x)}(0)$$
$$\frac{2A1C+n^2A^{\prime \prime }(x)A^{}(0)+(2nm)A1A^{}(0)+(m+n)A^{\prime \prime }(x)C}{2(A1+nA^{\prime \prime }(x))}$$
(4-7)
with
$$A1=(\frac{AW^{}}{W(x)})^{}$$
$$C=\frac{AW^{}}{W}(0).$$
The raising and lowering operators of the one dimensional shape invariant potentials, obtained from the master function, and also all other necessary information in constructing their minimum uncertainty state are given in Table I. The Hermitian quantum operators associated with the generalized Harmonic quantum phase variables $`\widehat{X}`$ and $`\widehat{P}`$ can be expressed in terms of the operators $`\stackrel{~}{A},\stackrel{~}{B}`$ and their Hermitian conjugates as
$$\widehat{X}=x_0[\stackrel{~}{A}+\stackrel{~}{A}^{}+\stackrel{~}{B}+\stackrel{~}{B}^{}]$$
$$\widehat{P}=\frac{p_0}{2i}[\stackrel{~}{A}+\stackrel{~}{B}^{}(\stackrel{~}{A}^{}+\stackrel{~}{B})],$$
(4-8)
where $`\stackrel{~}{A}^{}`$ and $`\stackrel{~}{B}^{}`$, the Hermitian conjugate of the raising and lowering operators are
$$\stackrel{~}{A}^{}=(\frac{f_1f_3A^{\prime \prime }(x)}{f_1+f_3})\stackrel{~}{A}+(\frac{2f_1A^{\prime \prime }(x)}{f_1+f_3})\stackrel{~}{B}+$$
$$\frac{(f_2A^{}(0))(f_1+f_3)(f_1A^{\prime \prime }(x))(f_2+f_4)+\frac{1}{2}(f_2f_4)(f_1+f_3)\frac{1}{2}(f_1f_3)(f_2+f_4)}{f_1+f_3}$$
$$\stackrel{~}{B}^{}=(\frac{2f_3+A^{\prime \prime }(x)}{f_1+f_3})\stackrel{~}{A}+(\frac{f_3f_1+A^{\prime \prime }(x)}{f_1+f_3})\stackrel{~}{B}+$$
$$\frac{(f_4+A^{}(0))(f_1+f_3)(f_3+A^{\prime \prime }(x))(f_2+f_4)+\frac{1}{2}(f_1f_3)(f_2+f_4)\frac{1}{2}(f_2f_4)(f_1+f_3)}{f_1+f_3}$$
with $`f_1,f_2,f_3`$ and $`f_4`$:
$$f_1=\frac{1}{2}A^{\prime \prime }(x)\frac{1}{2}(\frac{A(x)W^{}(x)}{W(x)})^{}\frac{1}{4}A^{\prime \prime }(x),$$
$$f_2=\frac{1}{2}\frac{A(x)W^{}(x)}{W(x)}(0)\frac{1}{4}A^{}(0)\frac{n^2A^{\prime \prime }(x)A^{}(0)+(2nm)(\frac{A(x)W^{}(x)}{W(x)})^{}A^{}(0)+(mn)A^{\prime \prime }(x)C}{2((\frac{A(x)W^{}(x)}{W(x)})^{}+nA^{\prime \prime }(x))},$$
$$f_3=\frac{1}{2}\frac{A(x)W^{}(x)}{W(x)}(0)\frac{1}{2}A^{\prime \prime }(x)+\frac{1}{4}A^{\prime \prime }(x),$$
and
$$f_4=\frac{1}{2}\frac{A(x)W^{}(x)}{W(x)}(0)+\frac{1}{4}A^{}(0)$$
$$\frac{n^2A^{\prime \prime }(x)A^{}(0)+(2nm)(\frac{A(x)W^{}(x)}{W(x)})^{}A^{}(0)+(m+n)A^{\prime \prime }(x)C+2(\frac{A(x)W^{}(x)}{W(x)})^{}C}{2((\frac{A(x)W^{}(x)}{W(x)})^{}+nA^{\prime \prime }(x))}.$$
(4-9)
Generally speaking, the number $`n`$ should not appear anywhere and it must be replaced by the Hamiltonian. This is done by expressing $`n`$ in terms of $`E(n,m)`$ and replacing $`E(n,m)`$ by the Hamiltonian. Since the set of eigenfunctions $`\psi _n^m`$ are complete and we can expand every function in our Hilbert space in terms of them, therefore it is sufficient to consider the effect of the operators on these base. Therefore, in order not to make things too complicated, we do not bother to replace $`n`$ in terms of Hamiltonian, as in reference, except for the operators $`\widehat{X}`$ and $`\widehat{P}`$ which are going to be functions of the Hamiltonian.
## 5 The Most General Minimum Uncertainty States
The coherent and squeezed states are generally the minimum uncertainty states of general harmonic phase variables $`\widehat{X}`$ and $`\widehat{P}`$, which can be obtained by solving the eigenfunction equation of the operators $`\widehat{X}+i\frac{<G>}{2(\mathrm{\Delta }p)^2}\widehat{P}`$ , that is, we solve
$$(\widehat{X}+i\frac{<G>}{2(\delta p)^2}\widehat{P})\psi _{MUCS}=C\psi _{MUCS}$$
(5-1)
where $`G`$ is proportional to the commutation of $`\widehat{X}`$ and $`\widehat{P}`$ as $`[\widehat{X},\widehat{P}]=iG`$, and where $`C=<\widehat{X}>+i\frac{<G>}{2(\mathrm{\Delta }p)^2}<\widehat{P}>`$.
In order to solve the Eq.$`(51)`$, we expand $`\psi _{MUCS}`$ in terms of $`\psi _n^m`$, the eigenstate of Schrodinger equation $`(24)`$,
$$\psi _{MUCS}=\underset{j=m}{\overset{\mathrm{}}{}}a_j\psi _j^m(x).$$
(5-2)
Inserting the above expansion in Eq.$`(51)`$ and using the independence of the eigenstates $`\psi _j^m`$, we get a recursion relation between the coefficients $`a_j`$. This is possible if we know how the adjoint operators $`\stackrel{~}{A}^{}`$ and $`\stackrel{~}{B}^{}`$ act over $`\psi _j^m`$ by writing them in terms of the raising and lowering operators $`\stackrel{~}{A}`$ and $`\stackrel{~}{B}`$ as
$$\stackrel{~}{A}^{}=(\frac{f_1f_3A^{\prime \prime }(x)}{f_1+f_3})\stackrel{~}{A}+(\frac{2f_1A^{\prime \prime }(x)}{f_1+f_3})\stackrel{~}{B}+$$
$$\frac{(f_2A^{}(0))(f_1+f_3)(f_1A^{\prime \prime }(x))(f_2+f_4)+\frac{1}{2}(f_2f_4)(f_1+f_3)\frac{1}{2}(f_1f_3)(f_2+f_4)}{f_1+f_3}$$
$$\stackrel{~}{B}^{}=(\frac{2f_3+A^{\prime \prime }(x)}{f_1+f_3})\stackrel{~}{A}+(\frac{f_3f_1+A^{\prime \prime }(x)}{f_1+f_3})\stackrel{~}{B}+$$
$$\frac{(f_4+A^{}(0))(f_1+f_3)(f_3+A^{\prime \prime }(x))(f_2+f_4)+\frac{1}{2}(f_1f_3)(f_2+f_4)\frac{1}{2}(f_2f_4)(f_1+f_3)}{f_1+f_3}$$
Therefore, the most general quantum harmonic phase coordinate $`\widehat{X}`$ and momentum $`\widehat{P}`$, given in Eq.$`(51)`$, can be expressed in terms of the raising and lowering operators.
$$\widehat{X}=2x_0(\stackrel{~}{A}+\stackrel{~}{B})$$
$$\frac{2i}{p_0}\widehat{P}=\frac{4f_3+2A^{\prime \prime }(x)}{f_1+f_3}\stackrel{~}{A}+\frac{4f_1+2A^{\prime \prime }(x)}{f_1+f_3}\stackrel{~}{B}+$$
$$\frac{2(f_1+f_3)(f_4f_2)+2(f_2+f_4)(f_1f_3)+2A^{}(0)(f_1+f_3)2A^{\prime \prime }(x)(f_2+f_4)}{f_1+f_3}$$
Substituting the operators $`\widehat{X}`$ and $`\widehat{P}`$, written only in terms of the raising and lowering operators as above, in Eq.$`(51)`$ and using the expansion of $`\psi _{MUCS}(x)`$ in terms of $`\psi _n^m(x)`$, we obtain the following recursion relation for $`a_{n+m}`$:
$$a_{n+m+1}=\frac{1}{2x_0+\frac{p_0<G>}{4(\mathrm{\Delta }p)^2}\frac{4f_3+2A^{\prime \prime }(x)}{f_1+f_3}}\frac{\mu _{n+1,m}}{E(n+1,m)}\times $$
$$[(C\frac{p_0<G>}{4(\mathrm{\Delta }p)^2}g)a_{n+m}(2x_0\frac{p_0<G>}{4(\mathrm{\Delta }p)^2}\frac{4f_12A^{\prime \prime }(x)}{f_1+f_3})\mu _{n,m}a_{n+m1}].$$
(5-3)
where
$$g=\frac{2(f_1+f_3)(f_4f_2)+2(f_2+f_4)(f_1f_3)+2A^{}(0)(f_1+f_3)2A^{\prime \prime }(x)(f_2+f_4)}{f_1+f_3}.$$
In principle, by iterating the recursion relation (5-3) we can determine $`a_{n+m}`$ in terms of $`a_0`$ and $`a_{m+1}`$. In general, due to the appearance of $`a_{n+m}`$ and $`a_{n+m1}`$ on the right hand side of this recursion relation, we are not able to obtain a closed form for the coefficient $`a_{n+m}`$ in terms of $`a_m`$ and $`a_{m+1}`$. Note that, by an appropriate choice of the parameters appearing in MUCS, as e.g. taking such as the averge of the quantum phase coordinates, the average of the commutator of the quatum phase coordinates or the squeezing percentage that is the ratio standard deviations, $`(2x_0\frac{p_0<G>}{4(\mathrm{\Delta }p)^2}\frac{4f_12A^{\prime \prime }(x)}{f_1+f_3})`$ vanishes. We iterate
$$\frac{a_{n+m+1}}{a_{n+m}}=\frac{(C\frac{p_0<G>}{4(\mathrm{\Delta }p)^2}g)}{2x_0+\frac{p_0<G>}{4(\mathrm{\Delta }p)^2}\frac{4f_3+2A^{\prime \prime }(x)}{f_1+f_3}}\frac{\mu _{n+1,m}}{E(n+1,m)}$$
to obtain $`a_{n+m}`$ only in terms of $`a_m`$:
$$\frac{a_{n+m}}{a_m}=k_0^n\underset{j=m}{\overset{n+m}{}}\frac{\mu _j}{E(j)}$$
(5-4)
with
$$k_0=\frac{(C\frac{p_0<G>}{4(\mathrm{\Delta }p)^2}g)}{2x_0+\frac{p_0<G>}{4(\mathrm{\Delta }p)^2}\frac{4f_3+2A^{\prime \prime }(x)}{f_1+f_3}}.$$
Substituting the coefficient $`a_{n+m}`$, given in Eq.$`(54)`$, into the expansion $`\psi _{MUCS}`$, we get
$$\psi _{MUCS}=\underset{n=m}{\overset{\mathrm{}}{}}a_mk_0^n\underset{j=m}{\overset{n+m}{}}\frac{\mu _j}{E(j)}\psi _n^m.$$
(5-5)
In the rest of this section we investigate some important special cases:
a- For $`A(x)=1`$ we get the harmonic oscillator with eigenfunction
$$\psi _n(x)=e^{\frac{\alpha x^2}{4}}H_n(x)$$
and the coherent states yield
$$\psi _{MUCS}(x)=e^{txt^2/2}$$
, where $`t=\frac{k_0}{\alpha }`$.
b- $`A(x)=x,\beta =1,m=0,and\alpha =\lambda +\frac{1}{2}`$ lead to
$$\psi _n(x)=x^{\frac{\alpha +1/2}{2}}e^{x/2}n!L_n^\alpha (x)$$
and the coherent states yield
$$\psi _{MUCS}(x)=x^{\frac{\alpha +1/2}{2}}e^{x/2}I_{\lambda +1/2}(2\sqrt{k_0x}),$$
in agreement with reference.
c- $`A(x)=x(1x),\alpha =\beta =0,m=1/2\lambda ,andn=n+\lambda 1/2`$ lead to
$$\psi _n(x)=\sqrt{\mathrm{cos}(t)}P_{n+\lambda 1/2}^{1/2\lambda }(\mathrm{sin}(t))$$
and the coherent states yield
$$\psi _{MUCS}(x)=$$
in agreement with reference.
d- $`A(x)=1+x^2,\alpha =1/2,and\beta =0,`$ lead to
$$\psi _n(x)=P_{m1}^n(\mathrm{tanh}(t))$$
and the coherent states yield
$$\psi _{MUCS}(x)=$$
in agreement with reference.
e- $`A(x)=x^2,\beta =1,\lambda =n\alpha /2`$ lead
$$\psi _n(x)=x^{\frac{\alpha +1}{2}}e^{x/2}L_n^{\alpha 1}(x)$$
and the coherent states yield
$$\psi _{MUCS}(x)=$$
in agreement with reference. If the coeffient of $`a_n`$ in the recursion relation (5-3) vanishes, that is for:
$$C=\frac{p_0<G>}{4(\mathrm{\Delta }p)^2}g$$
we encounter another interesting special case. This is again possible by an appropriate choice of the parameters of MUCS. In this case the MUCS can consist only of the odd or even associative orthogonal polynomials. In particular, if we choose even master function together with the corresponding even weight function, then the associative orthogonal function will be even for even $`n+m`$ and will be odd for odd $`n+m`$, provided that we choose a symmetric interval,that is $`[a,a]`$. Therefore, MUCS is either even or odd known as even or odd coherent (squeezed) states in the special case of the harmonic Hamiltonian (Schrodinger-Cat states) . Finally we can choose the parameters of MUCS such that the coefficients of $`a_{n+m}`$ and $`a_{n+m1}`$ on the right hand side of the recursion (5-3) become null. This means that in the expansion of MUCS, given in (5-2), only the ground state of the Hamiltonian will remain, that is, the ground state of the system belongs to its MUCS ensemble. At the end of this section it should be remined that, eventhough the closed form of MUCS is only available in seldom special case. But this is not a serious prblem, since one can iterate the recursion relation (5-3) numerically, to calculate the MUCS and related average quantities over these states which is under separate investigation.
## 6 Eigenstates of Annihilation Operators
It is well-known that the coherent states of harmonic oscillators are eigenstates of the annihilation operator $`a`$, that is, we have
$$a\alpha >=\alpha \alpha >.$$
Here we try to generalize this idea to the general shape invariant potentials ($`V_m(x(\xi ))`$). This is possible if we can find the eigenstate of annihilation operator ($`A\alpha >=\alpha \alpha >`$) associated with these potentials. Here in this section, by multiplying the annihilation operator by an appropriate function of Hamiltonian and using the result of reference , we get a closed form for their eigenstates. Denoting the eigenstates of the annihilation operator $`A`$ by $`\psi _{AOCS}(\beta _0,x)`$, we have the following eigen-equation
$$F(g^1(H))A\psi _{AOCS}(\beta _0,x)=\beta _0\psi _{AOCS}(\beta _0,x),$$
where $`F(g^1(H))`$ is an arbitrary function of Hamiltonian which is to be determined below. Expanding $`\psi _{AOCS}(\beta _0,x)`$ in terms of the associated orthogonal functions $`\varphi _{n,m}(x)`$ and using the relation
$$F(g^1(H))A\underset{n=0}{\overset{\mathrm{}}{}}C_n\varphi _{n,m}=\beta _0\underset{n=0}{\overset{\mathrm{}}{}}C_n\varphi _{n,m},$$
we get the following recursion relation
$$C_{n+1}\frac{F(n+1)E(n+1,m)}{\mu _{n+1,m}}=\beta _0C_n.$$
(6-1)
Substituting for $`\mu _{n+1,m}`$ we get
$$\frac{C_{n+1}}{C_n}=\beta _0\frac{\frac{1}{2}nA^{\prime \prime }(x)B_0}{E(n+1,m)F(n+1)}\frac{a_n}{a_{n+1}}$$
with
$$B_0=\frac{1}{2}(n+1)A^{\prime \prime }(x)+(\frac{A(x)W^{}(x)}{W(x)})^{}.$$
By defining $`F(g^1(H))`$ as
$$F(n+1)=\frac{(n+1)(\frac{1}{2}nA^{\prime \prime }(x)B_0)}{E(n+1,m)}$$
we can solve the recursive relation $`(61)`$ if we choose $`C_n=\frac{\lambda ^n}{n!a_n}`$ and $`\lambda =\beta _0`$. Substituting these in the expansion of $`\psi _{AOCS}(\beta _0,x)`$ we get
$$\psi _{AOCS}(\beta _0,x)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{\beta _0^n}{n!a_n}\varphi _{n,m}.$$
Now, using the result of reference we obtain
$$\psi _{AOCS}(\beta _0,x)=(1)^m(A(x))^{\frac{m}{2}}(\frac{d}{dx})^m(\frac{W(z)}{W(x)}\frac{dz}{dx})$$
with $`z=x+tA(z)`$. From the result thus obtained it is clear that in general, the $`\psi _{AOCS}`$ states are different from the corresponding $`\psi _{MUCS}`$ states and they only coincide in the case of harmonioc oscillator. To see their difference more explicitly, using the result of reference, we give in the rest of this section some of the $`\psi _{AOCS}`$ states.
I-The choice of $`A(x)=x,\beta =1,m=0,\alpha =\lambda +\frac{1}{2}`$ leeds
$$\psi _n(x)=x^{\frac{\alpha +1/2}{2}}e^{x/2}n!L_n^\alpha (x)$$
and the using coherent states
$$\psi _{MUCS}(x)=x^{\frac{\alpha +1/2}{2}}e^{x/2}I_{\lambda +1/2}(2\sqrt{k_0x})$$
which is in agreement with reference.
II-The choice of $`A(x)=x(1x),\alpha =\beta =0,m=1/2\lambda ,n=n+\lambda 1/2`$ leeds
$$\psi _n(x)=\sqrt{\mathrm{cos}(t)}P_{n+\lambda 1/2}^{1/2\lambda }(\mathrm{sin}(t))$$
and the using coherent states
$$\psi _{MUCS}(x)=$$
which is in agreement with reference.
III-The choice of $`A(x)=1+x^2,\alpha =1/2,\beta =0,`$ leeds
$$\psi _n(x)=P_{m1}^n(\mathrm{tanh}(t))$$
and the using coherent states
$$\psi _{MUCS}(x)=$$
which is in agreement with reference.
IV-The choice of $`A(x)=x^2,\beta =1,\lambda =n\alpha /2`$ leeds
$$\psi _n(x)=x^{\frac{\alpha +1}{2}}e^{x/2}L_n^{\alpha 1}(x)$$
and the using coherent states
$$\psi _{MUCS}(x)=$$
which is in agreement with reference.
## 7 Time Evolution of the Minimum Uncertainty States
In this section we investigate the time evolution of the generalized Harmonic quantum phase variable $`\widehat{X}`$ and $`\widehat{P}`$. Since they do not have any explicit time dependence, thus their time dependence can be written as
$$\widehat{X}(t)=e^{\frac{iHt}{\mathrm{}}}\widehat{X}e^{\frac{iHt}{\mathrm{}}},$$
(7-1a)
$$\widehat{P}(t)=e^{\frac{iHt}{\mathrm{}}}\widehat{P}e^{\frac{iHt}{\mathrm{}}},$$
(7-1b)
which follows from the Heisenberg equations of motion
$$\dot{\widehat{X}}=\frac{1}{i\mathrm{}}[\widehat{X},H]=\frac{\widehat{P}}{m}$$
$$\dot{\widehat{P}}=\frac{1}{i\mathrm{}}[\widehat{P},H]=\widehat{X}B_1(H)+i\widehat{P}B_0$$
with $`B_1(H)`$ and $`B_0`$ defined as:
$$B_0=\frac{1}{2m}\mathrm{}A^{\prime \prime }(x)$$
$$B_1=A^{\prime \prime }(H\gamma +\frac{2m1}{4}A^{\prime \prime })+\frac{12m}{2}A^{\prime \prime }(\frac{AW^{}}{W})^{}\frac{1}{2}((\frac{AW^{}}{W})^{})^2\frac{4m^21}{8}A^{\prime \prime }.$$
Now, using the Baker-Hausdorf formula in Eq.$`(71a)`$, we can calulate $`\widehat{X}(t)`$ :
$$\widehat{X}(t)=\underset{n=0}{\overset{\mathrm{}}{}}(\frac{it}{\mathrm{}})^n\frac{1}{n!}X_n$$
with
$$X_n=\widehat{X}f_n(H)+\widehat{P}g_n(H)$$
(7-2)
By iterating the following recursion relations
$$f_{n+1}(H)=\frac{\mathrm{}}{i}B_1(H)g_n(H)$$
(7-3)
$$g_{n+1}(H)=\frac{\mathrm{}}{i}(\frac{1}{m}f_n(H)+(iB_0)g_n(H))$$
(7-4)
with
we get the following expression for the function $`g_n(H)`$
$$g_n(H)=Ar_+^n+Br_{}^n$$
where $`r_{},r_+,A`$ and $`B`$ are
$$r_+=\frac{1}{2}\mathrm{}B_0+\frac{1}{2}\mathrm{}\sqrt{B_0^24B_1(H)/m}=\mathrm{}\omega _0+\mathrm{}\omega _H$$
$$r_{}=\frac{1}{2}\mathrm{}B_0\frac{1}{2}\mathrm{}\sqrt{B_0^24B_1(H)/m}=\mathrm{}\omega _0\mathrm{}\omega _H$$
$$A=B=\frac{i\mathrm{}}{2m\omega _H}$$
with $`\omega _0`$ and $`\omega _H`$ defined as
$$\omega _0=B_0/2\mathrm{a}\mathrm{n}\mathrm{d}\omega _H=\sqrt{B_0^24B_1(H)/m}/2,$$
respectively. Substituting the result thus obtained for $`g_n(H)`$ in (7-4a) we can determine $`f_n(H)`$.
Now, inserting all these in (7-3), we obtain the closed form for $`\widehat{X}(t)`$ and $`\widehat{P}(t)`$ as follows:
$$\widehat{X}(t)=\widehat{X}e^{i\omega _0t}[\mathrm{cos}(\omega _Ht)i\frac{\omega _0}{\omega _H}\mathrm{sin}(\omega _Ht)]$$
$$+\widehat{P}e^{i\omega _0t}2\frac{\omega _0}{\omega _H}\mathrm{sin}(\omega _Ht)$$
(7-5)
Starting from the formula (7-1b) and performing some calculation which is similar to the above calculation, we will obtain
$$\widehat{P}(t)=\widehat{P}e^{i\omega _0t}[\mathrm{cos}(\omega _Ht)+i\frac{\omega _0}{\omega _H}\mathrm{sin}(\omega _Ht)]$$
$$+\widehat{X}e^{i\omega _0t}(2\frac{\omega _0}{\omega _H})\mathrm{sin}(\omega _Ht)$$
(7-6)
Again the result thus obatained are in agreement with refernces for the special caes provided that we insert the corresponding eigen value of the operators of eigen frequency $`\omega _H`$, given in Table I. The time dependence of the generlized quantum phase coordinates given in (7-5) and (7-6) is almost similar to the time evolution of the phase coordinates of the qunantum oscillator except for the appearance of the $`\omega _0`$ and also the energy depenence of the $`\omega _H`$. We see that in case of the general shape invarint hamiltonian the frequency is a hamiltonian dependent operator and it is only constant in special case of oscillator.
## 8 CONCLUSION
In this paper a general algorithm has been given for the generation of the minimum uncertainty coherent and squeezed states in some one-dimensional hamiltonians with shape invariant potential, obtained from the master function. It looks like that the shape invariance symmetry of these hamiltonian might be the reason for the observation the MUCS. Since Solvability of these quantum systems are mainly due to the existance of this symmetry . But this not the only reason , Actually the main role belongs the existence of the lowering and raising operators or ladder ones , which map different energy eigenstates of a given hamiltonian into each other. As quoted in the introduction, the coherent and squeezed states generated by harmonic oscillator have already play such an importantat role in different branches of physics. Definetly the MUCS have been generated in refrence and here will soon play very important role in almost all branches of physics. Therefore, it deserve to find all other hamiltonian which can generate MUCS, which is under investigation.
ACKNOWLEDGEMENT
We wish to thank Dr. S. K. A. Seyed Yagoobi for his careful reading the article and for his constructive comments. |
warning/0003/astro-ph0003111.html | ar5iv | text | # Narrow Line Seyfert 1 Galaxies and the Evolution of Galaxies & Active Galaxies
## 1 Introduction
Years after their discovery (Osterbrock & Pogge 1985) the NLS1s have attracted attention of the AGN community at least partly due to their peculiar X-ray properties. The NLS1s are Seyfert 1 galaxies with relatively narrow widths of permitted optical emission lines (full width at half maximum (FWHM)$`\genfrac{}{}{0pt}{}{_<}{^{}}`$2000 km/s), strong optical FeII/H$`\beta `$ ratio, weak \[OIII\] emission, and as such were found to occupy one extreme end of the Boroson & Green (1992) “eigenvector 1”. An extremely strong anti-correlation was found between soft X-ray spectral slopes and H$`\beta `$ FWHM in Seyfert 1s (Boller, Brandt & Fink 1996) and quasars (Laor et al. 1997) meaning that a relation between the lines and continuum exists. Eigenvector 1 was later found (Brandt & Boller 1998) to correlate strongly with the soft X-ray power-law slope. Since the soft X-rays are formed in the vicinity of the central black hole, eigenvector 1 has probably a more fundamental physical meaning.
NLS1s, as a class, show peculiar continuum properties as well. They have very steep soft X-ray slopes ($`<\alpha >=2.13,F_\nu \nu ^\alpha `$, while for “normal” Sy1s $`<\alpha >=1.34\pm 0.03`$) and sometimes show rapid large amplitude variability. The hard X-ray spectra are steep as well (Brandt, Mathur & Elvis 1997). In the optical and UV range, most of NLS1s show a weak “big blue bump” (BBB), which is most likely due to the shift of the BBB (sometimes out of the optical/UV range) towards higher energies. The high energy tail of the BBB is apparent as the unusually strong and steep soft X-ray excess. NLS1s also show strong IR emission and some high polarization.
Many continuum properties of the NLS1s can be explained in terms of high accretion rate compared to the Eddington limit (ṁ = Ṁ/Ṁ<sub>Edd</sub>) and so, a small black hole (BH) mass for a given luminosity. The high accretion rate explanation for the X-ray properties of NLS1s was first proposed by Pounds, Done & Osborne (1995), in analogy with Galactic BH candidates whose soft X-ray spectra become steep in their high state. The narrow widths of the emission lines can be explained if the Broad Line Region (BLR) scales as L<sup>1/2</sup> and emission line clouds are virialized around the small mass BH (Laor et al. 1997). As an alternative, Wandel (1997) has argued that the continuum, with steeper X-ray slope, has stronger ionizing power, and hence the BLR is formed at a larger distance from the center. The resulting smaller velocity dispersion produces narrower lines. In general there is a reasonable consensus that large ṁ is the cause of the observed peculiar properties of NLS1 (an alternative being a pole-on view). A natural question to ask as a next step would be “what determines the accretion rate in an active galaxy?” Is it the age?
## 2 Are NLS1s the Active Galaxies in the making?
Here we present a number of arguments in support of our proposal that NLS1 might be Active Galaxies in early phase of their evolution.
1). Smaller BH mass. As per the well known correlation of Magorrian et al. (1998) smaller mass BHs reside in galaxies with smaller spheroids. Since NLS1s have relatively smaller mass BHs compared to normal Seyferts, the spheroids of their host galaxies might be smaller (see also Laor, 1998). Indeed, in the compilation of Wandel (1999), the NLS1 galaxy NGC4051 has the smallest black hole to bulge mass ratio. An accreting BH would also grow in mass with time \[the Salpeter time scale of growth is determined by t$`{}_{s}{}^{}=3\times 10^7(L_{Edd}/L_B)\eta _{0.1}`$ yr. where $`\eta _{0.1}`$ is the radiative efficiency in the units of 0.1 (see Fabian 1999)\]. Since NLS1 accrete at close to Eddington limit, their BHs would grow faster. So, smaller BHs in NLS1s are likely to be younger as well.
2). Super-solar gas phase metallicities. There are a couple of lines of evidence to suggest that NLS1s may have super-solar gas phase metallicities. One comes from the study of high ionization emission lines. Wills et al. (1999) found that the strength of NV $`\lambda 1240`$ emission line was systematically larger while the strength of the CIV $`\lambda 1549`$ was systematically smaller in AGN with narrow emission lines. The NV/CIV ratio serves as an abundance indicator as shown by Hamann & Ferland (1993). So, the NLS1s may have large nitrogen abundance. Absorption lines, being insensitive to density, serve as better indicators of metallicities. In the narrow line AGN PG1404+226, Ulrich et al. (1999) found that while the strengths of Ly$`\alpha `$ and CIV absorption lines were in reasonable agreement with those expected from the ionized X-ray absorber (See Mathur 1997 for X/UV absorber models), the NV absorption line was significantly stronger. This observation, again, can be understood in terms of high nitrogen abundance in this narrow line AGN (Mathur & Komossa, 2000).
Large nitrogen abundance is obtained when overall metallicities are high with N$`/`$H $``$ (Z$`/`$Z)<sup>2</sup> where Z is the solar abundance. Nitrogen is preferentially enhanced because of secondary CNO nucleosynthesis (see Hamann & Ferland 1999, HF99, for details on AGN metallicities). Thus the observations of emission as well as absorption lines in NLS1s imply super-solar gas phase metallicities. The strength of the fluorescent Fe-K alpha line in some NLS1s is also indicative of super-solar abundance (Fabian 1999).
Such metal enrichment is possible when the initial mass function of star formation is flat, favorable for high mass star formation, and the evolution is fast. Such a star formation scenario is likely to be present in deep potential wells like galactic nuclei and protogalactic clumps (HF99). Moreover, high metallicities are achieved while consuming less gas (HF99). The NLS1s may then represent that early phase in galactic evolution when rapid star formation is taking place in the nucleus.
3). IR brightness. Many NLS1s are observed to be bright in the infra-red (Moran, Halpern, & Helfand 1996). Young, star forming galaxies are also bright IR sources. It is possible that a part of the nuclear IR flux is from a nuclear star burst.
4) Analogy with high redshift quasars. In $`\mathrm{\S }3`$ we argue that NLS1s are analogous to very high redshift quasars. High redshift (z$`>`$4) quasars are believed to quasars in the early phase of evolution compared to the $`z1`$ quasars. By analogy, NLS1s may as well be in the early evolutionary phase compared to the normal Seyfert galaxies.<sup>1</sup><sup>1</sup>1We note that Grupe (1996) has also argued that the age since an AGN was born might be the underlying reason for some NLS1-type correlations he has studied.
## 3 Are NLS1s low redshift analogues of high redshift quasars?
That the AGN phenomenon was so much stronger at z$``$ 2–3 than today has long elicited the suspicion that there is a connection between the youth of a galaxy and the likelihood that an AGN forms inside it. The question then naturally arises, “what are the local counterparts to the young galaxies in early universe, in which local AGN may live?” (see, e.g, Krolik 1999). A standard answer to this question is “Starburst galaxies”. Heckman (1999) has argued that starburst galaxies are the low redshift analogues of Lyman break galaxies at high redshift. Similarly, we ask, what are the low redshift analogues of high redshift (z$`\genfrac{}{}{0pt}{}{_>}{^{}}`$4) quasars? We propose that they might be NLS1s.
It is interesting to note the similarity of the properties of NLS1s with high redshift (z$`\genfrac{}{}{0pt}{}{_>}{^{}}`$4) quasars.
1) Hamann & Ferland (1993) found high metallicities in high redshift quasars (Z$`\genfrac{}{}{0pt}{}{_>}{^{}}`$Z at z$`\genfrac{}{}{0pt}{}{_>}{^{}}`$4). Large metallicities in NLS1s may make them low redshift, low luminosity analogues of high redshift quasars.
2). NLS1s and BALQSOs. Similarities between the observed properties of low ionization Broad Absorption Lines Quasars (BALQSOs) and NLS1s have been reported in the literature (e.g. Lawrence et al. 1997, Leighly et al. 1997). Both these classes show strong FeII$`\lambda 4570`$ and AlIII$`\lambda 1857`$ and weak CIV$`\lambda 1549`$ and \[OIII\]$`5007`$ emission lines. Their continua are red in the optical and strong in the IR. Evidence of relativistic outflow is also reported in three NLS1s (Leighly et al. 1997). If these two classes are indeed related (Brandt 1999), then NLS1, at least those those with some evidence of outflow, might be low redshift, low luminosity cousins of BALQSOs. BALQSOs are tentatively identified with that phase in quasar evolution when the matter around the nuclear BH is being blown away, and a quasar emerges (see, e.g. Fabian 1999). NLS1s may then represent a similar early evolutionary phase at low redshift.
3). Optical spectra of a sample of z$`\genfrac{}{}{0pt}{}{_>}{^{}}`$4 quasars revealed that their emission lines are typically narrower than the low redshift quasars (FWHM $`\genfrac{}{}{0pt}{}{_<}{^{}}`$2000 km/s, Shields & Hamann 1997). The normal explanation of this observation is that these are type 2 quasars, where the broad emission lines are obscured from our line of sight. Alternatively, these high z quasars might be true “narrow” broad line objects.
4) We note here another interesting connection with high redshift. As discussed in $`\mathrm{\S }1`$, NLS1s have strong FeII emission lines. Quasars Q0014+813 and Q0636+680 at redshifts z=3.398 and z=3.195 respectively, were observed to have very strong FeII emission (Elston, Thompson, & Hill 1994). Are they also highly accreting objects at early evolutionary phase? Note also the narrow UV emission lines (FWHM $`\genfrac{}{}{0pt}{}{_<}{^{}}`$2150 km s<sup>-1</sup>) in the ultra strong UV FeII emitter Q2226-3905 (Graham, Clowes, & Campusano, 1996).
All these similarities point towards NLS1s being low redshift, low luminosity analogues of high redshift quasars.
## 4 Do NLS1s reside in rejuvenated galaxies?
In the previous section we have argued that NLS1s may represent an early phase in AGN evolution. Whether they reside in young galaxies is a separate question and a step further. That young galaxies are gas rich is helpful; they would have the large reservoir of gas necessary to sustain the close to Eddington rate accretion in NLS1s. But do we have any evidence that they indeed reside in young galaxies? There is no published systematic study of the properties of the host galaxies of NLS1s. However, some of the NLS1s are originally from Zwicky (e.g. I Zw 1) and Markarian (e.g. Mrk 766) sample of galaxies implying that they are blue. While the blue color might be due to big blue bumps in the active nuclei, as in normal Seyfert galaxies, NLS1s have weak blue bumps ($`\mathrm{\S }1`$) and so the blue colors might be a result of actively star forming galaxies. Some NLS1s are IRAS galaxies (e.g. IRAS 13349+2438), infra-red bright, and star forming. Using the catalogs of galaxies RC3 (de Vaucouleurs et al. 1991) and UGC (Nilson, 1973) we looked into the morphology of a small sample of NLS1s with X-ray absorption features and found information on seven of them. Three were found to be compact (I Zw 1, Mrk 507, and Mrk 1298), two with signatures of inner ring (NGC 4051 and Ark 564), and three with nuclear bars (NGC 4051, Mrk 776 and Ark 564). These are signatures of recent activity, most likely due to galaxy- galaxy interactions or mergers. In this scenario the galaxies are newly formed, or rejuvanated.
That NLS1s reside in young galaxies is also consistent with the hypothesis that the formation and evolution of galaxies and their active nuclei is intimately related (Rees 1997, Fabian 1999, Granato et al. 1999, Haehnett & Koffmann 1999). In this scenario, the process of formation of a massive BH and the active nucleus is the very process of galaxy formation. The active nucleus and the galaxy evolve together, with BH accreting matter and the galaxy making stars. At one stage the winds from the active nucleus blow away the matter surrounding it and a quasar emerges. This is not only the end of active evolution of the quasar but that of the galaxy as well, which is evacuated of its interstellar medium. The quasar then shines as long as there is fuel in the accretion disk (Fabian 1999). In this scenario, high redshift quasars represent early stage of galaxy evolution, BALQSOs at z$`2`$ represent the stage when the gas is being blown away and z$`1`$ quasars would be the passively evolving population. Massive ellipticals found today might be the dead remnants of what were once quasars.
The quasar phenomenon may thus be a result of galaxy formation due to primordial density fluctuations. At low redshift, when new galaxies are formed due to interactions or mergers, similar evolution may take place. As argued above, the NLS1s may represent a crucial early phase. (In our scenario, the accretion rate ṁ is large in the early stages of evolution and reduces later on. This is opposite to the proposal by Wandel (1999) in which ṁ increases with time.)
In fact, there might be some NLS1s with a star burst component (see $`\mathrm{\S }2`$). Soft X-ray spectra of NLS1 are steep and often variable. However, Page et al. (1999) report that while the power-law component in the NLS1 Mrk 766 varied, the thermal black-body component did not. This component might well be due to a nuclear starburst. Note also the strong CO emission in the prototype NLS1 I Zw 1 (Barvainis, Alloin & Antonucci 1989). Schinnerer, Eckart & Tacconi (1998) mapped I Zw 1 in Co and found a circumnuclear ring of diameter 1.8 kpc. The authors found strong evidence for a nuclear starburst. There is also a companion to I Zw 1, supporting an interpretation of starburst due to interaction. Similarly, AGN activity is known to exist in star burst galaxies (see Heckman 1999 for a review). Dennefeld et al. (1999) report observations of narrow optical emission lines in a sample of IR selected starburst galaxies.
## 5 Observational Tests
Here we propose several observations that could test the ideas presented above. (1) Emission line properties of high redshift quasars: Objects in the Shields & Hamann sample, for example, were selected on the basis of their colors, in particular very red B$``$R which results when Ly$`\alpha `$ is redshifted into the R band. This could lead to a selection bias in favor of objects with narrow, peaky profiles (Shields & Hamann 1997). It would be important to remove such selection bias before we can conclude that z$`>4`$ quasars have narrow emission lines. A broader wavelength coverage with more than just two bands would be useful. The “drop-out” technique (Steidel et al. 1995) used for finding Lyman break galaxies would be another way to remove emission line bias. (2) X-ray properties of high redshift quasar: Only about a dozen z$`>4`$ quasars are detected in X-rays (see the latest update by Kaspi, Brandt & Schneider, 2000). However, X-ray spectra of z$`>4`$ quasars are still not available. It will be interesting to see if they are steep, and highly variable like those of NLS1s. We will be studying X-ray spectra of radio-loud as well as radio quiet z$`>4`$ quasars with XMM. (3) Morphology and environment of NLS1s: Are NLS1s preferentially found in younger galaxies and/or in more disturbed environments? A systematic study of host galaxy properties of a well defined sample of NLS1s is needed. (4) Search for starburst components in NLS1s: Evidence for the starburst–NLS1 connection ($`\mathrm{\S }4`$) is suggestive, but not yet statistically sound. Do the nuclear components of starburst galaxies show narrow emission lines? Do NLS1s show evidence of a starburst component more often than normal Seyferts? Spectroscopic as well as high resolution imaging observations would help towards establishing the connection between the two.
## 6 On the correlation of FeII strength and X-ray spectral slope
While photoionization models reproduce the properties of optical and UV emission lines observed in AGN spectra with reasonable success, the case of FeII lines is different. Wills, Netzer & Wills (1985) found that standard photoionization models cannot explain the strength of the the observed FeII lines in AGN. Collin-Souffrin and collaborators (1988) as well as Kwan (1984) advocated that collisional ionization would be important in the production of FeII. Collin (1999) has again shown that strong FeII emission cannot be produced by photoionization with any set of parameters, and even by making iron abundance reasonably super-solar. The importance and necessity of collisional ionization of iron was, however, not appreciated at least in part due to the observed correlation between FeII$`\lambda 4570`$ equivalent width and the soft X-ray energy index (Wilkes, Elvis & McHardy, 1987, Shastri et al. 1993). Standard photoionization models for the line emission from the broad line regions of quasars (Krolik 1988) imply a close link between the energy index of the ionizing X-ray continuum and the strength of emission lines. So the Wilkes et al. and Shastri et al. correlation was interpreted as a result of photoionization. Note, however, that the correlation is in in the opposite sense to that predicted by the standard photoionization model in which FeII emission is generated deep within the cloud and thus is sensitive to harder X-rays. NLS1s may provide us with the clue to understand this observed, conflicting trend as discussed below.
As discussed in $`\mathrm{\S }1`$, there is a general consensus that large accretion rate, Ṁ/Ṁ<sub>Edd</sub>, is the likely driver of the many observed properties of NLS1s. Large strength of FeII emission may then be linked to the large accretion rate. In a model by Kwan et al. (1995) FeII line emission is produced in an accretion disk. The accretion disks with larger accretion rate may simply have more mass to produce stronger FeII. Thus we argue that the observed correlation of FeII strength and soft X-ray slope is a consequence of the correlation with the accretion rate and support collisional ionization as the origin.
## 7 Conclusions
We have argued that NLS1s are likely to be AGN in the making and reside in rejuvenated galaxies. As such they represent a crucial early phase in the evolution of galaxies and active galaxies. What we need now is a systematic study of host galaxy properties of a well defined sample of NLS1s and their comparison with the hosts of normal Seyferts. Some evidence presented above is based on a small number of objects and generalization may not be appropriate. Studies at high redshift also suffer from selection effects. Understanding and correcting for them is crucial in establishing the analogy with NLS1s on firm footing. X-ray spectra of high redshift radio-quiet quasars will provide an addition piece of information towards this goal. It would of great interest to find out whether star burst galaxies are parent population of NLS1s.
## Acknowledgments
I thought about NLS1s and their place in the cosmic “big picture” while preparing my talk for the NLS1 workshop at Bad Honnef, Germany. I am grateful to Th. Boller and the organizing committee for inviting me to the workshop “Observational and theoretical progress in the study of Narrow Line Seyfert 1 Galaxies”.
It is my pleasure to thank F. Hamann, C. Reynolds, D. Weinberg, A. Pradhan, R. Pogge, B. Peterson, M. Elvis, B. Ryden and P. Osmer for useful discussions and encouragement, and the referee Niel Brandt for useful comments.
This research has made use of the NASA/IPAC Extragalactic Database (NED) which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration. Support through NASA grant NAG5-3249 (LTSA) is gratefully acknowledged. |
warning/0003/hep-th0003088.html | ar5iv | text | # From cool pions to the chiral phase transition
## Abstract
Using the ideas of effective field theory and dimensional reduction, we relate the parameters of two low energy models of QCD: the $`O(N)`$ nonlinear sigma model in $`D=3+1`$, which describes the dynamics of cool pions, and the $`O(N)`$ Heisenberg magnet in $`D=3+0`$, which is commonly argued to reproduce the correct critical behaviour of the chiral phase transition. As a result, we obtain a generalized expression for the finite temperature pion decay constant which reproduces, in certain limits, the available expressions in the literature.
The thermodynamic properties of the low energy dynamics of QCD can be apropriately described in terms of different $`O(N)`$ $`\sigma `$ models, deppending on whether or not we are close to the critical temperature $`T_c150`$ Mev. At low temperatures, the only degrees of freedom which are excited are pions. One may then use the results of chiral perturbation theory, which in its lowest order (very lowest energies) is just the nonlinear $`\sigma `$ model. In fact, as first pointed out by Weinberg , a suitable effective field theory involving Goldstone fields automatically generates transition amplitudes which obey the low energy theorems of current algebra and PCAC. The interaction among the Goldstone bosons is described by an effective Lagrangian, which is invariant under global chiral transformations. Near the critical temperature, on the other hand, it has been demonstrated that the long wavelength fluctuation modes of the massless $`N_f=2`$ QCD belong to the same universality class of the $`O(4)`$ Heisenberg magnet in three spatial dimensions . The reasoning behind this proposal is based on counting the light degrees of freedom. The transition region is dominated by the longitudinal and transverse fluctuations of the order parameter, the $`\sigma `$ and $`\stackrel{}{\pi }`$ fields, which go soft at the transition temperature. Being bosonic, $`\sigma `$ and $`\stackrel{}{\pi }`$ have zero frequency Matsubara modes, $`\omega _n=0`$, which turn to be the only relevant degrees of freedom in the scaling region. The fermions themselves, even if they are massless at zero temperature, do not influence the nature of the phase transition at finite temperature.
Before proceeding, however, we should mention that although this is a very intuitive scenario, there still is some divergence concerning the nature of the critical properties of massless $`N_f=2`$ QCD. For example, when one uses lattice simulation to analyse the critical behavior of the order parameter (the $`q\overline{q}`$ condensate), for which universality arguments suggest
$$q\overline{q}\left|(TT_c)/T_c\right|^{0.38\pm 0.01},$$
(1)
one concludes that lattice data is indeed consistent with the $`O(4)`$ spin system . Conversely, when one tries to check universality arguments by looking at the behavior of global thermodynamical quantities, such as the specific heat, one finds that, while in the $`O(4)`$ spin system the singular contribution of the soft modes vanishes at the critical point, lattice data for the $`N_f=2`$ QCD actually show a huge peak in the specific heat around $`T_c`$. As pointed out by Shuriak , this certainly mean that near criticality many new degrees of freedom become available or are significantly changed. This discrepancy may be understood if we remember that, while critical correlation functions are dominated by the long wavelength fluctuation modes, global thermodynamical quantities, such as the free energy or the specific heat, receive contribution from all energy modes. In fact, in any asymptotically free theory, like QCD, the effect of raising the temperature is to populate the weakly coupled higher energy modes . In this sense, corrections to a free gas should fall with the temperature. The $`n=0`$ mode makes only a perturbatively small contribution and all other higher modes must be considered. The conclusion is that global thermodynamic quantities simply cannot be calculated using classical statistical mechanics.
In this brief report we will not enter into the above discussion but rather study further the relation between the parameters of the two effective models for the low energy dynamics of QCD, from the point of view of effective field theory and dimensional reduction. This will be done with the aid of a ($`3+1`$)-dimensional $`O(N)`$ linear $`\sigma `$ model which will serve as a link between the above two effective models. Let us begin by considering the Lagrangian density
$$(\stackrel{}{\varphi }(x))=\frac{1}{2}(_\mu \varphi ^a)^2+\frac{1}{2}\mu ^2\varphi ^a\varphi ^a+\frac{\lambda }{4}(\varphi ^a\varphi ^a)^2,$$
(2)
which describes the long distance properties of a Heisenberg ferromagnet, where $`\stackrel{}{\varphi }(x)`$ is a $`N`$-component real massive scalar field playing the role of a slowly varying order parameter. In the ordered phase ($`\mu <\mu _c`$), the field has a nonvanishing expectation value
$$\stackrel{}{\varphi }(x)\stackrel{}{\varphi }(x)=M^2(\mu ,\lambda )=\frac{\mu ^2}{\lambda },$$
(3)
breaking spontaneously the $`O(N)`$ rotational symmetry to $`O(N1)`$ invariance. In order to consistently quantize the theory a prefered direction in the internal space must be choosen. This choice is completely arbitrary, and for our purposes it will be convenient to choose the direction of the unitary $`N`$-component field $`\widehat{\varphi }(x)\stackrel{}{\varphi }(x)/|\stackrel{}{\varphi }(x)|`$. This can be implemented by the following change of variables in the functional integral
$$\stackrel{}{\varphi }(x)=\rho (x)\widehat{\varphi }(x),$$
(4)
where $`\widehat{\varphi }(x)\widehat{\varphi }(x)=1`$ by construction. In terms of the new variables $`\rho (x)`$ and $`\widehat{\varphi }(x)`$, the functional integral becomes
$$Z=\rho ^{N1}(x)𝒟\rho (x)𝒟\widehat{\varphi }(x)e^{S(\rho ,\widehat{\varphi })},$$
(5)
with
$`S(\rho ,\widehat{\varphi })`$ $`=`$ $`{\displaystyle }\mathrm{d}^4x\{{\displaystyle \frac{1}{2}}\rho ^2(x)\left(_\mu \widehat{\varphi ^a}(x)\right)^2`$ (6)
$`+`$ $`{\displaystyle \frac{1}{2}}\left(_\mu \rho (x)\right)^2+{\displaystyle \frac{1}{2}}\mu ^2\rho ^2+{\displaystyle \frac{1}{4}}\lambda \rho ^4\}.`$ (7)
The theory has a natural UV cutoff, namely the meson mass $`M`$ defined by (3). We can then perform integration over the $`\rho (x)`$ field in order to generate an effective local action $`S_{eff}(\widehat{\varphi })`$ for the long wavelength fluctuations of $`\widehat{\varphi }`$
$$e^{S_{eff}(\widehat{\varphi })}=\rho ^{N1}(x)𝒟\rho (x)e^{S(\rho ,\widehat{\varphi })}.$$
(8)
Although it is generally not possible to compute this integral exactly, due to the $`\rho `$ quartic self-interaction, one can nevertheless compute it perturbatively. Indeed, since $`M^2(\mu ,\lambda )=\mu ^2/\lambda `$ is the vacuum expectation value of the field $`\rho `$, we can write
$$\rho (x)=M+\rho ^{}(x),$$
(9)
where $`\rho ^{}`$ are fluctuations of the $`\rho `$ field around $`M`$. In terms of the new variable $`\rho ^{}`$ the action (7) reads
$`S(\rho ^{},\widehat{\varphi })`$ $`=`$ $`{\displaystyle }\mathrm{d}^4x\{{\displaystyle \frac{1}{2}}[M^2+2M\rho ^{}+\rho _{}^{}{}_{}{}^{2}](_\mu \widehat{\varphi ^a}(x))^2`$ (10)
$`+`$ $`{\displaystyle \frac{1}{2}}(_\mu \rho ^{}(x))^2+{\displaystyle \frac{\mu }{2}}[M+\rho ^{}]^2+{\displaystyle \frac{\lambda }{4}}[M+\rho ^{}]^4\}.`$ (11)
The zeroth order in a perturbative expansion of the above effective field theory is obtained by neglecting all $`\rho ^{}`$. As a result we obtain, after the identification $`M^2=f_\pi ^2`$, the expression
$$S_{eff}^0(\widehat{\varphi })=f_\pi ^2\mathrm{d}^4x\frac{1}{2}\left[_\mu \widehat{\varphi ^a}(x)\right]^2,$$
(12)
which is simply the action of a nonlinear $`\sigma `$ model. It must be emphasized that loop corrections coming from the $`\rho ^{}`$ integration renormalize the coeficient $`f_\pi ^2(\mu ,\lambda )`$ in front of (12) . Additional $`\widehat{\varphi }`$ interactions can be expanded in local terms, provided we are exploring momenta much smaller than any mass scale in the problem. In this sense, we conclude that the nonlinear $`\sigma `$ model completely describes the long distance properties of a Heisenberg ferromagnet at low temperature.
Let us now turn to the other extreme. Near the critical temperature it is the long wavelength fluctuation modes of QCD which dominates the scaling behavior of correlation functions. Thus, roughly speaking, at least these correlation functions should be described by statistical mechanics. In terms of the discrete frequency sum of finite temperature quantum field theory, this means that one needs to retain only the zero-frequency components of the fields. Keeping only these zero-frequency modes yields a field theory in a lower dimension with temperature dependent renormalized coupling constants. The reduced theory is an effective theory for the zero modes of the original fields and can be explicitly obtained from the original theory by integrating out the nonzero Matsubara frequencies of all fields, bosons and/or fermions.
Dimensional reduction of the $`O(N)`$ invariant linear $`\sigma `$ model was performed in , to one-loop order, using a modified minimal subtraction scheme at zero momenta and $`T_00`$. It was shown that, from the Lagrangian density (2), integration over nonstatic modes (nonzero Matsubara frequencies) give a dimensionally reduced effective free energy of the Landau-Ginzburg type for the bosonic zero mode $`\varphi _0`$, $`\stackrel{}{\varphi }(x)=(\varphi _0(x),0,\mathrm{},0)`$,
$$F(\varphi _0)=\mathrm{d}^3𝐱\left\{\frac{1}{2}(_i\varphi _0)^2+\frac{1}{2}\mu _R^2\varphi _0^2+\frac{\lambda _RT}{4}\varphi _0^4\right\},$$
(13)
with all possible nonrenormalizable interaction terms supressed by the temperature . In the above expression, $`\mu _R`$ and $`\lambda _R`$ are the thermally renormalized mass and coupling constant given by
$$\mu _R^2=\mu ^2+(N+2)\left\{\frac{\lambda (T^2T_0^2)}{12}\frac{\lambda \mu ^2}{4\pi ^2}\mathrm{ln}\left(\frac{T}{T_0}\right)\right\}$$
(14)
and
$$\lambda _R=\lambda \frac{(N+8)}{6}\frac{\lambda ^2}{16\pi ^2}\mathrm{ln}\left(\frac{T}{T_0}\right).$$
(15)
From the above expressions it is clear that at long distances the only effect of the nonstatic modes is to set the scale of the coupling constants to be the temeprature. The dependence of both $`\mu _R`$ and $`\lambda _R`$ on the choice of the thermal renormalization point $`T_0`$ is controlled by a homogeneous renormalization group equation
$$\left(T_0\frac{}{T_0}+\beta _{T_0}\frac{}{\lambda }+\gamma _{T_0}\mu ^2\frac{}{\mu ^2}\right)F(\varphi _0)=0,$$
(16)
with renormalization group functions
$$\beta _{T_0}=\frac{(N+8)}{6}\frac{\lambda ^2}{16\pi ^2}$$
(17)
and
$$\gamma _{T_0}=(N+2)\left[\frac{\lambda T_0^2}{12\mu ^2}+\frac{\lambda }{16\pi ^2}\right],$$
(18)
in accordance to .
Let us now define, in analogy to eq. (3), the quantity
$$f_\pi ^2(T,T_0)=\frac{\mu _R^2}{\lambda _R},$$
(19)
which relates the parameters of the two low energy models of QCD, namely $`f_\pi ^2`$ in (12) and $`\mu _R`$ and $`\lambda _R`$ in (13). It is not difficult to see that at $`T_0=0`$, and neglecting all the logarithms, expression (19) exactly reproduces the result of Bochkarev and Kapusta for the finite temperature pion decay constant <sup>*</sup><sup>*</sup>*There should be no reason to worry about the limit $`T_0=0`$ in expressions (14) and (15). If we had adopted a renormalization prescription based on subtractions at $`T_0=0`$ the logarithms in expressions (14) and (15) would be naturaly replaced by $`\mathrm{ln}(T/\sqrt{\mu })`$ .
$$f_\pi ^2(T,T_0=0)f_\pi ^2(T)=f_\pi ^2\left[1\frac{N+2}{12}\frac{T^2}{f_\pi ^2}\right].$$
(20)
We shall now give a physical interpretation to the quantity $`f_\pi ^2(T,T_0)`$. At zero temperature, $`f_\pi 93`$ Mev and is connected to the matrix elements of the axial vector current $`𝒜_\mu `$ by $`<0|𝒜_\mu ^a|\pi ^b(p)>=if_\pi p_\mu \delta ^{ab}`$, where the upper index in $`𝒜_\mu ^a`$ is related to isospin. At finite temperatures, on the other hand, it has been recently proposed by Bochkarev and Kapusta , using standard linear response theory, that the quantity $`f_\pi ^2(T)`$ would measure the strength of the coupling of the Goldstone bosons to the longitudinal part of the axial spectral density for $`𝒜_\mu `$, in the limit of zero external momentum. Alternatively, it was further suggested in that $`f_\pi ^2(T,0)/f_\pi ^2`$ would give the percentual of (or the probability of finding) Goldstone excitations in the ground state. For this reason, it is possible to speak about a finite temperature phase transition in the system described by (12), associated to the complete decoupling of the Goldstone bosons from the ground state, which happens at a critical temperature
$$T_c^2=\frac{12}{N+2}f_\pi ^2.$$
(21)
Following the above discussion, we are then lead to give similar interpretation for the quantity $`f_\pi ^2(T,T_0)/f_\pi ^2`$ as measuring the percentual of (or the probability of finding) Goldstone excitations between the ground state and the state corresponding to $`T_0`$. This is in fact evident from (14) and (15) from which we obtain $`f_\pi ^2(T,T_0=T)=f_\pi ^2`$, $`T`$. This trivially states that the percentual of Goldstone excitations between the ground state and the state corresponding to $`T_0=T`$ is unity (or $`100\%`$), as expected. In this sense, we conclude that expression (19) generalizes the results of by allowing us to have access to the dynamics of the Goldstone system in a wide energy region, ranging from the ground state $`T_0=0`$ to the state corresponding to $`T_0=T`$.
Expression (19) also reveals the existence of a critical line of second order phase transitions defined by
$$f_\pi ^2(T_c,T_0)=0.$$
(22)
In fact, by solving the above equation for $`T_c`$ we obtain, neglecting all the logarigthms,
$$T_c^2(T_0)=\frac{12}{N+2}f_\pi ^2+T_0^2,$$
(23)
which is consistent with the result (21) for $`T_0=0`$. That $`T_c^2(T_0)`$ is a monotonic growing function of $`T_0`$ should not be surprising. It simply means that it is more and more likely to find excitations populating higher energy states, than in the ground state.
We have obtained an expression for the finite temperature pion decay constant which somehow generalizes the results of . We believe that our expression (19) may be useful in studies in which one is interested in computing transition amplitudes between arbitrary energy states, other than the ground state. In fact, since the pion contribution to correlation functions must depend on the probability of the system to be in a given energy state, the quantity (19) will act as a weight in computing correlation functions.
The author is indebted with A. P. C. Malbouisson and N. F. Svaiter for inumerous discussions and comments in early stages of this work. FAPERJ is also acknowledged for finantial support. |
warning/0003/cond-mat0003132.html | ar5iv | text | # Conductivity sum rule; comparison of coherent and incoherent c-axis coupling
## Abstract
We calculate the $`c`$-axis kinetic energy difference between normal and superconducting state for coherent and for incoherent interlayer coupling between CuO<sub>2</sub> planes. For coherent coupling the ratio of the missing conductivity spectral weight to the superfluid density is equal to one and there is no violation of the conventional sum rule, but for the incoherent case we find it is always greater than one whatever the nature of the impurity potential may be. To model more explicitly YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-x</sub> around optimum doping, which is found to obey the sum rule, we consider a plane-chain model and show that the sum rule still applies. A violation of the sum rule of either sign is found even for coherent coupling when the in-plane density of electronic states depends on energy on a scale of the order of the gap.
It has been proposed that the interlayer coupling along the $`c`$ axis of a high transition temperature $`(T_c)`$ superconductor is incoherent, and the electronic kinetic energy along the $`c`$ axis changes when the system enters the superconducting state. Recently, Basov et al. have reported that there is a significant discrepancy between the superfluid density $`\rho _s`$ and the spectral weight missing from the real part of the $`c`$-axis conductivity $`N_nN_s=8_{0^+}^{\omega _c}𝑑\omega \left[\sigma _{1c}^n(\omega )\sigma _{1c}^s(\omega )\right]`$, where $`\omega _c`$ is a cutoff frequency of the order of a bandwidth, in several high-$`T_c`$ cuprate superconductors such as optimally doped Tl<sub>2</sub>Ba<sub>2</sub>Cu<sub>6+x</sub> (Tl$`2201`$). This implies that the conventional sum rule of Ferrel, Grover and Tinkham (FGT) is violated. However, the spectral discrepancy becomes vanishing for YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6.85</sub> and disappears for the optimally doped YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6.95</sub> (YBCO) crystal with $`T_c93K`$ as the $`c`$-axis response become more coherent with increasing oxygen content. Basov et al. also pointed out that there is no such discrepancy in the in-plane response for any cuprate. Moreover, for over-doped Tl$`2201`$, the sum rule discrepancy vanishes and a Drude-like peak develops in the conductivity for $`T>T_c`$. Theses observations, therefore, suggest that for coherent interlayer coupling in the cuprate superconductors the conventional sum rule is obeyed.
In this paper, we consider both coherent and incoherent $`c`$-axis coupling between CuO<sub>2</sub> planes. For the coherent case we find that the superfluid density remains equal to the missing optical spectral weight; in other words, it does not violate the FGT sum rule. The $`c`$-axis kinetic energies in the normal and superconducting state have the same value. For incoherent $`c`$-axis coupling the ratio of the missing area to the superfluid density is always larger than one in disagreement with some recent experiments. In YBCO the CuO chains play an important role in the electrodynamics and at the optimal doping a plane-chain model is needed to be complete. Here we use this model to investigate the $`c`$-axis conductivity sum rule. An algebraic calculation of the electronic kinetic energies is complicated and a numerical calculation is required although it can be reduced a lot in a special case, in which only the leading order in perturbation theory is kept. This case is particularly interesting because it has been shown to exhibit a pseudogap in the $`c`$-axis conductivity. Finally, we discuss the possibility that the FGT sum rule is violated in the plane-plane case when the in-plane density of states depends on energy even for coherent case.
The Hamiltonian $`H`$ for a cuprate superconductor with coherent $`c`$-axis coupling is $`H=H_0+H_c`$, where $`H_0`$ describes a $`d`$-wave superconductor in a plane and $`H_c=_{i\sigma }t_{}\left[c_{i1\sigma }^+c_{i2\sigma }+c_{i2\sigma }^+c_{i1\sigma }\right]`$ is a coherent interlayer coupling due to the overlap of electronic wave functions which is represented by $`t_{}`$; therefore, by coherent coupling we mean a tight binding like coupling along the $`c`$ axis. It will not be necessary to treat $`t_{}`$ as a constant in what follows. It can depend on an angle in the plane. For incoherent coupling the Hamiltonian is $`H_c^{}=_{i\sigma }V_i\left[c_{i1\sigma }^+c_{i2\sigma }+c_{i2\sigma }^+c_{i1\sigma }\right]`$, where $`V_i`$ is an impurity scattering potential, so that impurity scattering mediates the $`c`$-axis hopping and an impurity average is implied.
In the presence of an external vector potential $`A_z`$, $`H_c`$ is modified to $`H_c(A_z)`$ by the phase factor $`\mathrm{exp}(ieA_z)`$ for $`c_{i1\sigma }^+c_{i2\sigma }`$ and $`\mathrm{exp}(ieA_z)`$ for $`c_{i2\sigma }^+c_{i1\sigma }`$. For the response to an external field, $`H_c(A_z)`$ is expanded up to second order of $`A_z`$ to obtain the current $`j_c=\delta H_c(A_z)/\delta A_z=j_p+j_d`$, where $`j_p=ied_{i\sigma }t_{}\left[c_{i1\sigma }^+c_{i2\sigma }c_{i2\sigma }^+c_{i1\sigma }\right]`$ and $`j_d=e^2d^2H_cA_z`$ with $`d`$ the interlayer spacing. In linear response theory, $`j_c=\left[\mathrm{\Pi }+e^2d^2H_c\right]A_z`$, where $`\mathrm{\Pi }`$ is the current-current correlation function associated with $`j_p`$ and $`H_c`$ is the perturbation of $`j_d`$ due to $`H_c`$. The conductivity $`\sigma _c(𝐪,\omega )`$ is given by
$$\sigma _c(𝐪,\omega )=\frac{i}{\omega }\left[\mathrm{\Pi }(𝐪,\omega )e^2d^2H_c\right].$$
(1)
In the Matsubara formalism,
$$\mathrm{\Pi }(𝐪,\omega )=2(ed)^2T\underset{\omega ^{}}{}\underset{𝐤}{}t_{}^2\text{T}\text{r}\left[\widehat{\tau }_0\widehat{G}(𝐤,\omega ^{})\widehat{\tau }_0\widehat{G}(𝐤,\omega ^{}+\omega )\right],$$
(2)
and
$$H_c=2T\underset{\omega }{}\underset{𝐤}{}t_{}^2\text{T}\text{r}\left[\widehat{\tau }_3\widehat{G}(𝐤,\omega )\widehat{\tau }_3\widehat{G}(𝐤,\omega )\right],$$
(3)
where $`\widehat{\tau }_i`$ is the Pauli matrix in the spin space, and $`\widehat{G}(𝐤,\omega )`$ is the Green’s function in Nambu representation, namely, $`\widehat{G}(𝐤,\omega )=(i\omega \widehat{\tau }_0+\xi _𝐤\widehat{\tau }_3\mathrm{\Delta }_𝐤\tau _1)/(\omega ^2+\xi _𝐤^2+\mathrm{\Delta }_𝐤^2)`$ wiht $`\xi _𝐤`$ the in-plane energy and $`\mathrm{\Delta }_𝐤`$ the gap which has $`d_{x^2y^2}`$ symmetry in the cuprates.
The $`c`$-axis conductivity sum rule of the system is
$$\frac{2}{\pi }_0^{\mathrm{}}𝑑\omega \sigma _{1c}(\omega )=e^2d^2H_c.$$
(4)
We use the unit such that $`\mathrm{}=c=k_B=1`$ and set the volume of the system to be unity. From the sum rule, the superfluid density $`\rho _s`$ can be written as
$$\rho _s=8_{0^+}^{\omega _c}𝑑\omega \left[\sigma _{1c}^n(\omega )\sigma _{1c}^s(\omega )\right]4\pi e^2d^2\left[H_c^sH_c^n\right],$$
(5)
where $`\omega _c`$ is the cutoff frequency for interband transitions that $`H_c`$ does not account for.
Since the difference between the superfluid density and the missing spectral weight is proportional to the kinetic energy difference between normal and superconducting state as seen in Eq. (5), it is necessary to calculate $`H_c^sH_c^n`$ to see if the FGT sum rule is violated by coherent $`c`$-axis coupling. For the normal state,
$$H_c^n=4T\underset{\omega }{}\underset{𝐤}{}t_{}^2G_0(𝐤,\omega )^2,$$
(6)
where $`G_0(𝐤,\omega )`$ is a normal state Green’s function and $`t_{}`$ may depend on $`k_z`$ and $`\varphi =\mathrm{tan}^1(k_y/k_x)`$. We assume a cylindrical Fermi surface with $`\xi =k^2/2m\mu `$, where $`\mu `$ is a chemical potential in the plane, and a $`d`$-wave gap $`\mathrm{\Delta }_𝐤=\mathrm{\Delta }(T)\mathrm{cos}2\varphi _k`$. Then, we obtain
$`H_c^n`$ $`=4{\displaystyle \underset{k_z}{}}{\displaystyle \frac{d\varphi }{2\pi }t_{}^2_{\omega _c}^{\omega _c}𝑑\xi N(\xi )\frac{f(\xi )}{\xi }}`$ (8)
$`=4N(0){\displaystyle \underset{k_z}{}}{\displaystyle \frac{d\varphi }{2\pi }t_{}^2\mathrm{tanh}(\frac{\omega _c}{2T})},`$
where the integration range is limited by $`\omega _c`$, and the density of states, $`N(\xi )`$, is approximated by a constant value $`N(0)`$ around the Fermi energy. Later, we will discuss the effect of $`N(\xi )`$ on $`H_c`$ and will note the possibility that the FGT sum rule may be violated even for coherent $`c`$-axis coupling. Since $`f(\xi )/\xi =\delta (\xi )`$ at zero temperature ($`T=0`$), $`H_c^n`$ turns out to be $`4N(0)_{k_z}𝑑\varphi /(2\pi )t_{}^2`$. For a superconducting state with superconducting Green’s functions $`G(𝐤,\omega )`$ and $`F(𝐤,\omega )`$,
$`H_c^s`$ $`=4T{\displaystyle \underset{\omega }{}}{\displaystyle \underset{𝐤}{}}t_{}^2\left[G(𝐤,\omega )^2F(𝐤,\omega )^2\right]`$ (11)
$`=4T{\displaystyle \underset{\omega }{}}{\displaystyle \underset{k_z}{}}{\displaystyle \frac{d\varphi }{2\pi }t_{}^2_{\omega _c}^{\omega _c}𝑑\xi N(\xi )\frac{\omega ^2\xi ^2+\mathrm{\Delta }_𝐤^2}{(\omega ^2+\xi ^2+\mathrm{\Delta }_𝐤^2)^2}}`$
$`=4N(0){\displaystyle \underset{k_z}{}}{\displaystyle \frac{d\varphi }{2\pi }t_{}^2\frac{\omega _c}{\sqrt{\omega _c^2+\mathrm{\Delta }_𝐤^2}}\mathrm{tanh}(\frac{\sqrt{\omega _c^2+\mathrm{\Delta }_𝐤^2}}{2T})}.`$
The difference between $`H_c^s`$ and $`H_c^n`$ is of the order of $`(\mathrm{\Delta }(T)/\omega _c)^2`$; therefore, coherent $`c`$-axis coupling does not violate the FGT sum rule as long as $`\omega _c>>\mathrm{\Delta }(0)`$ even if $`t_{}`$ depends on $`\varphi `$. Note that the difference is largest at $`T=0`$ and vanishes as $`TT_c`$.
The calculations for incoherent (impurity mediated) $`c`$-axis coupling proceed in the same way as before. Note that in this case $`j_p=ied_{i\sigma }V_i\left[c_{i1\sigma }^+c_{i2\sigma }c_{i2\sigma }^+c_{i1\sigma }\right]`$ and $`j_d=e^2d^2H_c^{}A_z`$, and an impurity configuration average is required. We derive the normalized missing spectral weight $`(N_nN_s)/\rho _s`$ under assumption of a constant density of states and show that it is greater than one.
The penetration depth $`\lambda _c`$ can be calculated in two ways. Based on the Kramers-Kronig relation for the conductivity, we obtain $`\lambda _c`$, namely, $`1/4\pi \lambda _c^2=lim_{\omega 0}[\omega \text{I}\text{m}\sigma _c(0,\omega )]`$. Alternatively, using Eq.(5) we can also calculate $`\lambda _c(=1/\sqrt{\rho _s})`$. Equate these two expressions of $`\lambda _c`$, then after integration over energy we arrive the formula as follows:
$$\frac{(N_nN_s)}{\rho _s}=\frac{1}{2}+\frac{1}{2}\frac{\underset{\omega }{}𝑑\varphi _k𝑑\varphi _p|V(\varphi _k,\varphi _p)|^2\left[1\frac{\omega ^2}{\sqrt{\omega ^2+\mathrm{\Delta }_k^2}\sqrt{\omega ^2+\mathrm{\Delta }_p^2}}\right]}{_\omega 𝑑\varphi _k𝑑\varphi _p|V(\varphi _k,\varphi _p)|^2\frac{\mathrm{\Delta }_k}{\sqrt{\omega ^2+\mathrm{\Delta }_k^2}}\frac{\mathrm{\Delta }_p}{\sqrt{\omega ^2+\mathrm{\Delta }_p^2}}}.$$
(12)
The second term in Eq. (12) can easily be shown to be bigger than one half whatever the angular dependence of the impurity potential $`V(\varphi _k,\varphi _p)`$ may be. Thus the normalized missing spectral weight is always greater than one. In a simple model of impurity scattering, for which $`|V(\varphi _k,\varphi _p)|^2=|V_0|^2+|V_1|^2\mathrm{cos}2\varphi _p\mathrm{cos}2\varphi _k`$, we found that $`(N_nN_s)/\rho _s1.58`$. This incoherent coupling model, therefore, does not agree with recent findings. The sum rule is one for YBCO around optimum doping indicating coherent $`c`$-axis coupling and less than one for the underdoped case. To treat YBCO around optimum doping more realistically we need to include the complications introduced by existence of the chains along the $`b`$ axis.
Penetration depth ($`\lambda _{a(b)}`$) experiments in YBCO have shown that both $`\lambda _a`$ and $`\lambda _b`$ are linear $`T`$ at a low temperature and that a considerable amount of the condensate resides on the chains. To treat this case we need to consider a plane-chain coupling model. We assume the hybridization of Fermi surfaces between plane and chain arising through coherent coupling. For simplicity we also assume the gap in the chain has a $`d`$-wave symmetry and its magnitude is of the order of that in the plane. The Hamiltonian for a coupled plane-chain system is $`H=_𝐤\widehat{C}_𝐤^+\widehat{h}_𝐤\widehat{C}_𝐤`$, where $`\widehat{C}_𝐤^+=(C_{1𝐤}^+,C_{1𝐤},C_{2𝐤}^+,C_{2𝐤})`$ and
$$\widehat{h}_𝐤=\left(\begin{array}{cccc}\xi _{1𝐤}& \mathrm{\Delta }_{1𝐤}& t(k_z)& 0\\ \mathrm{\Delta }_{1𝐤}& \xi _{1𝐤}& 0& t(k_z)\\ t(k_z)& 0& \xi _{2𝐤}& \mathrm{\Delta }_{2𝐤}\\ 0& t(k_z)& \mathrm{\Delta }_{2𝐤}& \xi _{2𝐤}\end{array}\right)$$
(13)
where $`t(k_z)=t_0\mathrm{cos}(k_zd/2)`$ for coherent coupling between plane and chain, $`\xi _{1(2)}`$ is the energy dispersion in the plane (chain), and $`\mathrm{\Delta }_{1(2)}`$ is a gap of the plane (chain). We point out here that the conclusion we make later does not depends on the simple form of $`t(k_z)`$, and that $`\mathrm{\Delta }_{1(2)}`$ and $`\xi _{1(2)}`$ depend only on $`k_x`$ and $`k_y`$.
The Hamiltonian of the plane-chain coupling model is also decomposed into two parts, $`H=H_0+H_c`$. $`H_0`$ is for the superconductivity in the plane-chain coupling system and its eigenvalues can be reduced to $`\pm E_\pm =\pm \sqrt{ϵ_\pm ^2+\mathrm{\Delta }_𝐤^2}`$, where $`ϵ_\pm `$ are normal state energy dispersions $`ϵ_\pm =(\xi _1+\xi _2)/2\pm \sqrt{(\xi _1\xi _2)^2/4+t(k_z)^2}`$ with $`\mathrm{\Delta }_{1𝐤}=\mathrm{\Delta }_{2𝐤}=\mathrm{\Delta }_𝐤`$ for simplicity. Extensive work on this Hamiltonian can be found in Ref.
In order to calculate the linear response of the system to the external electromagnetic field, we modify $`H_c`$ with the phase factor mentioned before and follow the same procedure to derive the current $`j_c=j_p+j_d`$. Then, we obtain
$$j_p=\frac{edt_0}{2}\underset{𝐤}{}\mathrm{sin}(k_zd/2)\widehat{C}_𝐤^+\widehat{\sigma }_1\widehat{\tau }_0\widehat{C}_𝐤,$$
(14)
where $`d/2`$ is the distance between a plane and a chain and $`\widehat{\sigma }_1`$ is a Pauli matrix in the plane-chain space, and $`j_d=e^2(d/2)^2H_cA_z`$, where $`H_c=_𝐤t(k_z)\widehat{C}_𝐤^+\widehat{\sigma }_1\widehat{\tau }_3\widehat{C}_𝐤.`$ The $`c`$-axis conductivity for $`𝐪=0`$, $`\sigma _c(0,\omega )`$, of the sysyem is also derived to be $`\sigma _c(0,\omega )=(i/\omega )\left[\mathrm{\Pi }(0,\omega )e^2(d/2)^2H_c\right]`$, where
$$\mathrm{\Pi }(0,\omega )=(edt_0/2)^2T\underset{\omega ^{}}{}\underset{𝐤}{}\mathrm{sin}^2(k_zd/2)\text{T}\text{r}\left[\widehat{\sigma }_1\widehat{\tau }_0\widehat{G}(𝐤,\omega ^{})\widehat{\sigma }_1\widehat{\tau }_0\widehat{G}(𝐤,\omega ^{}+\omega )\right],$$
(15)
and
$$H_c=t_0^2T\underset{\omega }{}\underset{𝐤}{}\mathrm{cos}^2(k_zd/2)\text{T}\text{r}\left[\widehat{\sigma }_1\widehat{\tau }_3\widehat{G}(𝐤,\omega )\widehat{\sigma }_1\widehat{\tau }_3\widehat{G}(𝐤,\omega )\right],$$
(16)
with $`\widehat{G}(𝐤,\tau )=𝒯[\widehat{C}_𝐤(\tau )\widehat{C}_𝐤^+(0)]`$, which is a $`(4\times 4)`$ matrix. The Green’s function $`\widehat{G}(𝐤,\omega )`$ is given by $`\widehat{G}(𝐤,\omega )=(i\omega \widehat{h}_𝐤)^1`$. We emphasize that the Hamiltonian in this model is quite different from the usual macroscopic tunneling Hamiltonian, for which, for example, $`\widehat{G}_{13}(𝐤,\tau )=𝒯[C_{1𝐤}(\tau )C_{2𝐤}^+(0)]`$ is not allowed because each layer is independent (as is the case for the previous coupling model); however, it is possible in the present model because of the hybridization through the chain between the two Fermi surfaces of plane and chain. Introducing a unitary matrix $`U`$ which diagonalizes $`\widehat{h}_𝐤`$, one can show $`\widehat{G}_{ij}(𝐤,\omega )=_{m=1}^4U_{im}U_{mj}^+/(i\omega E_m)`$, where $`E_m=\pm E_\pm `$ if $`\mathrm{\Delta }_{1𝐤}=\mathrm{\Delta }_{2𝐤}`$.
$`H_c`$ becomes complicated and the energy dispersion in the chain is quite different from that in the plane so that a numerical calculation is required to see if the difference between $`H_c^s`$ and $`H_c^n`$ is negligible. However, since $`t_0`$ in Eq. (16) is assumed small we may expand $`H_c`$ in terms of $`t_0`$ and keep only the leading order, which is $`t_0^2`$. This case includes only interband Hamiltonian but is still very interesting as it can exhibit a $`c`$-axis pseudogap. In this approximation with $`\mathrm{\Delta }_{1𝐤}=\mathrm{\Delta }_{2𝐤}=\mathrm{\Delta }_𝐤`$ and for $`\mu ^{}=\mu `$ as a special case, $`H_c^s`$ becomes
$$H_c^s=4T\underset{\omega }{}\underset{𝐤}{}t(k_z)^2\frac{\omega ^2\xi _1\xi _2+\mathrm{\Delta }_𝐤^2}{(\omega ^2+\xi _1^2+\mathrm{\Delta }_𝐤^2)(\omega ^2+\xi _2^2+\mathrm{\Delta }_𝐤^2)}.$$
(17)
Note that $`H_c^s`$ in Eq. (17) is almost same as $`H_c^s`$ in Eq. (11) for the simple coherent coupling case except that now $`\xi _1\xi _2`$ and $`d/2`$ appears rather than $`d`$ in $`t(k_z)`$. One, therefore, may expect that $`\delta H_c`$ will vanish to order $`(\mathrm{\Delta }(T)/\omega _c)^2`$. It is obvious that $`\delta H_c`$ is identically zero along the nodal lines, and $`\delta H_c`$ is largest along the anti-nodal directions. Since $`\xi _1=k^2/2m\mu `$ and $`\xi _2=k_y^2/2m\mu `$, we introduce $`\xi =\xi _1`$, then $`\xi _2=\xi \mathrm{sin}(\varphi )^2\mu \mathrm{cos}(\varphi )^2`$. If $`\varphi =\pi /2`$, then $`\xi _2=\xi `$; therefore, it can be seen that $`\delta H_c_{\varphi =\pi /2}`$ is of the order of $`(\mathrm{\Delta }/\omega _c)^2`$. For $`\varphi =0`$, $`\xi _2=\mu `$ and it can be shown that
$`H_c_{\varphi =0}^s`$ $`2N(0){\displaystyle \underset{k_z}{}}t_{}^2{\displaystyle _{\omega _c}^{\omega _c}}d\xi [{\displaystyle \frac{\mu }{\xi ^2\mu ^2}}\mathrm{tanh}({\displaystyle \frac{\mu }{2T}})`$ (19)
$`{\displaystyle \frac{\xi ^2}{(\xi ^2\mu ^2)\sqrt{\xi ^2+\mathrm{\Delta }^2}}}\mathrm{tanh}({\displaystyle \frac{\sqrt{\xi ^2+\mathrm{\Delta }^2}}{2T}})].`$
Now, the leading order of $`\delta H_c_{\varphi =0}`$ changes to $`(\mathrm{\Delta }/\mu )^2`$. It is possible to show that for an arbitrary $`\varphi `$, as long as $`\mu `$ and $`\omega _c>>\mathrm{\Delta }`$, $`\delta H_c_\varphi `$ is negligible, and consequently, the FGT sum rule is not violated in the plane-chain coupling model.
In a numerical calculation for the general case without the above simplification, we have computed $`\delta H_c/H_c^s`$, which is the fractional change in kinetic energy. We have taken $`\xi _2=k_y^2/2m\mu ^{}`$ for the chain energy dispersion, where $`\mu ^{}`$ is a chemical potential in the chain. For simplicity, we also assume $`\mu ^{}=\mu `$. It has been shown that $`\mu ^{}\mu (<<\mu )`$ may correspond to the pseudogap seen in the $`c`$-axis response of over-doped YBCO; however, it makes no difference in the numerical evaluation of $`\delta H_c/H_c^s`$. We choose $`T=12K`$, $`\mathrm{\Delta }(T)=20`$meV, $`t_0=2`$meV, $`\mu =500`$ meV and $`\omega _c=400`$meV. We found that $`\delta H_c/H_c^s`$ becomes more negligible as we increase the summation range of the Matsubara frequency $`\omega `$. For $`|\omega |200\pi T`$, $`\delta H_c/H_c^s8.8\times 10^3`$, and for $`|\omega |2000\pi T`$, $`\delta H_c/H_c^s5.4\times 10^3`$.
One may consider a plane-plane coupling through a chain. In order to investigate the $`c`$-axis kinetic energy for such a coupling, one needs to replace $`\xi _2`$ and $`\mathrm{\Delta }_2`$ with $`\xi _1`$ and $`\mathrm{\Delta }_1`$, respectively. For the hopping amplitude, $`t(k_z)`$ can be simply changed to $`t(k_z)^2`$ because the plane-chain and chain-plane distances are the same and equal to $`d/2`$. Then, one can algebraically show that $`\delta H_c`$ is as negligible as before. It is also possible to see that $`\delta H_c`$ has a symmetry with respect to $`\xi _1\xi _2`$ and $`\mathrm{\Delta }_1\mathrm{\Delta }_2`$; in other words, $`\delta H_c`$ for the chain-plane coupling is the same as that for the plane-chain coupling. Therefore, it implies that $`\delta H_c`$ along the $`c`$ axis is conserved for coherent coupling.
So far we have taken the density of states as a constant: $`N(\xi )=N(0)`$ $`(\omega _c\xi \omega _c)`$ and concluded that the difference of the $`c`$-axis electronic kinetic energies between normal and supercondcuting state is negligible. Now we would like to consider the effect of $`N(\xi )`$ on the sum rule when it is a function of $`\xi `$ to illustrate possible changes. If it varies strongly with $`\xi `$, it clearly cannot be approximated by $`N(0)`$. We taylor-expand $`N(\xi )`$ up to $`\xi ^2`$ near $`\xi =0`$. Then, $`N(\xi )=N(0)+\xi N(\xi )/\xi |_0+(\xi /2)^2^2N(\xi )/\xi ^2|_0`$. At $`T=0`$, $`H_c^n`$ of Eq. (8) does not change; however, Eq. (11) for $`H_c^s`$ changes due to $`(\xi /2)^2N\mathrm{"}(0)`$, where $`N\mathrm{"}(0)=^2N(\xi )/\xi ^2|_0`$. Assuming $`t_{}`$ in Eqs. (8) and (11) does not depend on $`\varphi `$, we obtain
$$\delta H_c/H_c^n(8N(0))^1N\mathrm{"}(0)\mathrm{\Delta }(0)^2\mathrm{ln}(\omega _c/\mathrm{\Delta }(0)).$$
(20)
Note that this correction can have either sign depending on the sign of the second derivative. For $`N\mathrm{"}(0)/N(0)1/\omega _c^2`$, $`\delta H_c/H_c^nx^{1.65}`$, where $`x=\omega _c/\mathrm{\Delta }(0)`$, because $`\mathrm{ln}(x)/x^2x^{1.65}`$ when $`x>>1`$. If $`N\mathrm{"}(0)/N(0)1/(\mathrm{\Delta }(0)\omega _c)`$, then $`\delta H_c/H_c^nx^{0.65}`$; thus, $`\delta H_c`$ is considerable. For this to be the case $`N\mathrm{"}(0)`$ needs to exhibit variation on an energy scale of order $`\mathrm{\Delta }`$ rather than $`\omega _c`$. In a realistic model a Taylor expansion about $`\xi =0`$ may not be accurate but our calculations serve to illustrate the main point. Violation of the FGT sum rule of either sign can result from an energy dependence in the in-plane electronic density of states $`N(\xi )`$. The exact amount depends on details and cannot be known without a specific knowledge of the band structure involved. In-plane dynamics gets reflected in $`c`$-axis properties.
For coherent interlayer coupling between CuO<sub>2</sub> planes the superfluid density is equal to the missing optical weight; the FGT sum rule is satisfied. This applies even in more realistic model for YBCO around optimum doping such as the plane-chain model with two atoms per unit cell. On the other hand incoherent $`c`$-axis coupling mediated through impurity scattering gives a sum rule which is always larger than one and is in disagreement with experiment. To get the sum rule to be less than one as is observed in underdoped YBCO and other systems such as optimally doped Tl2201, it may be necessary to go to more exotic non-Fermi liquid pseudogap model for the in-plane motion as discussed recently by Ioffe and Millis. Their arguments however do not apply to optimally doped Tl2201 because this system does not show a pseudogap. Their pseudogap argument that leads to the cancellation of $`G(𝐤,\omega )`$ and $`G_0(𝐤,\omega )`$ contribution to the ratio of missing area to superfluid density making it one half instead of one for the preformed pair model was made for coherent c-axis coupling, but we find it also applies to the incoherent case. Another interesting model for the in-plane dynamics is the ”mode” model of Norman et al. introduced from consideration of ARPES data. In more conventional models a sum rule violation of either sign can also be obtained if there is a strong energy dependence to the density of states near the Fermi surface on the scale of a few times the gap.
W.K. is grateful to N. D. Whelan and C. Kallin for useful discussions and J.P.C. to D. Basov for discussions. This work was supported in part by the Natural Sciences and Engineering Research Council of Canada (NSERC) and by the Canadian Institute for Advanced Research (CIAR). |
warning/0003/gr-qc0003054.html | ar5iv | text | # The binary black-hole dynamics at the third post-Newtonian order in the orbital motion
## I Introduction
The motion of binary systems is one of the most important problems in general relativity, particularly in view of the future detection of gravitational waves from such systems. The simplest two-body problem is the one where the components of the system are objects which are as spherically symmetric and point-like as possible. In general relativity, those objects are black holes. In spite of the extension of black holes, in binary orbit approximate calculations of post-Newtonian (PN) type (which use expansions in powers of $`1/c`$, where $`c`$ denotes the velocity of light), the ansatz of Dirac delta functions to describe point-like sources turned out to be very successful up to the 2.5PN approximation , and at the 3.5PN order of approximation too . At the 3PN order, i.e. at the order $`(1/c^2)^3`$, only partial success was achievable . Two terms in the binary point-mass calculations came out ambiguously: a kinetic one, depending on the bodies’ linear momenta squared , and a static one, not depending on the momenta . The authors were able to show that the static ambiguity is intimately related to the fact that different exact solutions (Brill-Lindquist and Misner-Lindquist) of the time-symmetric and conformally flat initial value problem for two black holes do exist, which at the time reveal the emergence of black holes in binary point-mass calculations at the 3PN order . In this paper we confront the kinetic ambiguity with the center-of-mass motion and claim that it can be fixed through the standard relation between the center-of-mass velocity and the total linear momentum which results from global Lorentz invariance.
We employ the following notation: $`𝐱_a=\left(x_a^i\right)`$ ($`a=1,2`$; $`i=1,2,3`$) denotes the position of the $`a`$th point mass in the 3-dimensional Euclidean space endowed with a standard Euclidean scalar product (denoted by a dot). We also define $`𝐫_{12}𝐱_1𝐱_2`$, $`r_{12}|𝐫_{12}|`$, $`𝐧_{12}𝐫_{12}/r_{12}`$; $`||`$ stands here for the Euclidean length of a vector. The linear momentum of the $`a`$th body is denoted by $`𝐩_a=\left(p_{ai}\right)`$.
## II The generalized conservative Hamiltonian to 3PN order
The post-Newtonian approximation order 2.5PN (as well 3.5PN) is of purely dissipative character. Subtraction of the 2.5PN terms from the 3PN Hamiltonian results in the conservative 3PN Hamiltonian. This Hamiltonian is primarily of higher order as it depends on time derivatives of the positions and momenta of the bodies .
The conservative 3PN higher-order Hamiltonian for two-body point-mass systems was calculated in Ref. . We present here the explicit form of this Hamiltonian in the center-of-mass reference frame (where $`𝐩_1+𝐩_2=0`$). It is convenient to introduce the following reduced variables
$$𝐫\frac{𝐱_1𝐱_2}{GM},𝐩\frac{𝐩_1}{\mu }=\frac{𝐩_2}{\mu },\widehat{t}\frac{t}{GM},\widehat{H}^{\mathrm{NR}}\frac{H^{\mathrm{NR}}}{\mu },$$
(1)
where
$$\mu \frac{m_1m_2}{M},Mm_1+m_2,\nu \frac{\mu }{M}=\frac{m_1m_2}{(m_1+m_2)^2}.$$
(2)
In Eq. (1) the superscript NR denotes a ‘non-relativistic’ Hamiltonian, i.e. the Hamiltonian without the rest-mass contribution $`Mc^2`$. We also introduce the reduced-time derivatives:
$$\dot{𝐫}\frac{d𝐫}{d\widehat{t}}=\frac{d(𝐱_1𝐱_2)}{dt},\dot{𝐩}\frac{d𝐩}{d\widehat{t}}=\frac{G}{\nu }\frac{d𝐩_1}{dt}=\frac{G}{\nu }\frac{d𝐩_2}{dt}.$$
(3)
The reduced two-point-mass conservative Hamiltonian up to the 3PN order reads
$`\widehat{H}^{\text{NR}}(𝐫,𝐩,\dot{𝐫},\dot{𝐩})`$ $`=`$ $`\widehat{H}_\mathrm{N}(𝐫,𝐩)+{\displaystyle \frac{1}{c^2}}\widehat{H}_{1\mathrm{P}\mathrm{N}}(𝐫,𝐩)`$ (5)
$`+{\displaystyle \frac{1}{c^4}}\widehat{H}_{2\mathrm{P}\mathrm{N}}(𝐫,𝐩)+{\displaystyle \frac{1}{c^6}}\widehat{H}_{3\mathrm{P}\mathrm{N}}(𝐫,𝐩,\dot{𝐫},\dot{𝐩}),`$
where (here $`r|𝐫|`$ and $`𝐧𝐫/r`$)
$$\widehat{H}_\mathrm{N}(𝐫,𝐩)=\frac{𝐩^2}{2}\frac{1}{r},$$
(6)
$$\widehat{H}_{1\mathrm{P}\mathrm{N}}(𝐫,𝐩)=\frac{1}{8}(3\nu 1)(𝐩^2)^2\frac{1}{2}\left[(3+\nu )𝐩^2+\nu (𝐧𝐩)^2\right]\frac{1}{r}+\frac{1}{2r^2},$$
(7)
$`\widehat{H}_{2\mathrm{P}\mathrm{N}}(𝐫,𝐩)`$ $`=`$ $`{\displaystyle \frac{1}{16}}\left(15\nu +5\nu ^2\right)(𝐩^2)^3`$ (10)
$`+{\displaystyle \frac{1}{8}}\left[\left(520\nu 3\nu ^2\right)(𝐩^2)^22\nu ^2(𝐧𝐩)^2𝐩^23\nu ^2(𝐧𝐩)^4\right]{\displaystyle \frac{1}{r}}`$
$`+{\displaystyle \frac{1}{2}}\left[(5+8\nu )𝐩^2+3\nu (𝐧𝐩)^2\right]{\displaystyle \frac{1}{r^2}}{\displaystyle \frac{1}{4}}(1+3\nu ){\displaystyle \frac{1}{r^3}},`$
$`\widehat{H}_{3\mathrm{P}\mathrm{N}}(𝐫,𝐩,\dot{𝐫},\dot{𝐩})`$ $`=`$ $`{\displaystyle \frac{1}{128}}\left(5+35\nu 70\nu ^2+35\nu ^3\right)(𝐩^2)^4`$ (19)
$`+{\displaystyle \frac{1}{16}}\{(7+42\nu 53\nu ^26\nu ^3)(𝐩^2)^3`$
$`+(12\nu )\nu ^2[2(𝐧𝐩)^2(𝐩^2)^2+3(𝐧𝐩)^4𝐩^2]\}{\displaystyle \frac{1}{r}}`$
$`+{\displaystyle \frac{1}{48}}[3(27+140\nu +96\nu ^2)(𝐩^2)^2`$
$`+6(8+25\nu )\nu (𝐧𝐩)^2𝐩^2(35267\nu )\nu (𝐧𝐩)^4]{\displaystyle \frac{1}{r^2}}`$
$`+{\displaystyle \frac{1}{1536}}\{[48003(8944315\pi ^2)\nu 7808\nu ^2]𝐩^2`$
$`+9(2672315\pi ^2+448\nu )\nu (𝐧𝐩)^2\}{\displaystyle \frac{1}{r^3}}`$
$`+{\displaystyle \frac{1}{96}}\left[12+(87263\pi ^2)\nu \right]{\displaystyle \frac{1}{r^4}}`$
$`+\widehat{D}(𝐫,𝐩,\dot{𝐫},\dot{𝐩})+\widehat{\mathrm{\Omega }}(𝐫,𝐩).`$
In Eq. (19) $`\widehat{D}(𝐫,𝐩,\dot{𝐫},\dot{𝐩})`$ denotes that part of the 3PN Hamiltonian $`\widehat{H}_{3\mathrm{P}\mathrm{N}}`$ which depends on the time derivatives $`\dot{𝐫}`$ and $`\dot{𝐩}`$. Its explicit form reads
$$\widehat{D}=\widehat{D}_1r+\widehat{D}_0+\widehat{D}_1\frac{1}{r}+\widehat{D}_2\frac{1}{r^2}+\widehat{D}_3\frac{1}{r^3},$$
(20)
where
$`\widehat{D}_1`$ $`=`$ $`{\displaystyle \frac{1}{12}}\nu ^2[4(𝐧𝐩)(𝐧\dot{𝐩})(𝐩\dot{𝐩})5(𝐧\dot{𝐩})^2𝐩^2`$ (22)
$`5(𝐧𝐩)^2\dot{𝐩}^2(𝐧𝐩)^2(𝐧\dot{𝐩})^2+13𝐩^2\dot{𝐩}^2+2(𝐩\dot{𝐩})^2],`$
$`\widehat{D}_0`$ $`=`$ $`{\displaystyle \frac{1}{8}}\nu ^2\{(𝐧\dot{𝐫})(𝐩\dot{𝐩})[5𝐩^2+(𝐧𝐩)^2]`$ (25)
$`(𝐧𝐩)\left[𝐩^2(\dot{𝐩}\dot{𝐫})+2(𝐩\dot{𝐩})(𝐩\dot{𝐫})\right]{\displaystyle \frac{1}{3}}(𝐧𝐩)^3(\dot{𝐩}\dot{𝐫})`$
$`+(𝐧\dot{𝐩})[𝐩^2+(𝐧𝐩)^2][(𝐧𝐩)(𝐧\dot{𝐫})(𝐩\dot{𝐫})]\},`$
$`\widehat{D}_1`$ $`=`$ $`{\displaystyle \frac{1}{24}}\nu \{2(1710\nu )(𝐧𝐩)(𝐧\dot{𝐩})(𝐧\dot{𝐫})(15+22\nu )(𝐧𝐩)^2(𝐧\dot{𝐩})`$ (32)
$`(518\nu )(𝐧\dot{𝐩})𝐩^22(65\nu )(𝐧𝐩)(𝐩\dot{𝐩})`$
$`2(12\nu )(𝐧\dot{𝐫})(𝐩\dot{𝐩})2(72\nu )[(𝐧𝐩)(\dot{𝐩}\dot{𝐫})+(𝐧\dot{𝐩})(𝐩\dot{𝐫})]\}`$
$`+{\displaystyle \frac{1}{16}}\nu ^3[8(𝐧𝐩)(𝐧\dot{𝐫})𝐩^2(𝐩\dot{𝐫})+8(𝐧𝐩)^3(𝐧\dot{𝐫})(𝐩\dot{𝐫})`$
$`+2(𝐧𝐩)^2𝐩^2\dot{𝐫}^2+5(𝐩^2)^2\dot{𝐫}^2+(𝐧𝐩)^4\dot{𝐫}^25(𝐧\dot{𝐫})^2(𝐩^2)^2`$
$`6(𝐧𝐩)^2(𝐧\dot{𝐫})^2𝐩^25(𝐧𝐩)^4(𝐧\dot{𝐫})^2`$
$`4(𝐧𝐩)^2(\dot{𝐩}\dot{𝐫})^24𝐩^2(\dot{𝐩}\dot{𝐫})^2],`$
$`\widehat{D}_2`$ $`=`$ $`{\displaystyle \frac{1}{48}}\nu [5(57\nu )(𝐧𝐩)^3(𝐧\dot{𝐫})+10(35\nu )(𝐧𝐩)^2(𝐧\dot{𝐫})^2`$ (37)
$`+3(1735\nu )(𝐧𝐩)(𝐧\dot{𝐫})𝐩^228(38\nu )(𝐧𝐩)(𝐧\dot{𝐫})(𝐩\dot{𝐫})`$
$`15(23\nu )(𝐧𝐩)^2(𝐩\dot{𝐫})+2(2477\nu )(𝐧𝐩)^2\dot{𝐫}^2`$
$`+2(929\nu )(𝐧\dot{𝐫})^2𝐩^23(49\nu )𝐩^2(𝐩\dot{𝐫})`$
$`2(1237\nu )𝐩^2\dot{𝐫}^2+4(617\nu )(𝐩\dot{𝐫})^2],`$
$`\widehat{D}_3`$ $`=`$ $`{\displaystyle \frac{1}{1536}}\nu \{3[927\pi ^210832+36(485\pi ^2)\nu ](𝐧𝐩)(𝐧\dot{𝐫})`$ (41)
$`6\left[3(16+\pi ^2)2(45\pi ^2464)\nu \right](𝐧\dot{𝐫})^2`$
$`+\left[11600927\pi ^2+20(9\pi ^280)\nu \right](𝐩\dot{𝐫})`$
$`+2[176+3\pi ^2+6(17615\pi ^2)\nu ]\dot{𝐫}^2\}.`$
The last term in Eq. (19), $`\widehat{\mathrm{\Omega }}(𝐫,𝐩)`$, contains the parameters which parametrize the ambiguities of the 3PN Hamiltonian. Its explicit form is
$$\widehat{\mathrm{\Omega }}(𝐫,𝐩)=\omega _{\text{kinetic}}\left[𝐩^23(𝐧𝐩)^2\right]\frac{\nu ^2}{r^3}+\omega _{\text{static}}\frac{\nu }{r^4},$$
(42)
where $`\omega _{\text{kinetic}}`$ parametrizes the kinetic ambiguity, and the static ambiguity is parametrized by $`\omega _{\text{static}}`$.
The 3PN higher-order Hamiltonian $`\widehat{H}_{3\mathrm{P}\mathrm{N}}(𝐫,𝐩,\dot{𝐫},\dot{𝐩})`$, Eq. (19), can be reduced to a usual Hamiltonian depending only on $`𝐫`$ and $`𝐩`$. Details of the reduction can be found in Sec. II of Ref. .
## III The ambiguous part of the 3PN Hamiltonian
It was found in , that the ambiguous part $`\mathrm{\Omega }`$ of the 3PN Hamiltonian, Eq. (42), takes, in a non-center-of-mass reference frame and in non-reduced variables, the form ($`G`$, the Newtonian gravitational constant, is put equal to one)
$`\mathrm{\Omega }(𝐱_1,𝐱_2,𝐩_1,𝐩_2)`$ $`=`$ $`{\displaystyle \frac{1}{c^6}}{\displaystyle \underset{a}{}}{\displaystyle \underset{ba}{}}\{\omega _{\text{static}}{\displaystyle \frac{m_a^3m_b^2}{r_{ab}^4}}`$ (44)
$`+{\displaystyle \frac{1}{2}}\omega _{\text{kinetic}}{\displaystyle \frac{m_am_b}{r_{ab}^3}}[(𝐩_a)^23(𝐧_{ab}𝐩_a)^2]\}.`$
In the following we shall show how center-of-mass considerations can impose restrictions onto the “kinetic” ambiguous term. To make contact with the intended center-of-mass motion calculations we transform (44) to the Lagrangean level. On this level, $`\mathrm{\Omega }`$ appears as $`\mathrm{\Omega }`$ with
$`\mathrm{\Omega }(𝐱_1,𝐱_2,𝐯_1,𝐯_2)`$ $`=`$ $`{\displaystyle \frac{1}{c^6}}{\displaystyle \underset{a}{}}{\displaystyle \underset{ba}{}}\{\omega _{\text{static}}{\displaystyle \frac{m_a^3m_b^2}{r_{ab}^4}}`$ (46)
$`+{\displaystyle \frac{1}{2}}\omega _{\text{kinetic}}{\displaystyle \frac{m_a^3m_b}{r_{ab}^3}}[(𝐯_a)^23(𝐧_{ab}𝐯_a)^2]\},`$
where the bodies’ velocities are given by $`𝐯_a=𝐩_a/m_a`$.
In view of the relation between the total linear momentum $`𝐏`$ and the center-of-mass velocity $`𝐕=\dot{𝐗}`$ ($`𝐗`$ denotes the center-of-mass coordinate) demanded by global Lorentz invariance, ,
$$𝐏=\frac{H}{c^2}\frac{d𝐗}{dt},𝐏=𝐩_1+𝐩_2,$$
(47)
where $`H`$ is the conserved total energy, $`H=Mc^2+H^{\text{NR}}`$, the following expression (which is one of the contributions to $`𝐏=L/𝐯_1+L/𝐯_2`$)
$`𝐃{\displaystyle \frac{\mathrm{\Omega }}{𝐯_1}}{\displaystyle \frac{\mathrm{\Omega }}{𝐯_2}}=\omega _{\text{kinetic}}{\displaystyle \frac{1}{c^6}}{\displaystyle \underset{a}{}}{\displaystyle \underset{ba}{}}{\displaystyle \frac{m_a^3m_b}{r_{ab}^3}}\left[𝐯_a3(𝐧_{ab}𝐯_a)𝐧_{ab}\right]`$ (48)
has either to be a total time derivative (in which case no restrictions are imposed on $`\mathrm{\Omega }`$) or to be fixed by other terms in Eq. (47) (note that no static ambiguity is involved here).
The simplest way to show that $`𝐃`$ is not a total time derivative is to write it in the form
$$D^i=\omega _{\text{kinetic}}m_1m_2[(m_1^2+m_2^2)V^j+\frac{m_1m_2(m_1m_2)}{M}v^j]_{1j}_{1i}\frac{1}{r_{12}},$$
(49)
($`v^i`$ denotes the relative velocity, $`v^i=v_1^iv_2^i`$, and here, $`MV^i=m_1v_1^i+m_2v_2^i`$ holds) and to average $`D^i`$ over a circular orbit of the relative motion simply assuming that the center-of-mass motion is orthogonal to the relative motion. The result is (notice that the relative motion part in the Eq. (49) is a total time derivative)
$$<𝐃>=\omega _{\text{kinetic}}\frac{m_1m_2}{r_{12}^3}(m_1^2+m_2^2)𝐕.$$
(50)
In the center-of-mass frame $`D^i`$ is always a total time derivative. This fits with Ref. where the kinetic ambiguity, in the center-of-mass frame, was shown to be influenced by coordinate transformations which additionally only influence, unimportant in this context, the static potential at 3PN order.
## IV Conclusions
In this paper we have shown that the kinetic ambiguity obtained in binary point-mass calculations at the 3PN order should be getting fixed by center-of-mass motion considerations. The static ambiguity was known to be fixable by referring to e.g., the Brill-Lindquist ($`\omega _{\text{static}}=0`$) or the Misner-Lindquist solution ($`\omega _{\text{static}}=1/8`$). The energy content of the latter solution is influenced by the topology of the involved non-simply connected 3-space.
G.S. wishes to acknowledge useful discussions with Luc Blanchet about Lorentz invariance. The authors also thank Thibault Damour for constructive remarks. This work was supported in part by the Max-Planck-Gesellschaft Grant No. 02160-361-TG74 (G.S.) and the KBN Grant No. 2 P03B 094 17 (P.J.). |
warning/0003/astro-ph0003129.html | ar5iv | text | # Non-Linear Gravitational growth of large scale structures inside and outside standard Cosmology
## 1. Introduction
In Cosmology the standard picture of gravitational growth, and also many aspects of fundamental physics, are extrapolated many orders of magnitude, from the scales and times where our current theory of gravity (General Relativity, GR) has been experimentally tested, into the distant universe. In particular, current limits on the (parametrized) Post Newtonian formalism mostly restrict to our very local Universe (see Will 1993). It is important to evaluate how much our predictions and cosmological picture depend on the underlying hypothesis (see Peebles 1999 for insightful comments on the state of this subject). The other side of this argument is that cosmology can be used to test fundamental physics, such as our theory of gravity.
One aspect of GR that could be questioned or tested without modifying the basic structure or symmetry of the theory are Einstein’s field equations, relating the energy content ($`T_{\mu \nu }`$) to the curvature ($`R_{\mu \nu }`$). One such modification, which will be considered here, is scalar-tensor theories (STT), such as Brans-Dicke (BD) theory. A more generic, but also more vague, way of testing the importance of Einstein’s field equations is to model independently the geometry and the matter content, thus allowing for the possibility of other relations between them. Some simple aspects of this idea will be illustrated here by studying structure formation in a flat, matter dominated universe but with a more general growth law for the Hubble rate —see section 3.2 below. Similarly, we will also consider results for a generic equation of state: $`p=\gamma \rho `$, where $`\gamma `$ can be chosen independently of the cosmological parameters ($`\mathrm{\Omega }_M`$, $`\mathrm{\Omega }_k`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$).
Our aim in this paper is to explore certain variations of the standard model to see how they affect structure formation. The idea is to find a way to parameterize variations from GR that might produce differences large enough to be observable. The variations considered could have other observable consequences (eg in the local universe or in the radiation dominated regime) which might rule them out as a viable new theory. But even if this were the case, we still would have learn something about how structure formation depends on the underlying theory of Gravity or the assumptions about the equation of state. This aspect of the theory has hardly been explored and it therefore represents an important step forward in analyzing alternatives to the current paradigm, eg non-baryonic matter (see Peebles 1999), and could also help to set limits on variations of GR or the equation of state at high red-shifts.
Here we consider two main regimes for structure formation in non-standard gravity/cosmology: weakly non-linear and strongly non-linear large scale clustering. We study the shear-free or spherical collapse (SC) model, which corresponds to the spherically symmetric (or local) dynamics (see below). This approximation works very well at least in two different contexts, that will be explored here.
The first one is the growth of the smoothed 1-point cumulants of the probability distribution for large scale density fluctuations: the SC model turns out to reproduce exactly the leading order perturbation theory predictions (Bernardeau 1992), and turns out to be an excellent approximation for the exact dynamics as compared to N-body simulations both with Gaussian (Fosalba & Gaztañaga 1998a, 1998b) and non-Gaussian initial conditions (Gaztañaga & Fosalba 1998). The measured 1-point cumulants in galaxy catalogues have been compared with these predictions (eg Bouchet et al. 1993, Gaztañaga 1992,1994, Gaztañaga & Frieman 1994, Baugh, Gaztañaga & Efstathiou 1995, Gaztañaga 1995, Baugh & Gaztañaga 1996, Colombi etal 1997, Hui & Gaztañaga 1999).
The second one is the study of the epoch of formation and abundance of structures (such as galaxies and clusters), using the Press & Schechter (1974) formalism and its extensions (eg Bond et al. 1991, Lacey & Cole 1993, Sheth & Lemson 1999, Scoccimarro et al. 2000). Given some Gaussian initial conditions, this formalism can predict the number of structures (halos) of a given mass that will form at each stage of the evolution. One can use the SC model to predict the value of the critical linear over-density, $`\delta _c`$, that will collapse into virialized halos. It turns out that the analytical predictions for the halo mass function and formation rates are remarkably accurate as compared to N-body simulations (Lacey & Cole 1994). One can also use this type of modeling to predict clustering properties of halos (eg Mo & White 1996, Mo, Jing & White 1997), cluster abundances (White, Efstathiou & Frenk 1993, Bahcall & Fan 1998) or weak lensing through mass functions (Jain & Van Waerbeke 2000). The observed cluster abundances have been used as a strong discriminant for cosmological models and also as a way to measure the amplitude of mass fluctuations, $`\sigma _8`$ (see White, Efstathiou & Frenk 1993, Bahcall & Fan 1998).
In summary, we propose to address a very specific question here: how different are the above non-linear predictions when using a non-standard cosmology and non-standard theory of Gravity? To answer this question we will consider two non-standard variations: scalar-tensor models and some examples of a cosmology that do not obey Einstein’s field equations. The paper is organized as follows: In §2 we give a summary on how non-linear structure formation relates to the underlying theory of Gravity (see Weinberg 1972, Peebles 1993, Ellis 1999 and references therein, for a review on the relation between gravitational theory and cosmology). This section covers old ground with some detail as an introduction to later sections and for the reader that is not familiar with this subject or notation. We also present the more general case of an ideal (relativistic) fluid. As far as we know, some of the non-linear results presented here are new. In §3 we show how these predictions change in the two examples of non-standard gravity. Observational consequences are explored in §4. In §5 we present a discussion and the conclusions.
## 2. Gravitational Growth inside GR
The self-gravity of an over-dense region work against the expansion of the universe so that this region will expand at a slower rate that the background. This increases the density contrast so that eventually the region collapses. The details of this collapse depends on the initial density profile. Here we will focus in the spherically symmetric case. We will revise non-linear structure growth in the context of the fluid limit and the shear-free approximation. These turns out to be very good approximation for the applications that will be considered later (leading order and strongly non-linear statistics).
We start with the Raychaudhuri’s equation, which is valid for an arbitrary Ricci tensor $`R_{\mu \nu }`$. We use Einstein’s field equations and the continuity equation to turn Raychaudhuri’s equation into a second order differential equation for the density contrast. We first present the matter dominated (non-relativistic) case, with solutions for the linear and non-linear regimes. Later, in §2.5, we assess the more generic case of an ideal (relativistic) fluid and its corresponding solution.
### 2.1. Einstein’s and Raychaudhuri’s Equations
We start recalling that the metric tensor $`g_{\mu \nu }`$ defines the line element of space-time:
$$ds^2=g_{\mu \nu }dx^\mu dx^\nu $$
(1)
which in the homogeneous and isotropic model of the cosmological principle can be written as (see eg Weinberg 1972):
$$ds^2=dt^2a^2(t)\left[\frac{dr^2}{1+kr^2}+r^2\left(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2\right)\right]$$
(2)
As usual we will work in comoving coordinates x related to physical coordinates by $`\text{r}_p=a(t)\text{x}`$, where $`a(t)=(1+z)^1`$ is the cosmic scale factor, and $`z`$ the corresponding red-shift ($`a_01`$). Thus all geometrical aspects of this universal line element are determined up to the function $`a(t)`$ and the arbitrary constant $`k`$, which defines the usual open, Einstein-deSitter and closed universes. The function $`a(t)`$ can be found for each energy content by solving the corresponding equations of motion, eg the gravitational field equations.
In this section we consider Einstein’s equations:
$$R_{\mu \nu }+\mathrm{\Lambda }g_{\mu \nu }=8\pi G\left(T_{\mu \nu }\frac{1}{2}g_{\mu \nu }T\right)$$
(3)
where $`Tg^{\mu \nu }T_{\mu \nu }`$ is the trace of the energy-momentum tensor; we have included a cosmological constant term to keep the equations general at this stage. For an ideal fluid, we have:
$$T_{\mu \nu }=pg_{\mu \nu }+(p+\rho )u_\mu u_\nu $$
(4)
We can now use the field equations and the above energy-momentum to find the scale factor $`a(t)`$ in the metric:
$`{\displaystyle \frac{3\ddot{a}}{a}}`$ $`=`$ $`4\pi G\rho \left(1+{\displaystyle \frac{3p}{\rho }}\right)+\mathrm{\Lambda }`$ (5)
$`H^2{\displaystyle \frac{\dot{a}^2}{a^2}}`$ $`=`$ $`{\displaystyle \frac{8\pi G\rho }{3}}+{\displaystyle \frac{k}{a^2}}+{\displaystyle \frac{\mathrm{\Lambda }}{3}},\dot{}{\displaystyle \frac{d}{dt}}`$ (6)
In the fluid approximation, deviations from the mean background $`\overline{\rho }`$ are characterized by fluctuations in the density and velocity fields. The continuity equation for a non-relativistic fluid is (Peebles 1993):
$$\frac{}{\tau }\delta (\text{x},\tau )+\{\left[1+\delta (\text{x},\tau )\right]\text{v}(\text{x},\tau )\}=0$$
(7)
where $`\delta (\text{x},\tau )\rho (\text{x},\tau )/\overline{\rho }1`$ is the local density contrast, $`\text{v}(\text{x},\tau )`$ the peculiar velocity (see Eq. below), and $`\tau `$ the conformal time defined by
$$d\tau =\frac{dt}{a(t)}\frac{d}{dt}=\frac{1}{a}\frac{d}{d\tau }$$
(8)
The continuity equation can also be written
$$\frac{d\delta }{d\tau }+(1+\delta )\theta =0,\theta \text{v}$$
(9)
In order to find an equation of motion for the density contrast alone we shall resort to the Raychaudhuri equation (see eg Wald 1984)
$$\frac{d\mathrm{\Theta }}{ds}+\frac{1}{3}\mathrm{\Theta }^2=\sigma _{ij}\sigma ^{ij}+\omega _{ij}\omega ^{ij}+R_{\mu \nu }u^\mu u^\nu $$
(10)
where $`\mathrm{\Theta }_\mu u^\mu `$, $`\sigma _{ij}`$ is the shear tensor, $`\omega _{ij}`$ the vorticity tensor, and $`R_{\mu \nu }`$ the Ricci tensor, and $`s`$ the proper time parameter; $`u^\mu `$ is the fluid’s 4-velocity, $`u^0=1`$, and
$$\text{u}=\dot{a}(t)\text{x}+\text{v}(\text{x},t)$$
(11)
It is important to stress that Raychaudhuri’s equation, Eq. , is purely geometric: it describes the evolution in proper time of the dilatation coefficient $`\mathrm{\Theta }`$ of a bundle of nearby geodesics. There is no physics in this equation until a relationship between $`R_{\mu \nu }`$ and the matter contents of the universe is specified by means of a set of field equations. This makes it very useful for our purposes in this paper, as we shall later make reference to a different set of field equations.
If Einstein’s field equations, Eq. and , are assumed then it is readily verified that
$$R_{\mu \nu }u^\mu u^\nu =4\pi G\rho \left(1+\frac{3p}{\rho }\right)+\mathrm{\Lambda }$$
(12)
### 2.2. Shear free and matter domination
In a matter dominated regime ($`p=0`$), $`\rho a^3`$. Equation for the Hubble rate $`H`$, can be rewritten using the notation: $`\mathrm{\Omega }_M8\pi G\rho _0/(3H_0^2)`$, which is the ratio of the current matter density to the critical density, $`\mathrm{\Omega }_k=k/H_0^2`$ gives the global curvature, and $`\mathrm{\Omega }_\mathrm{\Lambda }=\mathrm{\Lambda }/(3H_0^2)`$ where $`\mathrm{\Lambda }`$ is the cosmological constant, so that $`\mathrm{\Omega }_M+\mathrm{\Omega }_k+\mathrm{\Omega }_\mathrm{\Lambda }=1`$:
$$H^2(z)=H_0^2\left[\mathrm{\Omega }_M(1+z)^3+\mathrm{\Omega }_k(1+z)^2+\mathrm{\Omega }_\mathrm{\Lambda }\right]$$
(13)
We can now replace equation into equation . In a matter-dominated regime, and for a shear free, non-rotating cosmic fluid we obtain:
$$\frac{d\mathrm{\Theta }}{dt}+\frac{1}{3}\mathrm{\Theta }^2=4\pi G\rho +\mathrm{\Lambda }$$
(14)
On making use of equation we can split $`\mathrm{\Theta }`$ as
$$\mathrm{\Theta }_\mu u^\mu =\frac{3\dot{a}}{a}+\frac{\theta }{a}$$
(15)
so that, taking into consideration the field equations for the expansion factor $`a(t)`$ (Eqs. and ), equation can be recast in the form
$$\frac{d\theta }{d\tau }+(\tau )\theta +\frac{1}{3}\theta ^2=4\pi Ga^2\overline{\rho }\delta $$
(16)
where $`(\tau )d(\mathrm{ln}a)/d\tau `$. We can now eliminate $`\theta `$ between eqs. and to find the following second order differential equation for the density contrast:
$`{\displaystyle \frac{d^2\delta }{d\tau ^2}}+(\tau ){\displaystyle \frac{d\delta }{d\tau }}{\displaystyle \frac{3}{2}}^2(\tau )\mathrm{\Omega }_M(\tau )\delta `$ (17)
$`=`$ $`{\displaystyle \frac{4}{3}}(1+\delta )^1\left({\displaystyle \frac{d\delta }{d\tau }}\right)^2+{\displaystyle \frac{3}{2}}^2(\tau )\mathrm{\Omega }_M(\tau )\delta ^2`$
where we have shifted to the rhs all non-linear terms, and used the notation
$$\mathrm{\Omega }_M(\tau )=\frac{\mathrm{\Omega }_M}{\mathrm{\Omega }_M+a\mathrm{\Omega }_k+a^3\mathrm{\Omega }_\mathrm{\Lambda }}$$
(18)
Equation reproduces the equation of the spherical collapse model (SC). In other words, the SC approximation is the exact dynamics when shear is neglected (see Fosalba & Gaztañaga 1998a). As one would expect, this yields a local evolution, in the sense that the evolved field at a point is just given by a local (non-linear) transformation of the initial field at the same point, with independence of the surroundings. This SC solution yields the exact perturbation theory predictions for the cumulants at tree-level (leading order with Gaussian initial conditions) and it also is an excellent approximation for next to leading orders, see below. As mentioned in the introduction, one can also use the SC model to predict the value of the critical linear overdensity, $`\delta _c`$, that will collapse into virialized halos.
### 2.3. Linear growth
We next do a perturbative expansion for $`\delta `$. The first contribution is the linear theory solution. For this, equation clearly simplifies to
$$\frac{d^2\delta _l}{d\tau ^2}+(\tau )\frac{d\delta _l}{d\tau }\frac{3}{2}^2(\tau )\mathrm{\Omega }_M(\tau )\delta _l=0$$
(19)
where $`\delta _l`$ stands for the “linear” solution. Because the coefficients of the above equation are time dependent only, the spatial and temporal part factorise:
$$\delta _l(\text{x},\tau )=\delta _0(\text{x})D(\tau )$$
(20)
where $`D`$ is usually called the linear growth factor. Thus initial fluctuations, no matter of what size, are amplified by the same factor, and the statistical properties of the initial field are just linearly scaled. For example, the $`N`$-point correlation functions are:
$$\xi _N(r_1,..,r_N,t)=D^N\xi _N(r_1,..,r_N,0)$$
(21)
To find the solution to equation it is expedient to change the time variable to $`\eta =\mathrm{ln}(a)`$, so that
$$\frac{d}{d\eta }=\frac{1}{(\tau )}\frac{d}{d\tau }=\frac{1}{H}\frac{d}{dt}$$
(22)
We then have
$$\frac{d^2D}{d^2\eta }+\left(2+\frac{\dot{H}}{H^2}\right)\frac{dD}{d\eta }\frac{3}{2}\mathrm{\Omega }_M(\eta )D=0$$
(23)
where we can write
$`{\displaystyle \frac{\dot{H}}{H^2}}`$ $`=`$ $`{\displaystyle \frac{3}{2}}\left({\displaystyle \frac{\mathrm{\Omega }_M+2/3e^\eta \mathrm{\Omega }_k}{\mathrm{\Omega }_M+e^\eta \mathrm{\Omega }_k+e^{3\eta }\mathrm{\Omega }_\mathrm{\Lambda }}}\right)`$ (24)
$`\mathrm{\Omega }_M(\eta )`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }_M}{\mathrm{\Omega }_M+e^\eta \mathrm{\Omega }_k+e^{3\eta }\mathrm{\Omega }_\mathrm{\Lambda }}}`$ (25)
where $`\mathrm{\Omega }_M`$, $`\mathrm{\Omega }_k`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$ are just constants (the current value at $`a=1`$).
In the Einstein-deSitter universe ($`\mathrm{\Omega }_k=\mathrm{\Omega }_\mathrm{\Lambda }=0`$) we have that $`\mathrm{\Omega }_M(\eta )=1`$ and $`\dot{H}/H^2=3/2`$, so the differential equation becomes
$$\frac{d^2D}{d^2\eta }+\frac{1}{2}\frac{dD}{d\eta }\frac{3}{2}D=0$$
(26)
whose solutions
$$D=C_1e^\eta +C_2e^{3/2\eta }=C_1a+C_2a^{3/2}$$
(27)
reproduce the usual linear growth $`Da`$ and the decaying solutions $`Da^{3/2}`$.
### 2.4. Non-linear growth
The exact (non-perturbative) solution for the SC Eq. for the density contrast in an Einstein-deSitter universe admits a well known parametric representation:
$`\delta (\phi )`$ $`=`$ $`{\displaystyle \frac{9}{2}}{\displaystyle \frac{(\phi \mathrm{sin}\phi )^2}{(1\mathrm{cos}\phi )^3}}1`$
$`\delta _l(\phi )`$ $`=`$ $`{\displaystyle \frac{3}{5}}\left[{\displaystyle \frac{3}{4}}(\phi \mathrm{sin}\phi )\right]^{2/3}`$ (28)
for $`\delta _l>0`$, linear overdensity, and
$`\delta (\phi )`$ $`=`$ $`{\displaystyle \frac{9}{2}}{\displaystyle \frac{(\mathrm{sinh}\phi \phi )^2}{(\mathrm{cosh}\phi 1)^3}}1`$
$`\delta _l(\phi )`$ $`=`$ $`{\displaystyle \frac{3}{5}}\left[{\displaystyle \frac{3}{4}}(\mathrm{sinh}\phi \phi )\right]^{2/3}`$ (29)
for $`\delta _l<0`$, linear under-density (see Peebles 1993), where the parameter $`\phi `$ is just a parametrisation of the time coordinate. There is also a solution for the $`\mathrm{\Omega }_M1`$ case (see Bernardeau 1992, Fosalba & Gaztañaga 1998b). The continuous line in Figure 1 illustrates the solution to the above equation (the other lines will be explained later). Note the singularity at $`\delta _l1.686`$, which corresponds to the gravitational collapse (see §3.1.3 below).
If we are only interested in the perturbative regime ($`\delta _l0`$), which is the relevant one for the description of structure formation on large scales, the above solution can be expressed directly in terms of the linear density contrast, $`\delta _l`$, which plays the role of the initial size of the spherical fluctuation in Eq.. This way, the evolved density contrast in the perturbative regime is given by a local-density transformation of the linear density fluctuation,
$$\delta =f(\delta _l)=\underset{n=1}{\overset{\mathrm{}}{}}\frac{\nu _n}{n!}[\delta _l]^n$$
(30)
Notice that all the non-linear dynamical information in the SC model is encoded in the $`\nu _n`$ coefficients. We can now introduce the above power series expansion in Eq. and determine the $`\nu _n`$ coefficients one by one. Before we do this, it is convenient to change again the time variable to $`\eta =\mathrm{ln}(a)`$ as we did in the linear case, Eq:
$`{\displaystyle \frac{d^2\delta }{d^2\eta }}+\left(2+{\displaystyle \frac{\dot{H}}{H^2}}\right){\displaystyle \frac{d\delta }{d\eta }}{\displaystyle \frac{3}{2}}\mathrm{\Omega }_M(\eta )\delta `$ (31)
$`=`$ $`{\displaystyle \frac{4}{3}}{\displaystyle \frac{1}{1+\delta }}\left({\displaystyle \frac{d\delta }{d\eta }}\right)^2+{\displaystyle \frac{3}{2}}\mathrm{\Omega }_M(\eta )\delta ^2`$
We can now use the expansion in Eq. with $`\delta _l`$ given by the linear growth factor $`D=a=e^\eta `$ and compare order by order. For the Einstein-deSitter universe they turn out to be:
$$\nu _2=\frac{34}{21};\nu _3=\frac{682}{189}$$
(32)
and so on (see eg Folsalba & Gaztañaga 1998b for other cases). Once we have these coefficients we can get the evolution of the non-linear variance and higher order moments in terms of the initial conditions (see §4.1 below).
### 2.5. Equation of state $`p=\gamma \rho `$
We will now consider a perfect fluid with equation of sate $`p=\gamma \rho `$. Not all values of $`\gamma `$ make physical sense. Here, in the spirit of going beyond the standard paradigm, we will ignore these restrictions and assume that $`\gamma `$ can take any real constant value, irrespective of other cosmological parameters.
The time-component of the energy conservation equations $`_\nu T^{\mu \nu }=0`$ gives us (for $`p=\gamma \rho `$) both the background density behavior
$$\overline{\rho }a^{3(1+\gamma )}=\text{const}$$
(33)
and the continuity equation for the density contrast
$$\frac{d\delta }{d\tau }+(1+\gamma )(1+\delta )\theta =\gamma \overline{\rho }(\text{v}\delta )$$
(34)
where, like before, $`\tau `$ is the conformal time, and $`\theta \text{v}`$. This is the generalization of equation for a relativistic fluid. Note that an additional (quadratic) term now appears in the rhs of . The magnitude of this term is assessed by resorting to the space-components of the energy conservation equations $`_\nu T^{\mu \nu }=0`$: these are identically satisfied when $`\gamma =0`$, and they show that $`\text{v}\delta |\text{v}|^2/c^2`$, plus higher order contributions. These can be safely neglected since peculiar velocities are always very small compared to the speed of light; in fact the approximation $`|\text{v}|^2/c^20`$ is always made, even in the more standard case when $`\gamma =0`$. We shall therefore consistently adopt the following equation for the density contrast:
$$\frac{d\delta }{d\tau }+(1+\gamma )(1+\delta )\theta =0$$
(35)
Also, Hubble’s equation, Eq. , now becomes
$$H^2=H_0^2\left[\mathrm{\Omega }_Ma^{3(1+\gamma )}+\mathrm{\Omega }_ka^2+\mathrm{\Omega }_\mathrm{\Lambda }\right]$$
(36)
We can combine equation with the Raychaudhuri equation for this case —cf Eqs. and
$$\frac{d\mathrm{\Theta }}{dt}+\frac{1}{3}\mathrm{\Theta }^2=4\pi G\rho \left(1+3\gamma \right)+\mathrm{\Lambda }$$
(37)
to obtain, after some algebra,
$`{\displaystyle \frac{d^2\delta }{d^2\eta }}+\left(2+{\displaystyle \frac{\dot{H}}{H^2}}\right){\displaystyle \frac{d\delta }{d\eta }}{\displaystyle \frac{3}{2}}(1+\gamma )(1+3\gamma )\mathrm{\Omega }(\eta )\delta =`$
$`{\displaystyle \frac{4+3\gamma }{3+3\gamma }}{\displaystyle \frac{1}{1+\delta }}\left({\displaystyle \frac{d\delta }{d\eta }}\right)^2+{\displaystyle \frac{3}{2}}(1+\gamma )(1+3\gamma )\mathrm{\Omega }(\eta )\delta ^2`$ (38)
where we have expediently redefined $`\mathrm{\Omega }(\eta )`$ in Eq. to
$$\mathrm{\Omega }_M(\eta )=\frac{\mathrm{\Omega }_M}{\mathrm{\Omega }_M+e^{\eta (1+3\gamma )}\mathrm{\Omega }_k+e^{3\eta (1+\gamma )}\mathrm{\Omega }_\mathrm{\Lambda }}$$
(39)
and we can write
$$\frac{\dot{H}}{H^2}=\frac{3}{2}\left(\frac{(1+\gamma )\mathrm{\Omega }_Me^{3\eta \gamma }+2/3e^\eta \mathrm{\Omega }_k}{\mathrm{\Omega }_Me^{3\eta \gamma }+e^\eta \mathrm{\Omega }_k+e^{3\eta }\mathrm{\Omega }_\mathrm{\Lambda }}\right)$$
(40)
In an Einstein-deSitter universe ($`\mathrm{\Omega }_k=\mathrm{\Omega }_\mathrm{\Lambda }=0`$), $`\mathrm{\Omega }(\eta )=1`$, and the linear regime is governed by
$$\frac{d^2D}{d^2\eta }+\frac{13\gamma }{2}\frac{dD}{d\eta }\frac{3}{2}(1+\gamma )(1+3\gamma )D=0$$
(41)
which has the usual solutions of the form $`D=a^\alpha `$, with
$$\alpha _1=1+3\gamma ,\alpha _2=3(1+\gamma )/2$$
(42)
Figure 2 shows these perturbative solutions. The shaded region corresponds to the case where linear evolution is suppressed, eg $`\alpha <0`$. In this case, as can be seen from Eq.-, $`\nu _2`$ and $`\nu _3`$ have a very rapid variation. The growing mode for $`\gamma >1/3`$ is:
$`\alpha _1`$ $`=`$ $`1+3\gamma `$ (43)
$`\nu _2`$ $`=`$ $`{\displaystyle \frac{2\left(17+48\gamma +27\gamma ^2\right)}{3\left(1+\gamma \right)\left(7+15\gamma \right)}}`$ (44)
$`\nu _3`$ $`=`$ $`[\mathrm{\hspace{0.17em}72}+540\gamma +324\gamma ^2+{\displaystyle \frac{16}{\left(1+\gamma \right)^2}}+{\displaystyle \frac{24}{1+\gamma }}`$
$``$ $`{\displaystyle \frac{\left(6+18\gamma \right)\left(17+48\gamma +27\gamma ^2\right)}{\left(1+\gamma \right)\left(7+15\gamma \right)}}]`$
$`\times `$ $`\left(\mathrm{\hspace{0.17em}27}+144\gamma +189\gamma ^2\right)^1`$ (46)
For $`\gamma <1`$ the dominant linear growth is $`\alpha _2`$ and the values of $`\nu _2`$ and $`\nu _3`$ are constant:
$`\alpha _2`$ $`=`$ $`{\displaystyle \frac{3\left(1+\gamma \right)}{2}}`$ (47)
$`\nu _2`$ $`=`$ $`{\displaystyle \frac{3}{2}}`$ (48)
$`\nu _3`$ $`=`$ $`3`$ (49)
For radiation ($`\gamma =1/3`$) we have that $`\alpha _1=2`$ which reproduces the well known results (see Peebles 1993) and $`\nu _2=\frac{3}{2}`$ and $`\nu _3=3`$, which are new results as far as we know. Note that these values are identical to the case of negative pressure, $`\gamma <1`$, the only difference being in the linear growth, but for $`\gamma =7/3`$ all $`\alpha `$, $`\nu _2`$ and $`\nu _3`$ are identical to the radiation case. In the limit of strong pressure $`\gamma \mathrm{}`$ we find: $`\nu _2=6/5`$ and $`\nu _3=12/7`$ . As can be seen in Figure 2, and also in the equations above, there are poles for $`\nu _2`$ at $`\gamma =1`$ and $`\gamma =7/15`$.
Figure 3 shows the corresponding variation in $`\delta _c`$, defined as the value of the linear overdensity where the corresponding non-linear value becomes infinity (see §3.1.3).
## 3. Gravitational Growth outside GR
### 3.1. Scalar-Tensor Theories
Here we investigate how a varying $`G`$ could change the above results. We parameterize the variation of $`G`$ using scalar-tensor theories (STT) of gravity such as Brans-Dicke (BD) theory or its extensions.
To make quantitative predictions we will consider cosmic evolution in STTs, where $`G`$ is derived from a scalar field $`\varphi `$ which is characterized by a function $`\omega =\omega (\varphi )`$ determining the strength of the coupling between the scalar field and gravity. In the simplest BD models, $`\omega `$ is just a constant and $`G\varphi ^1`$ —see below. However if $`\omega `$ varies then it can change with cosmic time, so that $`\omega =\omega (z)`$. The structure of the solutions to BD equations is quite rich and depends crucially on the coupling function $`\omega (\varphi )`$ (see Barrow & Parsons 1996).
Here we shall be considering the standard BD model with constant $`\omega `$; the field equations are (see eg Weinberg 1972):
$`R_{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{8\pi }{\varphi }}\left(T_{\mu \nu }{\displaystyle \frac{1+\omega }{3+2\omega }}g_{\mu \nu }T\right){\displaystyle \frac{\omega }{\varphi ^2}}_\mu \varphi _\nu \varphi `$ (50)
$``$ $`{\displaystyle \frac{1}{\varphi }}_\mu _\nu \varphi `$
$`\mathrm{}\varphi `$ $`=`$ $`{\displaystyle \frac{8\pi }{3+2\omega }}T,(Tg^{\mu \nu }T_{\mu \nu })`$ (51)
The Hubble rate $`H`$ for a homogeneous and isotropic background universe can be easily obtained from the above;
$$H^2\left(\frac{\dot{a}}{a}\right)^2=\frac{8\pi \rho }{3\varphi }+\frac{k}{a^2}+\frac{\mathrm{\Lambda }}{3}+\frac{\omega }{6}\frac{\dot{\varphi }^2}{\varphi ^2}H\frac{\dot{\varphi }}{\varphi }$$
(52)
These equations must be complemented with the equation of state for the cosmic fluid. In a flat, matter dominated universe ($`p=0`$), an exact solution to the problem can be found:
$$G=\frac{4+2\omega }{3+2\omega }\varphi ^1=G_0(1+z)^{1/(1+\omega )}$$
(53)
and
$$a(t)=(t/t_0)^{(2\omega +2)/(3\omega +4)}$$
(54)
This solution for the flat universe is recovered in a general case in the limit $`t\mathrm{}`$, and also arises as an exact solution of Newtonian gravity with a power law $`Gt^n`$ (Barrow 1996). For non-flat models, $`a(t)`$ is not a simple power law and the solutions get far more complicated. To illustrate the effects of a non-flat cosmology we will consider general solutions that can be parametrized as Eq. but which are not simple power-laws in $`a(t)`$. In this case, it is easy to check that the new Hubble law given by Eq. becomes
$$H^2=H_0^2\left[\widehat{\mathrm{\Omega }}_M(1+z)^{3+1/(1+\omega )}+\widehat{\mathrm{\Omega }}_k(1+z)^2+\widehat{\mathrm{\Omega }}_\mathrm{\Lambda }\right]$$
(55)
where $`\widehat{\mathrm{\Omega }}_M`$,$`\widehat{\mathrm{\Omega }}_k`$ and $`\widehat{\mathrm{\Omega }}_\mathrm{\Lambda }`$ follow the usual relation $`\widehat{\mathrm{\Omega }}_M+\widehat{\mathrm{\Omega }}_k+\widehat{\mathrm{\Omega }}_\mathrm{\Lambda }=1`$, and are related to the familiar local ratios ($`z0`$: $`\mathrm{\Omega }_M8\pi G_0\rho _0/(3H_0^2)`$, $`\mathrm{\Omega }_k=k/H_0^2`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }=\mathrm{\Lambda }/(3H_0^2)`$) by
$`\widehat{\mathrm{\Omega }}_M`$ $`=`$ $`\mathrm{\Omega }_M{\displaystyle \frac{3(1+\omega )^2}{(2+\omega )(4+3\omega )}}`$
$`\widehat{\mathrm{\Omega }}_\mathrm{\Lambda }`$ $`=`$ $`\mathrm{\Omega }_\mathrm{\Lambda }{\displaystyle \frac{6(1+\omega )^2}{(3+2\omega )(4+3\omega )}}`$ (56)
$`\widehat{\mathrm{\Omega }}_k`$ $`=`$ $`\mathrm{\Omega }_k{\displaystyle \frac{6(1+\omega )^2}{(3+2\omega )(4+3\omega )}}`$
Thus the GR limit is recovered as $`\omega \mathrm{}`$.
We now investigate the density fluctuations in the above theory. Like in section II, we shall make use of the continuity equation in combination with the Raychaudhuri equation . As mentioned above, cf section 2.1, both of these are still valid within the context of BD theory: it is only needed to replace the Ricci tensor in the rhs of Eq. according to BD’s field equations, Eq. . Considering again a non-rotating, shear-free cosmic fluid, we find:
$`{\displaystyle \frac{d\mathrm{\Theta }}{dt}}+{\displaystyle \frac{1}{3}}\mathrm{\Theta }^2=`$ (57)
$`=`$ $`{\displaystyle \frac{4+2\omega }{3+2\omega }}{\displaystyle \frac{4\pi \rho }{\varphi }}\left(1+{\displaystyle \frac{1+\omega }{2+\omega }}{\displaystyle \frac{3p}{\rho }}\right)\omega {\displaystyle \frac{\dot{\varphi }^2}{\varphi ^2}}{\displaystyle \frac{\ddot{\varphi }}{\varphi }}`$
We shall still make use of a gravitational “constant” parametrized as in equation above; this is justified insofar as the characteristic length for the variation of $`\varphi `$ is typically much greater than that of the density fluctuations in a matter dominated universe —see eg (Nariai 1969). In this approximation, the above equation gives
$$\frac{d\theta }{d\tau }+(\tau )\theta +\frac{1}{3}\theta ^2=\frac{4+2\omega }{3+2\omega }\frac{4\pi a^2\overline{\rho }\delta }{\varphi }$$
(58)
where $`\tau `$ is again the conformal time parameter, $`d\tau =a^1dt`$, and $`\theta `$ is defined in equations and . Remarkably, this equation is very similar to the GR equation : we only need to replace in it the gravitational constant $`G`$ by its expression as a multiple of the varying scalar field $`\varphi `$ given in equation . Combining with the continuity equation we immediately find
$$\frac{d^2\delta }{d\tau ^2}+(\tau )\frac{d\delta }{d\tau }\frac{4}{3(1+\delta )}\left(\frac{d\delta }{d\tau }\right)^2=\frac{4+2\omega }{3+2\omega }\frac{4\pi a^2\rho \delta }{\varphi }$$
(59)
Like in section II, we change the independent variable in to $`\eta =\mathrm{ln}a`$, whereby we obtain
$`{\displaystyle \frac{d^2\delta }{d\eta ^2}}+\left(2+{\displaystyle \frac{\dot{H}}{H^2}}\right){\displaystyle \frac{d\delta }{d\eta }}{\displaystyle \frac{4}{3}}(1+\delta )^1\left({\displaystyle \frac{d\delta }{d\eta }}\right)^2`$ (60)
$`=`$ $`{\displaystyle \frac{4+2\omega }{3+2\omega }}{\displaystyle \frac{4\pi a^2\rho \delta }{H^2\varphi }}`$
Using equation to calculate $`\dot{H}`$, and assuming further that $`\widehat{\mathrm{\Omega }}_k=\widehat{\mathrm{\Omega }}_\mathrm{\Lambda }=0`$, we finally get
$`{\displaystyle \frac{d^2\delta }{d^2\eta }}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\omega }{1+\omega }}{\displaystyle \frac{d\delta }{d\eta }}{\displaystyle \frac{1}{2}}{\displaystyle \frac{(2+\omega )(4+3\omega )}{(1+\omega )^2}}\delta `$ (61)
$`=`$ $`{\displaystyle \frac{4}{3}}{\displaystyle \frac{1}{1+\delta }}\left({\displaystyle \frac{d\delta }{d\eta }}\right)^2+{\displaystyle \frac{1}{2}}{\displaystyle \frac{(2+\omega )(4+3\omega )}{(1+\omega )^2}}\delta ^2`$
We next examine the solutions to this equation.
#### 3.1.1 Linear growth
Let us call $`D(\eta )`$ the solution to the linearized version of equation , i.e.,
$$\frac{d^2D}{d^2\eta }+\frac{1}{2}\frac{\omega }{1+\omega }\frac{dD}{d\eta }\frac{1}{2}\frac{(2+\omega )(4+3\omega )}{(1+\omega )^2}D=0$$
(62)
Again the solutions are given by the roots $`\alpha _1`$ and $`\alpha _1`$ of the corresponding characteristic functions:
$$D=C_1a^{\alpha _1}+C_2a^{\alpha _2}$$
(63)
with
$`\alpha _1`$ $`=`$ $`{\displaystyle \frac{2+\omega }{1+\omega }}1+{\displaystyle \frac{1}{\omega }}+𝒪\left({\displaystyle \frac{1}{\omega ^2}}\right)`$ (64)
$`\alpha _2`$ $`=`$ $`{\displaystyle \frac{43\omega }{2+2\omega }}{\displaystyle \frac{3}{2}}{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{\omega }}+𝒪\left({\displaystyle \frac{1}{\omega ^2}}\right)`$ (65)
which reproduces the usual linear growth $`Da`$ and $`Da^{3/2}`$ in the limit $`\omega \mathrm{}`$. Note that $`\alpha _1`$ corresponds to the growing mode only for large values of $`|\omega |`$, but the situation is more complicated when $`\omega `$ is not large.
Figure 4 shows the values of $`\alpha _1`$ and $`\alpha _2`$ as functions of $`\omega `$. The effective $`G`$ in BD decreases as the Universe expands if $`1<\omega <\mathrm{}`$, and the expansion factor $`a(t)`$ stops for $`\omega =1`$; the growing mode in this regime is controlled by $`\alpha _1`$, since this is the positive root. The growing mode for $`4/3<\omega <1`$ is $`\alpha _2`$, but the universe shrinks to an eventual collapse in this regime (see Eq. ). Between $`2<\omega <4/3`$ the Universe expands again, but there are no growing modes, as can be seen in Figure 4 (both $`\alpha _1`$ and $`\alpha _2`$ are negative). For $`\omega <2`$ the expansion factor grows with time and $`\alpha _1`$ becomes the growing mode again. Notice that in this regime of $`\omega <2`$, $`\alpha _1<1`$, so that it is slower than for $`\omega >0`$. As we will show below this is compensated in part by a stronger non-linear growth.
#### 3.1.2 Non-linear growth
In the non-linear case we consider the full version of equation . We can now proceed as before, using the expansion in Eq. with $`\delta _l`$ given by the linear growth factor $`D=a^{\alpha _1}=e^{\alpha _1\eta }`$, and compare order by order. We find
$`\nu _2`$ $`=`$ $`{\displaystyle \frac{34\omega +56}{21\omega +36}}={\displaystyle \frac{34}{21}}\left[1{\displaystyle \frac{8}{119}}{\displaystyle \frac{1}{\omega }}+𝒪\left({\displaystyle \frac{1}{\omega ^2}}\right)\right]`$ (66)
$`\nu _3`$ $`=`$ $`{\displaystyle \frac{2(944+1136\omega +341\omega ^2)}{3(12+7\omega )(16+9\omega )}}`$
$`=`$ $`{\displaystyle \frac{682}{189}}\left[1+{\displaystyle \frac{3452}{21483}}{\displaystyle \frac{1}{\omega }}+𝒪\left({\displaystyle \frac{1}{\omega ^2}}\right)\right]`$ (67)
Note how for positive $`\omega `$ non-linear effects tend to compensate the increase in linear effects, cf Figure 4, whereas for $`\omega <4/3`$, the linear effects are reduced ($`\alpha <1`$) while non-linearities get larger.
Figure 5 shows the variation in $`\nu _2`$ as a function of $`\omega `$ using Eq.. Negative values of $`\omega `$ produce almost symmetrical variations in the opposite direction when $`|\omega |`$ is large. For small $`\omega `$ there is a pole at $`\omega =12/7`$ where $`\nu _2`$ diverges. But note that there is no growing linear mode in this case, which means that fluctuations are rapidly suppressed.
#### 3.1.3 Strongly non-linear regime
Figure 1 shows the fully non-linear solution for the overdensity $`\delta `$ as a function of the linear one $`\delta _l`$. The continuous line shows the standard solution to Eq. as given in Eq.-. As can be seen in the Figure, there is a critical value of $`\delta _l=3/2(3\pi /2)^{3/2}1.6865`$ where the non-linear fluctuations become infinite. This corresponds to the point where the spherical collapse occurs (see Peebles 1993). Thus an initial fluctuations $`\delta _0`$ will collapse after evolving a time $`t`$, such that the growth factor is $`D(t)=\delta _c/\delta _0`$. For the standard GR, flat and matter dominated case, this time would correspond to a formation red-shift: $`z_f=\delta _0/\delta _c1`$ (if we use $`a=1`$ today). For the BD case both $`\delta _c`$ and $`D(t)`$ are different, so that formation times $`z_f`$ will be correspondingly different (see Eq.). The short-dashed lines in Figure 1 correspond to the same exact solution in the BD model with $`\omega =10`$ and $`\omega =1`$. Right panel in Figure 5 illustrates how $`\delta _c`$ changes in the BD model as a function of $`\omega `$.
### 3.2. Gravitational Growth with $`H^2a^{3(1+ϵ)}`$
Consider now the flat case with $`\mathrm{\Omega }_k=\mathrm{\Omega }_\mathrm{\Lambda }=0`$. To account for a simple variation on the standard Einstein’s field equations we will consider the case where fluctuations grow according to the matter dominated case (ie $`\gamma =0`$) but the background evolves in a different way. We will assume that the Hubble rate goes like $`H^2a^{3(1+ϵ)}`$ rather than $`H^2a^3`$. It might be possible to find some motivation for this model, but this is beyond the scope of this work. Here we just want to introduce some parametric variations around the standard field equations to see how things might change. In this case we have:
$$\frac{d^2\delta }{d^2\eta }+\frac{13ϵ}{2}\frac{d\delta }{d\eta }\frac{3}{2}\delta =\frac{4}{3}\frac{1}{1+\delta }\left(\frac{d\delta }{d\eta }\right)^2+\frac{3}{2}\delta ^2$$
(68)
The solutions for the linear growth factor index and the non-linear coefficient $`\nu _2`$ are
$`\alpha _1`$ $`=`$ $`{\displaystyle \frac{1+3ϵ+\sqrt{256ϵ+9ϵ^2}}{4}}`$ (69)
$`\nu _2`$ $`=`$ $`{\displaystyle \frac{13130ϵ+45ϵ^2+(13ϵ)\sqrt{256ϵ+9ϵ^2}}{8418ϵ+27ϵ^2}}`$ (70)
These solutions as a function of $`ϵ`$ are illustrated in Figure 6, which also shows $`\nu _3`$. As can be seen in the Figure, the higher the linear growth index $`\alpha _1`$ the lower the non-linear coefficients.
Right panel in Figure 6 shows the corresponding variation in $`\delta _c`$.
## 4. Observational consequences
We will focus here on Gaussian initial conditions. That is, our initial field for structure formation is a spatial realization of a (three-dimensional) Gaussian distribution with a given power spectrum shape, and a very small initial amplitude. As we are interested in the gravitational regime alone, this field will be smoothed over a large enough scale, corresponding to the distance beyond which non-gravitational forces (eg hydrodynamics) can be neglected. Thus at each point the overdensity $`\delta (\text{x})`$ grows according to gravity, which in the shear free approximation is just a local dynamics: the spherical collapse (eg Eq.).
### 4.1. Cumulants
Consider the $`Jorder`$ moments of the fluctuating field:
$$m_J\delta ^J.$$
(71)
Here the expectation values $`\mathrm{}`$ correspond to an average over realizations of the initial field. On comparing with observations we assume the fair sample hypothesis (§30 Peebles 1980), by which we can commute spatial integrals with expectation values. Thus, in practice $`\mathrm{}`$ is the average over positions in the survey area. In this notation the variance is defined as:
$$Var(\delta )\sigma ^2m_2m_1^2$$
(72)
More generally, we introduce the connected moments $`\overline{\xi }_J`$, which carry statistical information independent of the lower order moments, and are formally denoted by a bracket with subscript $`c`$:
$$\overline{\xi }_J\delta ^J_c$$
(73)
The connected moments are also called cumulants, reduced moments or irreducible moments. They are defined by just subtracting the lower order contributions:
$`\overline{\xi }_1`$ $`=`$ $`m_10`$
$`\overline{\xi }_2`$ $`=`$ $`\sigma ^2=m_2\overline{\xi }_1^2=m_2`$
$`\overline{\xi }_3`$ $`=`$ $`m_33\overline{\xi }_2\overline{\xi }_1\overline{\xi }_1^3=m_3`$ (74)
$`\overline{\xi }_4`$ $`=`$ $`m_44\overline{\xi }_3\overline{\xi }_13\overline{\xi }_2^26\overline{\xi }_2\overline{\xi }_1^2\overline{\xi }_1^4=m_43m_2^2`$
and so on. It is useful to introduce the hierarchical ratios:
$$S_J=\frac{\overline{\xi }_J}{\overline{\xi }_2^{J1}}$$
(75)
which are also called normalized one-point cumulants or reduced cumulants. We shall use the term skewness, for $`S_3=\overline{\xi }_3/\overline{\xi }_2^2`$ and kurtosis, for $`S_4=\overline{\xi }_4/\overline{\xi }_2^3`$.
#### 4.1.1 Linear Theory
As mentioned in §2.3, initial fluctuations, $`\delta _0`$, no matter of what amplitude, grow all by the same factor, $`D`$; thus the statistical properties of the initial field are just linearly scaled in the final (linear) field, $`\delta _l`$:
$$\delta _l^J_c=D^J\delta _0^J_c$$
(76)
Consider for example the linear rms fluctuations $`\sigma _l`$ or its variance $`\sigma _l^2`$. In the linear regime we have:
$$\sigma _l^2\delta ^2(t)=D(tt_0)^2\delta _0^2=D(tt_0)^2\sigma _0^2$$
(77)
where $`\sigma _0`$ refers to some initial reference time $`t_0`$. To give an idea of this effect, consider the growth of fluctuations since matter domination, when the universe was about 1100 times smaller. In General Relativity (GR) in the matter dominated Einstein-deSitter universe, $`\sigma `$ would grow by a factor $`D1100`$. While, if we take $`\omega 10`$ in the DB theory, eg Eq., we have that fluctuations increase instead by a factor $`D2079`$, which is about $`1.9`$ times larger in $`\sigma `$, so the variance nowadays would be about $`3.6`$ times larger if we fixed it around the COBE variance of Cosmic Microwave Background (CBM) temperature fluctuations. For $`\omega 100`$, the variance would only be $`14\%`$ larger than in GR. This latter result is small, but it could be relevant for future precision measurements (eg MAP or PLANCK satellites to map CMB and 2DF or SLOAN DIGITAL SKY galaxy surveys). Similar considerations can be made for the values of $`\alpha `$ with a different cosmic equation of state, eg Eq. or a different Hubble law, Eq.. In general we can write that a small change in $`\alpha `$ would produce a relative change in the linear rms of
$$\frac{\mathrm{\Delta }\sigma }{\sigma }=\mathrm{ln}(1+z)\mathrm{\Delta }\alpha $$
(78)
Thus, a change of only $`1\%`$ in the absolute value of the equation of state $`\gamma `$, would produce a relative change of $`20\%`$ in $`\sigma `$ between recombination $`z1100`$ and now, cf Eq..
The hierarchical ratios (see Eq.) will scale as, $`S_J=S_J(0)/D^{J2}`$, where $`S_J(0)`$ are the initial ratios. This implies that the linear growth erases the initial skewness and kurtosis, so that $`S_J0`$, as time evolves (and $`D\mathrm{}`$). Note that if we want to do a meaningful calculation of these ratios or the cumulants, in general we might need to consider more terms in the perturbative series, Eq.. For Gaussian initial conditions $`S_J(0)=0`$, and we need to consider higher order terms in the perturbation series to find the leading order prediction.
#### 4.1.2 Weakly non-linear
The next to leading order solutions for the cumulants of the evolved field given the expansion Eq., can be easily found by just taking expectation values of different powers of $`\delta `$ (see eg Fosalba & Gaztañaga 1998a). For leading order Gaussian initial conditions we have
$`S_3`$ $`=`$ $`3\nu _2+𝒪(\sigma _l^2)`$
$`S_4`$ $`=`$ $`4\nu _3+12\nu _2^2+𝒪(\sigma _l^2)`$ (79)
For non-Gaussian initial conditions see Fry & Scherrer (1994) Chodorowski & Bouchet (1996), Gaztañaga & Mahonen (1996), Gaztañaga & Fosalba (1998).
If we use for $`\nu _2`$ the solution in Eq., eg $`\nu _2=34/21`$, the skewness yields $`S_3=3\nu _2=34/7`$, which reproduces the exact perturbation theory (PT) result by Peebles (1980) in the matter dominated Einstein-deSitter universe. Thus the shear-free or SC model gives the exact leading order result for the skewness. This is also true for higher orders (see Bernardeau 1992 and Fosalba & Gaztañaga 1998a) and for other cosmologies (eg Bouchet et al. 1992, Bernardeau 1994a, Fosalba & Gaztañaga 1998b, Kamionkowski & Buchalter 1999). For smoothed fields, the exact leading order results are slightly different:
$`S_3`$ $`=`$ $`{\displaystyle \frac{34}{7}}+\gamma _1`$ (80)
$`S_4`$ $`=`$ $`{\displaystyle \frac{60712}{1323}}+{\displaystyle \frac{62}{3}}\gamma _1+{\displaystyle \frac{7}{3}}\gamma _1^2`$ (81)
where $`\gamma _1`$ is the logarithmic slope of the smoothed variance (see Juszkiewicz et al. 1993, Bernardeau 1994a, 1994b). These can also be reproduced in the shear-free approximation as shown by Gaztañaga & Fosalba (1998); this results in a smoothing correction:
$`\overline{\nu _2}`$ $`=`$ $`\nu _2+{\displaystyle \frac{\gamma _1}{3}}`$
$`\overline{\nu _3}`$ $`=`$ $`{\displaystyle \frac{1}{4}}(2\gamma _1+\gamma _1^2+6\gamma _1\nu _2+4\nu _3)`$ (82)
and replacing $`\nu _2`$ and $`\nu _3`$ by $`\overline{\nu _2}`$ and $`\overline{\nu _3}`$ in Eq. (see Fosalba & Gaztañaga 1998a for more details). There are also corrections to the above expressions when measurements are taken in red-shift space (eg Hivon et al 1995, Scoccimarro, Couchman and Frieman 1999). Next to leading order terms have been estimated by Scoccimarro & Frieman (1996) (see also Fosalba & Gaztañaga 1998a,b).
The smoothed values of $`S_3`$ and $`S_4`$ can be measured as traced by the large scale galaxy distribution (eg Bouchet et al. 1993, Gaztañaga 1992, 1994, Szapudi el at 1995, Hui & Gaztañaga 1999 and references therein), weak-lensing (Bernardeau, Van Waerbeke & Mellier 1997, Gaztañaga & Bernardeau 1998, Hui 1999) or the Ly-alpha QSO absorptions (Gaztañaga & Croft 1999). These measurements of the skewness and kurtosis can be translated into estimations of $`\nu _2`$ and $`\nu _3`$ which can be used to place constraints on $`\gamma `$, $`\omega `$ or $`ϵ`$ using Eq., and . For small values of these parameters the relationship is linear, so the uncertainties in $`S_3`$ and $`S_4`$ would directly translate into the corresponding uncertainties in $`\gamma `$, $`\omega `$ or $`ϵ`$.
The expressions above apply to unbiased tracers of the density field; since galaxies of different morphologies are known to have different clustering properties, at least some galaxy species must be biased tracers of the mass. As an example, suppose the probability of forming a luminous galaxy depends only on the underlying mean density field in its immediate vicinity. Under this simplifying assumption, the relation between the galaxy density field $`\delta _{gal}(\text{x})`$ and the mass density field $`\delta (\text{x})`$ can be written as
$$\delta _{gal}(\text{x})=f(\delta (\text{x}))=\underset{n}{}\frac{b_n}{n!}\delta ^n(\text{x}),$$
where $`b_n`$ are the bias parameters. Thus, biasing and gravity could produce comparable non-linear effects. To leading order in $`\overline{\xi }_2`$, this local bias scheme implies $`\overline{\xi }_2^{gal}=b_1^2\overline{\xi }_2`$, and (see Fry & Gaztañaga 1993)
$`S_3^{gal}`$ $`=`$ $`{\displaystyle \frac{S_3}{b_1}}+3{\displaystyle \frac{b_2}{b_1^2}}`$
$`S_4^{gal}`$ $`=`$ $`{\displaystyle \frac{S_4}{b_1^2}}+12{\displaystyle \frac{b_2S_3}{b_1^3}}+4{\displaystyle \frac{b_3}{b_1^4}}+12{\displaystyle \frac{b_2^2}{b_1^4}}`$ (83)
Gaztañaga & Frieman (1994) have used the comparison of $`S_3`$ and $`S_4`$ in PT with the corresponding values measured APM Galaxy Survey (Maddox et al. 1990), to infer that $`b_11`$, $`b_20`$ and $`b_30`$, but the results are degenerate due to the relative scale-independence of $`S_N`$ and the increasing number of biasing parameters. One could break this degeneracy by using the configuration dependence of the projected 3-point function, $`q_3(\alpha )`$, as proposed by Frieman & Gaztañaga (1994), Fry (1994), Matarrese, Verde & Heavens (1997), Scoccimarro et al. (1998). As shown in Frieman & Gaztañaga (1999), the configuration dependence of $`q_3(\alpha )`$ on large scales in the APM catalog is quite close to that expected in perturbation theory (see Fry 1984, Scoccimarro et al. 1998, Buchalter, Jaffe & Kamionkowski 2000), suggesting again that $`b_1`$ is of order unity (and $`b_20`$) for these galaxies. These agreement indicates that large-scale structure is driven by non-linear gravitational instability and that APM galaxies are relatively unbiased tracers of the mass on these large scales.
The values of $`S_3`$ and $`S_4`$ in the APM are measured to agree with the standard matter dominated Einstein-deSitter universe within about $`10\%20\%`$ (see Gaztañaga 1994; Gaztañaga & Frieman 1994; Baugh, Gaztañaga & Efstathiou 1995; Gaztañaga 1995, Hui & Gaztañaga 1999), also in agreement with the shape information in the 3-point function (see Frieman & Gaztañaga 1999). For example, using the projected APM catalogue Gaztañaga 1994 (Table 3) finds an average of $`S_3=3.2\pm 0.2`$ and $`S_4\pm 20.6\pm 2.6`$ scales between 7 and 30 $`h^1\text{Mpc}`$. For an average APM slope of $`\gamma _11.7`$, these values are in agreement with the PT predictions in Eq. yield: $`S_33.1`$ and $`S_418`$.
The 1–sigma error bar of $`10\%`$ on large scales quoted by Gaztañaga 1994 is mostly statistical (sampling error). Other systematics effects due to biasing, projection, or large scale errors in the building of the APM catalogue could be of the same order (see Frieman & Gaztañaga 1999 and Hui & Gaztañaga 1999). Thus given the current uncertainties it would be conservative to take a $`20\%`$ error bar. Unfortunately, with such large error bars we can not constraint much the values of $`\gamma `$ $`\omega `$ or $`ϵ`$. Stronger constraints can be found if we take the more optimistic 1–sigma $`10\%`$ error bars in the measurements of $`S_3`$ and $`S_4`$. This case is shown as horizontal dotted lines in Figures 2, 5 and 6. From $`\nu _2`$ the $`10\%`$ uncertainty translates into
$`0.2<`$ $`\gamma `$ $`<0.4`$
$`2.4>`$ $`\omega `$ $`>1.0`$
$`0.9<`$ $`ϵ`$ $`<0.9`$ (84)
Note that this is still of marginal interest. For example, the constraints on $`\gamma `$ include the possibility of a radiation ($`\gamma =1/3`$), matter ($`\gamma =0`$) or negative pressure $`\gamma <0`$. From $`\nu _3`$ we can obtain stronger constraints from a $`10\%`$ error (but obviously systematic effects could be larger for higher order cumulants):
$`0.1<`$ $`\gamma `$ $`<0.15`$
$`3.4>`$ $`\omega `$ $`>0.2`$
$`0.35<`$ $`ϵ`$ $`<0.35`$ (85)
These bounds are more interesting. It is clear that forthcoming surveys (such as the SLOAN Digital Sky Survey) will dramatically improve this situation (for errors on statistics see Szapudi, Colombi and Bernardeau 1999, and references therein).
Note that the above results are independent of the normalization of fluctuations.
### 4.2. Collapsed objects
Press & Schechter (1974) formalism and its extensions (eg Bond et al. 1991; Lacey & Cole 1993) predict the evolution of the mass function of halos and also their clustering properties. Comparison with N-body simulations show a very good agreement of these prescriptions for a wide range of statistical properties (eg see Lacey & Cole 1994 and references therein). For example, the comoving number density of collapsed objects (halos or clusters) of mass $`M`$ is
$$n(M)dM=\sqrt{\frac{2}{\pi }}\left(\frac{\delta _c}{\sigma }\right)\frac{d\mathrm{ln}\sigma }{d\mathrm{ln}M}\mathrm{exp}\left(\frac{\delta _c^2}{2\sigma ^2}\right)\frac{\overline{\rho }dM}{M^2}$$
(86)
where $`\sigma =\sigma (R)`$ is the current linear rms fluctuation at the scale $`R`$ corresponding to the mass $`M=4/3\pi R^3\overline{\rho }`$, and $`\overline{\rho }`$ is the mean background. The value of $`\delta _c`$ corresponds to the value of the linear overdensity at the time of collapse. The collapsing structure virializes when the (non-linear) overdensity becomes very large ($`\delta \begin{array}{c}>\hfill \\ \hfill \end{array}100`$). The actual definition is not very important, as once $`\delta \begin{array}{c}>\hfill \\ \hfill \end{array}100`$, the non-linear collapse is quite rapid, as can be seen in the plots of Figure 1, and the corresponding value of $`\delta _l`$ does not change much. Here we will take $`\delta _c`$ to be the critical value where $`\delta \mathrm{}`$; other prescriptions (eg the value of $`\delta _l`$ corresponding $`\delta 178`$) yield similar results. For the standard Einstein-de Sitter case we have $`\delta _c1.686`$. Note that the above abundance depends on the ratio
$$\nu \frac{\delta _c}{\sigma }$$
(87)
The time of collapse or formation is just given by the ratio of $`\delta _c`$ to the linear overdensity $`\delta _l`$ today
$$z_f=\left(\frac{\delta _l}{\delta _c}\right)^{1/\alpha }1$$
(88)
so that an object which has $`\delta _l=\delta _c`$ now, has a formation red-shift $`z_f=0`$, while a fluctuation 4 times larger collapses at $`z_f=3`$ if $`\alpha =1`$ or at $`z_f=1`$ if $`\alpha =2`$.
Non-standard parametrisation of the spherical collapse considered in the previous sections can change the above formalism in two ways. If we label objects by its initial overdensity $`\delta _0`$ then the corresponding $`\delta _l`$ today is
$$\delta _l=\delta _0a^\alpha $$
(89)
So a different value of $`\alpha `$, from the standard GR result ($`\mathrm{\Delta }\alpha \alpha \alpha _{GR}`$), as shown in Figures 2, 4, and 6, will produce a different amplitude of linear fluctuations today. Moreover, as shown in §3.1.3 and Figures 3, 5, and 6, the solution to the spherical collapse equation produces different values of $`\delta _c`$, and therefore different mass functions and formation times. Finally, for a directly measurable quantity, such as the surface density of objects, typically one needs the volume element, which is also a function of the cosmology.
For example if fluctuations are normalized at a given red-shift, $`z_n`$, then the change in $`\delta _l`$ today will be
$$\frac{\mathrm{\Delta }\delta _l}{\delta _l}=\mathrm{\Delta }\alpha \mathrm{log}(1+z_n)$$
(90)
For recombination, eg COBE normalization, we have $`z_n1100`$, and
$$\frac{\mathrm{\Delta }\delta _l}{\delta _l}3\mathrm{\Delta }\alpha $$
(91)
In the case of the BD theory, we can see in Figure 4 that for $`\omega >1`$, $`\mathrm{\Delta }\alpha \alpha \alpha _{GR}>0`$, which means that $`\mathrm{\Delta }\delta _l>0`$. This makes sense as the linear growth is faster and, for fixed initial fluctuations, the final linear overdensity will be larger. As shown in right panel of Figure 5, $`\delta _c`$ will also be larger. Thus in this case the effects tend to compensate each other. This is true for both the formation red-shift $`z_f`$ or for $`\nu `$ in Eq. above. For the formation red-shift $`z_f`$ we have
$$\frac{\mathrm{\Delta }z_f}{1+z_f}\frac{1}{\alpha }\left(\frac{\mathrm{\Delta }\delta _l}{\delta _l}\frac{\mathrm{\Delta }\delta _c}{\delta _c}\right)$$
(92)
which is only valid for small changes. In the BD example given above with $`\omega =10`$ (and COBE normalization) we have that $`\mathrm{\Delta }\delta _l/\delta _l0.9`$ while $`\mathrm{\Delta }\delta _c/\delta _c0.01`$, so the net effect is still quite large. In this case, a formation red-shift of $`z_f=1`$ will change to $`z_f=1.39`$. Thus, a positive finite $`\omega `$ (which corresponds to a larger $`G`$ at high red-shifts) tends to produce larger (earlier) formation red-shifts and higher densities (or larger abundances) at a given red-shift, than the standard model. This goes in the direction of some recent observations (eg see Bahcall & Fan 1998; Robinson, Gawiser & Silk 1998, Willick 1999), which seem to need larger abundances that expected in some standard cosmologies. This interpretation is degenerate with respect to initial conditions and cosmological parameters.
Figure 7 illustrates the large differences in the cluster counts that can be seen between different cosmological models at $`z>1`$ (see Holder et al. 1999 for details). Deviations from General Relativity in the BD models with $`\omega =100`$ and $`\omega =25`$ can be noticed even at low redshift, when models are normalized to CMB fluctuations.
A similar trend is found for the case of Hubble rate $`H^2=a^{3(1+ϵ)}`$ parametrisation. A change of $`|\mathrm{\Delta }ϵ|0.3`$ (allowed by the bounds in Eq.), when normalized to COBE, also produces $`|\mathrm{\Delta }\delta _l/\delta _l|0.9`$ and a smaller effect on $`\mathrm{\Delta }\delta _c/\delta _c`$. This translates into a similar change (of several tens to hundreds of percent) in $`z_f`$. Earlier (later) formation times and larger (smaller) abundances are found for $`ϵ>0`$ ($`ϵ<0`$).
The change in the equation of state $`p=\gamma \rho `$ could produce comparable effects. The allowed values in Eq. of $`|\mathrm{\Delta }\gamma |0.1`$ translate into $`|\mathrm{\Delta }\delta _l/\delta _l|0.3`$, which results in similar changes for $`z_f`$ in either direction, with earlier formation for $`\gamma >0`$.
If the normalization is not fixed, ie we do not quite know what is the value of the initial fluctuation that gave rise to an object we see today (eg a cluster), then all the relative change in the formation or abundance comes through $`\delta _c`$, which tends to produce smaller (later) formation red-shifts ($`\delta _c`$ is larger than the standard GR value) and lower densities (or smaller abundances) at a given red-shift.
## 5. Discussion and Conclusions
We have reconsidered the problem of non-linear structure formation in two different contexts that relate to observations: 1-point cumulants of large scale density fluctuations and the epoch of formation and abundance of structures using the Press & Schechter (1974) formalism. We have use the the shear-free or spherical collapse (SC) model, which is very good approximation for the above applications. We have addressed the question of how different are the predictions when using a non-standard theory of Gravity, such as BD model, or non-standard cosmological model (eg a different equation of state or Hubble law). Note that these are slight variations on the standard theme in the sense that they preserved the main ingredients of GR, such as the covariance and the geometrical aspects of the theory, including the same metric, with only slight changes in the field equations.
We have also presented some preliminary bounds on $`\gamma `$, $`\omega `$ and $`ϵ`$ from observations of the skewness and kurtosis in the APM Galaxy Survey, eg Eq.-. These bounds are optimistic given the current data, but the situation is going to change rapidly, and one can hope to find much better bounds form upcoming data (such as 2DF or SDSS projects). In terms of the equation of state the bounds in Eq. would indicate that our Universe is neither radiation ($`\gamma =1/3`$) or vacuum dominated ($`\gamma =1`$), but somewhere in between (eg matter dominated). In terms of the Gravitational constant, the bounds on $`\omega `$ from Eq. would say that $`G`$ has not changed by more than $`5\%`$ from $`z1.15`$, or by distances of $`400h^1\text{Mpc}`$. Clustering at higher red-shift would probe much larger scales and times. In terms of $`ϵ`$ the bounds Eq., would say that the Hubble law does not differ by more than $`7\%`$ from the standard result (assumed here to be $`ϵ=0`$). We have also shown how halo and cluster abundances and formation times could change in these non-standard cases. The above bounds on $`\gamma `$, $`\omega `$ and $`ϵ`$ from observations of the skewness and kurtosis in the APM still allow significant changes (of several tens to hundreds of percent) on formation red-shifts $`z_f`$ and the corresponding abundances (see §4.2).
In the context of BD models the limits we find for $`\omega `$ are less restrictive than the solar system limits $`\omega \begin{array}{c}>\hfill \\ \hfill \end{array}100`$. However, BD models allow $`\omega =\omega (\varphi )`$ so that $`\omega `$ can increase with cosmic time, $`\omega =\omega (z)`$, in such a way that it could approach the general relativity predictions ($`\omega \mathrm{}`$) at present time and still give significant deviations at earlier cosmological times. It is important to recall that our theory of gravity has only be tested on stellar distances (a.u.) while we want to use it on cosmological scales ($`Mpc`$). Our working example shows, for the first time, how non-linear effects are changed in such a model and sets the framework to study non-linear effects of more complicated (or realistic) Scalar-Tensor theories of gravity.
It is straightforward to combine several of the changes proposed here to explore more general situations. One could for example parameterize theories in the ($`\gamma ,\omega `$) plane, eg different equations of state with different BD parameters, or consider the whole ($`\gamma ,\omega ,\mathrm{\Omega }_M,\mathrm{\Omega }_\mathrm{\Lambda }`$) space. One could also consider a different equation of state for the $`\mathrm{\Lambda }`$component, as in quintessence cosmologies (Caldwell, Dave, Steinhardt 1998), such models have already been used to predict cluster abundances within the “standard” cosmology (see Haiman, Mohr, Holder 2000 and references therein). This would obviously allow for a wider set of possible solutions and degeneracies. One should also consider other observational consequences of these variations, in particular relating to BD theory, such as the age of the Universe, the effects on CMB (eg see Chen & Kamionkowski 1999), radiation-matter transition (Liddle, Mazumdar & Barrow 1998), or the constraints from nucleosynthesis (Santiago et al. 1997). These considerations could rule out some aspects of the proposed variations on the standard model, or might require more elaborate solutions (eg $`\omega =\omega (\varphi )`$ which implies $`\omega =\omega (z)`$). But even if this were the case, we still have learn a few new things about how structure formation depends on the underlying theory of Gravity, which is a first step towards further analysis of these issues.
Throughout this paper we have assumed Gaussian initial conditions and no biasing. Both biasing (eg Fry & Gaztañaga 1993) and non-Gaussianities in the initial conditions (Gaztañaga & Fosalba 1998) would provide an additional source of degeneracy as they might produce similar effects as the non-standard variations presented here. This is the case for example when we have non-zero initial skewness or kurtosis, which could produced quite different values of $`S_3`$ and $`S_4`$ (eg see Gaztañaga & Mahonen 1996; Peebles 1999a,b; White 1999; Scoccimarro 2000), and therefore to the inferred values of $`\nu _2`$ and $`\nu _3`$. Biasing can have a very similar effect (eg see Mo, Jing & White 1997). One would also expect some level of degeneracy with biasing and initial conditions for cluster abundances or formation times (see Robinson, Gawiser & Silk 1998, Willick 1999).
Rather than proposing an alternative theory of gravity or cosmological model, the aim of this paper was to show that some small deviations from the current paradigm have significant and measurable consequences for non-linear structure formation. This could eventually help explaining some of the current puzzles confronting the theory, such as the need of non-baryonic dark matter. Alternatively, current and upcoming observations of non-linear clustering and mass functions can be used to explore our assumptions and place limits on the theory of gravity at large ($`\begin{array}{c}>\hfill \\ \hfill \end{array}1h^1\text{Mpc}`$) scales. This provides an interesting test for gravity as the driving force for structure formation and for our knowledge of the cosmological equation of state. A more comprehensive comparison with particular scenarios is left for future work.
## Acknowledgments
One of us (JAL) gratefully acknowledges financial support from the Spanish Ministry of Education, contract PB96-0384, and also Institut d’Estudis Catalans. EG acknowledges support from CSIC, DGICYT (Spain), project PB96-0925. We would like to thank IEEC, where most of this work was carried out. |
warning/0003/cond-mat0003505.html | ar5iv | text | # The onset of exciton absorption in modulation doped GaAs quantum wells
## Abstract
We study the evolution of the absorption spectrum of a modulation doped GaAs/AlGaAs semiconductor quantum well with decreasing the carrier density. We find that there is a critical density which marks the transition from a Fermi edge singularity to a hydrogen-like behavior. At this density both the lineshape and the transitions energies of the excitons change. We study the density dependence of the singularity exponent $`\alpha `$ and show that disorder plays an important role in determining the energy scale over which it grows.
The absorption spectrum in the presence of a Fermi sea of electrons has been a subject of theoretical and experimental research for more than three decades. The interest in this problem was triggered by the pioneering work of Mahan, who showed that in metals and bulk semiconductors this spectrum should exhibit a power law singularity at the Fermi energy. This singularity, which became known as the Fermi edge singularity (FES), reflects the response of the Fermi sea electrons to the attractive potential of the valence band hole. Mahan’s work was followed by a bulk of theoretical works, which established the many body nature of the FES and provided the tools to treat it .
The FES was indeed observed in the X-ray absorption of metals but was never observed in bulk semiconductors. Its first observation in semiconductors was in modulation doped quantum wells. It is manifested in these systems as a pronounced enhancement of the absorption at the Fermi edge, with an asymmetrical lineshape: a fast rise at the low energy side and a slow fall at high energies. A signature of the FES is its strong dependence on temperature and electron density. The high energy slope becomes more and more gradual when these parameters are increased .
An intriguing aspect of the behavior of a system of many electrons and a hole is the existence of a bound state. Initially, the singularity was associated with a bound state (the so-called Mahan exciton), which is built out of empty conduction band states below the Fermi level, but soon after, it was realized that this exciton is unstable . Nevertheless, a bound state of electrons and a hole should exist in a two-dimensional (2D) system. It is well known that in 2D bound states occur for arbitrarily small attractive potentials. Thus, the attractive potential of a valence hole should have a bound state even after the many-body interactions have been accounted for.
The recent discovery of the charged exciton, X<sup>-</sup>, in semiconductor quantum wells, at low-electron densities created a renewed interest in the problem. The X<sup>-</sup>, which consists of two electrons with opposite spins bound to a valence band hole, appears at the absorption and emission spectra as a spectral line at some energy below the neutral exciton line, X. A natural question which then arises is how does the low density spectrum, of the X<sup>-</sup> and X bound states, transform into the FES. This issue has been addressed in a number of theoretical publications . It was shown that an X<sup>-</sup> like bound state should persist to a high electron density, and be manifested in the absorption spectrum as a spectral line with an asymmetric singular shape of the form $`(\omega \omega _0)^\alpha `$. Another singular peak, which is exciton like, was predicted to appear at some energy above that line. Indeed, this prediction was recently confirmed in CdTe modulation doped quantum wells , where the X-X<sup>-</sup> doublet was observed up to densities of a few$`\times `$10<sup>11</sup> cm<sup>-2</sup>. Attempts to observe the X-X<sup>-</sup> doublet in the absorption spectrum of GaAs quantum wells in the presence of a high density two-dimensional electron gas (2DEG) have yielded inconclusive results .
In this paper we investigate the evolution of the absorption spectrum of a single GaAs/AlGaAs quantum well with decreasing electron density, from $`1.5\times 10^{11}`$ to $`1.5\times 10^{10}`$ cm<sup>-2</sup>. We tune the electron density continuously by applying a voltage to a semi-transparent gate relative to an ohmic contact, which is made into the 2DEG. We find that there is a critical density which marks the transition from a FES to a hydrogen-like behavior. At this critical density both the lineshape and the transitions energies of the excitons exhibit a clear change: the singular asymmetric line becomes a symmetric resonance and the energy difference between X and X<sup>-</sup> becomes constant. We study the density dependence of the singularity exponent $`\alpha `$ and show that the energy scale over which it grows is determined by disorder.
The absorption spectrum is obtained using photocurrent spectroscopy. Electrons are excited from surface states by the incident photons to energies above the Schottkey barrier, and then drift into the doped region. From there they tunnel into the well and are collected by the ohmic contacts of the 2DEG, giving rise to a photocurrent in the gate circuit. This mechanism gives rise to a photocurrent at photon energies which extend well below the GaAs gap. At photon energies which can be absorbed in the quantum well there is an additional contribution to the photocurrent, resulting from electron-hole pairs created in the well. The characteristics of this photocurrent will be discussed in details elsewhere. There are, however, two important experimental points which should be mentioned:
\- In this gated structures light illumination gives rise to some depletion of the carrier density. However, since the capacitance of the structure is constant, this only implies that we need to add some positive gate voltage in order to recover the density in the dark. We have verified this behavior over a large light intensity range.
\- The photocurrent from the quantum well rides on a background signal, due to surface state absorption. The value of this background photocurrent depends on the sample structure and surface treatment. The spectra are displayed after this background was subtracted.
The sample structure that we investigated consists of the following sequence of layers: a $`300`$ nm-thick GaAs buffer, a superlattice with $`50`$ periods of $`100`$ nm Al<sub>0.37</sub>Ga<sub>0.67</sub>As and $`30`$ nm GaAs , a $`20`$ nm-thick GaAs quantum well, a $`50`$ nm-thick Al<sub>0.37</sub>Ga<sub>0.67</sub>As spacer layer, Si delta-doping of $`1\times 10^{12}`$ cm<sup>-2</sup>, a $`100`$ nm Al<sub>0.37</sub>Ga<sub>0.67</sub>As, a $`20`$ nm thick n-type Al<sub>0.37</sub>Ga<sub>0.67</sub>As, and a GaAs cap layer. The wafer is processed into a field-effect transistor structure, with a $`5`$ nm thick semi-transparent PaAu gate. Determination of the electron density $`N`$ under illumination is done by measuring the photoluminescence (PL) spectrum as a function of magnetic field, and finding the magnetic field values at which an integer number of Landau levels are filled. At these magnetic fields there is an abrupt change in the PL spectrum . This procedure works well up to $`N8\times 10^{10}`$ cm<sup>-2</sup>. To obtain the lower electron densities we use the capacitance of this structure to extrapolate the curve of $`N(V_\text{g})`$, where $`V_\text{g}`$ is the gate voltage.
The measurements are conducted in a liquid helium storage dewar at a temperature of $`4.2`$ K, using an optical fiber based system. We use a single-mode fiber at close proximity, $`100`$ $`\mu `$m, to illuminate the sample. A second multi-mode fiber is set at a distance of a few mm from the sample and collects the emitted PL. The sample is excited by a tunable Ti-sapphire laser, which is coupled into the single-mode fiber. The photocurrent is amplified by a sensitive current amplifier and measured using a lock-in amplifier. We have conducted the measurements using power levels from $`10`$ nW to $`1`$ mW. Except for a shift in the $`N(V_\text{g})`$ curve and some broadening at high power levels the results are power independent. After an absorption measurement at a given gate voltage $`V_\text{g}`$, the laser is set to $`1.541`$ eV and the PL is measured.
Figure 1 shows a set of photocurrent and emission spectra taken for various gate voltages. We label each spectrum with the corresponding Fermi energy E$`{}_{\text{F}}{}^{}=\pi `$h $`{}_{}{}^{2}N/m^{}`$, where $`N`$ is the electron density and $`m^{}`$ is the electron effective mass ($`m^{}=0.067m_0`$). At large electron densities (Figs. 1a) the absorption edge is a broad step which is shifted to high energies relative to the emission. In Fig. 1b, which is taken at E$`{}_{\text{F}}{}^{}=2.5`$ meV, we observe the formation of two broad peaks: one at the absorption edge and the other a few meV higher. As the density is further reduced (E$`_\text{F}`$ $`<1.1`$ meV) the two broad peaks acquire an asymmetric singular lineshape, characterized by a steep rise at low energies and a slow fall at high energies (Fig. 1c-d). Following the notation of Ref. we label the low energy peak as $`\omega _1`$ and the high energy one as $`\omega _2`$. The singularity in $`\omega _2`$ increases very fast, and at E$`{}_{\text{F}}{}^{}=0.75`$ meV it becomes a symmetric resonance, the heavy-hole exciton (Fig. 1e). The $`\omega _1`$ peak exhibits less pronounced changes. As the density is reduced it becomes weaker and evolves into the well known charged exciton, X<sup>-</sup> (Figs. 1e-f). A replica of that scenario appears at the light-hole exciton energy, 4 meV higher.
The evolution of the singularity is surprisingly fast. Figure 2 describes a fit of a power law singularity $`A(\omega )=(\omega \omega _0)^\alpha `$ to the high side energy of $`\omega _2`$ at E$`{}_{\text{F}}{}^{}=0.9`$ meV. Such a fitting procedure allows us to accurately extract the exponent $`\alpha `$ at each density. The inset shows the dependence of $`\alpha `$ on E$`_\text{F}`$. It is seen that $`\alpha `$ increases by nearly an order of magnitude (from $`0.05`$ to $`0.4)`$ in a very narrow range, $`\mathrm{\Delta }`$E$`{}_{\text{F}}{}^{}=0.25`$ meV, which corresponds to reducing the electron density by less than $`1\times 10^{10}`$ cm<sup>-2</sup>.
The exponent $`\alpha `$ is related to the phase shift of the electrons at the Fermi surface, when scattering off the valence hole potential . In that sense it measures the efficiency of the Fermi sea electrons in screening that potential: the smaller $`\alpha `$ is, the better is the screening. Previous studies have reported a significantly broader density range over which the singularity was observed . For example, in Ref. a similar change of $`\alpha `$ was observed on an energy scale of $`\mathrm{\Delta }`$E$`{}_{\text{F}}{}^{}10`$ meV, more than an order of magnitude larger than in our measurements. It was argued that the energy scale over which $`\alpha `$ changes is related to the many body nature of the problem . As we shall show in the following, disorder plays a critical role and determines the energy scale at which the singularity is observed.
Examining the PL spectra, we notice that at the density, at which $`\alpha `$ starts to grow the PL lineshape transforms from a broad single peak to two narrow resonances, associated with the X and X<sup>-</sup>. This behavior of the PL was extensively studied by some of us . We have shown that the appearance of excitons marks the onset of strong localization of the electrons in the potential fluctuations of the ionized donors. These donors, which are at a distance of 50 nm (the spacer width), are randomly distributed in the plane and induce a spatially fluctuating electrostatic potential at the 2DEG. At high electron density the 2DEG efficiently screens the fluctuations, but as the density is reduced the screening becomes less efficient and they grow considerably . Thus, the growth of $`\alpha `$ is related to the onset of strong disorder in the sample.
It should be noticed that this sample is of high quality. This is evidenced by the high mobility ( $`10^6`$ cm<sup>2</sup>V<sup>-1</sup>s<sup>-1</sup>), the narrow exciton linewidth ($`0.3`$ meV), and the absence of a Stokes shift between the PL and absorption ($`<0.1`$ meV) . Thus, at E$`{}_{\text{F}}{}^{}>1.1`$ meV, where the potential fluctuations are suppressed, our sample could be considered as close to ideal. Nevertheless, the singularity is significantly suppressed at that range. Only when disorder sets in, and electrons can not efficiently screen the valence-hole potential, $`\alpha `$ becomes large enough to make $`\omega _2`$ and $`\omega _1`$ observable as singular peaks. Indeed, it is indicated in Ref. that there is strong disorder present in the samples and the mobility is rather low. Large disorder is also present in the narrow quantum wells studied in Ref. , as is evident by the very broad exciton line $`8`$ meV. The wide range in which singular lineshapes are observed in these works is, therefore, related to the large disorder in the samples which were used and not to a fundamental energy scale. The use of a gated sample in this study allows us to control the onset and amount of disorder.
Let us now turn to examining the density dependence of the energies of the $`\omega _1`$ and $`\omega _2`$ peaks (Fig. 3). The data is presented for a limited density range, where the singularity edge is clear and the peak energy can be unambiguously determined. In Fig. 3a we show the $`\omega _1`$ and $`\omega _2`$ transition energies. It is seen that the $`\omega _2`$ peak shifts to higher energies with increasing density while the energy of $`\omega _1`$ remains nearly constant. In Fig. 3b we present the energy difference h $`(\omega _1`$ $`\omega _2)`$ as a function of E$`_\text{F}`$. It is evident that there exists a threshold density value, $`N_\text{c}2.5\times 10^{10}`$ cm<sup>-2</sup> (which corresponds to E$`_\text{F}`$ $`=0.85`$ meV), below which h $`(\omega _2\omega _1)`$ is constant and equals $`1.2`$ meV. Only above this threshold density the energy difference h $`(\omega _2\omega _1)`$ grows. This behavior is in clear contrast with that reported in Ref. , where it was argued that h $`(\omega _2\omega _1)`$ changes linearly with E$`_\text{F}`$, all the way to zero density. The value of $`1.2`$ meV is well documented experimentally as the X<sup>-</sup> binding energy in a $`20`$ nm GaAs quantum well . In fact, it is the energy difference between the two peaks in the PL spectra (Fig. 1). It should be emphasized that the optical spectra (Figs. 1e-f) clearly show that the density is indeed changing in this range, as evidenced by the exchange of oscillator strength between the X and the X<sup>-</sup>. This conclusion is supported by transport measurements through the sample at this range, which show that the conductivity decreases as the gate voltage becomes more negative.
To explain this behavior let us consider the energy spectrum of the system , which is schematically shown at the inset of Fig. 3b ($`\mu `$ is the chemical potential and E$`_\text{c}`$ is the single particle band gap). At the limit of infinitesimal density there are two bound states, X and X<sup>-</sup>, at a relatively large energy distance (the exciton binding energy, E$`_\text{X}`$ ) below E$`_\text{c}`$. The energy separation between the two bound states is E$`_\text{B}`$, the X<sup>-</sup> binding energy. On the other hand, at the limit of high electron density there is only one, X<sup>-</sup>-like, bound state. The absorption spectrum of the system was predicted to consist of two peaks, with an energy difference given by
$$\text{h}\text{ }\text{ }(\omega _2\omega _1)=\mu \epsilon _\mathrm{b}$$
(1)
where $`\epsilon _\mathrm{b}`$ is the binding energy of the X<sup>-</sup> like bound state. It is important to note that both $`\mu `$and $`\epsilon _\mathrm{b}`$ in Eq. 1 are measured with respect to the same level, which we take as the bottom of the conduction band at zero electron density, E$`{}_{}{}^{0}{}_{\text{c}}{}^{}`$. Equation 1 has a simple meaning: it describes the energy cost for ionizing the X<sup>-</sup>-like bound state by exciting one of the two electrons to the chemical potential level $`\mu `$. However, it is clear that this relation is valid only at the high density limit. At low densities $`\mu 0`$ and $`\epsilon _\mathrm{b}>`$ E$`_\text{X}`$. Hence, Eq. 1 would imply that h $`(\omega _2\omega _1)>`$ E$`_\text{X}`$, and that we should observe a huge increase in h $`(\omega _2\omega _1)`$ below a certain density. Thus, the use of Eq. 1 down to zero density is incorrect , and consequently the determination of the X<sup>-</sup> binding energy in CdTe has to be re-examined. Our experimental findings show that below $`N_\text{c}`$ the correct relation is h $`(\omega _2\omega _1)=`$ E$`_\text{B}`$. Only at large densities, where there is only one bound state, we are at the limit covered by Eq. 1.
It is remarkable that at the threshold density $`N_\text{c}`$ the exciton lineshape undergoes a drastic change. This change is clearly seen in Fig. 1: the exciton lineshape is symmetric below $`N_\text{c}`$ (Fig. 1e-f) and becomes a singular asymmetric line above it (Figs. 1b-d). The comparison between the exciton lineshape in Figs. 1e and 1d is particularly interesting. The difference is manifested not only in the high energy side but also at the low energy side, which becomes steeper above $`N_\text{c}`$. The fact that both the energy dispersion and lineshape change at $`N_\text{c}`$ implies that this density marks the transition from a hydrogen-like behavior to a FES. It is interesting to note that the value of $`N_\text{c}`$, which is found here, is very close to that reported in electroreflectance measurements , where quenching of the excitonic absorption was studied.
Finally, we wish to comment on the behavior at high densities. It can be seen at Fig. 3b that at densities above $`N_\text{c}`$ the energy separation between the X and X<sup>-</sup> grows monotonically. The dashed line describes the relation h $`(\omega _2\omega _1)=`$ E$`_\text{B}`$ $`+`$ E$`_\text{F}`$, (where E$`{}_{\text{B}}{}^{}=1.2`$ meV). This relation results from Eq. 1 and was used to fit the data in Ref. . It is evident that all the measured points in Fig. 3b are below this curve, but as the density gets higher the data points are approaching the curve. Unfortunately, we can not determine the behavior accurately at the high density limit. The X and X<sup>-</sup> lines become too broad, and one can not reliably obtain their energy separation. Nevertheless, it is obvious from our data at higher densities that the energy separation indeed increases monotonically with density.
This research was supported by the Minerva foundation. |
warning/0003/astro-ph0003452.html | ar5iv | text | # Non-conservative Evolution of Cataclysmic Variables
## 1 Introduction
Cataclysmic variables (CVs) are close binary systems consisting of a low-mass main-sequence (MS) star that fills its Roche lobe and a white dwarf. Main-sequence star (mass-donating star, usually referred to as star 2 or secondary) loses matter through the vicinity of the inner Lagrangian point $`L_1`$. White dwarf (accretor star, referred to as star 1 or primary) accretes at least a fraction of this matter via accretion disk or via accretion columns in the polar zones in the case when white dwarf has a strong magnetic field.
Physics and evolution of CVs as well as evolution of similar to them low-mass X-ray binaries were investigated starting from the late sixties (see, e.g., and references therein). These studies were motivated by the recognition of the fact that evolution of CVs is determined by losses of orbital angular momentum due to radiation of gravitational waves and magnetically coupled stellar wind . A number of studies investigated the influence of the loss of angular momentum associated with the material outflow from the system on the CV evolution (see, e.g., ). However, in the absence of gas dynamical simulations of the mass transfer in the binaries, these processes were considered under the parametric approach.
Recent three-dimensional (3D) gas dynamical simulations of the structure of gaseous flows in semi-detached binaries proved an important role of a circumbinary envelope (see, e.g. ). These calculations also show that during mass exchange a significant fraction of matter leaves the system. The present work is mainly devoted to the numerical investigation of the evolution of CVs using the data of 3D gas dynamical calculations on the losses of mass and angular momentum from close binaries. The main attention is paid to the stability of mass exchange against runaway mass loss and its dependence on the donor to accretor mass ratio $`q=M_2/M_1`$.
It is known that the loss of matter by the donor star results in the violation of its hydrostatic and thermal equilibrium. Hydrostatic equilibrium is restored adiabatically, i.e. in dynamical time scale, while thermal equilibrium is restored in thermal time scale. This transition to a new equilibrium state leads to a change of donor radius. The sign of radius variation depends on the convective and radiative stability of outer envelope of the star. For stars with masses $`M\stackrel{<}{}M_{}`$ with deep convective envelope as well as for white dwarfs mass loss results in increase of the radius, while for stars with radiative envelopes it leads to the shrinkage of the star. Mass exchange in close binary is unstable when in the course of evolution mass-losing star tends to overfill the Roche lobe. It can occur when radius of the donor $`R_2`$ increases faster (or decreases slower) than the effective radius of Roche lobe $`R_{RL}`$.<sup>1</sup><sup>1</sup>1The effective radius of the Roche lobe $`R_{RL}`$ is determined as the radius of a sphere with a volume equal to the volume of the Roche lobe. The case when the radius of the donor increases while the radius of the Roche lobe decreases is also possible. Thus, the question of stability of mass exchange against runaway is determined by the balance of derivatives $`R_2/M_2`$ and $`R_{RL}/M_2`$. It means that unstable mass exchange is possible even for stars which contract when losing matter.
The study of conditions of stable mass exchange is motivated by two related problems: i) it is necessary to explain why $``$ 10% of CVs with well-determined masses of components have combination of donor mass and mass ratio of components which is “forbidden” in the evolutionary models based on conventional assumptions on the conservation of angular momentum of a binary; ii) for population synthesis studies of CVs it is necessary to distinguish progenitors of CVs among all binaries containing a white dwarf and a low-mass companion (i.e., specify in which binaries of this type stable mass exchange is possible). A problem similar to the last one appears in the studies of low-mass X-ray binaries.
The paper is organized as follows. In Section 2 we describe major factors affecting the evolution of CVs and conditions of stable mass exchange in binaries. The reasons which motivated us to rule out the conservative approximation for mass exchange are considered in Section 3. Sections 4 and 5 describe some results of 3D gas dynamical simulations of the flow structure in semi-detached binaries and introduce the model for description of angular momentum losses due to non-conservative mass exchange. Results of evolutionary calculations for conservative and “non-conservative” models are compared in Section 6. Criteria of stable mass exchange for “non-conservative” model with mass and angular momentum loss are determined in Section 7. Section 8 summarizes the main results of the work.
## 2 Major factors affecting the evolution of CVs
In the context of the present work the term “stable mass exchange” means dynamical stability. In other words, in the course of mass exchange the donor star can be out of thermal equilibrium and mass transfer rate may exceed the rate corresponding to the Kelvin time scale ($`\stackrel{}{M}3\times 10^7RL/M`$, where radius $`R`$, luminosity $`L`$, and mass of the star $`M`$ are in solar units and mass transfer rate $`\stackrel{}{M}`$ is in $`M_{}\text{yr}^1`$). Let us define unstable mass exchange as the situation when the radius of the donor changes faster than the effective radius of Roche lobe: $`\underset{2}{\overset{}{R}}>\underset{RL}{\overset{}{R}}`$. Using the derivatives of the radii w.r.t. mass of the donor we can rewrite the condition of the stability of mass transfer as:
$$\zeta _{}\frac{\mathrm{ln}R_2}{\mathrm{ln}M_2}>\frac{\mathrm{ln}R_{RL}}{\mathrm{ln}M_2}\zeta _{RL}.$$
(1)
The effective radius of Roche lobe can be estimated, for instance, using interpolation formula given by Eggleton :
$$R_{RL}0.49A\frac{q^{2/3}}{0.6q^{2/3}+\mathrm{ln}(1+q^{1/3})},$$
where $`A`$ is the semimajor axis of the orbit. For $`q\stackrel{<}{}1`$ the approximation suggested by Paczyński is more convenient:
$$R_{RL}\frac{2}{3^{4/3}}A\left(\frac{q}{1+q}\right)^{1/3}.$$
(2)
Derivative of the Roche lobe radius can be written as
$$\frac{\mathrm{ln}R_{RL}}{\mathrm{ln}M_2}=\frac{\mathrm{ln}A}{\mathrm{ln}M_2}+\frac{\mathrm{ln}(R_{RL}/A)}{\mathrm{ln}M_2}.$$
(3)
It is seen from (2) and (3) that the variation of effective radius of Roche lobe is defined by the change of donor and accretor masses $`M_2`$ and $`M_1`$ as well as by the change in the separation $`A`$, and, hence, it is influenced by the loss of mass and angular momentum from the system. As a rule, only systemic angular momentum is considered as the orbital momentum, i.e. the sum of angular momenta of two stars, which are considered as point masses. Momenta of the spin of the components, momentum of the accretion disc (if it’s present), as well as momentum of gaseous flows inside the system as a rule are not included into the total momentum. Deviation of the momentum of the donor from that of the point mass object (which can be significant, see ) usually is not taken into account as well. Such a simplified approach is compelled, since it is very difficult to include all the abovementioned factors into calculations. Having in mind all these simplifications, one gets for the orbital momentum of a binary system with the circular orbit
$$J=\sqrt{\frac{GM_1^2M_2^2A}{M_1+M_2}},$$
(4)
where $`G`$ is the gravitational constant.
According to the widely accepted models, the evolution of cataclysmic binaries is determined by the loss of angular momentum from the system via gravitational waves radiation (GWR) and/or magnetic stellar wind (MSW) of the donor as well as the mass exchange between components. In the standard models of evolution it’s accepted that the mass exchange among components does not change the systemic angular momentum, and influences it indirectly – through a possible loss of mass from the system and angular momentum ablation by leaving matter (term with subscript LOSS in Eq. 5 below). Therefore, within conventional models, the equation for the change of orbital momentum can be written as:
$$\frac{dJ}{dt}=\left(\frac{J}{t}\right)_{\text{GWR}}+\left(\frac{J}{t}\right)_{\text{MSW}}+\left(\frac{J}{t}\right)_{\text{LOSS}}.$$
(5)
Let us consider the terms of expression (5):
1) Loss of the systemic angular momentum due to GWR
The change of the systemic orbital momentum as a result of GWR is given by a formula (see, e.g., ):
$$\left(\frac{\mathrm{ln}J}{t}\right)_{\text{GWR}}=\frac{32G^3}{5c^5}\frac{M_1M_2(M_1+M_2)}{A^4},$$
(6)
where $`c`$ is the speed of light. Intensity of GWR is characterized by a very strong dependence on the separation of components and respectively on the orbital period. This process is essential for short-period systems with $`P_{orb}\stackrel{<}{}10^\text{h}`$, since only then the time scale of angular momentum loss becomes shorter than the Hubble time.
2) Loss of the systemic angular momentum due to MSW of the donor
A mechanism for the loss of the angular momentum of the system by means of the MSW was suggested by Schatzman and by Mestel . If the star has a convective envelope and, subsequently, a surface magnetic field, its own axial rotation is braked by the magnetic stellar wind, and the angular momentum loss rate can be essential even for small mass loss rate. The subsequent synchronization of the spin of the donor and orbital rotation caused by tidal interaction between components results in the loss of the orbital momentum of the system and reduction of the separation of components $`A`$. To take this effect into account, Verbunt and Zwaan extrapolated the data on rotation of F-type stars to the K and M type components of CVs and suggested a widely used semi-empirical formula which allows to determine the temporal behavior of the orbital momentum:
$$\left(\frac{\mathrm{ln}J}{t}\right)_{\text{MSW}}=0.5\times 10^{28}Ck^2G\frac{(M_1+M_2)^2}{M_1}\frac{R_2^4}{A^5},$$
(7)
were $`k`$ is the radius of gyration of the donor star; $`C`$ is a numerical factor determined from the comparison of theoretical calculations and observational data. Well-defined reduction of magnetic field for stars with masses less than $`0.3M_{}`$ permits to suggest that vanishing of a radiative core of a star when its mass decreases to this threshold leads to an abrupt “switch-off” of the so-called “$`\alpha \mathrm{\Omega }`$”-dynamo mechanism that is responsible for the generation of stellar magnetic fields . At the instant of cessation (or sharp reduction) of MSW the time scale of angular momentum loss \[defined by Eq. (7)\] is shorter than the thermal time scale of a donor with $`M0.3M_{}`$. Thus, the latter is out of thermal equilibrium and its radius is larger than the radius of thermally equilibrium star of the same mass.<sup>2</sup><sup>2</sup>2For discussion on the response of star on mass loss in different time scales see, e.g., . When the action of MSW and corresponding loss of angular momentum stop, the rate of reduction of the semimajor axis declines, the donor shrinks down to the equilibrium radius and loses the contact with Roche lobe . Meanwhile, the system continues to lose angular momentum due to GWR, components continue to approach each other and after some time the optical star fills its Roche lobe again. Further evolution of the system is determined by the loss of angular momentum via GWR.
Cessation of MSW after the star becomes completely convective is widely adopted hypothesis explaining satisfactorily the so-called orbital period gap in CVs. The value of coefficient $`C`$ in Eq. (7) can be determined by comparison of theoretical width of period gap and observational one. In the present study we adopted $`C=3.0`$ in accordance to .
3) Loss of the systemic angular momentum due to mass loss
Mass loss from the system is mainly considered as a parameter fine tuning of which permits to get an agreement with observational minimal period of CVs (e.g. ). More definite assumptions are made only for specific cases when examining evolution of the systems, in which: i) the rate of accretion on white dwarf is limited by the hydrogen burning rate ($`10^7÷10^6`$$`M_{}\text{yr}^1`$), and the mass loss rate by the donor is much higher than this limit, but does not exceed Eddington limit for the dwarf (that is close to $`1.5\times 10^5`$$`M_{}\text{yr}^1`$ , see, e.g., ), the mass excess being lost by means of stellar wind ; ii) mass is lost due to outbursts occurring on white dwarf (see, e.g., )<sup>3</sup><sup>3</sup>3The latter two studies considered also the loss of momentum due to interaction of the donor with the envelope of a Nova.. In both cases it is usually assumed that the specific angular momentum of the matter leaving the system is equal to the specific momentum of the accretor.
To describe the loss of mass and momentum out of the system, it is convenient to introduce parameters describing the degree of non-conservativity of the evolution respective to the mass :
$$\beta =\frac{\underset{1}{\overset{}{M}}}{\underset{2}{\overset{}{M}}}=\frac{M_1}{M_2}$$
and to the angular momentum:<sup>4</sup><sup>4</sup>4Sometimes specific angular momentum lost from the system is measured in units of $`\mathrm{\Omega }A^2`$ and used instead of $`\psi `$:
$$\alpha =\frac{\stackrel{}{J}}{\stackrel{}{M}}/\mathrm{\Omega }A^2=\frac{\stackrel{}{J}}{\stackrel{}{M}}/\frac{J}{\mu }=\frac{q}{(1+q)^2}\psi ,\mu =\frac{M_1M_2}{M}.$$
$$\psi =\frac{\stackrel{}{J}}{\stackrel{}{M}}/\frac{J}{M}=\frac{\mathrm{ln}J}{\mathrm{ln}M},$$
(8)
where $`M=M_1+M_2`$.
With these parameters, the expression for the loss of momentum by the matter leaving the system gets the following form:
$$\left(\frac{J}{t}\right)_{\text{LOSS}}=(1\beta )\underset{2}{\overset{}{M}}\psi \frac{J}{M}.$$
In the calculations of the cataclysmic binary evolution an equation for the variation of the semiaxis of the orbit is used directly, instead of the equation for the variation of momentum over time. Let us differentiate the expression (4), substitute the result into equation (5), and use the relation $`\underset{1}{\overset{}{M}}=\beta \underset{2}{\overset{}{M}}`$. Then the equation for the variation of orbital semiaxis can be written as:
$$\left(\frac{dA}{dt}\right)=\left(\frac{A}{t}\right)_{\text{EXCH}}+\left(\frac{A}{t}\right)_{\text{LOSS}}+\left(\frac{A}{t}\right)_{\text{GWR}}+\left(\frac{A}{t}\right)_{\text{MSW}}.$$
(9)
Here, index ‘EXCH’ corresponds to the change of $`A`$ resulting from the mass transfer between components. Note that while the term $`(J/t)_{\text{EXCH}}`$ is absent in the equation (5), the appropriate term $`(A/t)_{\text{EXCH}}`$ is present in the equation (9), being a natural consequence of the dependence of the orbital momentum both on semimajor axis of the orbit and masses of components. When mass exchange is conservative, the evolution of the orbital separation is determined by the assumption that the mass transfer does not change the orbital momentum of the system. In the course of mass exchange the matter that initially had specific angular momentum of the donor is transferred onto accretor and finally gets the specific momentum of the accretor. The assumption of conservation of the orbital momentum implies that if the mass of the donor is lower than the mass of accretor, the excess momentum of accreted gas should transform into the orbital one, and the mass exchange acts in the direction of increasing the orbital semiaxis. If the donor is more massive than the accretor, the lacking momentum of the accreted gas should be taken from the orbital momentum, therefore, the mass exchange acts in the direction of shrinking of the binary system.
The numerical investigation of the evolution of cataclysmic binaries consists of simultaneous calculations of the evolution of the donor and of the variation of the orbital separation in time. Let us consider the processes that determine the evolution of the orbital semiaxis. Using parameters $`\beta `$ and $`\psi `$, we can write formulae for the terms of equation (9) (here masses and distances are given in solar units and time is given in years):
$$\left(\frac{A}{t}\right)_{\text{EXCH}}=2A\frac{M_2M_1}{M_1M_2}\beta \underset{2}{\overset{}{M}},$$
(10)
$$\left(\frac{A}{t}\right)_{\text{LOSS}}=(1\beta )A\frac{2M_1(1q\psi )+M_2}{M_2(M_1+M_2)}\underset{2}{\overset{}{M}},$$
(11)
$$\left(\frac{A}{t}\right)_{\text{GWR}}=1.65\times 10^9\frac{M_1M_2(M_1+M_2)}{A^3},$$
(12)
$$\left(\frac{A}{t}\right)_{\text{MSW}}=6\times 10^7C\frac{(M_1+M_2)^2}{M_1}\left(\frac{R_2}{A}\right)^4.$$
(13)
It is seen from Eqs (10)–(13) that changing of $`R_{RL}`$ is mainly determined by the value of $`q`$. On the other hand, changing of $`R_2`$ depends on $`M_2`$ and $`\underset{2}{\overset{}{M}}`$. Thus the most convenient way for studying of limits of stable mass exchange is analysis of “ratio of donor-star mass to accretor mass $`q`$ – donor-star mass $`M_2`$” diagram.
In a number of cases one may simplify expression (3) and obtain analytical formulae for the derivative of the Roche lobe radius. In particular, when the angular momentum changes due to GWR and MSW can be neglected, one may deduce from (10) and (11) the expression for the change of separation $`A`$ \[the first term in the formula (3)\] as a function of mass ratio $`q`$ and the parameters of the loss of matter and momentum $`\beta `$ and $`\psi `$:
$$\frac{\mathrm{ln}A}{\mathrm{ln}M_2}=\frac{2\psi (1\beta )q+\beta q+2\beta q^2q2}{1+q}.$$
(14)
Using the expression (2) for $`R_{RL}`$ we can express the second term of (3) as:
$$\frac{\mathrm{ln}(R_{RL}/A)}{\mathrm{ln}M_2}={}_{}{}^{1}/_{3}^{}{}_{}{}^{1}/_{3}^{}\frac{(1\beta )q}{1+q}.$$
(15)
Using (14) and (15), for the case of totally conservative mass exchange ($`\beta =1`$), we obtain:
$$\frac{\mathrm{ln}R_{RL}}{\mathrm{ln}M_2}=2q{}_{}{}^{5}/_{3}^{}.$$
(16)
From the same equations, for the so-called ‘Jeans mode’ of mass loss, when the mass exchange in the system does not occur, but the donor loses its mass due to stellar wind and matter leaves the system carrying away specific angular momentum of the donor ($`\beta =0`$, $`\psi =q`$), we obtain:
$$\frac{\mathrm{ln}R_{RL}}{\mathrm{ln}M_2}={}_{}{}^{1}/_{3}^{}\frac{13q}{1+q}.$$
(17)
For the case of stellar wind from the accretor when all matter lost by the donor is transferred onto accretor but then leaves the system carrying away specific angular moment of the accretor (formally, one should take here $`\beta =0`$ and $`\psi =q^1`$) we obtain:
$$\frac{\mathrm{ln}R_{RL}}{\mathrm{ln}M_2}=\frac{2q^2q{}_{}{}^{5}/_{3}^{}}{1+q}.$$
(18)
The comparison of the rates of variation of the radius of the donor and effective radius of the Roche lobe (3) determines the state of stability/instability of mass exchange in any phase of evolution. It is possible to determine the boundaries of stable mass exchange region by two different methods:
* It is possible to calculate derivatives of the radius of the donor w.r.t. its mass directly for different $`\underset{2}{\overset{}{M}}`$ and to compare it with derivatives of the effective radius of Roche lobe depending on $`M_1`$, $`M_2`$, and $`A`$. This method was used, for instance, in and it permits to restrict a region of $`qM_2`$ diagram where mass exchange is dynamically stable and also to detect the regions where mass exchange should occur in different time scales: thermal, nuclear, or in the time scale of angular momentum loss.
* Alternatively, a requirement that the rate of accretion should not exceed the Eddington limit for the dwarf can be accepted as a criterion of stable mass exchange. This assumption seems to be more reasonable since evolutionary calculations often suggest very high rate of mass loss during very initial moments of time after overfilling the Roche lobe, while after that $`\stackrel{}{M}`$ falls off rapidly. In fact, this approach involves a restriction on dynamical stability of mass transfer but permits loss of matter in thermal time scale.
## 3 Stable mass exchange region in conservative on mass model of the CV evolution
Figure 1 depicts the boundary of the stable mass exchange region in conventional conservative model of the evolution of CV. Boundary 1 was found in , and boundary 2 was found in . These boundaries coincide well for of donors of low mass but differ a little for large ones. Unfortunately, in the procedure of the calculation of the boundary isn’t described in detail. We may assume that this discrepancy is due to the differences in the codes, in a different treatment of possible weakening of donor MSW at masses greater than $`1M_{}`$, and different techniques for computation of stellar radii derivatives in this region.
Figure 1 also shows CVs with known masses of components from the catalog . The total number of binaries in the diagram is 80, and 11 of them are in the region of unstable mass exchange. For 8 of these 11 systems uncertainties of parameters are known, they are shown in Fig. 1 as error bars. It is seen that the location of stars in a region of unstable mass exchange can not be explained by uncertainties in masses of components. In Fig. 1a we plot CVs with the best determined masses , parameters for a number of systems being different compared to the data from . Nevertheless, even for well-determined systems the problem of “improper” location of CVs remains – 3 out of 22 systems lie in the “forbidden” region.
To explain this disagreement we suggest a model of CVs evolution which involves loss of mass and angular momentum from the system. This model is based on 3D gas dynamical simulations of the flow structure in close binary systems.
## 4 Calculations of the rate of an angular moment loss in 3D gas dynamical models
The above-stated analysis proves that the rate of change of the distance between components $`A`$ in evolutionary calculations depends heavily on two parameters: the degree of non-conservativity of mass exchange w.r.t. the matter $`\beta `$ and w.r.t. to the angular momentum $`\psi `$. These parameters can be determined only by means of 3D gas dynamical simulations of mass transfer in binary systems. These investigations were made in for a binary system with mass ratio of components $`q={}_{}{}^{1}/_{5}^{}`$, and in for a binary with $`q=1`$. Besides that, we have specially made a numerical simulation of a binary system with $`q=5`$ to study the influence of parameter $`q`$ on the flow structure.
Analysis of the obtained results proves that in the binary systems considered by us for different values of parameter $`q`$ mass transfer is non-conservative and the non conservativity degree w.r.t. the matter is $`\beta 0.4÷0.6`$. Simulations of flows in the 3D numerical models prove also that the matter leaves the system with a significant angular momentum. The non conservativity degree of mass exchange w.r.t. the matter can be easily determined directly from gas dynamical model, while the determination of non conservativity degree w.r.t. the angular momentum is a sophisticated problem. Let us consider the equations which determines the transport of angular momentum in the system. We should consider Navier-Stokes equations for a correct analysis because viscosity plays a significant role in the redistribution of angular momentum. From the steady-state gas dynamical equations in a rotating frame we can deduce the equation for the angular momentum transfer
$$u\frac{\lambda }{x}+v\frac{\lambda }{y}+w\frac{\lambda }{z}+\frac{1}{\rho }\left((𝒓𝒓_{CM})\times div𝒫\right)_z=\left((𝒓𝒓_{CM})\times grad\mathrm{\Phi }\right)_z,$$
where, as usually, $`𝒗=(u,v,w)`$ is the velocity vector, $`\rho `$ is the density, $`\mathrm{\Phi }`$ is the Roche potential, $`\mathrm{\Omega }`$ is the orbital angular velocity, $`𝒫`$ is the stress tensor, $`𝒓_{CM}`$ is the radius-vector of the center of masses of the system, and angular momentum $`\lambda `$ (in a laboratory, i.e. inertial or observer’s frame) is defined by expression:
$$\begin{array}{c}\lambda =(xx_{CM})vyu+\mathrm{\Omega }\left((xx_{CM})^2+y^2\right)\hfill \\ \\ =(xx_{CM})\left(v+\mathrm{\Omega }(xx_{CM})\right)y(u\mathrm{\Omega }y).\hfill \end{array}$$
Let us write down the equation for angular momentum transfer in a divergent form:
$$div(\rho \lambda 𝒖)+\left((𝒓𝒓_{CM})\times div𝒫\right)_z=\rho \left((𝒓𝒓_{CM})\times grad\mathrm{\Phi }\right)_z.$$
From this equation, we can obtain the integral law of change of angular momentum as follows:
$$\underset{\mathrm{\Sigma }}{}\rho \lambda 𝒖\text{d}𝒏+\underset{\mathrm{\Sigma }}{}\left((𝒓𝒓_{CM})\times 𝒫\text{d}𝒏\right)_z=\underset{V}{}\rho \left((𝒓𝒓_{CM})\times grad\mathrm{\Phi }\right)_zdV\mathrm{\Pi }.$$
The term $`\mathrm{\Pi }`$ determines the change of the angular momentum due to non central nature of the force field, defined by the Roche potential. By definition, the value
$$𝑭_\lambda =\rho \lambda 𝒖+𝒫𝒓^{\mathbf{}},\text{where}𝒓^{\mathbf{}}=(y,xx_{CM},0)$$
(19)
is the specific flux of angular momentum $`\lambda `$. As a result, for the steady-state case we obtain:
$$\underset{\mathrm{\Sigma }_1}{}𝑭_\lambda \text{d}𝒏=\underset{\mathrm{\Sigma }_2}{}𝑭_\lambda \text{d}𝒏+\underset{\mathrm{\Sigma }_3}{}𝑭_\lambda \text{d}𝒏\mathrm{\Pi },$$
where $`\mathrm{\Sigma }_1`$ and $`\mathrm{\Sigma }_2`$ are the boundaries of the donor and the accretor, respectively, and $`\mathrm{\Sigma }_3`$ is the outer boundary of the calculation domain. By analogy to the flux of matter
$$\stackrel{}{M}=\underset{\mathrm{\Sigma }_3}{}𝑭_m\text{d}𝒏=\underset{\mathrm{\Sigma }_3}{}\rho 𝒖\text{d}𝒏,$$
one may estimate the flux of angular momentum from the system as
$$\stackrel{}{J}=\underset{\mathrm{\Sigma }_3}{}𝑭_\lambda \text{d}𝒏,$$
(20)
and to use it in the expression (8), determining the non conservativity parameter w.r.t. the angular momentum.
Note that numerical simulations were conducted for Eulerian equations for in-viscid gas and not for Navier-Stokes equations. Respectively, to calculate integral (20) we use isotropic term describing gas dynamical pressure in the expression for the specific flux of angular momentum: $`𝒫_{\alpha \beta }=P\delta _{\alpha \beta }`$.<sup>5</sup><sup>5</sup>5Here, as usual, $`\delta _{\alpha \beta }`$ is Kronecker’s symbol:
$$\delta _{\alpha \beta }=\{\begin{array}{c}1,\alpha =\beta ,\hfill \\ 0,\alpha \beta .\hfill \end{array}$$
This substitution is correct because on the outer border where the integral in (20) is estimated, viscosity does not play a significant role.
Using expression (20) in the gas dynamical calculations we obtained $`\psi 6`$ corresponding to $`\alpha =0.83`$ (see definition of $`\alpha `$ in the footnote on page 4) for the binary system with $`q={}_{}{}^{1}/_{5}^{}`$ , and $`\psi 5`$ that corresponds to $`\alpha =1.25`$ for a binary with $`q=1.`$ The last value is in good agreement with the results of two-dimensional calculations of binary system with equal masses of components where $`\alpha =1.65`$ was obtained.
However, the use of these estimates for $`\psi `$ in evolutionary calculations proves that this rate of an angular moment loss is so great that in the majority of binary systems the mass exchange rapidly becomes a runaway process and the rate of mass loss by the donor increases without any limit. Apparently, the application of formula (20) to the estimation of the angular momentum loss in the evolutionary calculations is not quite correct. since conjunction of gas dynamical and evolutionary models is not completely consistent. The considered gas dynamical model does not take into account the change in the position of stars in time and, hence, the change in the angular momentum of the gas due to the non-central nature of the force field is not compensated by respective change in the orbital angular momentum.
In a more general gas dynamical model, in which the change in the position of stars is taken into account, one may correctly estimate the change of the systemic angular momentum in the form of angular momentum flux through the external boundary for the “donor + accretor + gas” system. However, to take into account such gas dynamical calculations appears rather complicated. Such estimates are possible only for a specific stage in the life of a binary system. Moreover, thus obtained gas dynamical results cannot be used directly in the standard evolutionary models that ignore the presence of circumbinary envelope in the system. Therefore, a simplified model for estimation of the angular momentum loss developed on the basis of gas dynamical calculations will be a promising approach in the nearest future.
## 5 A simplified model for estimation of angular momentum loss rate in semi-detached binaries
Gas dynamical simulations of mass transfer in semi-detached binaries confirm that outflow of matter from the donor occurs through a small vicinity of the inner Lagrangian point $`L_1`$. In this case specific angular momentum of lost matter (in corotating frame) can be estimated as $`\lambda _{L_1}=\mathrm{\Omega }\mathrm{\Delta }^2`$, where $`\mathrm{\Delta }=|x_cx_{L_1}|`$. Respectively, the flow of angular momentum from the donor is equal to
$$F_\lambda ^{(2)}=\lambda _{L_1}\underset{2}{\overset{}{M}}=\mathrm{\Omega }\mathrm{\Delta }^2\underset{2}{\overset{}{M}}.$$
Due to non-conservativity of mass exchange, only a fraction $`\beta `$ of the matter outflowing through $`L_1`$ is accreted. Specific angular momentum of accreted matter is equal to the angular momentum of accretor (here we neglect the finite radius of accretor and/or residual angular momentum of accreted matter, and this approximation appears to be reasonable for CVs). It means that accreted matter has zero angular momentum in the rotating frame, and, hence, does not accelerate the rotation of accretor. Thus the question on the efficiency of the angular momentum transfer from the spin rotation of accretor to the orbital rotation of the binary can be ignored. To summarize, we can describe the flow of accreted angular momentum as
$$F_\lambda ^{(1)}=\lambda _{accr}\beta \underset{2}{\overset{}{M}}=\mathrm{\Omega }\left(\frac{M_2}{M}\right)^2A^2\beta \underset{2}{\overset{}{M}}.$$
These expressions for the angular momentum flows do not contradict the general formula for specific flow of angular momentum \[see (19)\]. Indeed, viscosity does not play a significant role neither in the vicinity of the surface of the donor nor in the vicinity of the accretor where the matter has already lost its angular momentum and the flow has radial direction in the rotating frame. It means, that in these cases the stress tensor $`𝒫`$ is reduced to the isotropic pressure as well, and does not deposit into the integral of specific flow of angular momentum. Our assumptions permit to write the expression for the flow of angular momentum lost from the system in the following form:
$$F_\lambda ^{loss}=\eta (F_\lambda ^{(2)}F_\lambda ^{(1)})=\eta (\mathrm{\Omega }\mathrm{\Delta }^2\underset{2}{\overset{}{M}}\mathrm{\Omega }\left(\frac{M_2}{M}\right)^2A^2\beta \underset{2}{\overset{}{M}}),$$
where $`\eta `$ is a parameter defining the fraction of the circumbinary envelope angular momentum that is lost with the matter leaving the system. Then $`1\eta `$ is the fraction of the circumbinary envelope angular momentum that returns back into the system via tidal interaction. Henceforth we adopt for the parameter $`\eta `$ a value of 1 that corresponds to the case when the angular momentum of the circumbinary envelope is totally lost with the matter leaving the system. In this case the specific angular momentum of the matter leaving the system is equal to
$$\lambda _{loss}=F_\lambda ^{loss}/\underset{loss}{\overset{}{M}}=\frac{\mathrm{\Omega }\mathrm{\Delta }^2\underset{2}{\overset{}{M}}\beta \left(\frac{M_2}{M}\right)^2\mathrm{\Omega }A^2\underset{2}{\overset{}{M}}}{(1\beta )\underset{2}{\overset{}{M}}}$$
or in the units of systemic specific angular momentum ($`\lambda _{syst}=\mathrm{\Omega }A^2M_1M_2/M^2`$):
$$\psi =\lambda _{loss}/\lambda _{syst}=\frac{\left(f(q)\frac{1}{1+q}\right)^2\beta \left(\frac{q}{1+q}\right)^2}{1\beta }\frac{(1+q)^2}{q},$$
(21)
where $`f(q)=x_{L_1}/A`$ is dimensionless distance from the center of the donor to $`L_1`$.
We would like to stress that formula (21) was derived under the assumption of aligned synchronous rotation of components of the binary system. It is known that the loss of angular momentum due to outflow of the matter from the vicinity of $`L_1`$ can result in the development of a misalignment of the vector of spin rotation of the donor and the vector of the orbital rotation of the binary . The flow structure in the binaries with misaligned asynchronous rotation was studied by the authors in . Having in mind the simplified character of evolutionary models we don’t consider the effect of the misalignment in the present work.
The graph of $`\psi (q)`$ is shown in Fig. 2 for various values of non-conservativity degree w.r.t. the matter $`\beta `$. The values $`\psi =0`$ and $`\psi =1`$ are also shown in the Figure by dashed lines. In the region $`\psi >1`$ each unit of the matter leaving the system carries away an angular moment that is higher than the mean specific angular momentum of the system. This results in diminishing of systemic specific angular momentum:
$$\frac{d}{dt}\left(\frac{J}{M}\right)<0.$$
In the region $`0<\psi <1`$ each unit of the matter leaving the system carries away an angular moment that is lower than the mean specific angular momentum of the system, and, hence,
$$\frac{d}{dt}\left(\frac{J}{M}\right)>0.$$
Figure 2 shows that the value of $`\beta `$ becomes negative in a certain range of $`q`$ for the majority of cases of non-conservative mass exchange (except $`\beta 0`$). This means that the escape of the matter to infinity causes redistribution of the angular momentum in the “donor + accretor + gas” system. This results in the growth of the total angular momentum of the binary: $`\stackrel{}{J}>0`$. It is evident that such a “pumping” of angular momentum into the system can stabilize the mass exchange in CV. Indeed, consideration of orbit semiaxis change due to conservative mass exchange shows that the binary system widens for $`q<1`$ only. For the “non-conservative” model the range of $`q`$ corresponding to the widening of binary due to mass and angular momentum losses is larger – for $`\beta =0.5`$ it happens for $`0<q<2.8`$.
Note that $`\psi (q,\beta )`$ dependence has a curios feature \[see Fig. 2 and formula (21)\] – all curves $`\psi (q)`$ for different values of $`\beta `$ pass through the same point. This point corresponds to the case when the center of mass of the binary system is located exactly equidistantly between $`L_1`$ and accretor, and in this case $`\psi `$ does not depend on $`\beta `$. Fortunately, this specific feature does not influence the analysis of solution.
Three-dimensional gas dynamical simulation of mass transfer in CVs confirms conclusion that the outflow of matter from the system can lead both to decrease and to increase of the total angular momentum of a binary. The typical flow structures \[see Fig. 3a,b\] for binaries with $`M_2/M_1=1:5`$ and $`M_2/M_1=5:1`$ demonstrate that in an inertial frame gas can have different directions of rotation (in both figures orbital rotation corresponds to the counter-clockwise direction). Respectively, the change of angular momentum has opposite signs in the obtained solutions – for $`q={}_{}{}^{1}/_{5}^{}`$ the systemic angular momentum decreased, and for $`q=5`$ it increased.
It seems reasonable to explain the different types of gas dynamical solutions using the following qualitative speculation:
* The flow of the gas leaving the system consists of the matter of the circumbinary envelope that has an angular momentum large enough to overcome the gravitational attraction of the accretor. Initial velocity of this matter (in rotational frame) has the same direction as the orbital movement of the accretor.
* Velocity of the gas leaving the system changes under the action of the gravitational attraction of both components of the binary system, centrifugal force, Coriolis force, and pressure gradient. The change of azimuthal velocity is determined mainly by the Coriolis force and gravitation forces, since centrifugal force is axially-symmetrical and deviation of pressure gradient from the axial symmetry is small as well. Action of Coriolis force results in deflecting of the stream in the direction opposite to that of the orbital movement. Consequently, this permits to subdivide the whole flow leaving the system through the vicinity of $`L_2`$ into two parts: the first part retains the direction of original movement, and the direction of movement of the second one is changed.
* It is evident that in an axially-symmetrical gravitational field the action of Coriolis force can not produce a solution when the direction of gas movement in the inertial frame is opposite to the orbital one. The situation is changed drastically for the case of non-axially-symmetrical gravitational field of a binary system. In this case gravitational attraction of the donor can accelerate the gas of the stream and result in the solution with opposite direction of movement in the inertial frame.
* It follows from Eq. (21) and Fig. 3 that there is a range of values of $`q`$ where gas moves in the direction opposite to the orbital one in the inertial frame. These values of $`q`$ correspond to the cases when attraction of the donor is large enough to accelerate the gas properly (left end of the range), but not too large to produce accretion of gas on the surface of the donor (right end of the range). The values of $`q`$ out of this range correspond to the case when the flow directed along the orbital movement prevails.
Appearance of an additional range of values of $`q`$ where the “non-conservative” mass exchange is accompanied by increase of systemic angular momentum results in expansion of the stable mass exchange region. To illustrate this fact let us consider a graph that is often used in the qualitative studies of the stability of mass transfer in CV. This is a $`q\zeta `$ graph with curves corresponding to the analytical expressions for derivatives of effective radius of the Roche lobe w.r.t. the mass of the donor $`\zeta _{RL}`$, and straight lines corresponding to the values of derivatives of the radius of the star w.r.t. to its mass $`\zeta _{}`$, the latter being known from the investigations of the internal structure of star (see, e.g., ). Mass exchange is stable for the range of $`q`$ where curve $`\zeta _{RL}`$ passes below the straight line $`\zeta _{}`$ for the donor of corresponding type. This graph is presented in Fig. 4. It shows the derivative of effective radius of Roche lobe $`\zeta _{RL}=\mathrm{ln}R_{RL}/\mathrm{ln}M_2`$ as a function of $`q`$ for a scenario in which the loss of angular momentum is defined by formula (21) for various values of the parameter $`\beta `$. Dependencies $`\zeta _{RL}(q)`$ for the scenario of the completely conservative mass exchange and for the case when the outflowing matter carries away specific angular momentum of the donor or specific angular momentum of accretor are shown as well. The same Fig. 4 also shows the values of derivative $`\zeta _{}=\mathrm{ln}R_2/\mathrm{ln}M_2`$ for completely convective and degenerate stars ($`\zeta _{}={}_{}{}^{1}/_{3}^{})`$, for thermally equilibrium stars of MS with $`MM_{}`$ ($`\zeta _{}=0.6)`$, and for subgiants with degenerate low-mass helium cores ($`\zeta _{}=0)`$. The analysis of Fig. 4 proves that the outflow of matter and loss of angular momentum from the system in accordance to law (21) stabilize mass exchange in a wider range of $`q`$ compared to the conventional models of conservative mass exchange. It is also worth noting that for conservative scenario of mass exchange its time scale for stars with masses $`M_{}`$ is determined by the time scale of the loss of angular momentum or the nuclear time scale for $`q\stackrel{<}{}1.2`$ , while for above-considered “non-conservative” model this limit on $`q`$ is sufficiently higher.
We would like to stress also that in the case under consideration there exists a region of unstable mass exchange where value of $`q`$ is very small. The reason of instability \[i.e. the violation of criteria (1)\] is fast increase of the ratio of the lost angular momentum to the mean specific momentum of binary system \[see Eq. (14) and (21)\]. This circumstance can lead to the destruction of a donor of very low mass.<sup>6</sup><sup>6</sup>6Another scenario of development of mass exchange instability for low $`q`$ is also known. It may be caused by low efficiency of interaction between accretion disk and orbital movement (see, e.g., ). Thus, CVs that begin their evolution in a stable mode finish it by a catastrophic disruption of the donor when $`q`$ becomes lower than a certain threshold.
## 6 Evolution of CVs with losses of mass and angular momentum
To take into account results of 3D gas dynamical simulations of the mass transfer in CVs and to incorporate proper estimates of the loss of angular momentum we investigated the evolution of CVs on the basis of various assumptions on its conservativity.
Evolution of the donor was calculated using a modified version of the code for numerical modeling of low-mass stars which was used in the previous works of the authors (see, e.g., ). The code uses opacity tables of Huebner et al. with addition of data of Alexander et al. for low-temperature region. Equation of state was adopted from Fontaine et al. with modifications suggested by Denisenkov . Rates of nuclear reaction were adopted from .
All runs were carried out under assumption that the donor fills its Roche lobe immediately after arriving on ZAMS. Masses of the donors were adopted in the range from $`0.1M_{}`$ to $`1.2M_{}`$ and for various values of $`q`$. Chemical composition was adopted as follows: $`X=0.70`$, $`Y=0.28`$, $`Z=0.02`$. For the mixing length parameter a value of $`l/H_p=1.8`$ was adopted. We used the method of calculation of the donor mass loss rate for known stellar radius and effective radius of Roche lobe suggested by Kolb & Ritter . Based on the results of gas dynamical calculations, we assumed that the fraction of matter accreted by white dwarf is $`\beta =0.5`$, and that angular momentum carried away by the matter leaving the system is defined by expression (21). We didn’t include into the model mass and angular momentum loss due to outbursts.
Let us consider evolution of CV with parameters corresponding to the stable mass exchange in conservative model. In Fig. 5 (“orbital period – logarithm of donor-star mass loss rate” diagram) two evolutionary tracks are shown. One of the tracks is calculated in the conservative approximation (track 1), and another one – in the “non-conservative” approximation (track 2). Initial masses of both donor and accretor were equal to $`1.0M_{}`$. It is seen that these two tracks are almost identical. During the early phase evolution is driven by angular momentum loss due to the MSW of the donor. This results in decrease of the orbital semiaxis and orbital period in the course of evolution. Decrease of the mass of the donor as well as its radius lead to decrease of $`(A/t)_{\text{MSW}}`$ \[see Eq. (13)\], and consequently to decrease of mass loss rate by the donor $`\underset{2}{\overset{}{M}}`$. For the “non-conservative” track, the value of $`\underset{2}{\overset{}{M}}`$ is remarkably lower than for the “conservative” one, because in this phase the value of $`\psi `$ is negative (see Fig. 2) and net loss of orbital angular momentum in the “non-conservative” model is lower than in the conservative one. The situation is inverted in later phases of evolution as $`q`$ reaches the values corresponding to large angular momentum loss.
After the star becomes completely convective, MSW from the donor stops. Mass of the star at this instant is determined by the deviation from the thermal equilibrium. Single MS stars become completely convective at $`M0.36M_{}`$, while mass-losing components of binaries do at lower masses – $`M0.25M_{}÷0.30M_{}`$ (under adopted assumption on chemical composition and opacity of stellar matter).
In the “conservative” track mixing occurs at $`M_2=0.265M_{}`$ (when the orbital period is equal to $`P_{orb}=3^\text{h}.27`$), and in the “non-conservative” track – at $`M_2=0.249M_{}`$ ($`P_{orb}=3^\text{h}.54`$).<sup>7</sup><sup>7</sup>7Since we were interested in peculiarities of CV evolution that are related to its possible non conservativity, we formally continued our calculations beyond accretor mass of $`M_{Ch}=1.4M_{}`$ while in reality in this case the evolution of binary has to be terminated by a thermonuclear explosion of accretor. The continuation of calculation for $`M_1>M_{Ch}1.4M_{}`$ can be advocated by the spin-up of accretor due to accretion because in this case the critical mass can remarkable exceed $`M_{Ch}`$. Thus, values of donor mass and orbital period differ only a little for the “conservative” track and for the “non-conservative” one.
The rate of angular momentum loss decreases abruptly and contraction of binary systems becomes slower after the donor MSW is stopped. As a result donor contracts ceases to lose mass. Since at the moment of cessation of mass loss the radius of the donor is larger than the radius of thermally equilibrium MS star, the star contracts to the equilibrium radius. As a result the ratio of stellar radius to the effective radius of Roche lobe decreases even more. In this phase the evolution of binary system is determined by GWR only. During this detached phase the binary can not be identified as CV, and certain range of the orbital periods becomes devoid of stars. This explains the origin of so called period gap of CVs.
Observed period gap of CVs lies between $`2^\text{h}.1`$ and $`3^\text{h}.1`$ . The boundaries of the “theoretical” period gap depend on adopted rate of angular momentum loss due to MSW and on input physics parameters for evolutionary code that determine theoretical radii of stars. We used a code that gives somewhat overestimated value for the lower boundary. Gap boundaries for the “conservative” track are equal to $`2^\text{h}.5÷3^\text{h}.3`$, and for the “non-conservative” one they are equal to $`2^\text{h}.4÷3^\text{h}.5`$. This discrepancy is related to the different rate of angular momentum loss just prior to gap, and with different masses of accretors after termination of semi-detached phase of evolution as well.
Note that the period gap is characterized not by total absence of CVs but rather by deficiency of CVs. The presence of CVs in the period gap can be explained as follows: i) for CVs with masses of donors equal to $`0.25M_{}÷0.4M_{}`$ mass exchange begins immediately in gap; ii) in the case when the donor fills its Roche lobe not on ZAMS but later, when hydrogen is almost burn out, then as mass decreases down to $`0.3M_{}`$ complete mixing does not occur due to presence a helium core in the star (see, e.g., ). Though, it must be admitted that the number of such systems is very small.
The phase of evolution after period gap is characterized by lower rate of mass loss by the donor, because evolution now is driven by GWR. In this phase the rate of angular momentum loss from the system is sufficiently lower compared to the donor MSW driven phase of evolution. Another feature of this phase of evolution is existence of a minimum period of binary. Short-period cutoff is related to the onset of significant degeneracy of the material of the donor and consequently to the change of mass-radius dependence law : for small masses of donors $`M_20.05M_{}`$ their radii increase as mass decreases and the exponent of the mass-radius law tends to $`{}_{}{}^{1}/_{3}^{}`$. Minimal period corresponds to the value of exponent equal to $`+{}_{}{}^{1}/_{3}^{}`$.<sup>8</sup><sup>8</sup>8This follows from approximation (2) for the effective radius of Roche lobe $`R_{RL}A(M_2/M)^{1/3}`$, condition $`R_2R_{RL}`$, third Kepler’s law $`P_{orb}A^{3/2}M^{1/2}`$, and definition of $`\zeta _{}`$ (1). As a result we have:
$$\stackrel{}{P}/\stackrel{}{P}={}_{}{}^{3}/_{2}^{}\stackrel{}{A}/\stackrel{}{A}{}_{}{}^{1}/_{2}^{}\stackrel{}{M}/\stackrel{}{M},$$
$$\underset{2}{\overset{}{R}}/\underset{2}{\overset{}{R}}={}_{}{}^{1}/_{3}^{}\underset{2}{\overset{}{M}}/\underset{2}{\overset{}{M}}{}_{}{}^{1}/_{3}^{}\stackrel{}{M}/\stackrel{}{M}+\stackrel{}{A}/\stackrel{}{A},$$
$$\underset{2}{\overset{}{R}}/\underset{2}{\overset{}{R}}=\zeta _{}\underset{2}{\overset{}{M}}/\underset{2}{\overset{}{M}},$$
and, finally
$$\stackrel{}{P}/\stackrel{}{P}(\zeta _{}{}_{}{}^{1}/_{3}^{})\underset{2}{\overset{}{M}}/\underset{2}{\overset{}{M}}.$$
This conclusion is independent of the assumption of conservativity or non-conservativity of mass exchange. For “conservative” track the minimum period is equal to 75 minutes and for “non-conservative” one it is equal to 81 minutes. The difference in the values of $`P_{min}`$ is related to the difference in the total mass of the system and to the difference in the deviations of the radii of donor stars from thermal equilibrium values. The latter deviations are caused by different rates of mass loss. After passing the minimum the orbital period begins to increase, and in this phase the rate of mass loss by the donor quickly decreases. It is seen from Fig. 5 that difference between “conservative” and “non-conservative” tracks becomes larger in this stage. This can be naturally explained by increase of the rate of angular momentum loss with decrease of $`q`$ in the “non-conservative” case.
“Conservative” and “non-conservative” tracks have appreciably different values of mass transfer rates after passing the minimum of orbital period. This fact is of importance since the probability of discovery of a CV in a given range of $`\mathrm{log}P`$ is proportional to
$$p(\mathrm{log}P)\frac{(\underset{2}{\overset{}{M}})^\gamma }{\stackrel{}{P}/\stackrel{}{P}},$$
(22)
where $`\gamma `$ is a positive constant $`1`$ for visual magnitude limited sample .
After passing the minimum of the orbital period (i.e. in period range $`2^\text{h}÷2^\text{h}.5`$) the mean value of the probability of discovery of CV for the “non-conservative” track is 1.2 times larger than for the “conservative” one. Thus, despite a substantial difference in the rate of mass exchange, the probability of discovery of CV differs only a little in two cases under consideration. This is due to faster crossing of the corresponding period range in the “non-conservative” case. Nevertheless, predicted difference in $`\underset{2}{\overset{}{M}}`$ for “conservative” and “non-conservative” cases could be reflected in the distribution of CVs over types of variability since the character of variability of a CV depends on the rate of mass exchange.
Figure 5 also shows the evolutionary tracks for stars that would be dynamically unstable within the conservative model. Initial masses of the components were equal to $`1.0M_{}`$ and $`0.28M_{}`$. Of course, this combination of masses is rather extreme but we adopted these parameters for the sake of illustration.
Track 3 was calculated under the conventional assumptions on the loss of angular momentum due to GWR and MSW only. This is a typical case of unstable mass exchange with unlimited growth of mass exchange rate and fast convergence of the time scale of mass loss to the dynamical one, i.e. runaway. Strictly speaking this run should be considered as purely illustrative because in this case a common envelope has to form that increases the rate of angular momentum loss even more (it was not taken into account in calculations).
Track 4 was calculated within the model in which the binary loses matter and angular momentum in accordance with Eq. (21). It resembles the tracks for conventional case of stable evolution of CV and differs slightly from track 2 by a higher rate of mass loss in the initial phase of evolution. The latter is caused by some increase of $`(A/t)_{\text{MSW}}`$ at large $`q`$ \[see Eq. (13)\]. It results in more strong deviations of the donor from thermal equilibrium and corresponding change of the upper border of period gap: it is equal to $`4^\text{h}.4`$ for track 4, and $`3^\text{h}.5`$ for track 2. The minimum of period for track 4 is equal to 76 minutes, i.e. it is nearly the same as for track 1 but slightly less than for track 2.
## 7 Boundary of stable mass exchange region
Boundaries of stable mass exchange region were determined as follows: several sets of evolutionary tracks were calculated for systems with predefined masses of donors and values of $`q`$ that were varied with step 0.1. The value of $`q`$ was considered as critical if for a given evolutionary track the mass accretion rate didn’t exceed Eddington’s limit and if for the next evolutionary track (with $`q`$ higher by 0.1) it did. Note that consequent increasing of $`q`$ leads to unlimited growth of mass loss rate (at values of $`q`$ greater by 0.2 – 0.3 than critical). Note also, that results of such procedure are in good agreement with technique based on comparison of derivatives of stellar radius and effective Roche lobe radius.
The new boundary of the stable mass exchange region is shown in Fig. 6 (under assumption of “non-conservative” evolution). It is well seen that the new region of stable mass exchange contains all observed CVs with estimated masses of components. Thus, within the “non-conservative” model the problem of location of observed systems outside the region of stable mass exchange does not arise.
On the other hand, in the $`qM_2`$ diagram all observed CVs are concentrated near $`q\stackrel{<}{}1`$, while “non-conservative” model supposes sufficiently higher values of $`q`$. To investigate this phenomenon, we have calculated a set of “non-conservative” tracks with initial donor masses from $`M_2^{init}=0.1M_{}`$ up to $`M_2^{init}=1.2M_{}`$ with step $`0.1M_{}`$ and with such values of $`q`$ that tracks start in the vicinity of the new boundary of stable mass exchange region. These tracks are shown in Fig. 7 ($`qM_2`$ diagram). Based on calculations of these tracks, it is possible to get isochrones that are also shown in Fig. 7 for the moments of time equal to $`10^5`$, $`10^6`$, $`10^7`$, $`10^8`$, $`10^9`$, $`10^{10}`$ years (time is measured from the onset of mass exchange between components).<sup>9</sup><sup>9</sup>9Note, that the isochrones for $`10^9`$ years and $`10^8`$ years for tracks with initial donor mass $`M_2^{init}0.4M_{}`$ merge, since this interval of time binaries spend in the period gap. It is this “extreme” set of tracks that was chosen since these systems spent the majority of time between an old boundary of the stable mass exchange region and the new one. It is clear, that systems with smaller values of $`q`$ (tracks of which begin to the left from new boundary) spent there less time. Thus, CVs with parameters that are close to the “extreme” ones have most chances to be found out in this region.
With respect to the location of isochrones in Fig. 7 all tracks can be subdivided into two distinct groups: i) tracks for systems with massive donors (from $`0.4M_{}`$ up to $`1.2M_{}`$), for which the evolution prior to the period gap is defined mainly by the loss of systemic angular momentum due to the donor MSW; ii) tracks for systems of low-mass donors (from $`0.1M_{}`$ up to $`0.3M_{}`$), for which this stage is absent, because their donors are initially completely convective. For the tracks of the first group a high ($`10^6÷10^9`$$`M_{}\text{yr}^1`$) rate of mass loss by the donor prior the period gap is typical, and, hence, donor mass and mass ratio decrease rapidly. This results in a fast evolution in the $`qM_2`$ diagram. As a result, the part of the diagram near the new boundary is crossed very quickly. This explains the absence of observed CVs near new boundary. For the tracks of the second group (dropping a short initial phase with large rate of mass exchange) the evolution of binary system is characterized by low and nearly constant mass loss rate ($`10^{10}`$$`M_{}\text{yr}^1`$) for a long time until period minimum. Consequently, the mass of the donor and mass ratio of components change slowly. Hence, if beginning the evolution near the new boundary of stable mass exchange region, the system should remain there long enough. The absence of CVs in this region can be explained as follows: masses of dwarfs in zero-age CVs can not be higher than $`0.15M_{}`$ that is minimum mass of helium core for a star leaving MS.
## 8 Conclusions
Results of three-dimensional gas dynamical simulations of the flow structure in semi-detached binaries show that during mass exchange a significant fraction of the matter leaves the system. On the other hand, a number of observed CVs have a combination of donor mass and mass ratio of components forbidden in the evolutionary models based on conventional assumptions on the loss of angular momentum of binaries, since in these cases mass exchange should be unstable. This contradiction can be resolved in the framework of the model suggested here. This model takes into account losses of mass and angular momentum from the system according to the results of gas dynamical simulations. It is shown that approximation for outflowing angular momentum as the difference of specific angular momentum of the donor matter in $`L_1`$ and specific angular momentum of accreted matter (being multiplied by appropriate mass flows) permits to explain observations satisfactorily. Note, that the new “non-conservative” model practically does not change such parameters of evolutionary tracks as limits of period gap and minimal period.
Thus, the main conclusion of our study is: if we use a “non-conservative” approximation for CV evolution, we obtain a new boundary of the region of stable mass exchange in $`qM_2`$ plane which explains the distribution of observed CVs better than conventional models.
## ACKNOWLEDGMENTS
This work was supported by the Russian Foundation for Basic Research grants 99-02-17619 and 99-02-16037 and Russian Federation Presidential grant 99-15-96022.
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warning/0003/math0003152.html | ar5iv | text | # Theorem 1
Perturbation of $`l^1`$-copies and measure convergence in preduals of von Neumann algebras
H. Pfitzner
## Abstract
Let $`\mathrm{L}^1`$ be the predual of a von Neumann algebra with a finite faithful normal trace. We show that a bounded sequence in $`\mathrm{L}^1`$ converges to $`0`$ in measure if and only if each of its subsequences admits another subsequence which converges to $`0`$ in norm or spans $`l^1`$ ”almost isometrically”. Furthermore we give a quantitative version of an essentially known result concerning the perturbation of a sequence spanning $`l^1`$ isomorphically in the dual of a C-algebra.
§1 Introduction, main results
The present article deals with convergence in probability in $`\mathrm{L}^1`$-spaces from a functional analytic point of view. The $`\mathrm{L}^1`$-spaces in question are the preduals of von Neumann algebras with finite faithful normal traces. To consider an easy example we look at the commutative case: Let $`(\mathrm{\Omega },\mathrm{\Sigma },\mu )`$ be a finite measure space, let $`(f_n)`$ be a bounded sequence in $`\mathrm{L}^1(\mathrm{\Omega },\mathrm{\Sigma },\mu )`$. If (appropriately chosen representatives of) the $`f_n`$ have pairwise disjoint supports then clearly $`(f_n)`$ converges to $`0`$ in measure. From the functional analytic point of view such a sequence, up to normalization, is the canonical basis of an isometric copy of $`l^1`$. If one perturbes $`(f_n)`$ by a norm null sequence $`(g_n)`$ then $`(f_n+g_n)`$ still $`\mu `$-converges to $`0`$ and spans $`l^1`$ almost isometrically (in a sense to be made precise below in §2). It has been known \[10, Th. 2\] (see also \[19, Th. 3, Rem. 6bis\]) for quite a time that, roughly speaking, these are essentially the only examples of $`\mu `$-null sequences.
Theorem 1 contains the analogous statement for the predual of a von Neumann algebra with finite faithful normal trace. (For notation and definitions see §2.)
###### Theorem 1
Let $`(x_n)`$ be a bounded sequence in $`\mathrm{L}^1(𝒩,\tau )=𝒩_{}`$ where $`(𝒩,\tau )`$ is a von Neumann algebra with a finite normal faithful trace $`\tau `$. Then the following assertions are equivalent.
(i) $`x_n\stackrel{\tau }{}0`$.
(ii) For each subsequence $`(x_{n_k})`$ of $`(x_n)`$ there are a subsequence $`(x_{n_{k_l}})`$ and a sequence $`(y_l)`$ of pairwise orthogonal elements of $`\mathrm{L}^1(𝒩,\tau )`$ such that $`x_{n_{k_l}}y_l_10`$.
(iii) For each subsequence $`(x_{n_k})`$ of $`(x_n)`$ there is a subsequence $`(x_{n_{k_l}})`$ which tends to $`0`$ in $`_1`$ or spans $`l^1`$ almost isometrically.
(iv) For each subsequence $`(x_{n_k})`$ of $`(x_n)`$ there is a subsequence $`(x_{n_{k_l}})`$ which tends to $`0`$ in $`_1`$ or spans $`l^1`$ asymptotically.
The implications (i) $``$ (ii) $``$ (iii) $``$ (iv) hold also for unbounded sequences $`(x_n)`$, the implications (iii) $``$ (ii), (i) do not.
Implication (i)$``$(ii) has already appeard as a special case of a result of Sukochev \[20, Prop. 2.2\]. The other nontrivial implication (iii)$``$(ii) follows immediately from Theorem 2 which holds for the predual of any von Neumann algebra and is of independent interest:
###### Theorem 2
Let $`𝒩`$ be an arbitrary von Neumann algebra and $`(\varphi _m)`$ a bounded sequence in its predual $`𝒩_{}`$. If $`(\varphi _m)`$ spans $`l^1`$ almost isometrically then there are a subsequence $`(\varphi _{m_l})`$ of $`(\varphi _m)`$ and a sequence $`(\stackrel{~}{\varphi }_l)`$ of pairwise orthogonal functionals in $`𝒩_{}`$ such that $`\varphi _{m_l}\stackrel{~}{\varphi }_l0`$ as $`l\mathrm{}`$.
This amounts to saying that there are pairwise orthogonal projections $`s_l`$ and pairwise orthogonal projections $`t_l`$ in $`𝒩`$ such that $`\varphi _{m_l}t_l\varphi _{m_l}s_l0`$ as $`l\mathrm{}`$.
It is natural to ask what can be improved in Theorem 2 if one replaces the predual of the von Neumann algebra by the dual of a C-algebra. At the time of this writing this is not clear. What we have among other things is
###### Proposition 3
Let $`(\varphi _m)`$ be a bounded sequence that spans $`l^1`$ almost isometrically in the dual of an arbitrary C-algebra $`A`$. Then, given $`\epsilon >0`$, there are pairwise orthogonal positive normalized elements $`(a_n)`$ and pairwise orthogonal positive normalized elements $`(b_n)`$ in $`A`$ such that $`\varphi _{m_n}b_n\varphi _{m_n}a_n<\epsilon `$ for an appropriate subsequence $`\varphi _{m_n}`$ and all $`n\mathrm{I}N`$ .
For a more detailed discussion see §6.
As to the organization of the paper, after recalling some notation and definitions in the next section we gather some auxiliary results in §3 in order to prove Theorem 2 in §4. In §5 we prove Theorem 1 for the sake of completeness although, as already mentionned, it follows essentially from \[20, Prop. 2.2\] and Theorem 2. In §6 perturbations of $`l^1`$-copies in the dual of C-algebras are considered.
§2 Notation, definitions
Let $`(x_n)`$ be a sequence of nonzero elements in a Banach space $`X`$.
We say that $`(x_n)`$ spans $`l^1`$ $`r`$-isomorphically or just isomorphically if there exists $`r>0`$ (trivially $`r1`$) such that $`r(_{n=1}^{\mathrm{}}|\alpha _n|)_{n=1}^{\mathrm{}}\alpha _n\frac{x_n}{x_n}_{n=1}^{\mathrm{}}|\alpha _n|`$ for all scalars $`\alpha _n`$ (the second inequality being trivial).
We say that $`(x_n)`$ spans $`l^1`$ almost isometrically if there is a sequence $`(\delta _m)`$ in $`[0,1[`$ tending to $`0`$ such that $`(1\delta _m)_{n=m}^{\mathrm{}}|\alpha _n|_{n=m}^{\mathrm{}}\alpha _n\frac{x_n}{x_n}_{n=m}^{\mathrm{}}|\alpha _n|`$ for all $`m\mathrm{I}N`$.
Trivially the property of spanning $`l^1`$ almost isometrically passes to subsequences. Recall that James’ distortion theorem (see or ) for $`l^1`$ says that every isomorphic copy of $`l^1`$ contains an almost isometric copy of $`l^1`$. To be more precise, let $`r>0`$, $`[0,1[\delta _n0`$, and let $`(x_n)`$ be a normalized basis spanning $`l^1`$ $`r`$-isomorphically. Then it follows from the proof of that there is a sequence $`(\lambda _i)`$ of scalars and a sequence $`(F_n)`$ of pairwise disjoint finite subsets of $`\mathrm{I}N`$ such that $`(1\delta _m)_{n=m}^{\mathrm{}}|\alpha _n|_{n=m}^{\mathrm{}}\alpha _ny_n_{n=m}^{\mathrm{}}|\alpha _n|`$ for all scalars $`\alpha _n`$ and all $`m\mathrm{I}N`$ where $`y_n=_{iF_n}\lambda _ix_i`$ and where $`_{iF_n}|\lambda _i|\frac{1}{r}`$ for all $`n\mathrm{I}N`$.
Finally $`(x_n)`$ is said to span $`l^1`$ asymptotically isometrically or just to span $`l^1`$ asymptotically if there is a sequence $`(\delta _n)`$ in $`[0,1[`$ tending to $`0`$ such that $`_{n=1}^{\mathrm{}}(1\delta _n)|\alpha _n|_{n=1}^{\mathrm{}}\alpha _n\frac{x_n}{x_n}_{n=1}^{\mathrm{}}|\alpha _n|`$ for all scalars $`\alpha _n`$. We say that a Banach space is isomorphic (respectively almost isometric, respectively asymptotically isometric) to $`l^1`$ if it has a basis with the corresponding property. Clearly a sequence spanning $`l^1`$ asymptotically spans $`l^1`$ almost isometrically. The main result of states that the converse does not hold because there are almost isometric copies of $`l^1`$ which do not contain $`l^1`$ asymptotically. However, it follows from that this cannot happen in the predual of a von Neumann algebra because each sequence spanning $`l^1`$ almost isometrically in a von Neumann predual contains a sequence spanning $`l^1`$ asymptotically (cf. (iii) $``$ (iv) in the proof of Theorem 1). Note that the present definitions of almost and asymptotically isometric differ slightly from those in , by the term $`x_n/x_n`$ but that, of course, for normalized sequences the definitions are the same. Note also the technical detail that because of this term one might have $`x_n0`$ for a sequence spanning $`l^1`$ isomorphically (or almost or asymptotically isometrically) whereas sequences that are equivalent to the canonical $`l^1`$-basis (\[4, p. 43\]) are uniformly bounded away from $`0`$.
The dual of a Banach space $`X`$ is denoted by $`X^{}`$. We work with complex scalars. Two elements $`a,b`$ of a C-algebra are called orthogonal \- $`ab`$ in symbols - if $`ab^{}=0=a^{}b`$. Two elements $`\varphi ,\psi `$ of the predual of a von Neumann algebra are called orthogonal - $`\varphi \psi `$ in symbols - if they have orthogonal right and orthogonal left support projections. It is well know that $`\varphi \psi `$ if and only if the linear span of $`\varphi `$ and $`\psi `$ is isometrically isomorphic to the two-dimensional $`l_2^1`$; if $`\varphi `$ and $`\psi `$ are positive they are orthogonal if and only if $`\varphi \psi =\varphi +\psi `$.
Let $`𝒩`$ be a von Neumann algebra, $`a𝒩`$, $`\varphi 𝒩_{}`$ then $`a\varphi `$ denotes the normal functional $`𝒩x\varphi (xa)`$ and $`\varphi a`$ denotes the normal functional $`𝒩x\varphi (ax)`$.
Let $`\tau `$ be a finite faithful normal trace on a von Neumann algebra $`𝒩`$. The set $`I=\{x𝒩|\tau (|x|)<\mathrm{}\}`$ is an ideal in $`𝒩`$, can be normed by $`x\tau (|x|)=:x_1`$ and its Banach space completion is denoted by $`\mathrm{L}^1=\mathrm{L}^1(𝒩,\tau )`$. It is well-known that $`\mathrm{L}^1`$ is isometrically isomorphic to the predual $`𝒩_{}`$ via the map $`\mathrm{L}^1x\varphi _x𝒩_{}`$ where $`\varphi _x(y)=\tau (xy)`$ for $`y𝒩`$ and where $`\tau `$ is understood as the (well-defined) extension of $`\tau `$ from $`I`$ to $`\mathrm{L}^1`$ \[21, V.2.18\]. In particular, the multiplication on $`𝒩\times I`$ can be extended to $`𝒩\times \mathrm{L}^1`$, the map $`x\varphi _x`$ respects orthogonality and one has $`|\tau (xy)|x_1y_{\mathrm{}}`$ for $`x\mathrm{L}^1`$, $`y\mathrm{L}^{\mathrm{}}=\mathrm{L}^{\mathrm{}}(𝒩,\tau ):=𝒩`$. More generally one can define $`\mathrm{L}^p(𝒩,\tau )`$-spaces, $`1p<\mathrm{}`$, as the sets of those $`x\mathrm{L}^0`$ for which $`x_p:=\tau (|x|^p)^{1/p}<\mathrm{}`$ where $`\mathrm{L}^0=\mathrm{L}^0(𝒩,\tau )`$ is the space of $`\tau `$-measurable densely defined (in general unbounded) operators affiliated with $`𝒩`$ and where $`\tau `$ is understood as the extension of $`\tau `$ from $`𝒩`$ to $`\mathrm{L}^0`$. On $`\mathrm{L}^0`$ one defines the measure topology as the translation invariant topology in which the sets $`\{x\mathrm{L}^0|p𝒩_{proj}:xp𝒩,xp_{\mathrm{}}\epsilon ,\tau (p^{})\delta \}`$, $`\epsilon ,\delta >0`$, form a base of the zero neighborhoods. ($`𝒩_{proj}`$ denotes the set of projections of $`𝒩`$.) In this topology, $`\mathrm{L}^0`$ becomes a (well-defined) metrizable complete Hausdorff topological vector -algebra and all $`\mathrm{L}^p`$ embed injectively in $`\mathrm{L}^0`$. In particular, sum and product are well-defined in $`\mathrm{L}^0`$. All this (for the more general case of a semifinite trace) can be found for example in , \[22, Ch. 1\] or .
If a sequence $`(x_n)`$ in $`\mathrm{L}^0`$ converges to $`x\mathrm{L}^0`$ this is denoted by $`x_n\stackrel{\tau }{}x`$. In this context Chebyshev’s inequality reads $`\tau (\chi _{]\epsilon ,\mathrm{}[}(|x|))\tau (\frac{1}{\epsilon }|x|)=\frac{1}{\epsilon }x_1`$ for $`x\mathrm{L}^1`$ \- which means in particular that the norm topology is finer than the measure topology induced by $`\tau `$ \- and from \[8, A48\] we know that in accordance with the commutative case, $`x_n\stackrel{\tau }{}0`$ if and only if $`\tau (\chi _{]\epsilon ,\mathrm{}[}(|x_n|))0`$ as $`n\mathrm{}`$ for all $`\epsilon >0`$.
Basic properties and definitions which are not explained here can be found in or in - for Banach spaces and in , for C-algebras.
§3 Some auxiliary results
Let us first state an easy lemma which says that almost isometric and asymptotically isometric $`l^1`$-copies are stable with respect to perturbations by norm null sequences.
###### Lemma 4
Let $`(x_n)`$, $`(y_n)`$ be two sequences in a Banach space $`X`$ such that $`infx_n>0`$, $`y_n0`$ and $`x_n+y_n0`$.
If $`(x_n)`$ spans $`l^1`$ almost isometrically then so does $`(x_n+y_n)`$.
If $`(x_n)`$ spans $`l^1`$ asymptotically then so does $`(x_n+y_n)`$.
Proof: Suppose that $`(x_n)`$ spans $`l^1`$ almost isometrically. For all scalar sequences $`(\alpha _n)`$ one has
$`{\displaystyle \underset{n=m}{\overset{\mathrm{}}{}}}\alpha _n{\displaystyle \frac{x_n+y_n}{x_n+y_n}}`$
$``$ $`{\displaystyle \underset{n=m}{\overset{\mathrm{}}{}}}\alpha _n{\displaystyle \frac{x_n}{x_n}}{\displaystyle \underset{n=m}{\overset{\mathrm{}}{}}}\alpha _n\left(1{\displaystyle \frac{x_n}{x_n+y_n}}\right){\displaystyle \frac{x_n}{x_n}}{\displaystyle \underset{n=m}{\overset{\mathrm{}}{}}}\alpha _n{\displaystyle \frac{y_n}{x_n+y_n}}`$
$``$ $`((1\delta _m){\displaystyle \underset{n=m}{\overset{\mathrm{}}{}}}|\alpha _n|)(\underset{nm}{sup}|1{\displaystyle \frac{x_n}{x_n+y_n}}|{\displaystyle \underset{n=m}{\overset{\mathrm{}}{}}}|\alpha _n|)(\underset{nm}{sup}{\displaystyle \frac{y_n}{x_n+y_n}}{\displaystyle \underset{n=m}{\overset{\mathrm{}}{}}}|\alpha _n|)`$
$`=`$ $`(1\delta _m^{}){\displaystyle \underset{n=m}{\overset{\mathrm{}}{}}}|\alpha _n|`$
where $`\delta _m^{}=\delta _m+sup_{nm}|1\frac{x_n}{x_n+y_n}|+sup_{nm}\frac{y_n}{x_n+y_n}0`$ as $`m\mathrm{}`$. Hence $`(x_n+y_n)`$ spans $`l^1`$ almost isomorphically. The asymptotic case is proved similarly.
Lemmas 5 \- 7 seem to be known and are proved mainly for lack of suitable reference. (In part they overlap with \[17, Lem. 3-5\].)
###### Lemma 5
Let $`A`$ be a C-algebra, $`\omega `$ a positive functional on $`A`$ and $`a,b`$ elements of the unit ball of $`A`$. Then
$`a\omega \omega `$ $``$ $`(2\omega )^{1/2}|\omega \omega (a)|^{1/2}`$ (1)
$`\omega a\omega `$ $``$ $`(2\omega )^{1/2}|\omega \omega (a)|^{1/2}`$ (2)
$`b\omega a\omega `$ $``$ $`(2\omega )^{1/2}(|\omega \omega (a)|^{1/2}+|\omega \omega (b)|^{1/2})`$ (3)
Proof: Let $`xA`$ and $`x1`$. Set $`\gamma =\omega \omega (a)`$, thus $`\omega (a^{})=\omega \overline{\gamma }`$. Without loss of generality we assume $`\omega =1`$. The inequality of Cauchy-Schwarz yields
$`|\omega (x)a\omega (x)|^2`$ $`=`$ $`|\omega (x(1a))|^2\omega (xx^{})\omega ((1a)^{}(1a))`$
$``$ $`\omega ((1a)^{}(1a))=\omega (1a)\omega (a^{}a^{}a)`$
$`=`$ $`\gamma (1\overline{\gamma })+\omega (a^{}a)2\mathrm{R}\mathrm{e}\gamma 2|1\omega (a)|`$
whence (1); (2) follows analogously; (3) follows from (1), (2) and from $`\omega b\omega a\omega b\omega +b(\omega \omega a)\omega b\omega +\omega \omega a`$
###### Lemma 6
Let $`A`$ be a C-algebra, $`\varphi `$ a functional on $`A`$ and $`a,b`$ in the unit ball of $`A`$. Then
$`\varphi a|\varphi |`$ $``$ $`(2\varphi )^{1/2}|\varphi \varphi ^{}(a)|^{1/2}`$ (4)
$`|\varphi |a\varphi `$ $``$ $`(2\varphi )^{1/2}|\varphi \varphi (a)|^{1/2}`$ (5)
$`b\varphi a\varphi `$ $``$ $`(2\varphi )^{1/2}\left(|\varphi |\varphi |(a)|+|\varphi |\varphi ^{}|(b)|\right)^{1/2}.`$ (6)
Proof: Let $`\varphi =u|\varphi |`$ be the polar decomposition of $`\varphi `$. Then the polar decomposition of $`\varphi ^{}`$ is $`\varphi ^{}=u^{}|\varphi ^{}|`$ (cf. the proof of \[21, III.4.2\]), we have $`\varphi =|\varphi ^{}|u`$, $`|\varphi |=u^{}\varphi =|\varphi |^{}=\varphi ^{}u`$. Without loss of generality we assume $`\varphi =1`$.
Inequality (5) follows from
$`a\varphi |\varphi |=au|\varphi ||\varphi |\stackrel{(\text{1})}{}|2(1|\varphi |(au))|^{1/2}=|2(1\varphi (a))|^{1/2}.`$
Replacing $`\varphi `$ by $`\varphi ^{}`$ we have $`a\varphi ^{}|\varphi ^{}||2(1\varphi ^{}(a))|^{1/2}`$ whence (4) by
$`a|\varphi |\varphi =(a\varphi ^{}|\varphi ^{}|)ua\varphi ^{}|\varphi ^{}|`$. (6) follows from
$`\varphi b\varphi a`$ $``$ $`\varphi b\varphi +b\varphi b\varphi a`$
$`=`$ $`|\varphi ^{}|ub|\varphi ^{}|u+bu|\varphi |bu|\varphi |a`$
$``$ $`|\varphi ^{}|b|\varphi ^{}|+|\varphi ||\varphi |a`$
$`\stackrel{(\text{1})(\text{2})}{}`$ $`(2\varphi )^{1/2}(|\varphi |\varphi ^{}|(b)|+|\varphi |\varphi |(a)|)^{1/2}.`$
###### Lemma 7
Let $`𝒩`$ be a von Neumann algebra with predual $`𝒩_{}`$. If a functional $`\sigma `$ in the unit ball of $`𝒩_{}`$, projections $`r,l𝒩`$ and a number $`\beta ]0,1[`$ are such that $`r(|\sigma |)1\beta `$ and $`l(|\sigma ^{}|)1\beta `$ then $`\sigma \tau <5\sqrt{\beta }`$ where $`\tau =\frac{l\sigma r}{l\sigma r}`$.
Proof: $`l\sigma r\sigma 2\sqrt{\beta }`$ by (6) and $`\frac{l\sigma r}{l\sigma r}l\sigma r=\frac{1l\sigma r}{l\sigma r}l\sigma r\beta +\sigma l\sigma r\beta +\sigma l\sigma r\beta +2\sqrt{\beta }`$ thus $`\sigma \tau <5\sqrt{\beta }`$.
We recall some more definitions and notation. Let $`A`$ be a C-algebra. A projection $`pA^{\prime \prime }`$ is called open if it is the limit of an increasing net of positive elements of $`A`$ (\[16, 3.11\], ). If $`pA^{\prime \prime }`$ is open then $`B^{\prime \prime }=pA^{\prime \prime }p`$ where $`B=pA^{\prime \prime }pA`$ is a hereditary subalgebra. A projection $`qA^{\prime \prime }`$ is called closed if there is an open projection $`pA^{\prime \prime }`$ such that $`q=p^c`$ where $`p^c`$ denotes the complement $`1p`$ of $`p`$. (This makes sense also if $`A`$ is not unital because one always has $`1A^{\prime \prime }`$.) By definition the closure $`\overline{p}`$ of a projection $`pA^{\prime \prime }`$ is the infimum of all closed projections majorizing $`p`$. $`\chi _M`$ denotes the characteristic function of a set $`M`$. By functional calculus $`\chi _{]\epsilon ,1]}(x)`$ (respectively $`\chi _{[\epsilon ,1]}(x)`$) is an open (respectively closed) projection in $`A^{\prime \prime }`$ if $`1>\epsilon >0`$, $`xA`$, $`0x1`$, because $`\chi _{]\epsilon ,1]}`$ (respectively $`\chi _{[\epsilon ,1]}`$) is the pointwise limit of an increasing (respectively decreasing) sequence of continuous functions on $`[0,1]`$ (cf. \[1, II.3\]). As to $`\chi _{]\epsilon ,1]}(x)`$ this is easy but as to $`\chi _{[\epsilon ,1]}(x)`$ a bit more attention must be paid to the case where $`A`$ is not unital; in this case one works with the unitisation $`A_1`$ of $`A`$. Therfore, in order to avoid complications in the non-unital case , we state Lemma 8 only for unital $`A`$ (although the lemma holds also in the non-unital case).
The following Lemma 8 is a natural generalisation of \[17, Lem. 5\].
###### Lemma 8
Let $`A`$ be a unital C-algebra. For each $`\epsilon >0`$ and each $`n\mathrm{I}N`$ there is $`\delta =\delta (n,\epsilon )>0`$ with the following property.
If there are functionals $`\varphi _1,\mathrm{},\varphi _n`$ in the unit ball of $`A^{}`$ and open projections $`s,tA^{\prime \prime }`$ such that
$`(1\delta ){\displaystyle \underset{1}{\overset{n}{}}}|\alpha _k|{\displaystyle \underset{1}{\overset{n}{}}}\alpha _kt\varphi _ks{\displaystyle \underset{1}{\overset{n}{}}}|\alpha _k|(\alpha _k)\mathrm{C}`$ (7)
then there are open projections $`p_1,\mathrm{},p_nsA^{\prime \prime }s`$ with pairwise orthogonal closures in $`sA^{\prime \prime }s`$ and open projections $`q_1,\mathrm{},q_ntA^{\prime \prime }t`$ with pairwise orthogonal closures in $`tA^{\prime \prime }t`$ such that
$`p_k(|\varphi _k|)`$ $`>`$ $`(1\epsilon )\varphi _k`$ (8)
$`q_k(|\varphi _k^{}|)`$ $`>`$ $`(1\epsilon )\varphi _k`$ (9)
for $`k=1,\mathrm{},n`$.
In particular the $`\varphi _k`$ are close to normalized orthogonal elements $`\psi _k`$ on $`sAt`$ in the sense that $`\varphi _k\psi _k<5\sqrt{\epsilon }`$ where $`\psi _k=q_k\varphi _kp_k/q_k\varphi _kp_k`$ are normalized and pairwise orthogonal with left (right) supports majorized by $`t`$ (by $`s`$).
Proof: (a) First we suppose $`s=t=1`$ and deal only with the special case
(a1) of positive functionals $`\varphi _k`$.
Let $`\epsilon >0`$. For $`n=1`$ choose an $`x0`$ in the unit ball of $`A`$ such that $`\varphi _1(x)>1\epsilon `$ and set $`p_1=q_1=\chi _{]0,1]}(x)`$, $`\delta (1,\epsilon )=\epsilon `$.
Suppose now that the assertion holds true (for positive functionals, for $`s=t=1`$ and) for $`n\mathrm{I}N`$. By hypothesis on $`n`$ we choose $`\delta _n=\delta (n,\epsilon )`$ We define $`\delta _{n+1}>0`$ such that $`\delta _{n+1}+(32n\delta _{n+1})^{1/2}<\delta _n`$. Consider positive functionals $`\varphi _k`$, $`k=1,\mathrm{},n+1`$, in the unit ball of $`A`$ such that $`_1^{n+1}\alpha _k\varphi _k(1\delta _{n+1})_1^{n+1}|\alpha _k|`$. Set $`\sigma =\frac{1}{n}_1^n\varphi _k`$ and $`\tau =\varphi _{n+1}`$. Then $`(1\delta _{n+1})(|\alpha |+|\beta |)\alpha \sigma +\beta \tau |\alpha |+|\beta |`$ for all scalars $`\alpha ,\beta `$. In particular $`\sigma \tau 2(1\delta _{n+1})`$. There is a selfadjoint normalized element $`xA`$ such that $`(\sigma \tau )(x)>2(12\delta _{n+1})`$. Decompose $`x=x^++x^{}`$ in its negative and positive parts. Then $`(\sigma \tau )(x)=(\sigma (x^+)+\tau (x^{}))(\sigma (x^{})+\tau (x^+))>2(12\delta _{n+1})`$ whence $`\sigma (x^+)>2(12\delta _{n+1})\tau (x^{})>14\delta _{n+1}`$ and similarly $`\varphi _{n+1}(x^{})>14\delta _{n+1}`$. Together with $`\varphi _k(x^+)1`$ this gives $`\varphi _k(x^+)>14n\delta _{n+1}`$ for all $`k=1,\mathrm{},n`$ because otherwise one would have $`n\sigma (x^+)14n\delta _{n+1}+(n1)=n(14\delta _{n+1})`$. Define $`p=\chi _{]\eta ,1]}(x^+)`$, $`p_{n+1}=\chi _{]\eta ,1]}(x^{})`$ where $`\eta >0`$ is such that
$`p(\varphi _k)`$ $`>`$ $`14n\delta _{n+1}\text{ for }k=1,\mathrm{},n`$ (10)
$`\text{and }p_{n+1}(\varphi _{n+1})`$ $`>`$ $`14\delta _{n+1}>1\epsilon `$
By functional calculus the projections $`p`$ and $`p_{n+1}`$ are open and have orthogonal closures. $`B=pA^{\prime \prime }pAA`$ is a hereditary subalgebra of $`A`$. This explains the equality sign in the following formula:
$`{\displaystyle \underset{1}{\overset{n}{}}}\alpha _k\varphi _k|_B_B`$ $`=`$ $`{\displaystyle \underset{1}{\overset{n}{}}}\alpha _k(p\varphi _kp){\displaystyle \underset{1}{\overset{n}{}}}\alpha _k\varphi _k{\displaystyle \underset{1}{\overset{n}{}}}\alpha _k(\varphi _kp\varphi _kp)`$
$`\stackrel{(\text{3})(\text{10})}{>}`$ $`(1\delta _{n+1}){\displaystyle \underset{1}{\overset{n}{}}}|\alpha _k|\sqrt{32n\delta _{n+1}}{\displaystyle \underset{1}{\overset{n}{}}}|\alpha _k|`$
$`>`$ $`(1\delta _n){\displaystyle \underset{1}{\overset{n}{}}}|\alpha _k|.`$
By induction hypothesis applied to $`B`$ and to $`\varphi _k|_B`$ one gets $`n`$ open projections $`p_1,\mathrm{},p_nB^{\prime \prime }`$ with pairwise orthogonal closures in $`B^{\prime \prime }`$ \- whence in $`A^{\prime \prime }`$ \- such that (8) holds for $`k=1,\mathrm{},n+1`$. Furthermore, (9) holds with $`q_k=p_k`$ because we have supposed $`\varphi _k0`$. This proves the case where $`s=t=1`$ for positive functionals $`\varphi _k`$.
(a2) For the case of arbitrary functionals (but still with $`s=t=1`$) suppose that the lemma is false. Then there are $`\epsilon >0`$, a sequence $`(A_i)`$ of C-algebras, and $`\varphi _{k,i}(A_i)_1`$ such that for each $`i\mathrm{I}N`$,
$`(1{\displaystyle \frac{1}{i}}){\displaystyle \underset{k=1}{\overset{n}{}}}|\alpha _k|<{\displaystyle \underset{k=1}{\overset{n}{}}}\alpha _k\varphi _{k,i}{\displaystyle \underset{k=1}{\overset{n}{}}}|\alpha _k|(\alpha _k)\mathrm{C}`$ (11)
but for each $`i\mathrm{I}N`$ the $`\varphi _{k,i}`$ are far from orthogonal functionals, more precisely
$`\underset{kn}{\mathrm{min}}p_{k,i}(|\varphi _{k,i}|)1\epsilon \text{or}\underset{kn}{\mathrm{min}}q_{k,i}(|\varphi _{k,i}^{}|)1\epsilon `$ (12)
for all sequences $`(p_{k,i})_{k=1}^n`$ and $`(q_{k,i})_{k=1}^n`$ of open projections with orthogonal closures in $`A_i^{\prime \prime }`$.
We recall some basic facts on ultraproducts (see e.g. ). If $`𝒰`$ is an ultrafilter on an index set $`I`$ the ultraproduct $`X=(X_i)/𝒰`$ of a family $`(X_i)_{iI}`$ of Banach spaces is defined as the quotient $`l^{\mathrm{}}(X_i)/c_0(X_i)`$ where $`l^{\mathrm{}}(X_i)=\{(x_i)_{iI}|(x_i);=sup_𝒰x_i<\mathrm{}\}`$ and $`c_0(X_i)=\{(x_i)_{iI}l^{\mathrm{}}(X_i)|lim_𝒰x_i=0\}`$. With the quotient norm $`X`$ becomes a Banach space. By $`[x_i]_𝒰`$ we denote the equivalence class represented by $`(x_i)_{iI}l^{\mathrm{}}(X_i)`$. One has $`[x_i]_𝒰=lim_𝒰x_i`$ independently of the representative of $`[x_i]_𝒰`$. The ultraproduct $`(X_i^{})/𝒰`$ of the duals can be identified isometrically with a closed subspace of the dual $`X^{}`$ via $`[x_i^{}]_𝒰([x_i]_𝒰)=lim_𝒰x_i^{}(x_i)`$. An ultraproduct $`A=(A_i)/𝒰`$ of a family of C-algebras $`A_i`$ is canonically a C-algebra with pointwise multiplication and involution because in this case the null space $`c_0(X_i)`$ is an ideal in $`A`$.
Let now $`I=\mathrm{I}N`$ and $`𝒰`$ be a free ultrafilter on $`\mathrm{I}N`$ and set $`A=(A_i)/𝒰`$. For each element $`\psi A^{}`$ of the form $`\psi =[\psi _i]_𝒰`$ we have $`|\psi |=[|\psi _i|]_𝒰`$ and $`|\psi ^{}|=[|\psi _i^{}|]_𝒰`$. \[To see this choose $`a=[a_i]_𝒰`$ in the unit ball of $`A`$ such that $`a_i=1`$ and $`\psi (a)=lim_𝒰\psi _i(a_i)=lim_𝒰\psi _i=1`$. Then $`|\psi |=a\psi =[a_i\psi ]_𝒰`$ and $`|\psi _i|a_i\psi _i(2|\psi _i\psi _i(a_i)|)^{1/2}0`$ by (5) of Lemma 6 hence $`|\psi |=[|\psi _i|]_𝒰`$. For $`|\psi ^{}|=[|\psi _i^{}|]_𝒰`$ the proof is analogous.\]
The $`n`$ functionals $`|\varphi _k|=[|\varphi _{k,i}|]_𝒰`$ and the $`n`$ functionals $`|\varphi _k^{}|=[|\varphi _{k,i}^{}|]_𝒰`$ are pairwise orthogonal because the $`\varphi _k=[\varphi _{k,i}]_𝒰`$ are so by (11). Since the statement of the Lemma has been proved for positive functionals (and for $`s=t=1`$) there are, for each $`iI`$, two finite sequences $`(p_{k,i})_{k=1}^n`$, $`(q_{k,i})_{k=1}^n`$ of open projections with pairwise orthogonal closures in $`A_i^{\prime \prime }`$ such that $`lim_𝒰p_{k,i}(\varphi _{k,i})=[p_{k,i}]_𝒰(\varphi _k)=1`$ and $`lim_𝒰q_{k,i}(\varphi _{k,i}^{})=[q_{k,i}]_𝒰(\varphi _k^{})=1`$ for $`k=1,\mathrm{},n`$ which contradicts (12) and thus proves the lemma for the case where $`s=t=1`$.
(b) Now we turn to the general case of arbitrary open projections $`s,tA^{\prime \prime }`$. We assume without loss of generality that $`\varphi _k=1`$. After what has been proved in (a), we further assume without loss of generality that the $`\varphi _k`$ are pairwise orthogonal because (7) remains valid if $`t\varphi _ks`$ is replaced by $`\varphi _k`$. Then the $`|\varphi _k|`$ are orthogonal and so are the $`|\varphi _k^{}|`$. Thus we may further assume that
$`{\displaystyle \underset{1}{\overset{n}{}}}\alpha _k|\varphi _k|={\displaystyle \underset{1}{\overset{n}{}}}|\alpha _k|\text{and}{\displaystyle \underset{1}{\overset{n}{}}}\alpha _k|\varphi _k^{}|={\displaystyle \underset{1}{\overset{n}{}}}|\alpha _k|`$ (13)
for all scalars $`\alpha _k\mathrm{C}\text{ }`$. Let $`\varphi _k=u_k|\varphi _k|`$ be the polar decomposition then
$`1\delta \stackrel{(\text{7})}{}t\varphi _ks=tu_k|\varphi _k|s=|tu_k|\varphi _k|s|=s(|(tu_k|\varphi _k|s)|)s(|\varphi _k|)`$
where the last inequality follows from \[21, III.4.9\]. Analogously $`t(|\varphi _k^{}|)1\delta `$. Hence by (6) of Lemma 6, if $`\delta `$ is small enough, $`t\varphi _ks`$ is a small perturbation of $`\varphi _k`$ and since the absolut value is norm continuous on von Neumann preduals \[21, III.4.10\], $`|t\varphi _ks|`$ is a small perturbation of $`|\varphi _k|`$, and $`|s\varphi _k^{}t|`$ is a small perturbation of $`|\varphi _k^{}|`$. Finally, in view of (13) we may without loss of generality replace (7) by
$`(1\delta ){\displaystyle \underset{1}{\overset{n}{}}}|\alpha _k|{\displaystyle \underset{1}{\overset{n}{}}}\alpha _k|t\varphi _ks|{\displaystyle \underset{1}{\overset{n}{}}}|\alpha _k|(\alpha _k)\mathrm{C}`$ (14)
$`(1\delta ){\displaystyle \underset{1}{\overset{n}{}}}|\alpha _k|{\displaystyle \underset{1}{\overset{n}{}}}\alpha _k|s\varphi _k^{}t|{\displaystyle \underset{1}{\overset{n}{}}}|\alpha _k|(\alpha _k)\mathrm{C}`$ (15)
in the statement of the lemma. It remains to apply part (a) to the hereditary subalgebra $`sA^{\prime \prime }sA`$ and to $`tA^{\prime \prime }tA`$ (because the support projection of $`|t\varphi _ks|`$ (of $`|s\varphi _k^{}t|`$) is majorized by $`s`$ (by $`t`$)). This yields the desired open projections $`p_k(sA^{\prime \prime }sA)^{\prime \prime }=sA^{\prime \prime }s`$ and $`q_ktA^{\prime \prime }t`$ satisfying (8) and (9) if $`\delta `$ is small enough.
The last assertion of the lemma is immediate from Lemma 7.
###### Corollary 9
Let $`𝒩`$ be a von Neumann algebra. For each $`\epsilon >0`$ and each $`n\mathrm{I}N`$ there is $`\delta =\delta (n,\epsilon )>0`$ with the following property.
If there are functionals $`\varphi _1,\mathrm{},\varphi _n`$ in the unit ball of $`𝒩_{}`$ and (arbitrary) projections $`s,t𝒩`$ such that
$`(1\delta ){\displaystyle \underset{1}{\overset{n}{}}}|\alpha _k|{\displaystyle \underset{1}{\overset{n}{}}}\alpha _kt\varphi _ks{\displaystyle \underset{1}{\overset{n}{}}}|\alpha _k|(\alpha _k)\mathrm{C}`$ (16)
then there are pairwise orthogonal projections $`p_1,\mathrm{},p_ns𝒩s`$ and pairwise orthogonal projections $`q_1,\mathrm{},q_nt𝒩t`$ such that
$`\varphi _k\psi _k<\epsilon `$ $`\text{for }k=1,\mathrm{},n`$ (17)
where $`\psi _k=q_k\varphi _kp_k/q_k\varphi _kp_k`$.
Proof: For $`s=t=1`$ the assertion is immediate from Lemma 8 and Lemma 7. For arbitrary projections $`s,t𝒩`$ we proceede as in part (b) of the proof of Lemma 8 in order to show that (16) can be replaced by (14) and (15) and to apply this to the subalgebras $`s𝒩s`$ and $`t𝒩t`$.
§4 Proof of Theorem 2
Without loss of generality we assume that $`\varphi _m=1`$ for all $`m\mathrm{I}N`$. Let $`(\eta _n)`$ be a sequence of positive numbers such that $`\eta _n`$ converges.
By induction on $`n=1,2,\mathrm{}`$ we construct an increasing sequence $`(m_n)`$ in $`\mathrm{I}N`$, functionals $`\psi _{m_k}^{(n)}𝒩_{}`$ for $`k=1,\mathrm{},n`$, such that for all $`n\mathrm{I}N`$:
$`|\psi _{m_k}^{(n)}|`$ $``$ $`|\psi _{m_l}^{(n)}|k,l=1,\mathrm{},n,kl,`$ (18)
$`\psi _{m_k}^{(n)}`$ $`=`$ $`1k=1,\mathrm{},n`$ (19)
$`\psi _{m_k}^{(n)}\psi _{m_k}^{(n1)}`$ $`<`$ $`\eta _nk=1,\mathrm{},n1,`$ (20)
$`\psi _{m_n}^{(n)}\varphi _{m_n}`$ $`<`$ $`\eta _n.`$ (21)
For $`n=1`$ one may simply set $`\psi _{m_1}^{(1)}=\varphi _1`$; (18, $`n=1`$) and (20, $`n=1`$) are void, (19, $`n=1`$) and (21, $`n=1`$) are trivial.
Induction step $`nn+1`$.
Suppose $`m_k`$ and $`\psi _{m_k}^{(n)}`$ to be constructed for $`k=1,\mathrm{},n`$ according to (18) - (21).
Choose $`\delta _1=\delta (n,\eta _{n+1}/2)>0`$ according to Corollary 9 such that furthermore $`\delta _1<\eta _{n+1}/2`$. Let $`j\mathrm{I}N`$ be such that $`(2/j)^{1/2}<\delta _1`$. Now, again according to Corollary 9, choose $`\delta _0=\delta (nj,\eta _{n+1})`$.
Since $`(\varphi _m)`$ spans $`l^1`$ almost isometrically there is an index $`m_0`$ such that $`(\varphi _m)_{mm_0}`$ spans $`l^1`$ $`(1\delta _0)`$-isomorphically. By Corollary 9 (with $`s=t=1`$, $`\delta =\delta _0`$) we find a finite set $`N\mathrm{I}N`$ of cardinality $`nj`$ (for example $`N=\{m_0+1,\mathrm{},m_0+nj\}`$), a finite sequence of orthogonal projections $`(p_m)_{mN}`$ in $`𝒩`$ such that
$`\varphi _m{\displaystyle \frac{\varphi _mp_m}{\varphi _mp_m}}`$ $`<`$ $`\eta _{n+1}mN_{n+1}.`$ (22)
Set $`\varphi =_{k=1}^n|\psi _{m_k}^{(n)}|`$. $`\varphi `$ is positive. We have $`(_{mN}p_m)(\varphi )\varphi n`$. Thus there is an index $`m_{n+1}N`$ such that $`0p_{m_{n+1}}(\varphi )1/j`$ and $`0p_{m_{n+1}}(|\psi _{m_k}^{(n)}|)1/j`$ for $`k=1,\mathrm{},n`$. We set $`s=1\overline{p}_{m_{n+1}}`$ and define $`\stackrel{~}{\psi }_{m_k}^{(n+1)}=\psi _{m_k}^{(n)}s`$ for $`k=1,\mathrm{},n`$ and
$$\psi _{m_{n+1}}^{(n+1)}=\frac{\varphi _{m_{n+1}}p_{m_{n+1}}}{\varphi _{m_{n+1}}p_{m_{n+1}}}.$$
Then (19, $`n+1`$) holds for $`k=n+1`$ and (21, $`n+1`$) holds by (22). We have $`s(|\psi _{m_k}^{(n)}|)=\psi _{m_k}^{(n)}p_{m_{n+1}}(|\psi _{m_k}^{(n)}|)11/j`$ by (19). From this and (6) one gets that
$`\stackrel{~}{\psi }_{m_k}^{(n+1)}\psi _{m_k}^{(n)}(2/j)^{1/2}<\delta _1<{\displaystyle \frac{\eta _{n+1}}{2}},k=1,\mathrm{},n.`$ (23)
Thus, up to $`\delta _1`$ the $`\stackrel{~}{\psi }_{m_k}^{(n+1)}`$ are near to an isometric copy of $`l_n^1`$ because
$`{\displaystyle \underset{k=1}{\overset{n}{}}}|\alpha _k|`$ $``$ $`{\displaystyle \underset{k=1}{\overset{n}{}}}\alpha _k\stackrel{~}{\psi }_{m_k}^{(n+1)}={\displaystyle \underset{k=1}{\overset{n}{}}}\alpha _k\stackrel{~}{\psi }_{m_k}^{(n+1)}s`$
$``$ $`{\displaystyle \underset{k=1}{\overset{n}{}}}\alpha _k\psi _{m_k}^{(n)}{\displaystyle \underset{k=1}{\overset{n}{}}}\alpha _k(\stackrel{~}{\psi }_{m_k}^{(n+1)}\psi _{m_k}^{(n)})`$
$`\stackrel{(\text{23})}{}`$ $`{\displaystyle \underset{k=1}{\overset{n}{}}}\alpha _k\psi _{m_k}^{(n)}(2/j)^{1/2}{\displaystyle \underset{k=1}{\overset{n}{}}}|\alpha _k|`$
$`\stackrel{(\text{18},\text{19})}{=}`$ $`\left(1(2/j)^{1/2}\right){\displaystyle \underset{k=1}{\overset{n}{}}}|\alpha _k|`$
$`>`$ $`(1\delta _1){\displaystyle \underset{k=1}{\overset{n}{}}}|\alpha _k|.`$
It remains to apply Corollary 9 another time (with $`t=1`$, $`\delta =\delta _1`$) in order to get small normalized orthogonal perturbations $`\psi _{m_k}^{(n+1)}`$ \- whence (19, $`n+1`$) for $`kn`$ \- of the $`\stackrel{~}{\psi }_{m_k}^{(n+1)}`$ whose right supports are majorized by $`s`$ and thus orthogonal to the right support of $`\psi _{m_k}^{(n+1)}`$ such that $`\psi _{m_k}^{(n+1)}\stackrel{~}{\psi }_{m_k}^{(n+1)}<\eta _{n+1}/2`$ for $`k=1,\mathrm{},n`$. Together with (23) this gives (20, $`n+1`$). Finally one verifies (18, $`n+1`$) by observing that the support projections of the $`|\psi _{m_k}^{(n+1)}|`$ are the right supports of the $`\psi _{m_k}^{(n+1)}`$. This ends the induction.
By construction, $`(\psi _{m_k}^{(n)})_{n\mathrm{I}N}`$ is a Cauchy sequence for each $`k`$ because $`\psi _{m_k}^{(n)}\psi _{m_k}^{(i)}_{l=i+1}^n\eta _l0`$ as $`n>i\mathrm{}`$. Let $`\psi _k=lim_n\psi _{m_k}^{(n)}`$ be its limit. Then $`\psi _k\varphi _{m_k}\varphi _{m_k}\psi _{m_k}^{(k)}+\psi _{m_k}^{(k)}lim_n\psi _{m_k}^{(n)}\eta _k+_{l=k+1}^{\mathrm{}}\eta _l0`$ as $`k\mathrm{}`$. The $`\psi _k`$ have pairwise orthogonal right supports because by continuity of the absolute value (\[21, III.4.10\]), if $`kl`$ one has
$`|\psi _k||\psi _l|`$ $`=`$ $`\underset{n\mathrm{}}{lim}|\psi _{m_k}^{(n)}||\psi _{m_l}^{(n)}|`$
$`\stackrel{(\text{18})}{=}`$ $`\underset{n\mathrm{}}{lim}|\psi _{m_k}^{(n)}|+|\psi _{m_l}^{(n)}|=\psi _k+\psi _l.`$
So far we have proved that if $`(\varphi _m)`$ spans $`l^1`$ almost isometrically then there is a subsequence $`(\varphi _{m_k})`$ and there are pairwise orthogonal projections $`s_k𝒩`$ (namely the right support projections of the $`\psi _k`$) such that $`\varphi _{m_k}\varphi _{m_k}s_k\varphi _{m_k}\psi _k+\psi _ks_k\varphi _{m_k}s_k2\varphi _{m_k}\psi _k0`$. Since $`(\varphi _{m_k}^{})`$ spans $`l^1`$ almost isometrically, too, there are pairwise orthogonal projections $`t_l𝒩`$ such that $`\varphi _{m_{k_l}}^{}\varphi _{m_{k_l}}^{}t_l0`$ for an appropriate sequence $`(m_{k_l})`$ in $`\mathrm{I}N`$. Set $`\stackrel{~}{\varphi }_l=t_l\varphi _{m_{k_l}}s_{k_l}`$. Then $`\varphi _{m_{k_l}}\stackrel{~}{\varphi }_l\varphi _{m_{k_l}}\varphi _{m_{k_l}}s_{k_l}+(\varphi _{m_{k_l}}t_l\varphi _{m_{k_l}})s_{k_l}\varphi _{m_{k_l}}\varphi _{m_{k_l}}s_{k_l}+\varphi _{m_{k_l}}^{}\varphi _{m_{k_l}}^{}t_l0`$.
The second statement of the theorem is trivial by the definiton of the $`\stackrel{~}{\varphi }_l`$. This ends the proof.
From Remark 2 after the proof of Theorem 1 at the end of the next section it follows that Theorem 2 does not hold for unbounded sequences $`(\varphi _m)`$.
§5 Proof of Theorem 1
(i) $``$ (ii). Let $`(x_{n_k})`$ be a subsequence of $`(x_n)`$. If $`(x_{n_k})`$ contains a sequence $`(x_{n_{k_l}})`$ such that $`x_{n_{k_l}}=0`$ for all $`l\mathrm{I}N`$ then we simply choose $`y_l=0`$ for $`l\mathrm{I}N`$. Otherwise we may (pass to another subsequence and) suppose that $`x_{n_k}_10`$ for all $`k\mathrm{I}N`$. By norm density of $`\mathrm{L}^1\mathrm{L}^{\mathrm{}}`$ in $`\mathrm{L}^1`$ and the fact that the norm topology is finer than the measure topology we may suppose without loss of generality that $`0x_{n_k}_{\mathrm{}}<\mathrm{}`$ for all $`k\mathrm{I}N`$. We set $`\epsilon _l=2^l/\tau (1)`$ for $`l\mathrm{I}N`$.
By induction over $`l\mathrm{I}N`$ we construct a strictly increasing subsequence $`(n_{k_l})`$ of $`(n_k)`$, projections $`p_l𝒩`$ and positive numbers $`\delta _l`$ such that for all $`l\mathrm{I}N`$
$`\tau (p_l)<\delta _l\text{where}p_l=\chi _{]\epsilon _l,\mathrm{}[}(|x_{n_{k_l}}|)`$ (24)
and where
$`\delta _l`$ $`=`$ $`{\displaystyle \frac{2^l}{\mathrm{max}_{\mathrm{\hspace{0.17em}1}ml1}x_{n_{k_m}}_{\mathrm{}}}},\text{if}l2.`$ (25)
For $`l=1`$ we choose $`n_{k_1}=n_1`$ and any $`\delta _1>\tau (p_1)`$. For the induction step $`ll+1`$ we suppose $`n_{k_m}`$, $`p_m`$, and $`\delta _m`$ to be constructed for $`m=1,\mathrm{},l`$, we define $`\delta _{l+1}`$ by (25) and choose $`n_{k_{l+1}}`$ such that
$$\tau (\chi _{]\epsilon _{l+1},\mathrm{}[}(|x_{n_{k_{l+1}}}|)<\delta _{l+1}$$
which is possible because $`x_n\stackrel{\tau }{}0`$. We define $`p_{l+1}`$ by (24). This settles (24, $`l+1`$) and ends the induction.
By (25) we have
$`\delta _{l+1+r}={\displaystyle \frac{2^{(l+1+r)}}{\mathrm{max}_{ml+r}x_{n_{k_m}}_{\mathrm{}}}}2^{(r+1)}{\displaystyle \frac{2^l}{x_{n_{k_l}}_{\mathrm{}}}}`$
for $`r\mathrm{I}N\{0\}`$ which gives
$`{\displaystyle \underset{ml+1}{}}\delta _m={\displaystyle \underset{r0}{}}\delta _{l+1+r}{\displaystyle \frac{2^l}{x_{n_{k_l}}_{\mathrm{}}}}.`$ (26)
Put $`q_l=1_{ml+1}p_m`$ and $`\stackrel{~}{y}_l=x_{n_{k_l}}(p_lq_l)`$. By construction the $`\stackrel{~}{y}_l`$ have pairwise orthogonal right support projections and their left support projections are majorized by the ones of the $`x_{n_{k_l}}`$. We show that $`x_{n_{k_l}}\stackrel{~}{y}_l_10`$. In order to save indices we use the abbreviations $`x=x_{n_{k_l}}`$, $`p=p_l`$, $`q=q_l`$, $`\stackrel{~}{y}=\stackrel{~}{y}_l`$ until the end of formula (27):
$`x\stackrel{~}{y}_1`$ $``$ $`xxp_1+xp\stackrel{~}{y}_1`$ (27)
$`=`$ $`x(1p)_1+x(p(pq))_1`$
$``$ $`x(1p)_{\mathrm{}}\tau (1)+x_{\mathrm{}}\tau (p(pq))`$
$`\stackrel{()}{=}`$ $`|x|\chi _{[0,\epsilon _l]}(|x|)_{\mathrm{}}\tau (1)+x_{\mathrm{}}\tau \left((pq)q\right)`$
$``$ $`\epsilon _l\tau (1)+x_{\mathrm{}}\tau (1q)`$
$``$ $`\epsilon _l\tau (1)+x_{\mathrm{}}\left({\displaystyle \underset{ml+1}{}}\tau (p_m)\right)`$
$`\stackrel{(\text{24}),(\text{26})}{}`$ $`2^{(l1)}.`$
For $`()`$ we used that $`p(pq)`$ and $`(pq)q`$ are equivalent projections for any two projections $`p,q`$ (\[21, V.1.6\]) hence $`\tau (p(pq))=\tau ((pq)q)`$.
So far we have proved that given a $`\tau `$-null subsequence $`(x_{n_k})`$ there are $`x_{n_{k_l}}`$ and there are $`\stackrel{~}{y}_l`$ whose right supports are orthogonal and whose left supports are majorized by the left supports of the $`x_{n_{k_l}}`$ such that $`x_{n_{k_l}}\stackrel{~}{y}_l_10`$. In particular, $`\stackrel{~}{y}_l\stackrel{\tau }{}0`$ whence $`\stackrel{~}{y}_l^{}\stackrel{\tau }{}0`$. Thus we can apply the same reasoning (up to passing to appropriate subsequences) in order to find perturbations $`y_l^{}`$ of the $`\stackrel{~}{y}_l^{}`$ which have both orthogonal right and orthogonal left supports such that $`\stackrel{~}{y}_ly_l_1=\stackrel{~}{y}_l^{}y_l^{}_10`$ hence $`x_{n_{k_l}}y_10`$. This ends the proof of (i) $``$ (ii).
(ii) $``$ (i) Since $`\tau `$ is finite and the $`y_l`$ are pairwise orthogonal we have that $`y_l\stackrel{\tau }{}0`$. And $`x_{n_{k_l}}y_l0`$ entails $`x_{n_{k_l}}y_l\stackrel{\tau }{}0`$ hence $`x_{n_{k_l}}\stackrel{\tau }{}0`$. Thus each subsequence of $`(x_n)`$ contains a subsequence which converges to $`0`$ in measure whence $`x_n\stackrel{\tau }{}0`$.
(ii) $``$ (iii) follows from Lemma 4: Suppose (ii) holds and $`infx_{n_k}_1>0`$ for a subsequence $`(x_{n_k})`$ of $`(x_n)`$. Then by (ii), there are orthogonal $`y_l`$ and there is $`(x_{n_{k_l}})`$ such that $`x_{n_{k_l}}y_l_10`$. One may suppose that $`infy_l_1>0`$ hence $`(y_l)`$ spans $`l^1`$ isometrically. Thus by Lemma 4, the sequence $`(x_{n_{k_l}})=(y_l+(x_{n_{k_l}}y_l))`$ spans $`l^1`$ almost isometrically.
(iii) $``$ (iv): Von Neumann preduals are L-embedded spaces \[6, IV.1.1\], thus by each sequence spanning $`l^1`$ almost isometrically admits a subsequence spanning $`l^1`$ asymptotically.
(iv) $``$ (iii) is trivial.
(iii) $``$ (ii) follows from Theorem 2.
See the following Remark 2 for an example which shows that in general (iii) does not imply (i), (ii) for unbounded sequences $`(x_n)`$.
Remarks:
1. As an illustration of how to get an orthogonal subsequence consider the sequence $`x_n=n\mathrm{\hspace{0.17em}1}_{[0,\mathrm{\hspace{0.17em}1}/n]}`$ in $`\mathrm{L}^1([0,1])`$. One may take, for example, $`y_l=x_{n_l}1_{]1/n_{l+1},1]}=n_l1_{]1/n_{l+1},1/n_l]}`$ where $`n_l=2^{(2^l)}`$.
2. In general (iii) does not imply (i), neither (ii), if the sequence $`(x_n)`$ is unbounded. Take the bounded sequence $`x_n=n^21_{[1/n+1,\mathrm{\hspace{0.33em}1}/n[}+\frac{1}{n}`$ in $`\mathrm{L}^1([0,1])`$. It converges to zero in measure and does not contain a norm null sequence. Hence by (i)$``$(iii) an appropriate subsequence $`(x_{n_k})`$ spans $`l^1`$ almost isometrically. Thus the unbounded sequence $`(n_k^2x_{n_k})`$ satisfies (iii) but not (i). It cannot satisfy (ii) either because (ii) $``$ (i) holds also for unbounded sequences. This means in particular that Theorem 2 does not hold for unbounded sequences $`(\varphi _m)`$.
3. A few straightforward modifications show that (i)$``$(ii) holds accordingly also for $`\mathrm{L}^p(𝒩,\tau )`$, $`1p<\mathrm{}`$. (Cf. .)
§6 $`l^1`$-copies in the dual of C-algebras, proof of Proposition 3
The proof of the main result of gives the following: Let $`(\varphi _m)A^{}`$ be a bounded sequence of selfadjoint functionals on a C-algebra $`A`$, let $`\epsilon >0`$. If $`(\varphi _m)`$ spans $`l^1`$ $`r`$-isomorphically ($`0<r<1`$) then there is a subsequence $`(\varphi _{m_n})`$ and there is a sequence $`(x_n)`$ of pairwise orthogonal normalized selfadjoint elements of $`A`$ such that $`\varphi _{m_n}(x_n)>(1\epsilon )r\varphi _{m_n}`$. This amounts to saying that $`|\varphi _{m_n}|(|x_n|)>(1\epsilon )r\varphi _{m_n}`$ (to see this it is enough to decompose both $`\varphi _{m_n}`$ and $`x_n`$ in their positive and negative parts) or that, via Lemma 6, $`\varphi _{m_n}a_n\varphi _{m_n}a_n0`$ where $`a_n=|x_n|`$. This is Lemma 10 for selfadjoint $`\varphi _m`$ with the better factor $`r`$ instead of $`r^2`$ in $`(\text{29})`$ and $`(\text{30})`$.
With Lemma 8 at one’s disposal, the proof of Lemma 10 \- and thus of Proposition 3 \- is a straightforward modification of and gives a kind of quantitative version of which holds for arbitrary functionals, not only selfadjoint ones. (We give the entire proof of Proposition 3 not only for the sake of completeness but also because it is quite lengthy wherefore the usual argument ”The details are left to the reader” would be exaggerated.) Yet, it does not complete the subject ”perturbations of $`l^1`$-copies in C-algebras” as at least two questions remain open.
Firstly, is it necessary in Theorem 2 or in Proposition 3 to pass to subsequences? In the commutative case it is not, as a result of Dor shows that, if $`𝒩_{}=\mathrm{L}^1([0,1])`$ contains a $`(1\delta )`$-isomorphic copy of $`l^1`$ then the whole canonical basis of this copy can be perturbed in norm so to span $`l^1`$ isometrically with the perturbation smaller than $`\delta ^{}`$ and $`\delta ^{}0`$ as $`\delta 0`$. Furthermore Arazy proved that if the predual of an arbitrary von Neumann algebra $`𝒩`$ contains a $`(1\delta )`$-copy of $`l^1`$ then the whole copy is complemented by a projection whose norm is majorized by $`1+\delta ^{}`$ \- a result which has recently been generalized by N. Ozawa to the category of operator spaces.
Secondly, can the $`(m_n)`$, $`(a_n)`$ and $`(b_n)`$ in Proposition 3 be arranged such that $`\varphi _{m_n}b_n\varphi _{m_n}a_n0`$ as $`n\mathrm{}`$? \[Let us sketch in passing why this would generalize Lemma 10. If $`(\varphi _m)A^{}`$ is normalized and spans $`l^1`$ $`r`$-isomorphically then by James’ distortion theorem there are blocks $`\psi _n=_{iF_n}\lambda _i\varphi _i`$ spanning $`l^1`$ almost isometrically such that $`_{iF_n}|\lambda _i|1/r`$. Now, if there are appropriate $`a_n,b_nA`$ such that $`\psi _nb_n\psi _na_n0`$ (after passing, if necessary, to an appropriate subsequence of $`(\psi _n)`$), one deduces that $`|\psi _n|(a_n)>\sqrt{1\epsilon _n}`$ with $`0<\epsilon _n0`$. Thus for each $`n`$ there is $`i_nF_n`$ such that $`|\varphi _{i_n}|(a_n)>(1\epsilon _n)r^2`$ because otherwise by \[21, III.4.7\] one would have the contradiction
$`\sqrt{1\epsilon _n}`$ $`<`$ $`|\psi _n|(a_n)=|{\displaystyle \underset{iF_n}{}}\lambda _i\varphi _i|(a_n)\left({\displaystyle \underset{F_n}{}}|\lambda _i|\varphi _i\right)^{1/2}\left({\displaystyle \underset{F_n}{}}|\lambda _i|(|\varphi _i|(a_n^2))\right)^{1/2}`$
$``$ $`{\displaystyle \frac{1}{r}}\left(\underset{F_n}{\mathrm{max}}|\varphi _i|(a_n^2)\right)^{1/2}\sqrt{1\epsilon _n}.`$
Similarly one obtains $`|\varphi _{i_n}^{}|(b_n)>(1\epsilon _n)r^2`$.\]
Proposition 3 follows immediately from Lemma 10 (and (6) of Lemma 6) with $`s=1=t`$. The technical part concerning $`s,t`$ is added because it might be usefull for answering the second question just mentionned above.
###### Lemma 10
Let $`A`$ be a C-algebra (unital or not), $`r>0`$, let $`(\varphi _m)`$ be a normalized sequence in $`A^{}`$ spanning $`l^1`$ $`r`$-isomorphically that is such that
$`r{\displaystyle |\alpha _m|}{\displaystyle \alpha _m\varphi _m}{\displaystyle |\alpha _m|(\alpha _m)}\mathrm{C}.`$ (28)
Then, given $`\epsilon >0`$, there are a sequence $`(m_n)`$ in $`\mathrm{I}N`$ and a sequence $`(a_n)`$ of pairwise orthogonal positive normalized elements in $`A`$ and another sequence $`(b_n)`$ of pairwise orthogonal positive normalized elements in $`A`$ such that
$`|\varphi _{m_n}|(a_n)`$ $`>`$ $`(1\epsilon )r^2`$ (29)
$`|\varphi _{m_n}^{}|(b_n)`$ $`>`$ $`(1\epsilon )r^2`$ (30)
for each $`n\mathrm{I}N`$.
Moreover, if $`s`$ and $`t`$ are open projections in $`A^{\prime \prime }`$ such that $`s`$ ($`t`$) majorizes the right (left) supports of all $`\varphi _m`$ (that is $`t\varphi _ms=\varphi _m`$ for all $`m\mathrm{I}N`$) then one can obtain in addition that $`a_nsA^{\prime \prime }s`$ and $`b_ntA^{\prime \prime }t`$.
Moreover, if the $`\varphi _m`$ are selfadjoint one can obtain in addition $`|\varphi _{m_n}|(a_n)=|\varphi _{m_n}^{}|(a_n)>(1\epsilon )r`$ instead of $`(\text{29})`$ and $`(\text{30})`$.
Proof: First we suppose that $`A`$ is unital.
It is enough to construct a sequence $`(p_n)`$ of orthogonal open projections in $`sA^{\prime \prime }s`$ such that
$`|\varphi _{m_n}|(p_n)`$ $`>`$ $`(1\epsilon )r^2`$ (31)
for an appropriate subsequence $`(\varphi _{m_n})`$ because then, by the definition of open projections, for all $`n\mathrm{I}N`$ positive elements $`a_np_n`$ can be choosen so to be pairwise orthogonal (since the $`p_n`$ are) and so to satisfy (29); finally, since (28) remains valid if $`\varphi _n^{}`$ is substituted for $`\varphi _n`$ the same reasoning that leads to (29) shows the existence of a sequence $`(b_n)`$ in $`tA^{\prime \prime }t`$ as desired in (30).
Let $`0<\epsilon <1`$ and choose a sequence $`(\epsilon _n)`$ of positive numbers such that $`\epsilon _n=\epsilon `$ and $`\epsilon _n\frac{3}{4}`$ for all $`n\mathrm{I}N`$.
By induction over $`n=1,2,\mathrm{}`$ we construct a sequence $`(p_n)`$ of open projections in $`sA^{\prime \prime }s`$, a sequence of indices $`(m_n)`$, a decreasing sequence $`(N_n)`$ of infinite subsets of $`\mathrm{I}N`$, i.e. $`\mathrm{}N_{n+1}N_n\mathrm{}N_1N_0=\mathrm{I}N`$, such that we have for all $`n\mathrm{I}N`$:
$`\overline{p_n}`$ $``$ $`sA^{\prime \prime }s`$ (32)
$`\overline{p_i}\overline{p_n}`$ $`=`$ $`0i<n`$ (33)
$`\overline{p_n}(|\varphi _m|)`$ $`<`$ $`{\displaystyle \frac{1}{72}}r^2\epsilon _n^4mN_n`$ (34)
$`p_n(|\varphi _{m_n}|)`$ $`>`$ $`r^2(1{\displaystyle \underset{1}{\overset{n}{}}}\epsilon _i)`$ (35)
$`m_n`$ $``$ $`N_{n1}.`$ (36)
We start the induction with $`n=1`$.
Choose $`j_1\mathrm{I}N`$ with $`1/j_1<r^2\epsilon _1^4/72`$. For $`j_1`$ and $`\epsilon _1/4`$ Lemma 8 yields a number $`\delta _1=\delta _1(j_1,\epsilon _1/4)>0`$ and without loss of generality we assume $`\delta _1\epsilon _1/4`$. By James’ distortion theorem applied to (28) there are pairwise disjoint finite sets $`F_k^{(1)}N_0=\mathrm{I}N`$, a finite sequence $`(\lambda _i^{(1)})_{iF_k^{(1)}}\mathrm{C}\text{ }`$ and functionals $`\tau _k^{(1)}=_{iF_k^{(1)}}\lambda _i^{(1)}\varphi _i`$ for $`k\mathrm{I}N`$, such that
$`{\displaystyle \underset{F_k^{(1)}}{}}|\lambda _i^{(1)}|`$ $``$ $`{\displaystyle \frac{1}{r}},`$ (37)
$`(1\delta _1){\displaystyle \underset{k1}{}}|\alpha _k|`$ $``$ $`{\displaystyle \underset{k1}{}}\alpha _k\tau _k^{(1)}{\displaystyle \underset{k1}{}}|\alpha _k|(\alpha _k)_{kN_0}\mathrm{C}.`$ (38)
Again by Lemma 8 there are pairwise orthogonal open projections $`p_k^{(1)}sA^{\prime \prime }s`$, $`kj_1`$, such that
$`p_k^{(1)}(|\tau _k^{(1)}|)>(1\epsilon _1/4)\tau _k^{(1)}kj_1,`$ (39)
and since the projections can be chosen to have orthogonal closures in $`sA^{\prime \prime }s`$ we have
$`\left({\displaystyle \underset{1}{\overset{j_1}{}}}\overline{p_k^{(1)}}\right)(|\varphi _m|)1mN_0.`$
Therefore there exist a $`k_1j_1`$ and an infinite set $`N_1N_0`$ such that
$`\overline{p_{k_1}^{(1)}}(|\varphi _m|){\displaystyle \frac{1}{j_1}}<{\displaystyle \frac{r^2\epsilon _1^4}{72}}mN_1.`$ (40)
Set $`p_1=p_{k_1}^{(1)}`$, $`\tau _1=\tau _{k_1}^{(1)}`$, $`F_1=F_{k_1}^{(1)}`$. Then (32) holds for $`n=1`$. Now we infer that
$`p_1(|\tau _1|)\stackrel{(\text{39})}{>}\tau _1(1{\displaystyle \frac{\epsilon _1}{4}})\stackrel{(\text{38})}{}(1\delta _1)(1{\displaystyle \frac{\epsilon _1}{4}})(1{\displaystyle \frac{\epsilon _1}{4}})^2>\sqrt{1\epsilon _1},`$
which in turn yields the existence of an index $`m_1F_1N_0`$ as desired in (35) and (36) for $`n=1`$, because otherwise we would have
$`p_1(|\tau _1|)`$ $`=`$ $`p_1\left(|{\displaystyle \underset{F_1}{}}\lambda _i^{(1)}\varphi _i|\right)`$
$`\stackrel{()}{}`$ $`\left({\displaystyle \underset{F_1}{}}\lambda _i^{(1)}\varphi _i\right)^{1/2}\left({\displaystyle \underset{F_1}{}}p_1\left(|\lambda _i^{(1)}\varphi _i|\right)\right)^{1/2}`$
$``$ $`\left({\displaystyle \underset{F_1}{}}|\lambda _i^{(1)}|\right)^{1/2}\left(\underset{iF_1}{\mathrm{max}}p_1(|\varphi _i|){\displaystyle \underset{F_1}{}}|\lambda _i^{(1)}|\right)^{1/2}`$
$`\stackrel{(\text{37})}{}`$ $`{\displaystyle \frac{1}{r}}\left(r^2(1\epsilon _1)\right)^{1/2}=\sqrt{1\epsilon _1}.`$
Here inequality $`()`$ follows from \[21, III.4.7\]. For (33, n=1) nothing needs to be proved. Inequality (34, n=1) corresponds to (40). The first induction step is done.
Induction step $`nn+1`$:
Suppose $`p_k`$, $`N_k`$, $`m_k`$ to be constructed for $`kn`$ according to (33) – (36).
Since the $`\overline{p_k}`$ are orthogonal in $`sA^{\prime \prime }s`$, $`_1^n\overline{p_k}`$ is closed by \[1, Th. II.7\]. Therefore $`s_n=s_1^n\overline{p_k}sA^{\prime \prime }s`$ is open. Set $`\stackrel{~}{\varphi }_m=\varphi _ms_n`$.
Claim:
The normalized functionals $`\left(\frac{\stackrel{~}{\varphi }_m}{\stackrel{~}{\varphi }_m}\right)_{mN_n}`$ form an $`l^1`$-basis with
$`r\left(1{\displaystyle \underset{1}{\overset{n}{}}}\epsilon _i^2\right)^{1/2}{\displaystyle \underset{mN_n}{}}|\alpha _m|`$ $``$ $`{\displaystyle \underset{mN_n}{}}\alpha _m{\displaystyle \frac{\stackrel{~}{\varphi }_m}{\stackrel{~}{\varphi }_m}}`$ (41)
$``$ $`{\displaystyle \underset{mN_n}{}}|\alpha _m|(\alpha _m)_{mN_n}\mathrm{C}.`$
Set $`\eta =\frac{r^2}{72}_1^n\epsilon _i^4`$. Then
$`(ss_n)(|\varphi _m|)=\left({\displaystyle \underset{1}{\overset{n}{}}}\overline{p_k}\right)(|\varphi _m|)\stackrel{(\text{34})}{<}\eta mN_n,`$ (42)
thus since $`s(|\varphi _m|)=\varphi _m=1`$
$`\varphi _ms_n\varphi _m`$ $`\stackrel{(\text{6})}{}`$ $`|2(\varphi _ms_n(|\varphi _m|))|^{1/2}`$ (43)
$`=`$ $`\left(2{\displaystyle \underset{1}{\overset{n}{}}}\overline{p}_k(|\varphi _m|)\right)^{1/2}\stackrel{(\text{42})}{<}\sqrt{2\eta }mN_n;`$
further we note that for all $`mN_n`$
$`\stackrel{~}{\varphi }_m`$ $``$ $`\varphi _m=1,`$ (44)
$`01\stackrel{~}{\varphi }_m`$ $`=`$ $`\varphi _m\stackrel{~}{\varphi }_m\varphi _m\stackrel{~}{\varphi }_m\stackrel{(\text{43})}{}\sqrt{2\eta }`$ (45)
$`\stackrel{~}{\varphi }_m`$ $`\stackrel{(\text{45})}{}`$ $`1\sqrt{2\eta },`$ (46)
hence
$`{\displaystyle \frac{\stackrel{~}{\varphi }_m}{\stackrel{~}{\varphi }_m}}\stackrel{~}{\varphi }_m`$ $`\stackrel{(\text{44})}{}`$ $`{\displaystyle \frac{1}{\stackrel{~}{\varphi }_m}}1=\left(1\stackrel{~}{\varphi }_m\right){\displaystyle \frac{1}{\stackrel{~}{\varphi }_m}}`$ (47)
$`\stackrel{(\text{45})(\text{46})}{}`$ $`{\displaystyle \frac{\sqrt{2\eta }}{1\sqrt{2\eta }}}`$
and
$`{\displaystyle \frac{\stackrel{~}{\varphi }_m}{\stackrel{~}{\varphi }_m}}\varphi _m`$ $``$ $`{\displaystyle \frac{\stackrel{~}{\varphi }_m}{\stackrel{~}{\varphi }_m}}\stackrel{~}{\varphi }_m+\stackrel{~}{\varphi }_m\varphi _m`$ (48)
$`\stackrel{(\text{47})(\text{43})}{}`$ $`\sqrt{2\eta }\left(1+{\displaystyle \frac{1}{1\sqrt{2\eta }}}\right)<3\sqrt{2\eta }`$
because $`\epsilon <1`$, $`r1`$, thus $`\sqrt{2\eta }<1/2`$. Then (41) follows from
$`{\displaystyle \underset{mN_n}{}}\alpha _m{\displaystyle \frac{\stackrel{~}{\varphi }_m}{\stackrel{~}{\varphi }_m}}`$ $``$ $`{\displaystyle \underset{N_n}{}}\alpha _m\varphi _m{\displaystyle \underset{N_n}{}}\alpha _m{\displaystyle \frac{\stackrel{~}{\varphi }_m}{\stackrel{~}{\varphi }_m}}{\displaystyle \underset{N_n}{}}\alpha _m\varphi _m`$
$`\stackrel{(\text{28})}{}`$ $`r\left(1\underset{N_n}{sup}{\displaystyle \frac{\stackrel{~}{\varphi }_m}{\stackrel{~}{\varphi }_m}}\varphi _m\right){\displaystyle \underset{N_n}{}}|\alpha _m|`$
$`\stackrel{(\text{48})}{}`$ $`r(13\sqrt{2\eta }){\displaystyle \underset{N_n}{}}|\alpha _m|`$
$`=`$ $`r\left(1{\displaystyle \frac{1}{2}}({\displaystyle \underset{1}{\overset{n}{}}}\epsilon _i^4)^{1/2}\right){\displaystyle \underset{N_n}{}}|\alpha _m|`$
$`>`$ $`r\left(1({\displaystyle \underset{1}{\overset{n}{}}}\epsilon _i^4)^{1/2}\right)^{1/2}{\displaystyle \underset{N_n}{}}|\alpha _m|`$
$`>`$ $`r\left(1{\displaystyle \underset{1}{\overset{n}{}}}\epsilon _i^2\right)^{1/2}{\displaystyle \underset{mN_n}{}}|\alpha _m|`$
and the Claim is established.
Choose a number $`j_{n+1}\mathrm{I}N`$ such that $`1/j_{n+1}<r^2\epsilon _{n+1}^2/4`$. Further choose a number $`\delta _{n+1}=\delta _{n+1}(j_{n+1},\epsilon _{n+1}^2/4)>0`$ according to Lemma 8 and such that moreover $`\delta _{n+1}\epsilon _{n+1}^2/4`$. Now we apply James’ distortion theorem. By (41) there are pairwise disjoint finite sets $`F_k^{(n+1)}N_n`$, a finite sequence $`(\lambda _i^{(n+1)})_{iF_k^{(n+1)}}\mathrm{C}\text{ }`$ and functionals $`\tau _k^{(n+1)}=_{iF_k^{(n+1)}}\lambda _i^{(n+1)}\frac{\stackrel{~}{\varphi }_i}{\stackrel{~}{\varphi }_i}`$ for each $`k\mathrm{I}N`$ such that
$`{\displaystyle \underset{iF_k^{(n+1)}}{}}|\lambda _i^{(n+1)}|`$ $``$ $`{\displaystyle \frac{1}{r(1_1^n\epsilon _i^2)^{1/2}}}k\mathrm{I}N,`$ (49)
$`(1\delta _{n+1}){\displaystyle \underset{k1}{}}|\alpha _k|`$ $``$ $`{\displaystyle \underset{k1}{}}\alpha _k\tau _k^{(n+1)}s_n{\displaystyle \underset{k1}{}}|\alpha _k|.`$ (50)
Again by Lemma 8, applied to the open projections $`s_n`$ and $`1`$, to the functionals $`\tau _k^{(n+1)}A^{}`$, and to (50), there exist open projections $`p_k^{(n+1)}s_nA^{\prime \prime }s_n`$, $`kj_{n+1}`$, with pairwise orthogonal closures in $`s_nA^{\prime \prime }s_n`$ such that
$`p_k^{(n+1)}(|\tau _k^{(n+1)}|)>\tau _k^{(n+1)}\left(1{\displaystyle \frac{\epsilon _{n+1}^2}{4}}\right)`$ (51)
for $`kj_{n+1}`$. Since the projections have orthogonal closures we have
$`\left({\displaystyle \underset{1}{\overset{j_{n+1}}{}}}\overline{p_k^{(n+1)}}\right)(|\varphi _m|)1mN_n.`$
Therefore there exist an index $`k_{n+1}j_{n+1}`$ and an infinite subset $`N_{n+1}N_n`$ such that
$`\overline{p_{k_{n+1}}^{(n+1)}}(|\varphi _m|){\displaystyle \frac{1}{j_{n+1}}}<{\displaystyle \frac{r^2\epsilon _{n+1}^4}{72}}mN_{n+1}.`$
Set $`p_{n+1}=p_{k_{n+1}}^{(n+1)}`$, $`\tau _{n+1}=\tau _{k_{n+1}}^{(n+1)}`$, $`F_{n+1}=F_{k_{n+1}}^{(n+1)}`$. Then (32), (33) and (34) hold for $`n+1`$. Now we infer that
$`p_{n+1}(|\tau _{n+1}|)`$ $`\stackrel{(\text{51})}{>}`$ $`\tau _{n+1}\left(1{\displaystyle \frac{\epsilon _{n+1}^2}{4}}\right)`$ (52)
$`\stackrel{(\text{50})}{}`$ $`(1\delta _{n+1})(1{\displaystyle \frac{\epsilon _{n+1}^2}{4}})>\sqrt{1\epsilon _{n+1}^2},`$
hence there is an index $`m_{n+1}F_{n+1}N_n`$ as desired in (36) such that
$`p_{n+1}\left({\displaystyle \frac{\stackrel{~}{\varphi }_{m_{n+1}}}{\stackrel{~}{\varphi }_{m_{n+1}}}}\right)>r^2(1\epsilon _{n+1}^2)\left(1{\displaystyle \underset{1}{\overset{n}{}}}\epsilon _i^2\right),`$ (53)
because otherwise the following estimates would contradict (52):
$`p_{n+1}(|\tau _{n+1}|)`$ $`=`$ $`p_{n+1}\left(|{\displaystyle \underset{F_{n+1}}{}}\lambda _i^{(n+1)}{\displaystyle \frac{\stackrel{~}{\varphi _i}}{\stackrel{~}{\varphi _i}}}|\right)`$
$`\stackrel{()}{}`$ $`\left({\displaystyle \underset{F_{n+1}}{}}\lambda _i^{(n+1)}{\displaystyle \frac{\stackrel{~}{\varphi _i}}{\stackrel{~}{\varphi _i}}}\right)^{1/2}\left({\displaystyle \underset{F_{n+1}}{}}p_{n+1}(|\lambda _i^{(n+1)}{\displaystyle \frac{\stackrel{~}{\varphi _i}}{\stackrel{~}{\varphi _i}}}|)\right)^{1/2}`$
$``$ $`\left({\displaystyle \underset{F_{n+1}}{}}|\lambda _i^{(n+1)}|\right)^{1/2}\left(\underset{iF_{n+1}}{\mathrm{max}}p_{n+1}(|{\displaystyle \frac{\stackrel{~}{\varphi _i}}{\stackrel{~}{\varphi _i}}}|){\displaystyle \underset{F_{n+1}}{}}|\lambda _i^{(n+1)}|\right)^{1/2}`$
$`\stackrel{(\text{49})}{}`$ $`{\displaystyle \frac{1}{r(1_1^n\epsilon _i^2)^{1/2}}}\left(r^2(1\epsilon _{n+1}^2)(1{\displaystyle \underset{1}{\overset{n}{}}}\epsilon _i^2)\right)^{1/2}=\sqrt{1\epsilon _{n+1}^2}.`$
Here inequality $`()`$ follows from \[21, III.4.7\]. Note that for a functional $`\varphi `$ with polar decomposition $`\varphi =u|\varphi |`$ one has $`|\varphi s|=|u|\varphi |s||\varphi |`$ by \[21, III.4.9\] which explains inequality $`()`$ below; now (35, $`n+1`$) follows from
$`p_{n+1}(|\varphi _{m_{n+1}}|)`$ $`\stackrel{()}{}`$ $`p_{n+1}(|\varphi _{m_{n+1}}s_n|)`$
$`\stackrel{(\text{53})}{>}`$ $`\varphi _{m_{n+1}}s_nr^2(1\epsilon _{n+1}^2)\left(1{\displaystyle \underset{1}{\overset{n}{}}}\epsilon _i^2\right)`$
$`\stackrel{(\text{46})}{}`$ $`\left(1\sqrt{2\eta }\right)r^2(1\epsilon _{n+1}^2)\left(1{\displaystyle \underset{1}{\overset{n}{}}}\epsilon _i^2\right)`$
$`>`$ $`r^2\left(1{\displaystyle \underset{1}{\overset{n+1}{}}}\epsilon _i\right)`$
where the last inequality follows from the following completely elementary estimates:
$`\left(1\sqrt{2\eta }\right)(1\epsilon _{n+1}^2)\left(1{\displaystyle \underset{1}{\overset{n}{}}}\epsilon _i^2\right)=`$
$`=`$ $`\left(1{\displaystyle \frac{r}{6}}({\displaystyle \underset{1}{\overset{n}{}}}\epsilon _i^4)^{1/2}\right)(1\epsilon _{n+1}^2)\left(1{\displaystyle \underset{1}{\overset{n}{}}}\epsilon _i^2\right)`$
$`>`$ $`\left(1{\displaystyle \frac{r}{3}}{\displaystyle \underset{1}{\overset{n}{}}}\epsilon _i^2\right)(1\epsilon _{n+1}^2)\left(1{\displaystyle \underset{1}{\overset{n}{}}}\epsilon _i^2\right)`$
$`=`$ $`\left(1{\displaystyle \underset{1}{\overset{n+1}{}}}\epsilon _i^2\right){\displaystyle \frac{r}{3}}\left[1(1\epsilon _{n+1}^2){\displaystyle \underset{1}{\overset{n}{}}}\epsilon _i^2\left(1+{\displaystyle \frac{3}{r}}\right)\epsilon _{n+1}^2\right]{\displaystyle \underset{1}{\overset{n}{}}}\epsilon _i^2`$
$`>`$ $`\left(1{\displaystyle \underset{1}{\overset{n+1}{}}}\epsilon _i^2\right){\displaystyle \frac{r}{3}}{\displaystyle \underset{1}{\overset{n}{}}}\epsilon _i^2>\left(1{\displaystyle \underset{1}{\overset{n+1}{}}}\epsilon _i^2\right){\displaystyle \frac{1}{3}}{\displaystyle \underset{1}{\overset{n+1}{}}}\epsilon _i^2`$
$`=`$ $`\left(1{\displaystyle \underset{1}{\overset{n+1}{}}}\epsilon _i\right)+{\displaystyle \underset{1}{\overset{n+1}{}}}\epsilon _i(1\epsilon _i{\displaystyle \frac{\epsilon _i}{3}})1{\displaystyle \underset{1}{\overset{n+1}{}}}\epsilon _i`$
since we assumed $`\epsilon _i\frac{3}{4}`$ for all $`i\mathrm{I}N`$. Thus (35, $`n+1`$) is proved. This ends the induction and the proof if $`A`$ is unital.
If $`A`$ is not unital we consider its unitisation $`A_1=A\mathrm{C}\text{ }\mathrm{\hspace{0.33em}1}`$ on which multiplication is defined by $`(a,\lambda )(b,\mu )=(ab+\lambda b+\mu a,\lambda \mu )`$. Note that if $`\widehat{\varphi }A_1^{}`$ is a norm preserving Hahn-Banach extension of $`\varphi A^{}`$ then $`|\widehat{\varphi }|A_1^{}`$ is an extension of $`|\varphi |A^{}`$. \[To see this let $`u_nA`$, $`u_n1`$ be such that $`\widehat{\varphi }=\varphi =lim\varphi (u_n)`$. Then both $`|\varphi |u_n\varphi 0`$ and $`|\widehat{\varphi }|(u_n,0)\widehat{\varphi }0`$ by (5) of Lemma 6. Hence, $`|\widehat{\varphi }|((x,0))=lim\widehat{\varphi }((x,0)(u_n,0))=lim\widehat{\varphi }((xu_n,0))=lim\varphi (xu_n)=|\varphi |(x)`$ for any $`x=(x,0)A`$.\] Let $`\widehat{\varphi }_m`$ be a norm preserving Hahn-Banach extension of $`\varphi _m`$. Then (28) remains valid and by what has been proved for the unital case we get normalized pairwise orthogonal positive $`(a_n,\lambda _n)A_1`$ such that $`(a_n,\lambda _n)s`$ and $`|\widehat{\varphi }_{m_n}|((a_n,\lambda _n))>(1\epsilon )r^2`$ for an appropriate subsequence $`(\widehat{\varphi }_{m_n})`$. Since the $`(a_n,\lambda _n)`$ are pairwise orthogonal all (but possibly one) of them have $`\lambda _n=0`$ whence $`|\varphi _{m_n}|(a_n)=|\widehat{\varphi }_{m_n}|((a_n,0))>(1\epsilon )r^2`$. Likewise, we get (30).
The last statement concerning selfadjoint functionals has been discussed in the beginning of this section.
Acknowledgement I thank Dirk Werner for several helpful discussions.
Hermann Pfitzner
Université d’Orléans
BP 6759
F-45067 Orléans Cedex 2
France
e-mail: pfitzner@labomath.univ-orleans.fr |
warning/0003/nlin0003046.html | ar5iv | text | # Variational principles in the analysis of traffic flows. (Why it is worth to go against the flow.)
## 1 Introduction
During the last decade problems related to transport in complex systems attracted a huge amount of interest in particular due to their evident practical importance. In this paper we deal with theoretical aspects of phenomena arising in the modeling of highway traffic flow. Previously much of the effort in the construction and analysis of such models was concentrated on discrete (on time and on space) stochastic models introduced in and later studied by many authors (see for review and further references). All these models were based on the idea to describe the dynamics in terms of cellular automata and to a large extent were studied by means of numerical simulation (especially because of low computational cost of the numerical realization of cellular automata rules, which made it possible to realize large-scale real-time simulations of urban traffic ).
My own interest to this type of problems is mainly due to the following practical observation. Going by foot in a slowly moving crowd it is faster to go against the “flow” than in the same direction as other people go. This effect is especially pronounced if one is moving near a boundary between two “flows” of people going in opposite directions. A standard probabilistic model of a diffusion of a particle against/along the flow clearly contradicts to this observation, which very likely indicates a very special (nonrandom) intrinsic structure of the flow in this case. The main aim of the present paper is to study how this structure emerges from arbitrary (random) initial configurations in some simple models of the traffic flow.
The main quantity of interest in traffic models is the average velocity of cars $`\mathrm{V}`$ and its dependence on the density of cars $`\rho `$ (called a fundamental diagram) is typically studied in the steady state. Various approaches starting from the mean-field approximation to combinatorial techniques and statistical mechanics methods were used in the analysis of this type of models. In what follows we shall restrict the consideration to deterministic discrete time and space traffic models. The simplest model (among that we consider) describes the dynamics of cars moving along a one-row motor road and is defined as follows.
The road is associated to a one-dimensional ordered lattice of size $`N`$ with periodic boundary conditions and each position on the lattice is either occupied by a particle, or empty. Denote the number of particles in the configuration by $`m`$. Then the density of particles $`\rho `$ is equal to $`\frac{m}{N}`$. On the next time step each particle remains on its place if the next position is occupied, and moves forward by one place otherwise. We shall call this model the traffic model with slow particles to distinguish it from other ones.
Quite recently in this model (and some its generalizations) was studied by means of cellular automata techniques. One of the most intriguing phenomenon related to this model is a drastic change of the shape of the curve describing the dependence $`\mathrm{V}(\rho )`$ of the average velocity of particles on their density in the steady state when the density passes the $`1/2`$ value (see Fig. 1). It was found numerically and confirmed analytically in the limit of large $`N`$ (and for a typical initial configuration) that $`\mathrm{V}(\rho )`$ is equal to $`1`$ while $`\rho <1/2`$ and then goes down to zero as $`\frac{1}{\rho }1`$ for $`\rho 1/2`$.
Despite an apparent simplicity of the model with slow particles its dynamics (especially during the transient period and/or in the high density case) is rather nontrivial. The description of the dynamics in terms of cellular automata makes it possible (using a rather complicated combinatorial techniques) to derive an asymptotic description in the limit of large $`N`$ . Another approach based on the ultradiscrete limit of the Burgers equation was proposed in , where the dependence $`\mathrm{V}(\rho )`$ was proven for a lattice of size $`N`$ but with the estimate of the transient period $`N/2`$, which rules out the generalization for the case of the infinite lattice and nonperiodic initial configurations. In this paper we consider this model from a bit more general point of view as a discrete time dynamical system (map $`T`$) acting in the space of all possible configurations $`𝐗`$ – collections of zeros and ones (describing the positions of particles). We derive the variational approach based on the observation that the average velocity of any configuration does not decrease in time (Proposition 2.4). Simultaneously to the dynamics of particles one can study the dynamics of empty places. Observe that each of these dynamics determines the other one in the unique way. To make use of this observation we introduce a dual map $`T^{}`$ corresponding to the dynamics of empty places. By means of these two basic ideas (variational approach and dual maps techniques) we first prove the formula for the dependence of the average velocity on the density for any (may be small) finite lattice lengths $`N`$ and any (not necessary typical) initial configurations. We find also that steady state configurations demonstrate certain periodic in time patterns, whose features are described in Theorem 2.1 as well as the convergence to the steady state and the duration of the transient period. Qualitatively our main result about this model is that the following alternative takes place: either the flow consists of only free particles (i.e. there are no clusters of particles), or there are no clusters of empty places.
The paper is organized as follows. In Section 2 we study in detail the above formulated simplest model, which we call 1D periodic model with slow particles. In Section 3 we consider the same model but on the unbounded lattice, which significantly changes the dynamics and cannot be obtained in the limit of large $`N`$ from the previous one. To demonstrate the power of our dual maps techniques we introduce in Section 4 the model of the one-row motor road with speedy cars. The latter means that instead of the moving by at most 1 position, a particle moves forward until the next occupied position. We study both periodic and unbounded 1D cases. From the mathematical point of view the main difference of this model from the previous one is that the dynamics is not local, which makes it impossible to describe this model in terms of cellular automata. In Section 5 we generalize the traffic models with slow particles for a more practically interesting case of a multi-row motor road and study its properties. Finally in Section 6 we discuss a model of a passive tracer in the flow generated by the traffic model with slow particles confirming our practical observation above.
It is worth mention that to the best of our knowledge only the simplest 1D periodic traffic model with slow particles (among the models considered in the paper) was discussed in the literature previously. For some missing definitions related to dynamical systems theory (especially for systems acting on discrete phase spaces) we refer the reader to books on ergodic theory of dynamical systems (see, for example, ).
## 2 1D traffic model with slow particles on the finite lattice
Let $`𝐗=\{0,1\}^N`$ be the set of all possible configurations – collections $`X`$ of $`N`$ elements from the alphabet of 2 letters $`0`$ and $`1`$. We consider a map $`T:𝐗𝐗`$ defined as follows:
$$TX(x):=\{\begin{array}{cc}1\hfill & \text{if }X(x)=0\text{ and }X(x1)=1\hfill \\ 1\hfill & \text{if }X(x)=1\text{ and }X(x+1)=1\hfill \\ 0\hfill & \text{otherwise}.\hfill \end{array}$$
(2.1)
We assume here periodic boundary conditions, i.e. $`N+11`$ and $`0N`$. Observe that this map is not one-to-one and thus the backward (in time) dynamics cannot be defined in a unique way. See examples of the dynamics under the action of the map $`T`$ on Fig. 2.
We shall say that there is a particle at a position $`x`$ on the lattice if $`X(x)=1`$ and that this position is empty otherwise. A particle at a position $`x`$ is called free if $`X(x+1)=0`$. A group (more than 1) of consecutive particles (empty places) we call a cluster of particles (empty places). Observe that the number of particles is preserved under dynamics. Under the action of the map $`T`$ on the next time step each particle will either go forward by 1 place if this place is empty (occupied by 0) or will remain on the same place otherwise. Introducing the notion of a local velocity:
$$v(X,x):=\{\begin{array}{cc}1\hfill & \text{if }X(x)=1\text{ and }X(x+1)=0\hfill \\ 0\hfill & \text{otherwise},\hfill \end{array}$$
we define the average (in space) velocity of (particles in) the configuration $`X`$ as
$$\mathrm{V}(X):=\frac{1}{m(X)}\underset{x}{}v(X,x),$$
where $`m(X)_xX(x)`$ is the total number of particles in the configuration. Observe that $`\mathrm{V}(X)`$ is equal to the number of free particles divided by the total number of particles in the configuration.
###### Theorem 2.1
For any $`N`$ and any initial configuration with $`mN`$ particles after at most $`\mathrm{min}(m,Nm)`$ iterations the configuration will become periodic (in time) with the period $`N`$ and the average velocity $`\mathrm{V}=\mathrm{min}(1,\frac{N}{m}1)`$.
###### Remark 2.1
The smallest period (in time) of the periodic configuration above might be smaller than $`N`$, indeed, for $`N=2m`$ the smallest period is equal to $`2`$, i.e. there exists $`X𝐗`$ such that $`T^2X=X`$.
The proof of the theorem consists of the following lemmas.
###### Lemma 2.2
The length of any cluster of particles cannot increase, and the number of free particles cannot decrease.
Proof. Fix a cluster of particles. On the next time step its first particle (with the largest $`x`$-coordinate) goes out of the cluster and at most one particle can join the cluster from behind. If the number of particles is 2 and no particle will join the cluster from behind on the next time step, only one particle will remain in it, and according to our definition the cluster disappears.
Consider now the number of free particles. On the next time step each cluster of particles loses the first its particle which becomes a free one and at most the same number of free particles can join clusters coming from behind. Therefore the number of free particles cannot decrease.
###### Corollary 2.3
Traffic jams cannot appear from nothing.
On the other hand, since the average velocity of a configuration is equal to the number of free particles in it and by Lemma 2.2 this number grows in time we come the the following variational principle.
###### Proposition 2.4
(Variational principle) The functional $`\mathrm{V}(X)`$ (average velocity) increases under the dynamics up to the moment when it takes its maximal possible value.
###### Lemma 2.5
Any configuration after a finite number of time steps gets into a periodic one.
Proof. The phase space $`𝐗`$ of the considered dynamical system $`(𝐗,T)`$ is finite and therefore any its trajectory $`T^tX`$ will begin repeating after a finite number of iterations. Thus any limit set of the map $`T`$ consists of periodic configurations.
The map $`T`$ describes the dynamics of particles. It turns out that often it is simpler to study the dynamics of empty places instead. To make use of this observation for a configuration $`X`$ we introduce a dual one $`X^{}(x):=1X(x)`$ for all $`x`$ and define a dual map $`T^{}`$ acting on the same space $`𝐗`$:
$$T^{}X(x):=\{\begin{array}{cc}0\hfill & \text{if }X(x)=1\text{ and }X(x+1)=0\hfill \\ 0\hfill & \text{if }X(x)=0\text{ and }X(x1)=1\hfill \\ 1\hfill & \text{otherwise},\hfill \end{array}$$
(2.2)
assuming again periodic boundary conditions. One can easily see that the map $`T^{}`$ describes the motion of empty places under the action of the map $`T`$, which in this case satisfies the same rules as the the motion of particles except it goes in the backward direction in space. A direct computation gives the following representation, which can be considered as a definition of the dual map in the general case (not only for this specific model).
###### Lemma 2.6
$`TX=(T^{}X^{})^{}`$.
###### Corollary 2.7
All results for the map $`T`$ hold true also for the dual map and vice versa.
Proof of Theorem 2.1. Assume first that the density of particles $`\rho (X):=\frac{m}{N}`$ in the initial configuration $`X`$ is less or equal to $`1/2`$. If all the particles in this configuration are free then each particle moves with the velocity $`1`$ and the trajectory starting from this configuration is periodic (in time) with (may be not minimal) period $`N`$. Indeed, after $`N`$ iterations each particle will return to its initial position.
If there are clusters of particles in $`X`$ we need to show that after at most $`m`$ iterations all particles will become free. Indeed, by Lemma 2.2 the number of free particles does not decrease. Denote by $`\stackrel{~}{m}`$ the length of the largest cluster of empty places. Let us prove by induction on $`m`$ that for any pair of positive integers $`m,n`$ such that $`mn/2`$ and for any initial configuration after at most $`t_c:=\mathrm{min}(m,\stackrel{~}{m})`$ iterations all particles will become free. This statement is trivial for $`m=1`$ and $`n2`$. Assuming that it holds for some $`m`$ let us prove it for $`m+1`$. Let $`X`$ be a configuration with $`m+1`$ particle on the lattice of size $`n2(m+1)`$ and let $`x_0`$ be the position of the first particle after (one of) the largest clusters of empty spaces of length $`\stackrel{~}{m}`$. Fix this particle and consider the dynamics of others. By the induction hypothesis during the first $`\mathrm{min}(m,\stackrel{~}{m})`$ iterations other particles do not collide with the chosen one. Therefore their dynamics is the same as if there are only $`m`$ particles. Again by the induction hypothesis all these particles will become free after $`\mathrm{min}(m,\stackrel{~}{m})`$ iterations and at that moment either the chosen particle is free also, or it forms a cluster of size 2. Therefore after the next iterations it becomes free as well.
It remains to analyze initial configurations with the density of particles greater than $`1/2`$. This situation is much more complex, because any configuration of this type contains clusters, which are interchanging particles between themselves and never disappear completely. To overcome this difficulty we consider a dual map $`T^{}`$. The density of empty places is equal to $`\frac{Nm}{N}<1/2`$. Therefore by Corollary 2.7 and the first part of the proof under the action of the map $`T^{}`$ a dual configuration $`X^{}`$ after at most $`t_0=Nm`$ iterations will get into a periodic configuration $`(T^{})^{t_0}X^{}`$ consisting of $`Nm`$ free particles with the period (in time) $`N`$ for the map $`T^{}`$. Clearly for the dual to it configuration we have the following identity $`((T^{})^{t_0}X^{})^{}=T^{t_0}X`$ and the latter configuration is $`T`$-periodic with the same period. To finish the proof we calculate the average velocity:
$$\mathrm{V}(T^{t_0}X)=\frac{Nm}{m}=\frac{N}{m}1.$$
It is worth notice that our estimate of the duration of the transient period $`\mathrm{min}(m,Nm)`$ is exact: if $`m<Nm`$ and the initial configuration consists of only one cluster of length $`m`$ the duration of the transient period is equal to $`m`$.
###### Corollary 2.8
Qualitatively our main result about this model is that the following alternative takes place: either the flow consists of only free particles (i.e. there are no clusters of particles), or there are no clusters of empty places.
On the other hand, there is no a priori information about the distribution of lengths of clusters of particles because this distribution depends on the initial configuration and can differ between various periodic limiting configurations.
## 3 1D traffic model with slow particles on the unbounded lattice
Consider now the unbounded one-dimensional case, i.e. $`𝐗:=\{0,1\}^{\text{ZZ}^1}`$ and the map $`T`$ is defined by the formula (2.1) in the same way as in the previous section, except for the periodic boundary conditions. For a configuration $`X𝐗`$ we define the notion of a subconfiguration $`X_k^n:=\{X(k),X(k+1),\mathrm{},X(n)\}`$, i.e a collection of entries of $`X`$ between the pair of given indices $`k<n`$, and introduce the corresponding density and the average velocity:
$$\rho (X_k^n):=\frac{m(X_k^n)}{nk+1},\mathrm{V}(X_k^n):=\frac{1}{m(X_k^{n1})}\underset{x=k}{\overset{n1}{}}v(X,x),$$
where $`m(X_k^n)`$ stays for the number of particles in the subconfiguration $`X_k^n`$.
By the density and the average velocity (of particles) of a entire configuration $`X𝐗`$ we mean the following limits (if they are well defined):
$$\rho (X):=\underset{n\mathrm{}}{lim}\rho (X_n^n),\mathrm{V}(X):=\underset{n\mathrm{}}{lim}\mathrm{V}(X_n^n),$$
otherwise one can consider the corresponding upper and lower limits, which we denote by $`\rho _\pm (X)`$ and $`\mathrm{V}_\pm (X)`$.
Notice that in distinction to the finite case these quantities make sense not for all possible configurations. We shall say that a configuration $`X`$ satisfies the regularity assumption (or simply regular) if there exists a number $`\rho `$ and a monotonous one-to-one function $`\phi (n)0`$ as $`n\mathrm{}`$, such that for any $`n\text{ZZ}^1`$, $`N\text{ZZ}_+^1`$ and any subconfiguration $`X_{n+1}^{n+N}`$ of length $`N`$ the number of particles in this subconfiguration $`m(X_{n+1}^{n+N})`$ satisfies the inequality
$$\left|\frac{m(X_{n+1}^{n+N})}{N}\rho \right|\phi (N).$$
(3.1)
It is clear that at least for a configuration $`X`$ satisfying the regularity assumption the density $`\rho (X)`$ is well defined and is equal to the value $`\rho `$ in the formulation of the assumption.
###### Theorem 3.1
Let the initial configuration $`X`$ satisfies the regularity assumption with $`\rho 1/2`$. Then after a finite number of iterations the average velocity of particles becomes equal to $`\mathrm{min}(1,\frac{1}{\rho }1)`$.
Proof. Notice that Lemma 2.2 holds true in this case also. Moreover it can be applied for the dual map as well and thus the length of any cluster of empty places cannot decrease. This shows that the variational principle (Proposition 2.4) holds in this case also. Moreover it can be reformulated to take care about configurations for which neither the density, nor the average velocity are well defined.
###### Proposition 3.1
(Variational principle) For any configuration $`X`$ any pair of its particle denote by $`x^{}(t)<x^{\prime \prime }(t)`$ their positions at the moment $`t`$. Then the average velocity $`\mathrm{V}(X_{x^{}(t)}^{x^{\prime \prime }(t)})`$ of the subconfiguration $`X_{x^{}(t)}^{x^{\prime \prime }(t)}`$ increases monotonically with $`t`$ up to the moment (may be infinite) when it takes its maximal possible value.
The basic technical step of the proof of Theorem 3.1 is given by the following statement.
###### Lemma 3.2
If $`X`$ satisfies the regularity assumption, then the same holds true for $`T^tX`$ for any $`t>0`$.
Proof. Clearly it is enough to prove this statement for $`t=1`$. Assume that this is not true and for some $`N>\phi ^1(\rho )`$ there is a subconfiguration $`(TX)_{n+1}^{n+N}`$ of length $`N`$ in the configuration $`TX`$ such that
$$m((TX)_{n+1}^{n+N})<(\rho \phi (N))N.$$
This can happen only if the following equalities are satisfied simultaneously
$$m(X_{n+1}^{n+N})=(\rho \phi (N))N,X(n)+X(n+N+1)=0,X(n+N)=1,$$
i.e. on the next time step no particle from behind will come to this interval and there is a free particle in the last position of the considered interval, which leaves it. On the other hand, from these equalities we immediately deduce that
$$m(X_n^{n+N1})=(\rho \phi (N))N1,$$
which contradicts the regularity assumption. In the same way one can prove that the inequality
$$m(X_n^{n+N1})(\rho +\phi (N))N$$
cannot break as well.
Now we can return to the proof of Theorem 3.1. Assume first that $`\rho <1/2`$. Denoting by $`[]`$ the integer part of a number, we get that for any $`NN_c:=[\phi ^1(\frac{1}{2}\rho )]`$ the number of particles in any subconfiguration $`(T^tX)_{n+1}^{n+N}`$ can be estimated as
$$m((T^tX)_{n+1}^{n+N})(\rho +\phi (N))N(\rho +\phi (N_c))NN/2.$$
Applying the same machinery as in the previous section one can show that after at most
$$t_c:=(\rho +\phi (N_c))N_c\frac{N_c}{2}\frac{1}{2}\phi ^1(\frac{1}{2}\rho )$$
iterations all particles will become free ones, which implies the unit average velocity.
In the remaining case $`\rho >1/2`$ we follow the same idea as in the proof of Theorem 2.1 and pass to the dual map and the dual configuration. Observe that if a configuration $`X`$ satisfies the regular assumption with the density $`\rho `$ and the rate function $`\phi `$, then the dual configuration $`X^{}`$ also satisfies it with the density $`1\rho `$ and the same rate function $`\phi `$. Indeed $`m((X^{})_{n+1}^{n+N})=Nm(X_{n+1}^{n+N})`$ for any $`n,N`$ and thus
$$|\frac{m((X^{})_{n+1}^{n+N})}{N}(1\rho )|=|\frac{m(X_{n+1}^{n+N})}{N}\rho )|\phi (N).$$
On the other hand, the density of the dual configuration $`1\rho 1/2`$ which, having in mind that the only difference between the maps $`T`$ and $`T^{}`$ is the direction of motion, gives us the possibility to apply the first part of the proof.
Observe that in the proof of Theorem 3.1 we actually derived an estimate of the length of the transient period as $`t_c:=(\rho +\phi (N_c))N_c`$ which goes to infinity as $`\rho 1/2`$. This is the reason why Theorem 3.1 does not cover the boundary case $`\rho =1/2`$, which we discuss below.
###### Theorem 3.2
Let the initial configuration $`X`$ satisfies the regularity assumption with $`\rho =1/2`$ and let $`x^{}(t)<x^{\prime \prime }(t)`$ be positions of two fixed particles at the arbitrary moment $`t`$. Then the average velocity of the subconfiguration $`X_{x^{}(t)}^{x^{\prime \prime }(t)}`$ converges to $`1`$ as $`t\mathrm{}`$.
Proof. Choose an integer $`\widehat{N}`$ and consider a configuration $`\widehat{X}`$ obtained from the configuration $`X`$ by the following operation: for each $`k`$ we remove from the configuration $`X`$ the closest from behind particle to the position $`k\widehat{N}`$.
$$\left|\frac{m(\widehat{X}_{n+1}^{n+N})}{N}(\rho \frac{1}{\widehat{N}})\right|\left|\frac{m(X_{n+1}^{n+N})}{N}\rho \right|+\left|\frac{m(\widehat{X}_{n+1}^{n+N})}{N}\frac{m(X_{n+1}^{n+N})}{N}+\frac{1}{\widehat{N}}\right|\phi (N)+\frac{1}{\widehat{N}}.$$
Thus the configuration $`\widehat{X}`$ is also regular but with the density $`\frac{1}{2}\frac{1}{\widehat{N}}<\frac{1}{2}`$. Thus according to Theorem 3.1 after a finite number of iterations $`t_c`$ the average velocity of the configuration $`T^{t_c}\widehat{X}`$ becomes equal to $`1`$.
Making an opposite operation, namely inserting a particle to the configuration $`X`$ to the closest from behind to $`k\widehat{N}`$ empty position for each $`k`$, we obtain a regular configuration with the density $`\frac{1}{2}+\frac{1}{\widehat{N}}>1/2`$. Again by Theorem 3.1 after a finite number of iterations the average velocity of this configuration becomes equal to
$$\frac{1}{\frac{1}{2}\frac{1}{\widehat{N}}}1=1+\frac{4}{\widehat{N}2}1\mathrm{as}\widehat{N}\mathrm{}.$$
Thus both (arbitrary close as $`\widehat{N}\mathrm{}`$) approximations to the configuration $`X`$ have after a finite number of iterations (depending on $`\widehat{N}`$) the average velocity deviating from $`1`$ by $`O(\widehat{N})`$. It remains to show that the average velocity of a subconfiguration of the configuration $`X`$ can be estimated from above and from below by those from above approximations. Let $`X`$ and $`Y`$ are two configurations such that $`X(x)Y(x)`$ for all $`x`$ and let $`x^{}(t)<x^{\prime \prime }(t)`$ be positions of two fixed particles in the configuration $`X`$ at the arbitrary moment $`t`$. Denote by $`y^{}(t)<y^{\prime \prime }(t)`$ positions of the same particles in the configuration $`Y`$. Then
$$\mathrm{V}(X_{x^{}(t)}^{x^{\prime \prime }(t)})\mathrm{V}(X_{x^{}(t)}^{x^{\prime \prime }(t)})$$
for any moment of time $`t`$. Indeed, additional particles in the configuration $`Y`$ present only obstacles to the motion of other particles, thus making the average velocity slower (or at least not faster).
## 4 1D traffic model with speedy particles
The dual map $`T^{}`$ that we made use of in the previous sections was mirror symmetric with respect to the map $`T`$. In the model considered in this section we show that the dual map might have a different structure as well.
Let us start with the one-dimensional finite periodic case, but with speedy particles instead of particles moving with the velocity $`1`$ as above. This means that if in front of a particle there are exactly $`n`$ empty places it moves by $`n`$ places forward and remains in the same place if the next place is occupied. By the average velocity of (particles in) the configuration $`X`$ we mean the total distance covered by the particles from $`X`$ during the next time step divided by the number of particles. The corresponding map of the space $`𝐗=\{0,1\}^N`$ into itself we again denote by $`T`$. See examples of the dynamics under the action of the map $`T`$ on Fig. 3.
###### Theorem 4.1
The dual map in this case satisfies the relation $`T^{}X^{}=TX^{}`$, and for any nontrivial (having both zeros and ones) initial configuration $`X𝐗`$ the average velocity does not depend on time and is equal to $`\frac{N}{m(X)}1`$, where $`m(X)`$ is the number of particles in the configuration $`X`$.
Proof. We consider this result as an illustration of the dual maps techniques and and a base for the introduction of the passive tracer model considered in Section 6. Therefore we do not give explicit representations for the maps $`T,T^{}`$ and present only a sketch of the proof in this case.
The relation between the map $`T`$ and the dual one $`T^{}`$ can be checked by a direct computation (which we leave for the reader). Notice that unlike the dual map in Section 2 the dynamics of empty places is exactly the same as the dynamics of particles and occurs in the same direction in space.
According to the definition each particle on the next time step moves forward by the number of empty places in front of it. Therefore the total number of places to which the particles in the configuration $`X`$ move forward is equal to the number of empty places. Thus the average velocity is equal to $`\mathrm{V}(X)=\frac{Nm(X)}{m(X)}=\frac{N}{m(X)}1`$.
Using the same argument as in the proof of Theorem 2.1 we deduce that any trajectory of the map $`T`$ gets periodic after a finite number of iterations and a trivial observation that the number of empty places is invariant under dynamics finishes the proof.
Consider now the model of speedy particles on the unbounded 1D lattice $`\text{ZZ}^1`$ in analogy to the model in Section 3. The map $`T`$ is defined as above and acts in the space $`𝐗=\{0,1\}^{\text{ZZ}^1}`$. We shall say that a configuration $`X𝐗`$ satisfies the Law of Large Numbers if for any integer $`n`$ the limit
$$\underset{N\mathrm{}}{lim}\frac{m(X_{n+1}^{n+N})}{N}$$
is well defined and does not depend on $`n`$. Clearly the above limit (if it exists) coincides with the density of particles $`\rho (X)`$ in the configuration $`X`$.
The dual map in the considered case satisfies the same relation $`T^{}X^{}=TX^{}`$ as in Theorem 4.1 and the dynamics is described by the following statement.
###### Theorem 4.2
Let an initial configuration $`X𝐗`$ satisfy the Law of Large Numbers and let $`\rho (X)(1\rho (X))0`$. Then the average velocity does not depend on time and is equal to $`\frac{1}{\rho (X)}1`$.
Proof. Let $`x^{}<x^{\prime \prime }`$ be positions of two particles in the configuration $`X`$. Then the average velocity of the subconfiguration $`X_x^{}^{x^{\prime \prime }}`$ is equal to the number of empty places in this subconfiguration, divided by the number of particles in it, i.e.
$$\mathrm{V}(X_x^{}^{x^{\prime \prime }})=\frac{x^{\prime \prime }x^{}m(X_x^{}^{x^{\prime \prime }})}{m(X_x^{}^{x^{\prime \prime }})}=\frac{x^{\prime \prime }x^{}}{m(X_x^{}^{x^{\prime \prime }})}.$$
Now taking into account that this relation holds for arbitrary positions of particles, we immediately get the desired statement.
###### Remark 4.1
It is of interest that the average velocity in both discussed speedy particles models coincides with the average velocity in the models with slow particles in the case of high density (i.e. when $`\rho (X)>1/2`$). The explanation is that in the models with slow particles with the density $`\rho (X)>1/2`$ the typical distance between particles (in the steady state) is $`0`$ or $`1`$, which makes no difference between the dynamics of speedy and slow particles.
Observe that in the models with speedy particles the dynamics is richer than in the models with slow particles, for example the statement of Lemma 2.2 does not hold, i.e. clusters of particles can grow and traffic jams become typical even for the case of low density of particles.
## 5 2D traffic model
From the point of view of real traffic problems a clear shortage of all models considered above is the absence of the possibility to go around particles staying in a traffic jam. Indeed in a one-row motor road (which was the main considered example) this is not possible. To overcome this restriction we consider a model describing the motion of slow particles on a multi-row motor road.
A lattice in our model is a $`N\times K`$-strip, describing a one-way cyclic road of length $`N`$ consisting of $`K`$ parallel rows. We denote the first coordinate corresponding to the spread of the road by $`x\{1,\mathrm{},N\}`$, and the second one (the row number) by $`y\{1,\mathrm{},K\}`$, and assume periodic boundary conditions on $`x`$. In terms of particles the dynamics is defined as follows. If there is a particle at a position $`(x,y)`$ (i.e. $`X(x,y)=1`$) then this particle moves forward by one place to the position $`(x+1,y)`$ if this place is not occupied, else moves to the left to the position $`(x,y+1)`$ if this place and the place before it are empty, else moves to the right to the position $`(x,y1)`$ if this place and the place before and the next place to the right are not occupied, and remains on its place otherwise. To be consistent we assume that the $`0`$-th, $`(K+1)`$-th and $`(K+2)`$-th (virtual) rows are completely occupied. Recall that the space is $`N`$-periodic on the $`x`$-coordinate.
Observe that this model is nothing more than the simplest formulation of the standard traffic rules.
The space of all possible configurations of this model is $`𝐗=\{0,1\}^{NK}`$ and the corresponding map describing the dynamics of the configurations in this space can be written as follows:
$$TX(x,y):=\{\begin{array}{cc}1\hfill & \text{if }X(x,y)=0\text{ and }X(x1,y)=1\hfill \\ 1\hfill & \text{if }X(x,y)=1\text{ and }X(x+1,y)=1\hfill \\ 1\hfill & \text{if }X(x,y)+X(x1,y)=0\hfill \\ & \text{ and }X(x,y+1)X(x+1,y+1)=1\hfill \\ 1\hfill & \text{if }X(x,y)+X(x1,y)+X(x,y+1)=0\hfill \\ & \text{ and }X(x,y1)X(x+1,y+1)=1\hfill \\ 0\hfill & \text{otherwise}.\hfill \end{array}$$
(5.1)
The dual map for this model can be easily defined and is described by the following statement. Again to be consistent we need to add an additional $`(1)`$-th completely occupied virtual row to the lattice.
###### Lemma 5.1
The dual map $`T^{}`$ satisfies the general statement of Lemma 2.6 and the explicite formula for $`T^{}X(x,y)`$ can be obtained from the relation (5.1) by changing everywhere $`x+1`$ by $`x1`$, $`x1`$ by $`x+1`$, and $`y+1`$ by $`y1`$.
This statement can be checked by a direct computation.
###### Corollary 5.2
The dynamics of empty places under this model is the same as the dynamics of particles except it occurs in the opposite direction in space along the $`x`$-coordinate, i.e. the correspondence between the dynamics is the mirror symmetry.
We are interested in the average velocity along the $`x`$-coordinate, which we define in the same way as in the previous sections. From the first sight it seems that the presence of additional possibilities to go to the right or to the left (if the way forward is blocked) should improve the traffic. The following statement shows that this is not the case and moreover in the 2D case the average velocity might be much slower than the velocity predicted by the 1D model with slow particles.
###### Theorem 5.1
For any initial configuration $`X`$ in the case of the periodic bounded lattice and any regular initial configuration $`X`$ with $`\rho 1/2`$ in the case of the unbounded (on the $`x`$-coordinate) lattice after a sufficiently large number of iterations the upper limit for the average velocity satisfies the inequality
$$\mathrm{V}_+(T^tX)\mathrm{min}(1,\frac{1}{\rho (X)}1),$$
while the lower limit can be arbitrary close to $`0`$ even in the case of low density of particles ($`\rho <1/2`$).
Notice that in the bounded case exactly as in the 1D model on the bounded lattice any initial configuration becomes periodic after a finite number of iterations. Moreover in this case both limits $`\mathrm{V}_\pm `$ coincide.
Proof. We start with the case $`\rho 1/2`$ and the periodic (on the $`x`$-coordinate) bounded lattice of size $`NK`$. Our first aim is to show that for any number $`\rho [0,1/KN,\mathrm{},1]`$ a configuration of the density $`\rho `$ cannot have the average velocity larger than the one predicted by the $`1D`$ model with slow particles. If $`\rho N[N/2]`$, then one can construct a configuration $`X_+`$ consisting of exactly $`\rho N`$ free particles on each of $`K`$ rows. Clearly the configuration $`X_+`$ is $`N`$-periodic and has the maximal possible average velocity $`1`$. Otherwise at least in one of the rows there is a cluster of particles and hence the average velocity is strictly smaller than $`1`$.
To get the lower bound we construct the following example. Assuming that $`N`$ is even, we choose a number $`k\{K/3,\mathrm{},K/2\}`$. Consider the configuration $`X`$ such that $`X(x,l)=1`$ for all $`x`$ and $`l=\{1,\mathrm{},k\}`$, while in the $`(k+1)`$-th row we assume that $`X(l,k+1)=0`$ for odd $`l`$ and is equal to $`1`$ otherwise. Then independently on the positions of other particles, neither particle from the first $`k`$ rows can move to the rows with numbers larger than $`k`$, no particles from other rows can get into the first $`k`$ rows. On the other hand, this configuration $`X_k`$ is time periodic with the period 2 and its average velocity is equal to
$$\mathrm{V}(X_k)=\frac{N/2}{kN+N/2}=\frac{1}{2k+1}\frac{1}{2K/3+1}.$$
Therefore $`\mathrm{V}(X_k)0`$ as $`K\mathrm{}`$.
Moreover choosing $`k=1`$ in the above example we get the average velocity $`\frac{N/2}{N+N/2}=1/3`$ for the configuration of density $`\rho =\frac{N+N/2}{NK}=\frac{3}{2K}0`$ as $`K\mathrm{}`$. Thus for an arbitrary low density we may have the average velocity less than $`1`$.
If the density $`\rho >1/2`$ we use the same trick as usual and pass to the dual map and dual configuration, which permits to use the first part of the proof to get the needed estimates.
The proof in the unbounded case is basically the same and we leave it to the reader.
There are other possible ways to model multi-row traffic flows, for example one might consider different rules for the row changing, the presence of high velocity particles, periodic boundary conditions between the boundary rows, etc. However if the model respects the law that on the next step a particle can follow its way along the same row (assuming that the next place is not occupied) independently on particles on neighboring rows, the example in the above proof demonstrating the phenomenon of the low average velocity under a low density of particles remains generic.
## 6 Passive tracer in the traffic flow
In the Introduction we have mentioned the practical observation that it is beneficial in some cases to go against the flow than in the same direction as it goes. To study this phenomenon we consider in this section a simple model of a passive tracer imitating the behavior of a person moving in a hurry in the traffic flow. As usual we assume that the tracer have its own (forward or backward) chosen direction of motion and does not make any impact to the flow.
Let $`X(t)`$ describes the 1D flow of particles and let at time $`t`$ the passive tracer occupies the position $`x`$. Then before the next time step of the model of the flow the tracer moves in its chosen direction to the closest (in this direction) position of a particle of the configuration $`X(t)`$. For example, if the going forward tracer occupies the position 2 and the closest particle in this direction occupies the position 5, then the tracer moves to the position 5. Then the next iteration of the flow occurs, the tracer moves to its new position, etc.
To be precise for a fixed configuration $`X𝐗`$ with $`m(X)>1`$ we introduce two maps $`\tau _X^+`$ and $`\tau _X^{}`$ of the ordered periodic lattice $`𝐋`$ into itself defined as follows:
$$\tau _X^+y:=\mathrm{min}(x𝐋:𝐲<𝐱,𝐗(𝐱)=\mathrm{𝟏}),$$
$$\tau _X^{}y:=\mathrm{max}(x𝐋:𝐲>𝐱,𝐗(𝐱)=\mathrm{𝟏})$$
for any $`y𝐋`$, where the order relation $`y<x`$ is induced by the order on the lattice $`𝐋`$. Then the simultaneous dynamics of the configuration of particles (describing the flow) and the tracer is defined by the skew product of two maps – the map $`T`$ and one of the maps $`\tau _{}^\pm `$, i.e.
$$(X,y)𝒯_\pm (X,y):=(TX,\tau _X^\pm y),$$
acting on the extended phase space $`\overline{𝐗}:=𝐗\times 𝐋`$.
Let $`S(t)`$ denotes the total distance covered by the tracer up to the moment $`t`$ with the positive sign if the tracer moves forward, and the negative sign otherwise. Then we define the average (in time) velocity of the tracer $`v(t)`$ as $`S(t)/t`$.
###### Theorem 6.1
For an arbitrary configuration $`X`$ with $`m(X)>1`$ of the 1D model with slow particles on the finite periodic lattice and for an arbitrary initial configuration $`X`$ satisfying the regularity assumption (3.1) with the density $`\rho \{0,\frac{1}{2},1\}`$ in the 1D unbounded case the average velocity of the passive tracer converges to $`1`$ if the tracer moves forward (along the flow)) and $`\rho 1/2`$, and to $`\mathrm{max}(1,\frac{1}{\rho }1)`$ if it moves backward (against the flow).
Proof. We start with the finite periodic lattice and the tracer moving in the forward direction. According to our definition after the first iteration of the map $`𝒯`$ the tracer will occupy the same position as some particle. Therefore it is enough to consider only extended configurations $`(X,y)`$, such that $`X(y)=1`$. Since $`\rho 1/2`$, then after a finite number of iterations the flow will consist of only free particles. Therefore the tracer will run down one of them and will follow it, but cannot outstrip. Indeed after each iteration of the flow this free particle occurs just one position behind the tracer.
We do not discuss the case $`\rho >1/2`$ because the average velocity of the tracer in this case sensitively depends on the choice of the initial configuration. For example in the case $`n=8`$ and $`m=6`$ (i.e. $`\rho =3/4`$) consider two initial configurations $`01011111`$ and $`01111011`$. In the first case the tracer after several iterations obtains the constant velocity $`1`$, while in the second case the velocity of the tracer is periodic with the period $`3`$ and consists of the repeating groups of $`(2+2+1)`$, i.e. the average speed is equal to $`5/3`$.
Consider now the case when the tracer is moving backward. Then each time the tracer encounters a particle, on the next time step this particle moves in the opposite direction and does not disturb the movement of the tracer until the collision with the next particle. Assume that the density of particles is less than $`1/2`$. Then by Theorem 2.1 after a finite number of iterations only free particles are present in the flow, which results in the convergence of the average velocity of the passive tracer to $`(\frac{1}{\rho }1)`$. Indeed on the spread of length $`N`$ there are $`m`$ particles, i.e. $`m`$ obstacles for the tracer, which gives the average velocity $`\frac{Nm}{m}`$. On the other hand, if the density of particles is greater or equal to $`1/2`$ again by Theorem 2.1 there are no clusters of empty places in the flow and thus after each iteration the tracer can move only by one position, which finishes the proof for the model of the flow with periodic boundary conditions.
The proof in the unbounded case is practically the same with the only difference that one should use Theorem 3.1 instead of Theorem 2.1. Notice that additionally to the regularity assumption we need to assume that $`\rho \{0,1\}`$ to be consistent with the definition of the dynamics of the passive tracer.
Observe that the motion against the flow is efficient only in the case of low density of particles. Certainly, this model is oversimplified and probably its predictions are unrealistic for the case of very high density of particles. However we believe that while the density is not low and not very high our description of the passive tracer is reasonable at least on the qualitative level.
## 7 Conclusion
This paper represents one of the first steps in the mathematical foundation of the analysis of traffic flows and we restrict ourselves here to the pure deterministic settings. The next step should describe ergodic (statistical) properties of the considered models with initial conditions chosen at random and random versions of these models as well. This circle of questions is especially interesting in the case of models on infinite lattices, where the dynamics of typical configurations cannot be obtained in the limit of infinitively large lattice sizes from our results about regular configurations even in the deterministic setting. Indeed, for a reasonable choice of the class of random initial conditions their realizations do not satisfy our regular assumption with the probability $`1`$. |
warning/0003/cond-mat0003398.html | ar5iv | text | # Density of States of Disordered Two-Dimensional Crystals with Half-Filled Band
## Abstract
A diagrammatic method is applied to study the effects of commensurability in two-dimensional disordered crystalline metals by using the particle-hole symmetry with respect to the $`nesting`$ vector $`𝑷_0=\{\pm \frac{\pi }{a},\frac{\pi }{a}\}`$ for a half-filled electronic band. The density of electronic states (DoS) is shown to have nontrivial quantum corrections due to both $`nesting`$ and elastic impurity scattering processes, as a result the van Hove singularity is preserved in the center of the band. However, the energy dependence of the DoS is strongly changed. A small offset from the middle of the band gives rise to disappearence of quantum corrections to the DoS .
73.20.Fz; 73.50.-h; 73.20.Dx; 73.50.Bk; 73.20.Jc
Great advance has been made in the theory of disordered metals after the pioneering work of Abrahams et al.. According to this paper all electronic states in one- and two-dimensional (1d and 2d) disordered systems are localized irrespective of the degree of randomness,. Electron-electron correlations in disordered metals have been shown to result in nontrivial corrections to the density of electronic states,(DoS), and conductivity,. The corrections to the conductivity are similar to the localization corrections obtained for a noninteracting electron gas, and quantum corrections to the DoS were shown to reduce it near the Fermi level.
Disordered metals in all abovementioned papers are modeled as a free electron gas moving in the random field of rigid impurities. However, at low concentrations of impurities the crystal usually has a periodical structure and the impurity atoms in most cases substitute the host atoms of the lattice. Then, the effects of commensurability of the electron wavelength, $`\lambda `$ , and the lattice constant, $`a`$, become essential in the scattering processes. The commensurability is known to exist at all ’rational’ points of the electron band,however it appears to be important for a half-filled band.
In this Letter we present our study concerning the effects of weak disorder on the electronic DoS of a 2d crystalline metal.The Hamiltonian of the model can be written as $`\widehat{H}=\widehat{H_o}+V(𝐫)`$, where $`\widehat{H_o}`$ is the Hamiltonian of noninteracting electrons in the perfect square lattice with nearest- neighbor hopping and $`V(𝐫)=_iU(𝐫𝐑_𝐢)`$ is the impurity potential with $`𝐑_𝐢`$ being the positional vector of an impurity randomly located on the $`ith`$ lattice site.
The one-particle DoS of the regular lattice can be expressed as $`\rho _0^{(d)}=\frac{2}{(2\pi \mathrm{})^d}\frac{d𝑺}{|ϵ(𝒌)|}`$ , where $`d𝑺`$ is the element of an isoenergetical surface in d-dimensional space. It can be shown that $`\rho _0^{(d)}(ϵ)`$ has a van Hove singularity at the points where the group velocity of the electron wave packet $`𝑽_𝒌=ϵ(𝒌)`$ vanishes,. For a three-dimensional (3d) regular lattice $`\rho _0^{(3)}(ϵ)`$ has integrable singularities.
For a pure 2d lattice with nearest-neighbor hopping the van Hove singularity has logarithmic character. The simplest electron spectrum for a 2d square lattice can be written in the tight-binding approximation as
$$ϵ(𝒌)=t[2\mathrm{cos}(k_xa)\mathrm{cos}(k_ya)];$$
(1)
where,$`k_{x,y}=\frac{2\pi }{aN_{x,y}}n_{x,y}`$ with $`\frac{N_{x,y}}{2}<n_{x,y}\frac{N_{x,y}}{2}`$ and only electron tunneling between nearest-neighboring sites with the tunneling integral $`t`$ is involved. The bandwidth is $`W=4t`$; and for a half-filled band case the Fermi energy becomes $`ϵ_F=2t`$. The DoS of a 2d square lattice with nearest-neighbor hopping is expressed by the elliptic integral of the first kind, an asymptotic expression of which has a logarithmic singularity in the middle of the energy band as:
$$\rho _0^{(2)}(ϵ)=\{\begin{array}{cc}\frac{1}{(\pi a)^2\sqrt{ϵ_F^2|\stackrel{~}{ϵ}|^2}}\mathrm{ln}(\frac{4t^2|\stackrel{~}{ϵ}|^2}{|\stackrel{~}{ϵ}|^2})\hfill & |\stackrel{~}{ϵ}|2t,\hfill \\ \frac{2}{(\pi a)^2|\stackrel{~}{ϵ}|},\hfill & \stackrel{~}{ϵ}\pm 2t\hfill \end{array}$$
(2)
where $`\stackrel{~}{ϵ}`$ is an electron energy measured from the Fermi level $`\stackrel{~}{ϵ}=ϵ2t`$.(Hereafter the tilde on $`ϵ`$ will be dropped).
The DoS of a noninteracting electron gas moving in the random field of impurities has no essential singularities near the Fermi surface. Inclusion of even short-range correlations in the 2d disordered metal gives rise to a decreasing DoS near the Fermi level,. The DoS of a $`2d`$ disordered crystal with substitutional impurities turns out to have a singularity near the middle of the band even for the noninteracting electron gas.
Notice that the commensurability effect for a $`1d`$ disordered crystal near the middle of the band has been studied by many authors,. Dyson first pointed out that the phonons’ DoS of a $`1d`$ disordered chain has a singularity as $`\rho ^{(1)}(ϵ)|ϵ|^1\mathrm{ln}^3|ϵ|`$ near the middle of the band. Later an analogous singularity has been found in the electronic DoS of many $`1d`$ models,. However, there exist a few computational studies of the DoS of a $`2d`$ disordered crystal,. By studying the averaged Green’s function for a disordered system with $`n`$ orbitals per site,the expansion coefficients in power of $`1/n`$ for $`d2`$ were shown in to diverge for energies approaching the band center. This fact was interpreted in as a existence of a van Hove singularity in the DoS.
The technical difficulties in analytical calculations of physical parameters are connected with the cosine energy spectrum and, furthermore, with an absence of perturbative parameters for half-filling. It is appropriate to notice that the recent attempt to explain both the linear resistivity and the nearly temperature independence of the thermopower in the cuprate superconductors is based on the existence of a van Hove singularity in the DoS of these materials. To study the problem analytically, the idealized model $`ϵ(𝒌)=k_xk_y`$ for the energy spectrum, which gives a logarithmic DoS at half-filling, was used. There exists also the $`nested`$ Fermi Liquid scattering approach to study the susceptibility of high-$`T_c`$ superconductors.
As it is known, the Fermi surface of an infinite $`2d`$ lattice with nearest-neighbor hopping is changed with band-filling and it is flat for the half-filling. In this case there exists a $`nesting`$ vector $`𝑷_o=\{\pm \frac{\pi }{a},\frac{\pi }{a}\}`$ that maps an entire section of the Fermi surface onto another, i.e. the Fermi surface is perfectly $`nested`$ for the half-filled band case. There, the following particle-hole symmetry of the electron dispersion with respect to the vector $`𝑷_o`$ for a half-filled band holds:
$$ϵ(𝒑+𝑷_o)ϵ_F=[ϵ(𝒑)ϵ_F]$$
(3)
The one-electron DoS can be calculated according to the following expression:
$$\rho (ϵ)=\frac{2}{\pi }\mathrm{Im}\frac{d^2p}{(2\pi )^2}G_R(𝒑,ϵ)$$
(4)
where $`G_R(𝒑,ϵ)`$ is the retarded Green’s function.
The new class of diagrams which gives an essential contribution to the DoS is drawn in Fig.1. Thin solid and dashed lines in Fig.1 correspond to the ’bare’ retarded Green’s functions $`G_{R}^{}{}_{}{}^{o}(𝒑,ϵ)`$ and $`G_{R}^{}{}_{}{}^{o}(𝒑+𝑷_o,ϵ)`$ .
Averaging over impurity realization is performed according to the crossed diagram technique in the Born approximation described in , i.e. the white-noise impurity potential is used which corresponds to the off- diagonal disorder in the problem. The main diagrams are selected according to the condition of $`ϵ_F\tau 1`$ or $`k_Fl1`$ with $`l=v_F\tau `$ and $`\tau `$ being an elastic scattering time, which means $`l\lambda `$.
The scattering processes in the problem are described by two relaxation times $`\tau _o`$ and $`\tau _\pi `$, (Fig.2). $`\tau _o`$ corresponds to the normal scattering process contribution to which gives the diagram shown in Fig.2a and $`\frac{1}{\tau _o}=\frac{C_{imp}}{(2\pi )^2}\frac{d𝑺}{|𝑽_𝒌|}|U(S)|^2`$,. Here, $`C_{imp}`$ and $`U`$ are the impurity concentration and potential, respectively. The diagrams which give a contribution to Green’s functions contain new impurity vertices, shown in Figs.2b,c. These vertices characterize simultaneous reflections of electrons on the Brillouin zone boundary in the process of scattering on an impurity. The following expression represents these vertices: $`\frac{1}{\tau _\pi }=\frac{C_{imp}}{(2\pi )^2}\frac{d𝐒}{|𝐕_𝐤|}|U(P_o,S)|^2`$. It can be seen from Figs.2b,c that the total momentum, generally speaking, is not conserved for these impurity vertices. Therefore they represent an $`Umklapp`$ process.
Bare retarded and advanced Green’s functions $`G_{R,A}^o(𝒑,ϵ)=[ϵ(ϵ(𝒑)ϵ_F)+\frac{i}{2\tau }\mathrm{sign}ϵ]^1`$ contain a total relaxation time, which is equal to $`\frac{1}{\tau }=\frac{1}{\tau _o}+\frac{1}{\tau _\pi }`$.
The maximally crossed diagrams are redrawn in Fig.1a so that the self-energy part contains the particle-particle propagator $`C_\pi (𝒌,ϵ)`$which we call $`\pi `$-cooperon, (Fig.3a). In contrast to the cooperon for an isotropic system $`C_\pi (𝒌,ϵ)`$ has a diffusion pole at a small total energy and a large total momentum $`𝐏_𝐨`$. In the limit of small $`ϵ`$ and $`k`$ , satisfying the conditions $`|ϵ|\tau 1`$ and $`kl1`$, $`C_\pi (𝒌,ϵ)`$ has the following form
$$C_\pi (𝒌,ϵ)=\frac{C_{imp}|U(P_o)|^2}{14i\tau |ϵ|+(kl)^2(\tau /\tau _\pi )^2};$$
(5)
which has a diffusion pole for $`\tau _\pi \tau `$. Notice that $`ϵ`$ is an energy measured from the Fermi level.
The self-energy parts in Figs.1 $`b(b^{})`$ and $`c(c^{})`$ contain a particle- hole propagator $`D_\pi (𝒌,ϵ)`$ which can be written for $`|ϵ|\tau 1`$ and $`kl1`$ as
$$D_\pi (𝒌,ϵ)=\frac{C_{imp}|U(P_o)|^2}{14i\tau |ϵ|+(kl)^2(\tau /\tau _\pi )^2}$$
(6)
$`\pi `$\- diffuson $`D_U(𝒌,ϵ)`$ has a diffusion pole in the particle- hole channel at a small total energy and large momenta difference $`𝐏_𝐨`$ for $`\tau \tau _\pi `$.
According to the Dyson equation, the retarded Green’s function $`G_R(𝒑,ϵ)`$ is expressed in the following form:
$$G_R(𝒑,ϵ)=\frac{1}{(G_R^o(𝒑,ϵ))^1\mathrm{\Sigma }(𝒑,ϵ)}=\underset{n=0}{\overset{\mathrm{}}{}}(G_R^o(𝒑,ϵ))^{n+1}(\mathrm{\Sigma }(𝒑,ϵ))^n$$
(7)
where, $`\mathrm{\Sigma }(𝒑,ϵ)=_{i=a}^c^{}\mathrm{\Sigma }_i(𝒑,ϵ)`$, and each of $`\mathrm{\Sigma }_i(𝒑,ϵ)`$ corresponds to the one of self-energy parts $`ac^{}`$ in Fig.1, respectively. By summing all diagrams in Fig.1 the self-energy part $`\mathrm{\Sigma }(𝐩,ϵ)`$ is reduced to the form:
$$\begin{array}{cc}\hfill \mathrm{\Sigma }(𝐩,ϵ)=& \frac{d^2k}{(2\pi )^2}\{C_\pi (𝒌,ϵ)\frac{2\tau }{\tau _\pi }(1\frac{\tau }{\tau _\pi })D_\pi (𝒌,ϵ)\}\hfill \\ & G_{}^{o}{}_{R}{}^{}(𝐩+𝐤+𝐏_𝐨,ϵ);\hfill \end{array}$$
(8)
Substituting Eqs.(7) and (8) into Eq.(4) the following expression for $`\rho (ϵ)`$ can be obtained after integrating the bare Green’s functions over $`𝐩`$ and taking a sum over $`n`$ :
$$\rho (ϵ)=\rho _o^{(2)}\{1\mathrm{Re}\frac{4\tau ^2\alpha (ϵ)}{\sqrt{1+4\tau ^2\alpha (ϵ)}(1+\sqrt{1+4\tau ^2\alpha (ϵ)})}\}$$
(9)
where
$$\alpha (ϵ)=\frac{d^2k}{(2\pi )^2}\{C_\pi (𝒌,ϵ)\frac{2\tau }{\tau _\pi }(1\frac{\tau }{\tau _\pi })D_\pi (𝒌,ϵ)\};$$
(10)
The obtained expression for the DoS contains high order logarithmically divergent contributions. So, infinite order impurity blocks are summed up according to Eqs.(7)-(8).
The $`nesting`$ processes are weakened away from the middle of the band. In this case the normal scattering processes become more probable and $`\tau _o\tau _\pi `$ . As a result the diffusion poles in the propagators $`C_\pi (𝒌,ϵ)`$ and $`D_\pi (𝒌,ϵ)`$ disappear. Therefore a small offset from the middle of the band gives rise to the disappearence of the quantum corrections to the DoS and $`\rho (ϵ)=\rho _o^{(2)}`$. When filling reaches the center of the band the Fermi surface becomes flat. In this case the scattering with $`nesting`$ seems to have preference, i.e. $`\tau _\pi \tau _o`$ and $`\tau =\tau _\pi `$. The impurity blocks $`C_\pi (𝒌,ϵ)`$ and $`D_\pi (𝒌,ϵ)`$ have a diffusion pole under this condition and we get from Eq.(10):
$$\alpha (ϵ)=\frac{1}{8\pi ϵ_F\tau _{\pi }^{}{}_{}{}^{3}}\mathrm{ln}(\frac{1}{4i\tau _\pi |ϵ|})$$
(11)
Substitution this expression for $`\alpha (ϵ)`$ into Eq.(9) gives the following expression for the DoS:
$$\rho (ϵ)=\rho _o^{(2)}(ϵ)\{1\frac{\alpha _o(ϵ)}{\sqrt{1+\alpha _o(ϵ)}[1+\sqrt{1+\alpha _o(ϵ)}]}\};$$
(12)
where,$`\alpha _o(ϵ)=\frac{1}{2\pi ϵ_F\tau _\pi }\mathrm{ln}(\frac{1}{4\tau _\pi |ϵ|})`$. Near the vicinity of the Fermi level Eq.(12) is approximated as
$$\rho (ϵ)=\rho _o^{(2)}(ϵ)[\frac{1}{2\pi ϵ_F\tau _\pi }\mathrm{ln}(\frac{1}{4\tau _\pi |ϵ|})]^{1/2}$$
(13)
$`\rho _o^{(2)}(ϵ)`$ in Eqs.(12) and (13) is the DoS of a pure 2d square crystal, which has the van Hove singularity expressed by Eq.(2). Elastic scatterings in the crystal with substitutional impurities preserve the central peak of the DoS, however the energy dependence of $`\rho (ϵ)`$ is changed from logarithmic dependence to the square root of the logarithm in the close vicinity of the middle of the band:
$$\rho (ϵ)\frac{2}{(\pi a)^2ϵ_F}(2\pi ϵ_F\tau _\pi )^{1/2}\mathrm{ln}^{1/2}(\frac{1}{4\tau _\pi |ϵ|})\text{as}|ϵ|0$$
(14)
The peak also becomes narrower than that of the van Hove one due to impurity scattering.
It is worth to compare here the result presented in the Letter with that obtained for a 1d lattice with half filled energy band, containing off-diagonal disorder,. The DoS of the pure 1d lattice is a smooth function of the energy within the band and it has a singularity only at the boundary of the energy band. Therefore the Dyson peak of the DoS in the middle of the half-filled band of the 1d disordered lattice is a result of strong Bragg reflection of the electrons in the process of scattering on impurities with a consequent interference of scattered waves. The calculation of density-density and current-current correlators for 1d disordered chains shows that the localization length diverges and static conductivity saturates to a constant value at $`T=0`$ with the approaching of half filling. In contrast to the 1d case the singular enhancement of the DoS in the 2d disordered lattice is due to a van Hove singularity in the bare DoS. Effect of impurity tends to decrease the DoS. The obtained result is also confirmed by the computational study of the DoS of a 2d square lattice with off-diagonal disorder . Numerical analysis of the participation numbers in shows that localization becoms stronger close to the band center.
Notice that a nested Fermi surface with a nesting vector $`𝐏^{}=\xi 𝐏_𝐨`$ ($`\xi <1`$) arise for a tight-binding models which include both nearest neighbor and next nearest neighbor hopping terms and correspond to a non-half filled band. In this case the van Hove singularity in the DoS of pure system disappears and the DoS vanishes on the Fermi surface.
In conclusion, we proposed a new diagrammatic approach to study $`nesting`$ effects on the DoS of 2d disordered square lattice. Calculations were performed in the diffusion approximation. However, unlike the “conventional” localization theory, the infinite order logarithmically divergent diagrams were summed up here according to the Dyson equation. The quantum corrections tend to decrease the DoS in the whole energy interval. As a result, the van Hove singularity in the middle of the band is preserved, however, its energy dependence is changed from logarithmic to square root of logarithm in the close vicinity of the band center. The calculations are valid for a half-filled energy band. A small shift from the band center results in the vanishing of the quantum corrections to the DoS.
The authors are indebted to a referee who pointed out the difference between impurity effects in 1d and 2d systems. One of the authors (E.P.N.) thanks A.Erzan and all members of the Physical Department of Istanbul Technical University, where the part of this work was performed, for long time hospitality. E.P.N. also thanks H.Feldmann for discussion. |
warning/0003/quant-ph0003142.html | ar5iv | text | # Conditional teleportation using optical squeezers and photon counting
## I Introduction
In quantum teleportation, an unknown state of a system is destroyed and created on another, distant system of the same type. The method was first suggested in and realized in for discrete variables, namely photonic (polarization) qubits. Subsequently, the concept has been extended to continuous variables , and then realized experimentally to teleport a coherent state by means of parametrically entangled (squeezed) optical beams and quadrature-component measurements . The concept of teleportation of continuous quantum variables has been further elaborated in .
The basic requirement of quantum teleportation is that the two parties share an entangled state with each other. In continuous-variable teleportation of quantum states of optical field modes, a two-mode squeezed vacuum is suited for playing the role of the entangled state. The quadrature components $`\widehat{q}_k`$ and $`\widehat{p}_k`$ ($`[\widehat{q}_k,\widehat{p}_k]`$ $`=`$ $`i`$, $`k`$ $`=`$ $`1,2`$) are correlated and anti-correlated, respectively, such that $`\mathrm{\Delta }(\widehat{q}_1`$ $``$ $`\widehat{q}_2)`$ $`<`$ $`1`$ and $`\mathrm{\Delta }(\widehat{p}_1`$ $`+`$ $`\widehat{p}_2)`$ $`<`$ $`1`$. For large squeezing, the correlations approach the original Einstein-Podolsky-Rosen (EPR) correlations (for EPR correlations in optical fields, see, e.g., ).
The first scheme of teleportation that uses an optical two-mode squeezed vacuum is based on (single-event) quadrature-component measurements exploiting the above mentioned quadrature-component correlations . Later on, it has been realized that there are photon-number and phase correlations in a two-mode squeezed vacuum which could also be used for a potential teleportation protocol . In the scheme in it is assumed that a measurement of the photon-number difference and the phase sum of the two modes on Alice’s side is performed. The obtained information is then sent to Bob who has to transform the quantum state of his mode by appropriate phase and photon-number shifting, thus creating the resulting teleported state. The scheme is conditional, as for some measured photon-number differences the state that is desired to be teleported cannot be re-created by Bob.
Unfortunately, the scheme in requires phase measurements for which no methods have been known so far. Is there any hope to realize teleportation based on such a scheme or a related one? In this paper we suggest a viable modification of the scheme proposed in which is based on (single-event) photon-number measurements on the output of a parametric amplifier (squeezer). The scheme is also conditional, and it applies to certain classes of quantum states. Even though the scheme is not universal, it can produce for some states and some measurement events higher teleportation fidelities than the scheme based on quadrature-component measurements .
The paper is organized as follows. In Sec. II we present the scheme and derive the expression for the teleported quantum state. In Sec. III we illustrate the method presenting numerical results, and we conclude in Sec. IV.
## II Theory
Let us consider the scheme sketched in Fig. 1. The entangled state is a two-mode squeezed vacuum produced by the first parametric amplifier from the vacuum state, $`\alpha `$ being the squeezing parameter. One of the two output modes of the first parametric amplifier is then used as one of the input modes of the second parametric amplifier (squeezing parameter $`\beta `$), and the mode whose quantum state $`|\psi _{\mathrm{in}}`$ is desired to be teleported is the other input mode. Alice measures the photon numbers $`n`$ and $`n`$ $`+`$ $`d`$ ($`d`$ $``$ $`n`$) at the output of the second parametric amplifier and communicates the result to Bob. Owing to Alice’s measurement, the state of the mode that was sent to Bob from the first parametric amplifier has been projected onto the state $`|\psi _{\mathrm{out}}`$. Bob now reproduces the input state by means of the transformation $`\widehat{A}|\psi _{\mathrm{out}}`$ $`=`$ $`|\psi _{\mathrm{tel}}`$, where the operator $`\widehat{A}`$ shifts the photon number according to the measured photon-number difference $`d`$.
The three modes are initially (i.e., before they enter any of the parametric amplifiers) prepared in the states $`|\psi _{\mathrm{in}}`$, $`|0`$, and $`|0`$, where the state $`|\psi _{\mathrm{in}}`$ that is desired to be teleported can be written in the Fock basis as
$$|\psi _{\mathrm{in}}=\underset{k}{}|kk|\psi _{\mathrm{in}}.$$
(1)
After passing the parametric amplifiers and detecting $`n`$ and $`n^{}`$ photons in the outgoing modes (modes $`0`$ and $`1`$) on Alice’s side, the state of Bob’s mode (mode $`2`$) is
$$|\psi _{\mathrm{out}}_2=P^{\frac{1}{2}}{}_{0}{}^{}n|{}_{1}{}^{}n^{}|\widehat{S}_{01}(\beta )\widehat{S}_{12}(\alpha )|\psi _{\mathrm{in}}_{0}^{}|0_{1}^{}|0_2,$$
(2)
where $`P`$ is the probability of that measurement event. The two-mode squeeze operator $`\widehat{S}_{kl}(\alpha )`$ is given by
$`\widehat{S}_{kl}(\alpha )=\mathrm{exp}\left(\alpha ^{}\widehat{a}_k\widehat{a}_l\alpha \widehat{a}_k^{}\widehat{a}_l^{}\right),`$ (3)
with $`\widehat{a}_k`$ ($`\widehat{a}_k^{}`$) being the photon destruction (creation) operator of the $`k`$th mode. It can be written in the Fock basis as
$`{}_{k}{}^{}m|{}_{l}{}^{}m^{}|\widehat{S}_{kl}(\alpha )|n_{k}^{}|n^{}_{l}^{}=\delta _{mm^{},nn^{}}e^{i(m^{}n^{})\phi _\alpha }`$ (6)
$`\times (1)^n^{}\sqrt{m!m^{}!n!n^{}!}{\displaystyle \frac{(\mathrm{sinh}|\alpha |)^n^{}(\mathrm{tanh}|\alpha |)^m^{}}{(\mathrm{cosh}|\alpha |)^{n+1}}}`$
$`\times {\displaystyle \underset{j=\mathrm{max}\{0,n^{}n\}}{\overset{\mathrm{min}\{m^{},n^{}\}}{}}}{\displaystyle \frac{\left(\mathrm{sinh}^2|\alpha |\right)^j}{j!(m^{}j)!(n^{}j)!(nn^{}+j)!}},`$
where $`\alpha `$ = $`|\alpha |e^{i\phi _\alpha }`$. For the following it is useful to introduce the coefficients
$`S_m^{}^m(d;\alpha )={}_{k}{}^{}m+d|{}_{l}{}^{}m|\widehat{S}_{kl}(\alpha )|m^{}+d_{k}^{}|m^{}_{l}^{}`$ (10)
$`={}_{k}{}^{}m|{}_{l}{}^{}m+d|\widehat{S}_{kl}(\alpha )|m^{}_{k}^{}|m^{}+d_{l}^{}`$
$`=e^{i(mm^{})\phi _\alpha }(1)^m^{}\sqrt{m!m^{}!(m+d)!(m^{}+d)!}`$
$`\times {\displaystyle \frac{(\mathrm{tanh}|\alpha |)^{m+m^{}}}{(\mathrm{cosh}|\alpha |)^{d+1}}}{\displaystyle \underset{j=0}{\overset{\mathrm{min}\{m,m^{}\}}{}}}{\displaystyle \frac{\left(\mathrm{sinh}^2|\alpha |\right)^j}{j!(mj)!(m^{}j)!(d+j)!}}.`$
The properties of the conditional quantum state $`|\psi _{\mathrm{out}}`$, Eq. (2), in which the mode $`2`$ is prepared after the detection of $`n`$ and $`n^{}`$ photons in the modes $`0`$ and $`1`$ respectively, are qualitatively different for different sign of the observed difference $`d`$ $`=`$ $`n^{}`$ $``$ $`n`$. In the case when $`d`$ $``$ $`0`$ is valid, then from Eq. (2) together with Eq. (10) it follows that ($`|\psi _{\mathrm{out}}_2`$ $``$ $`|\psi _{\mathrm{out}}`$)
$$m|\psi _{\mathrm{out}}=P^{\frac{1}{2}}S_m^{n+d}(d;\beta )S_0^m(0;\alpha )md|\psi _{\mathrm{in}},$$
(12)
where the detection probability $`P`$ is given by
$$P=\underset{m}{}|S_m^{n+d}(d;\beta )|^2|S_0^m(0;\alpha )|^2|md|\psi _{\mathrm{in}}|^2.$$
(13)
In the second case when $`d`$ $`>`$ $`0`$ is valid, we derive
$$m|\psi _{\mathrm{out}}=\{\begin{array}{cc}P^{\frac{1}{2}}S_{md}^n(d;\beta )S_0^m(0;\alpha )md|\psi _{\mathrm{in}},\hfill & md,\hfill \\ 0,\hfill & m<d,\hfill \end{array}$$
(14)
where
$$P=\underset{md}{}|S_{md}^n(d;\beta )|^2|S_0^m(0;\alpha )|^2|md|\psi _{\mathrm{in}}|^2.$$
(15)
From an inspection of Eqs. (12) and (14) we see that when the coefficients $`S_m^{n+d}S_0^m`$ and $`S_{md}^nS_0^m`$, respectively, change sufficiently slowly with $`m`$, then the state $`|\psi _{\mathrm{out}}`$, Eq. (2), imitates the state $`|\psi _{\mathrm{in}}`$, Eq. (1), but with a shifted Fock-state expansion, where the shift parameter is just given by the measured photon-number difference $`d`$. Obviously, if $`d`$ $`<`$ $`0`$ then the state $`|\psi _{\mathrm{out}}`$ does not contain any information about the Fock-state expansion coefficients $`m|\psi _{\mathrm{in}}`$ for $`m`$ $`<`$ $`|d|`$. With regard to teleportation, this means that the method is conditional. Successful teleportation of a quantum state whose Fock-state expansion starts with the vacuum can only be achieved if the number of photons detected in the mode $`1`$ is not smaller than the number of photons detected in the mode $`0`$. This limitation is exactly of the same kind as in the scheme in : the teleportation fidelity tends sharply to zero as the photon-number difference exceeds some (state-dependent) threshold value. Examples of the coefficients $`S_m^{n+d}S_0^m`$ are plotted in Fig. 2 for $`d`$ $`=`$ $`0`$.
To complete the teleportation procedure, Bob transforms the state $`|\psi _{\mathrm{out}}`$ applying on it photon-number shifting. Thus, the teleported state is
$$|\psi _{\mathrm{tel}}=\{\begin{array}{cc}\widehat{E}^d|\psi _{\mathrm{out}}\text{if}\hfill & d<0,\hfill \\ \widehat{E}^d|\psi _{\mathrm{out}}\text{if}\hfill & d>0,\hfill \end{array}$$
(16)
where
$$\widehat{E}=\underset{n}{}|nn+1|$$
(17)
(i.e., the operator $`\widehat{A}`$ in Fig. 1 is a power of $`\widehat{E}`$ or $`\widehat{E}^{}`$). The teleportation fidelity is then given by
$$F=|\psi _{\mathrm{in}}|\psi _{\mathrm{tel}}|^2.$$
(18)
For $`d`$ $``$ $`0`$, the teleportation scheme requires a realization of the transformations $`\widehat{E}`$ and $`\widehat{E}^{}`$. Unfortunately, there has been no exact implementation of these transformations in quantum optics so far. Photon adding and subtracting are transformations that are very close to the required ones. They are based on conditional measurement and could be realized using presently available experimental techniques . Their use of course reduces the efficiency of the scheme. Thus, the scheme may be presently confined to the case where $`d`$ $`=`$ $`0`$.
## III Results
From Eqs (12) and (14) together with Eqs. (16) – (18), the main results can be summarized as follows. $`(i)`$ Fock states can perfectly be teleported, i.e., the fidelity, Eq. (18), is equal to unity, which follows from the fact that parametric amplifiers conserve the photon-number difference. Therefore, high teleportation fidelities can also be expected for states with small photon number dispersion. For such states our method may be more suitable than the method in , where teleportation via measurement of conjugate quadrature components is realized. On the other hand, high teleportation fidelities are not expected for states with large mean photon number and large photon-number dispersion. In particular, for teleportation of highly excited coherent states or phase squeezed states the method in may be more suitable. $`(ii)`$ In comparison to the method in , our scheme does not require phase shifting of the output state $`|\psi _{\mathrm{out}}`$. The squeezing parameters $`\alpha `$ and $`\beta `$ can be chosen such that the coefficients $`S_m^{n+d}S_0^m`$ and $`S_{md}^nS_0^m`$ in Eqs. (12) and (14), respectively, are real, so that the Fock-state expansion coefficients $`m|\psi _{\mathrm{out}}`$ have the same phase as the coefficients $`md|\psi _{\mathrm{in}}`$. $`(iii)`$ A high teleportation fidelity can be expected, provided that the values of the coefficients $`S_m^{n+d}S_0^m`$ and $`S_{md}^nS_0^m`$ vary sufficiently slowly with $`m`$ in the relevant range of the Fock-state expansion of the input state $`|\psi _{\mathrm{in}}`$. On the other hand, in ranges where the coefficients change rapidly, reliable teleportation cannot be achieved. From Fig. 2 it is seen that the $`m`$-range in which $`S_{md}^nS_0^m`$ slowly varies with $`m`$ increases with the strength of squeezing, and thus the class of states that can be teleported reliably extends.
In order to illustrate the method, we have calculated the teleported state, assuming that input state is a superposition of two Fock states, $`|\psi _{\mathrm{in}}`$ $`=`$ $`2^{1/2}(|1`$ $`+`$ $`i|3)`$ and equal squeezing parameters $`\alpha `$ and $`\beta `$ are used. Figure 3 presents the dependence on the detected photon numbers $`n,n^{}`$ of the teleportation fidelity $`F`$, Eq. (18), and the success probability $`P`$, Eqs. (13) and (15). We observe that close to the diagonal (but not always directly on it) the fidelity reaches high values close to unity. For the values of $`n,n^{}`$ with $`n`$ $`=`$ $`n^{}+2`$ and $`n`$ $`=`$ $`n^{}+3`$ the fidelity is exactly $`0.5`$, which indicates that the Fock state $`|3`$ was in the input of the second squeezer and has therefore been re-created in the teleportation. For the values of $`n,n^{}`$ with $`n`$ $`>`$ $`n^{}+3`$ the fidelity drops to zero and so does the probability: such events do not occur for the input state under consideration.
To quantify the performance of the method, we have calculated the probability of events which yield teleportation fidelities larger than or equal to some upper value $`F_\mathrm{u}`$,
$$P_\mathrm{u}=\underset{\genfrac{}{}{0pt}{}{n,n^{}}{F(n,n^{})F_\mathrm{u}}}{}P(n,n^{}),$$
(19)
where $`P(n,n^{})`$ is the success probability of detecting $`n`$ and $`n^{}`$ photons \[Eqs. (13) and (15)\], and $`F(n,n^{})`$ is the corresponding teleportation fidelity. In the example, we find that $`P_\mathrm{u}`$ $``$ $`33\%`$ for $`F_\mathrm{u}`$ $`=`$ $`90\%`$. Measuring (in place of $`n`$ and $`n^{}`$ in our scheme) the quadrature components $`X_0`$ and $`P_1`$ in the scheme in would yield (for the same input state and the same squeezing parameter of the entangled state) $`P_\mathrm{u}`$ $``$ $`23\%`$.
Let us consider the more realistic case where $`n`$ $`=`$ $`n^{}`$, so that no photon-number shifting is necessary. From Fig. 4 we see that with increasing strength of squeezing a higher fidelity can be realized. However, the corresponding success probability decreases. In the figure, the overall probability of realizing a fidelity higher than 90% is $`P`$ $``$ $`1.97\%`$ for the squeezing parameters $`\alpha `$ $`=`$ $`\beta `$ $`=`$ $`1.5`$, whereas for $`\alpha `$ $`=`$ $`\beta `$ $`=`$ $`2`$ the probability reduces to $`P`$ $``$ $`1.06\%`$.
Clearly, if the input states are completely unknown, it cannot be predicted with what fidelity a state would be teleported. The scheme applies if the input states can be confined to a certain class of states, so that from an estimated fidelity it can be decided which photodetection results would represent a successful teleportation.
## IV Summary and Conclusions
We have suggested a viable modification of the teleportation scheme proposed in . Our scheme avoids the phase sum measurement that is not realizable at present. It uses instead the property of a nondegenerate parametric amplifier that the photon-number difference of the output beams is equal to that of the input beams. However, the price of avoiding phase measurements is a relatively low success probability of teleportation.
Our method and the method in , which is based on quadrature-component measurements, may complement one another. So, our method is better suited to teleportation of states with small photon number dispersion (Fock states can be teleported with fidelity equal to unity in principle). The method in is more suitable for teleportation of states with smooth quadrature-component distributions.
Although our method is realizable in principle, there are several non-trivial experimental challenges. First, precise photodetection is needed, i.e., detectors are required that are able to distinguish between different photon numbers. This does not only concern Alice’s measurement but also Bob’s photon-number shifting, e.g., by means of photon adding and subtracting. Second, the photodetection should be sufficiently mode-selective, i.e, one must be able to distinguish whether an incident photon comes from the mode under study or from another part of the spectrum generated by the parametric amplifiers. A central problem in any scheme that exploits quantum coherence is that of decoherence due to unavoidable losses. The effect of decoherence may be reduced, if the squeezing strengths are reduced. However, using smaller squeezing decreases the available teleportation fidelity, so that one has to find an optimum regime for the teleportation of a given class of states, the losses in the scheme, and the required fidelity.
###### Acknowledgements.
This work was supported by the Deutsche Forschungsgemeinschaft. We thank S. Braunstein and P. Kok for useful comments. |
warning/0003/math-ph0003034.html | ar5iv | text | # A Crossed Module giving the Godbillon-Vey Cocycle
## Introduction
There is a well known correspondence in homological algebra between (equivalence classes of) exact sequences $`\mathrm{\Lambda }`$-modules, starting in $`M`$ and ending in $`N`$ with $`n`$ modules in between, and elements of $`Ext_\mathrm{\Lambda }^n(N,M)`$. For $`\mathrm{\Lambda }=U(𝔤)`$ the universal enveloping algebra of a Lie algebra $`𝔤`$, this gives for example a correspondence between $`H^2(𝔤,V)`$ and short eaxct sequences
$$0V\widehat{g}=V\times 𝔤𝔤0$$
where $`\widehat{g}`$ is the semi-direct product of the abelian Lie algebra $`V`$ and $`𝔤`$. The Lie algebra law on $`\widehat{g}`$ is given by a 2-cocycle of $`𝔤`$ with values in $`V`$. Note that the short exact sequence in uniquely determined on the vector space level by $`V`$ and $`𝔤`$.
In the same way, there is a correspondence between $`H^3(𝔤,V)`$ and certain 4 term exact sequences called crossed modules. In these sequences, only the first and the last term are specified, leaving (du to exactness) a choice of one $`𝔤`$-module to complete the sequence.
In this article, we exhibit crossed modules corresponding to the Godbillon-Vey 3-cocycle for $`W_1`$, $`Vect(S^1)`$ and the 2-dimensional analogue of $`Vect(S^1)`$ $`Vect_{0,1}(\mathrm{\Sigma })`$ (or $`Vect_{1,0}(\mathrm{\Sigma })`$), $`\mathrm{\Sigma }`$ being a compact Riemann surface.
Acknoledgements
The author thanks J.-L. Loday for his question stimulating the present article and C. Kassel for making him aware of reference .
## 1 Crossed modules
In the same way as extensions, i.e. 3 term exact sequences, are related to Lie algebra cohomology in degree 2, crossed modules , i.e. certain 4 term exact sequences, correspond to cohomology in degree 3.
###### Definition 1
A crossed module of a Lie algebra is a homomorphism of Lie algebras $`\mu :𝔪𝔫`$ together with a $`𝔫`$-module structure $`\eta `$ on $`𝔪`$ such that
(a) $`\mu (\eta (n)m)=[n,\mu (m)]`$ for all $`n𝔫`$ and all $`m𝔪`$,
(b) $`\eta (\mu (m))m^{}=[m,m^{}]`$ for all $`m,m^{}𝔪`$.
One shows that $`ker(\mu )=:V`$ is an abelian Lie algebra and that the action of $`𝔫`$ on $`𝔪`$ induces a structure of a $`𝔤:=coker(\mu )`$-module on $`V`$. Note that in general $`𝔪`$ and $`𝔫`$ are not $`𝔤`$-modules.
An equivalence of crossed modules is defined in a natural way such that the restrictions of the maps on the kernel and the cokernel of $`\mu `$ are identical. Let us denote by $`\mathrm{crmod}(𝔤,V)`$ the set of equivalence classes of crossed modules. We have then the fundamental correspondence theorem A.2, p. 138,
$$\mathrm{crmod}(𝔤,V)H^3(𝔤,V).$$
(1)
Let us briefly review how to associate to a crossed module a 3 cocycle of $`𝔤`$ with values in $`V`$:
The exact sequence associated to a crossed module reads as follows:
$$0V\stackrel{i}{}𝔪\stackrel{\mu }{}𝔫\stackrel{\pi }{}𝔤0$$
The first step is to take a (continuous) section $`\sigma `$ of $`\pi `$ and to calculate the default of $`\sigma `$ to be a Lie algebra homomorphism, i.e.
$$\alpha (x_1,x_2):=\sigma ([x_1,x_2])[\sigma (x_1),\sigma (x_2)].$$
We have obviously $`\pi \alpha (x_1,x_2)=0`$, because $`\pi `$ is a Lie algebra homomorphism, so $`\alpha (x_1,x_2)im(\mu )=ker(\pi )`$. This means that there is a $`\beta (x_1,x_2)𝔪`$ such that
$$\mu (\beta (x_1,x_2))=\alpha (x_1,x_2).$$
Now, one can easily calculate that $`\mu (d\beta (x_1,x_2,x_3))=0`$ where $`d`$ is the Lie algebra cohomology boundary of cohomology of $`𝔤`$ with values in $`𝔪`$ where $`𝔤`$ acts on $`𝔪`$ by $`\eta \sigma `$.
This means that $`d\beta (x_1,x_2,x_3)ker(\mu )=im(i)=V`$, i.e. there is a $`\gamma (x_1,x_2,x_3)V`$ such that $`d\beta (x_1,x_2,x_3)=i(\gamma (x_1,x_2,x_3))`$. It is fairly obvious that $`\gamma `$ is a 3-cocycle of $`𝔤`$ with values in $`V`$.
## 2 A crossed module giving $`\theta H^3(W_1)`$
In the following, we want to exhibit the 4 term exact sequence related to the Godbillon-Vey class in $`H^3(W_1,)`$.
Recall that there is a semi-direct product of $`W_1`$ by its module of $`\lambda `$-densities. It is represented by a short exact sequence
$$0F_\lambda F_\lambda \times W_1W_10,$$
where the Lie algebra structure on the product is given by
$$[(f,a),(g,b)]=([f,g],L_fbL_ga+c(f,g)).$$
Here, $`c`$ is a 2-cocycle of $`W_1`$ with values in $`F_\lambda `$, the module of $`\lambda `$-densities on the line $``$. We take $`\lambda =1`$ and
$$c(f,g)=\left|\begin{array}{cc}f^{}& g^{}\\ f^{\prime \prime }& g^{\prime \prime }\end{array}\right|.$$
Now consider the short exact sequence of $`W_1`$-modules
$$0F_0\stackrel{d}{}F_10$$
Since $`W_1`$ acts by the Lie derivative $`L_X=di_X+i_Xd`$ on the density modules, the action commutes with the exterior differential $`d`$.
We can glue these 2 sequences together to get
$$0F_0\stackrel{d}{}F_1\times W_1W_10,$$
(2)
$`d`$ is trivially a Lie algebra homomorphism from the abelian Lie algebra $`F_0`$ to the semidirect product $`F_1\times W_1`$. There is an action of $`F_1\times W_1`$ on $`F_0`$: it is given by the action of $`W_1`$ on $`F_0`$. It gives a structure of a crossed module to the sequence (2). The second compatibility condition is trivial. The first one reads
$$d((f,a)b)=[(f,a),db]$$
The left hand side is just $`d(L_fb)=di_{f\frac{d}{dt}}db`$. The right hand side gives $`(0,L_f(db))`$ and commutation of $`d`$ and $`L_f`$ shows condition $`(a)`$. In conclusion, the 4 term exact sequence gives a crossed module. We claim that the crossed module corresponds to the Godbillon-Vey class via the isomorphism (1).
In fact, this is easy: one just has to move up the arrows in the wrong direction. First, one has to choose a section of the arrow $`F_1\times W_1W_1`$. But calculating the default of this section to be a Lie algebra map just gives the cocycle of the semi-direct product. Then we have to move up the de Rham differential $`d`$, i.e. we have to choose a primitive. Afterwards, we have to take the Lie algebra coboundary. So, if we denote by a hat the primitive, we have to calculate
$$d(\widehat{c})(f,g,h)=d(\widehat{\left|\begin{array}{cc}^{}& ^{}\\ ^{\prime \prime }& ^{\prime \prime }\end{array}\right|})(f,g,h).$$
But it is obvious that the Lie algebra coboundary commutes with taking the primitive. Furthermore, it is well known that we have the following relation
$$d(\left|\begin{array}{cc}^{}& ^{}\\ ^{\prime \prime }& ^{\prime \prime }\end{array}\right|)(f,g,h)=\left(\left|\begin{array}{ccc}f& g& h\\ f^{}& g^{}& h^{}\\ f^{\prime \prime }& g^{\prime \prime }& h^{\prime \prime }\end{array}\right|\right)^{}$$
This expresses that fact that the Gelfand-Fuks cocycle is obtained from the Godbillon-Vey cocycle by “integration over the manifold” (and is easily checked directly).
Thus, we can choose the primitive such that the 3-cocycle associated to the above 4 term exact sequence is the Godbillon-Vey cocycle which is the generator of $`H^3(W_1,)`$.
Remarks:
1) This reasoning also shows that $`H^2(W_1,F_1)`$ which is known to be 1-dimensional (consequence of Goncharova’s theorem, cf p. 120), is generated by the cocycle
$$c(f,g)=\left|\begin{array}{cc}f^{}& g^{}\\ f^{\prime \prime }& g^{\prime \prime }\end{array}\right|dx.$$
2) We owe to C. Roger the remark that the crossed module is obtained under Yoneda product $`Ext_{U(W_1)}^1(F_1,)\times Ext_{U(W_1)}^1(,F_1)Ext_{U(W_1)}^2(,)H^3(W_1,)`$. This possibility bears perhaps deep homological properties.
## 3 A crossed module giving $`\theta H^3(Vect(S^1),)`$
Now we want to generalize the preceeding situation to the circle. There is an obvious problem to do so: the sequence
$$0_0\stackrel{d}{}_10$$
where $`_0=𝒞^{\mathrm{}}(S^1)`$ and $`_1=\mathrm{\Omega }^1(S^1)`$ is not exact, but has a cokernel $`H^1(S^1,)=`$.
In order to get around this problem, let us recall some evident, but somewhat strange facts:
There is a chain of inclusions of Lie algebras
$$W_1^{pol}Vect(S^1)Vect()$$
Here, $`W_1^{pol}=_{n1}x^{n+1}\frac{d}{dx}`$ is the Lie algebra of (complexified) polynomial vector fields on the line, $`Vect(S^1)=\widehat{_ne^{int}\frac{d}{dt}}`$ is the Lie algebra of complexified vector fields on the circle (the hat meaning that $`e^{int}\frac{d}{dt}`$ for $`n`$ is a topological basis, because periodic functions can be approximated by Fourier polynomials), and $`Vect()=\{f(x)\frac{d}{dx}|f𝒞^{\mathrm{}}(,)\}`$ is the Lie algebra of complexified vector fields on the real line $``$.
The maps $`x^{n+1}\frac{d}{dx}ie^{int}\frac{d}{dt}`$ (i.e. one sets $`x=e^{it}`$) and $`f(t)\frac{d}{dt}\stackrel{~}{f}(x)\frac{d}{dx}`$ (where the field $`f(t)\frac{d}{dt}`$ is lifted to the universal covering to the unique 1-periodic field $`\stackrel{~}{f}(x)\frac{d}{dx}`$) are easily seen to be Lie algebra homomorphisms.
The strange fact is that $`W_1^{pol}`$ and $`Vect()`$ have the same continuous cohomology, but $`Vect(S^1)`$ does not. This is elucidated by the fact that the isomorphism in cohomology between $`W_1^{pol}`$ and $`Vect()`$ is not induced by the above inclusion, but by the “Taylor expansion at 0”-map $`Vect()W_1^{pol}`$ (which is a continuous surjection of Fréchet spaces by Borel’s lemma). One sees that $`W_1^{pol}`$ and $`Vect()`$ correspond to the 2 ways of associating $``$ to $`S^1`$: as its tangent space, or as its universal covering.
I recall all this just to motivate the fact that $`Vect(S^1)`$ acts on $`F_\lambda `$ (the $`\lambda `$-densities on $``$) and $`W_1^{pol}`$ acts on $`_\lambda `$ (the $`\lambda `$-densities on $`S^1`$). In the first case, we embed $`Vect(S^1)`$ into $`Vect()`$, and in the second case, we embed $`W_1^{pol}`$ into $`Vect(S^1)`$.
This gives us the possibility to consider the exact sequence
$$0F_0\stackrel{d}{}F_1\times Vect(S^1)Vect(S^1)0.$$
By the same arguments as above, this is a crossed module giving the Godbillon-Vey cocycle.
## 4 A crossed module giving $`\theta H^3(Vect_{1,0}(\mathrm{\Sigma }),)`$
Now, let us recall that $`H^3(Vect_{1,0}(\mathrm{\Sigma }),)`$, where $`Vect_{1,0}(\mathrm{\Sigma })`$ is the Lie algebra of differentiable vector fields of type $`(1,0)`$ (i.e. locally, such a field is written $`f(z,\overline{z})\frac{}{z}`$) on a compact Riemann surface $`\mathrm{\Sigma }`$, is also of dimension 1 and generated by a Godbillon-Vey type cocycle, see , . This cocycle is the integral over $`\mathrm{\Sigma }`$ over a determinant as in the usual Godbillon-Vey cocycle, all derivatives being with respect to $`z`$.
We propose a crossed module corresponding to this generator.
There are two obvious problems: first of all, there is the same problem as in the preceeding section, so we have to lift our fields to the universal covering, and second, there is the problem that the derivatives involved are only $``$, instead of being $`d=+\overline{}`$. For the second problem, we define a new action of $`Vect_{1,0}(\mathrm{\Sigma })`$ on $`_{i=0,1}\mathrm{\Omega }^{i,0}(\mathrm{\Sigma })`$ by the modified Cartan formula:
$$L_X\omega =i_X+i_X.$$
This means explicitly
$$L_X\omega =\{\begin{array}{cc}i_X((\omega ))& \mathrm{if}deg(\omega )=0\\ (i_X(\omega ))& \mathrm{if}deg(\omega )=1\end{array}$$
One calculates easily that this defines a Lie algebra action of $`Vect_{1,0}(\mathrm{\Sigma })`$ on $`_{i=0,1}\mathrm{\Omega }^{i,0}(\mathrm{\Sigma })`$ and that
$$0\mathrm{\Omega }^{0,0}(\mathrm{\Sigma })\stackrel{}{}\mathrm{\Omega }^{1,0}(\mathrm{\Sigma })$$
is an exact sequence of $`Vect_{1,0}(\mathrm{\Sigma })`$-modules, unfortunately the last map is not surjective. Now, let us lift these modules to the universal covering:
Suppose $`g>0`$ in order to have as the universal covering of $`\mathrm{\Sigma }`$ either $``$ or $``$, and let $`U`$ denote $``$ in case the genus is 1 and $``$ in case the genus is greater than 1. We have a short exact sequence of $`Vect_{1,0}(\mathrm{\Sigma })`$-modules
$$0\mathrm{\Omega }^{0,0}(U)\stackrel{}{}\mathrm{\Omega }^{1,0}(U)0$$
Now we can form the 4 term exact sequence
$$0\mathrm{\Omega }^{0,0}(U)\stackrel{}{}\mathrm{\Omega }^{1,0}(U)\times Vect_{1,0}(\mathrm{\Sigma })Vect_{1,0}(\mathrm{\Sigma })0$$
This sequence gives obviously by the same arguments as before a crossed module with corresponding 3-cohomology class represented by the Godbillon-Vey cocycle.
Note that the same construction works for the Lie algebra of holomorphic vector fields on an open Riemann surface $`Hol(\mathrm{\Sigma }_r)`$ where $`\mathrm{\Sigma }_r=\mathrm{\Sigma }\{p_1,\mathrm{},p_r\}`$ acting on the module of holomorphic $`\lambda `$-densities $`_\lambda (\mathrm{\Sigma }_r)`$ by the modified Cartan formula. This gives a crossed module describing $`H^3(Hol(\mathrm{\Sigma }_r),)`$ which is also known to be 1-dimensional . |
warning/0003/physics0003056.html | ar5iv | text | # 1 Problem
## 1 Problem
Deduce an axicon solution for a Gaussian laser beam in vacuum, i.e., a beam with radial polarization of the electric field.
## 2 Solution
If a laser beam is to have radial transverse polarization, the transverse electric must vanish on the symmetry axis, which is charge free in vacuum. However, we can expect a nonzero longitudinal electric field on the axis, noting that the projections onto the axis of the electric field vectors of rays all have the same sign, as shown in Fig. 1a. This contrasts with the case of linearly polarized Gaussian laser beams for which rays at $`0^{}`$ and $`180^{}`$ azimuth to the polarization direction have axial electric field components of opposite sign, as shown in Fig. 1b. The longitudinal electric field of axicon laser beams may be able to transfer net energy to charged particles that propagate along the optical axis, providing a form of laser acceleration .
Although two of the earliest papers on Gaussian laser beams discuss axicon modes (without using that term, and without deducing the simplest axicon mode), most subsequent literature has emphasized linearly polarized Gaussian beams. We demonstrate that a calculation that begins with the vector potential (sec. 2.1) leads to both the lowest-order linearly polarized and axicon modes. We include a discussion of Gaussian laser pulses as well as continuous beams, and find in sec. 2.2 that the temporal pulse shape must obey condition (8). The paraxial wave equation and its lowest-order, linearly polarized solutions are reviewed in secs. 2.3-4. Readers familiar with the paraxial wave equation for linearly polarized Gaussian beams may wish to skip directly to sec. 2.5 in which the axicon mode is displayed. In sec. 2.6 we find an expression for a guided axicon beam, i.e., one that requires a conductor along the optical axis.
### 2.1 Solution via the Vector Potential
Many discussions of Gaussian laser beams emphasize a single electric field component such as $`E_x=f(r,z)e^{i(kz\omega t)}`$ of a cylindrically symmetric beam of angular frequency $`\omega `$ and wave number $`k=\omega /c`$ propagating in vacuum along the $`z`$ axis. Of course, the electric field must satisfy the free-space Maxwell equation $`𝐄=0`$. If $`f(r,z)`$ is not constant and $`E_y=0`$, then we must have nonzero $`E_z`$. That is, the desired electric field has more than one vector component.
To deduce all components of the electric and magnetic fields of a Gaussian laser beam from a single scalar wave function, we follow the suggestion of Davis and seek solutions for a vector potential A that has only a single component. We work in the Lorentz gauge (and Gaussian units), so that the scalar potential $`\mathrm{\Phi }`$ is related to the vector potential by
$$𝐀+\frac{1}{c}\frac{\mathrm{\Phi }}{t}=0.$$
(1)
The vector potential can therefore have a nonzero divergence, which permits solutions having only a single component. Of course, the electric and magnetic fields can be deduced from the potentials via
$$𝐄=\mathrm{\Phi }\frac{1}{c}\frac{𝐀}{t},$$
(2)
and
$$𝐁=\times 𝐀.$$
(3)
For this, the scalar potential must first be deduced from the vector potential using the Lorentz condition (1).
The vector potential satisfies the free-space wave equation,
$$^2𝐀=\frac{1}{c^2}\frac{^2𝐀}{t^2}.$$
(4)
We seek a solution in which the vector potential is described by a single component $`A_j`$ that propagates in the $`+z`$ direction with the form
$$A_j(𝐫,t)=\psi (r_{},z)g(\phi )e^{i\phi },$$
(5)
where the spatial envelope $`\psi `$ is azimuthally symmetric, $`r_{}=\sqrt{x^2+y^2}`$, $`g`$ is the temporal pulse shape, and the phase $`\phi `$ is given by
$$\phi =kz\omega t.$$
(6)
Inserting trial solution (5) into the wave equation (4) we find that
$$^2\psi +2ik\frac{\psi }{z}\left(1\frac{ig^{}}{g}\right)=0,$$
(7)
where $`g^{}=dg/d\phi `$.
### 2.2 A Condition on the Temporal Pulse Shape $`g(\phi )`$
Since $`\psi `$ is a function of r while $`g`$ and $`g^{}`$ are functions of the phase $`\phi `$, eq. (7) cannot be satisfied in general. Often the discussion is restricted to the case where $`g^{}=0`$, i.e., to continuous waves. For a pulsed laser beam, $`g`$ must obey
$$\left|\frac{g^{}}{g}\right|1$$
(8)
for eq. (7) to be consistent.
It is noteworthy that a “Gaussian” laser beam cannot have a Gaussian temporal pulse. That is, if $`g=\mathrm{exp}[(\phi /\phi _0)^2]`$, then $`\left|g^{}/g\right|=2\left|\phi \right|/\phi _0^2`$, which does not satisfy condition (8) for $`|\phi |`$ large compared to the characteristic pulsewidth $`\phi _0=\omega \mathrm{\Delta }t`$, i.e., in the tails of the pulse.
A more appropriate form for a pulsed beam is a hyperbolic secant (as arises in studies of solitons):
$$g(\phi )=\text{sech}\left(\frac{\phi }{\phi _0}\right).$$
(9)
Then, $`\left|g^{}/g\right|=(1/\phi _0)\left|\mathrm{tanh}(\phi /\phi _0)\right|`$, which is less than one everywhere provided that $`\phi _01`$.
### 2.3 The Paraxial Wave Equation
In the remainder of this paper, we suppose that condition (8) is satisfied. Then, the differential equation (7) for the spatial envelope function $`\psi `$ becomes
$$^2\psi +2ik\frac{\psi }{z}=0.$$
(10)
The function $`\psi `$ can and should be expressed in terms of three geometric parameters of a focused beam, the diffraction angle $`\theta _0`$, the waist $`w_0`$, and the depth of focus (Rayleigh range) $`z_0`$, which are related by
$$\theta _0=\frac{w_0}{z_0}=\frac{2}{kw_0},\text{and}z_0=\frac{kw_0^2}{2}=\frac{2}{k\theta _0^2}.$$
(11)
We therefore work in the scaled coordinates
$$\xi =\frac{x}{w_0},\upsilon =\frac{y}{w_0},\rho ^2=\frac{r_{}^2}{w_0^2}=\xi ^2+\upsilon ^2,\text{and}\varsigma =\frac{z}{z_0},$$
(12)
Changing variables and noting relations (11), eq. (10) takes the form
$$_{}^2\psi +4i\frac{\psi }{\varsigma }+\theta _0^2\frac{^2\psi }{\varsigma ^2}=0,$$
(13)
where
$$_{}^2\psi =\frac{^2\psi }{\xi ^2}+\frac{^2\psi }{\upsilon ^2}=\frac{1}{\rho }\frac{}{\rho }\left(\rho \frac{\psi }{\rho }\right),$$
(14)
since $`\psi `$ is independent of the azimuth $`\varphi `$.
The form of eq. (13) suggests the series expansion
$$\psi =\psi _0+\theta _0^2\psi _2+\theta _0^4\psi _4+\mathrm{}$$
(15)
in terms of the small parameter $`\theta _0^2.`$ Inserting this into eq. (13) and collecting terms of order $`\theta _0^0`$ and $`\theta _0^2`$, we find
$$_{}^2\psi _0+4i\frac{\psi _0}{\varsigma }=0,$$
(16)
and
$$_{}^2\psi _2+4i\frac{\psi _2}{\varsigma }=\frac{^2\psi _0}{\varsigma ^2},$$
(17)
etc.
Equation (16) is called the the paraxial wave equation, whose solution is well-known to be
$$\psi _0=fe^{f\rho ^2},$$
(18)
where
$$f=\frac{1}{1+i\varsigma }=\frac{1i\varsigma }{1+\varsigma ^2}=\frac{e^{i\mathrm{tan}^1\varsigma }}{\sqrt{1+\varsigma ^2}}.$$
(19)
The factor $`e^{i\mathrm{tan}^1\varsigma }`$ in $`f`$ is the so-called Guoy phase shift , which changes from 0 to $`\pi /2`$ as $`z`$ varies from 0 to $`\mathrm{}`$, with the most rapid change near the $`z_0`$.
The solution to eq. (17) for $`\psi _2`$ has been given in , and that for $`\psi _4`$ has been discussed in .
With the lowest-order spatial function $`\psi _0`$ in hand, we are nearly ready to display the electric and magnetic fields of the corresponding Gaussian beams. But first, we need the scalar potential $`\mathrm{\Phi }`$, which we suppose has the form
$$\mathrm{\Phi }(𝐫,t)=\mathrm{\Phi }(𝐫)g(\phi )e^{i\phi },$$
(20)
similar to that of the vector potential. Then,
$$\frac{\mathrm{\Phi }}{t}=\omega \mathrm{\Phi }\left(1\frac{ig^{}}{g}\right)\omega \mathrm{\Phi },$$
(21)
assuming condition (8) to be satisfied. In that case,
$$\mathrm{\Phi }=\frac{i}{k}𝐀,$$
(22)
according to the Lorentz condition (1). The electric field is then given by
$$𝐄=\mathrm{\Phi }\frac{1}{c}\frac{𝐀}{t}ik\left[𝐀+\frac{1}{k^2}(𝐀)\right],$$
(23)
in view of condition (8). Note that $`(1/k)/x=(\theta _0/2)/\xi `$, etc., according to eqs. (11)-(12).
### 2.4 Linearly Polarized Gaussian Beams
Taking the scalar wave function (18) to be the $`x`$ component of the vector potential,
$$A_x=\frac{E_0}{ik}\psi _0g(\phi )e^{i\phi },A_y=A_z=0,$$
(24)
the corresponding electric and magnetic fields are found from eqs. (3), (23) and (24) to be the familiar forms of a linearly polarized Gaussian beam,
$`E_x`$ $`=`$ $`E_0\psi _0ge^{i\phi }+𝒪(\theta _0^2)E_0fe^{f\rho ^2}ge^{i\phi }`$
$`=`$ $`{\displaystyle \frac{E_0e^{\rho ^2/(1+\varsigma ^2)}g(\phi )}{\sqrt{1+\varsigma ^2}}}e^{i[kz+\varsigma \rho ^2/(1+\varsigma ^2)\omega t\mathrm{tan}^1\varsigma ]},`$
$`=`$ $`{\displaystyle \frac{E_0e^{r_{}^2/w^2(z)}g(\phi )}{\sqrt{1+z^2/z_0^2}}}e^{i\{kz[1+r_{}^2/2(z^2+z_0^2)]\omega t\mathrm{tan}^1(z/z_0)\}},`$
$`E_y`$ $`=`$ $`0,`$ (25)
$`E_z`$ $`=`$ $`{\displaystyle \frac{i\theta _0E_0}{2}}{\displaystyle \frac{\psi _0}{\xi }}ge^{i\phi }+𝒪(\theta _0^3)i\theta _0f\xi E_x,`$
$`B_x`$ $`=`$ $`0,`$
$`B_y`$ $`=`$ $`E_x,`$ (26)
$`B_z`$ $`=`$ $`{\displaystyle \frac{i\theta _0E_0}{2}}{\displaystyle \frac{\psi _0}{\upsilon }}ge^{i\phi }=i\theta _0f\upsilon E_x,`$
where
$$w(z)=w_0\sqrt{1+z^2/z_0^2}$$
(27)
is the characteristic transverse size of the beam at position $`z`$. Near the focus (r
<
w0,|z|<z0formulae-sequence
<
subscript𝑟perpendicular-tosubscript𝑤0𝑧subscript𝑧0r_{\perp}\mathrel{\vbox{\kern 0.0pt\hbox{$<$} \kern 0.0pt\hbox{$\sim$} }}w_{0},\left|z\right|<z_{0}), the beam is a plane wave,
$$E_xE_0e^{r_{}^2/w_0^2}e^{i(kz\omega tz/z_0)},E_z\theta _0\frac{x}{w_0}E_0e^{r_{}^2/w_0^2}e^{i(kz\omega t2z/z_0\pi /2)},$$
(28)
For large $`z`$,
$$E_xE_0e^{\theta ^2/\theta _0^2}\frac{e^{i(kr\omega t\pi /2)}}{r},E_z\frac{x}{r}E_x,$$
(29)
where $`r=\sqrt{r_{}^2+z^2}`$ and $`\theta r_{}/r`$, which describes a linearly polarized spherical wave of extent $`\theta _0`$ about the $`z`$ axis. The fields $`E_x`$ and $`E_z`$, i.e., the real parts of eqs. (29), are shown in Figs. 2 and 3.
The fields (25)-(26) satisfy $`𝐄=0=𝐁`$ plus terms of order $`\theta _0^2`$.
Clearly, a vector potential with only a $`y`$ component of form similar to eq. (24) leads to the lowest-order Gaussian beam with linear polarization in the $`y`$ direction.
### 2.5 The Lowest-Order Axicon Beam
An advantage of our solution based on the vector potential is that we also can consider the case that only $`A_z`$ is nonzero and has the form (18),
$$A_x=A_y=0,A_z=\frac{E_0}{k\theta _0}fe^{f\rho ^2}ge^{i(kz\omega t)}.$$
(30)
Then,
$$𝐀=\frac{A_z}{z}ikA_z\left[1\frac{\theta _0^2}{2}f(1f\rho ^2)\right],$$
(31)
using eqs. (11)-(12) and the fact that $`df/d\varsigma =if^2`$, which follows from eq. (19). Anticipating that the electric field has radial polarization, we work in cylindrical coordinates, $`(r_{},\varphi ,z)`$, and find from eqs. (3), (23), (30) and (31) that
$`E_{}`$ $`=`$ $`E_0\rho f^2e^{f\rho ^2}ge^{i\phi }+𝒪(\theta _0^2),`$
$`E_\varphi `$ $`=`$ $`0,`$ (32)
$`E_z`$ $`=`$ $`i\theta _0E_0f^2(1f\rho ^2)e^{f\rho ^2}ge^{i\phi }+𝒪(\theta _0^3).`$
The magnetic field is
$$B_{}=0,B_\varphi =E_{},B_z=0.$$
(33)
The fields $`E_x`$ and $`E_z`$ are shown in Figs. 4 and 5. The dislocation seen in Fig. 5 for $`\rho \varsigma `$ is due to the factor $`1f\rho ^2`$ that arises in the paraxial approximation, and would, I believe, be smoothed out on keeping higher-order terms in the expansion (15).
The transverse electric field is radially polarized and vanishes on the axis. The longitudinal electric field is nonzero on the axis. Near the focus, $`E_zi\theta _0E_0`$ and the peak radial field is $`E_0/\sqrt{2e}=0.42E_0`$. For large $`z`$, $`E_{}`$ peaks at $`\rho =\varsigma /\sqrt{2}`$, corresponding to polar angle $`\theta =\theta _0/\sqrt{2}`$. For angles near this, $`\left|E_{}\right|\rho \left|f\right|^21/z`$, as expected in the far zone. In this region, the ratio of the longitudinal to transverse fields is $`E_z/E_{}i\theta _0f\rho r_{}/z`$, as expected for a spherical wave front.
The factor $`f^2`$ in the fields implies a Guoy phase shift of $`e^{2i\mathrm{tan}^1\varsigma }`$, which is twice that of the lowest-order linearly polarized beams.
It is noteworthy that the simplest axicon mode (32)-(33) is not a member of the set of Gaussian modes based on Laguerre polynomials in cylindrical coordinates (see, for example, sec. 3.3b of ).
### 2.6 Guided Axicon Beam
We could also consider the vector potential
$$A_r_{}\psi _0ge^{i\phi },A_\varphi =A_z=0,$$
(34)
which leads to the electric and magnetic fields
$$E_r=E_0fe^{f\rho ^2}ge^{i\phi },E_\varphi =0,E_z=i\theta _0f\rho E_r,B_r=0,B_\varphi =E_r,B_z=0,$$
(35)
and the potential
$$A_r_{}=0,A_\varphi \psi _0ge^{i\phi },A_z=0,$$
(36)
which leads to
$$E_r=0,E_\varphi =E_0fe^{f\rho ^2}ge^{i\phi },E_z=0,B_r=E_\varphi ,B_\varphi =0,B_z=i\theta _0\frac{12f\rho ^2}{2\rho }E_\varphi .$$
(37)
The case of eqs. (36)-(37) is unphysical due to the blowup of $`B_z`$ as $`r_{}0`$.
The fields of eqs. (34)-(35) do not satisfy $`𝐄=0`$ at $`r_{}=0`$, and so cannot correspond to a free-space wave. However, these fields could be supported by a wire, and represent a TM axicon guided cylindrical wave with a focal point. This is in contrast to guided plane waves whose radial profile is independent of $`z`$ . Guided axicon beams might find application when a focused beam is desired at a point where a system of lenses and mirrors cannot conveniently deliver the optical axis, or in wire-guided atomic traps . Figures 2 and 3 show the functional form of the guided axicon beam (35), when coordinate $`x`$ is reinterpreted as $`r_{}`$. |
warning/0003/hep-ph0003040.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The aim of this paper is to find the experimentally observable consequences of collective phenomena in the high energy hadrons inelastic collision. We will pay the main attention on the phase transitions, living out other possible interesting collective phenomena.
First of all, the statistics experience dictate that we should prepare the system to the phase transition. The temperature in a critical domain and the equilibrium media are just this conditions. It is evident, they are not trivial requirement considering the hadrons inelastic collision at high energies.
The collective phenomena by definition suppose that the kinetic energy of particles of media are comparable, or even smaller, than the potential energy of theirs interaction. It is a quite natural condition noting that, for instance, the kinetic motion may destroy, even completely at a given temperature $`T`$, necessary for the phase transition long-range order. This gives, more or less definitely, the critical domain.
The same idea as in statistics seems natural in the multiple production physics. We will assume (A) that the collective phenomena should be seen just in the very high multiplicity (VHM) events, where, because of the energy-momentum conservation laws, the kinetic energy of created particles can not be high.
We will lean at this point on the $`S`$-matrix interpretation of statistics . It based on the $`S`$-matrix generalization of the Wigner function formalism of Carruzers and Zachariazen and the real-time finite temperature field theory of Schwinger and Keldysh .
First of all, the $`n`$-particle partition function in this approach coincide with the $`n`$ particles production cross section $`\sigma _n(s)`$ (in the appropriate normalization condition). Secondly, $`\sigma _n(s)`$ can be calculated from the $`n`$-point Wigner function $`W_n(X_1,X_2,\mathrm{},X_n)`$. In the relativistic case $`X_k=(u,q)_k`$ are the 4-vectors. So, the external particles are considered as the ‘probes’ to measure the interacting fields state, i.e. the low mean energy of probes means that the system is ‘cold’.
The multiple production phenomena may be considered also as the thermalization process of incident particles kinetic energy dissipation into the created particles mass. From this point of view the VHM processes are highly nonequilibrium since the final state of this case is very far from initial one. It is known in statistics that such process aspire to be the stationary Markovian with high level of entropy production. In the case of complete thermalization the final state is equilibrium.
The equilibrium we will classify as the condition in frame of which the fluctuations of corresponding parameter are Gaussian. So, in the case of complete thermalization the probes should have the Gauss energy spectra. In other terms, the necessary and sufficient condition of the equilibrium is smallness of mean value of energy correlators . From physical point of view, absence of this correlators means depression of the macroscopic energy flows in the system.
The multiple production experiment shows that created particles energy spectrum is far from Gauss law, i.e. the final states are far from equilibrium. The natural explanation of this phenomena consist in presence of (hidden) conservation laws in the interacting Yang-Mills fields system: it is known that presence of sufficient number of first integrals in involution prevents thermalization completely.
But nevertheless the VHM final state may be equilibrium (B) in the above formulated sense. This means that the forces created by the non-Abelian symmetry conservation laws may be frozen during thermalization process (remembering its stationary Markovian character in the VHM domain). We would like to take into account that the entropy $`𝒮`$ of a system is proportional to number of created particles and, therefore, $`𝒮`$ should tend to its maximum in the VHM region .
One may consider following small parameter $`(\overline{n}(s)/n)<<1`$, where $`\overline{n}(s)`$ is the mean value of multiplicity $`n`$ at given CM energy $`\sqrt{s}`$. Another small parameter is the energy of fastest hadron $`\epsilon _{max}`$. One should assume that in the VHM region $`(\epsilon _{max}/\sqrt{s})0`$. So, the conditions:
$$\frac{\overline{n}(s)}{n}<<1,\frac{\epsilon _{max}}{\sqrt{s}}0$$
(1.1)
would be considered as the mark of the processes under consideration. We can hope to organize the perturbation theory over them having this small parameters. In this sense VHM processes may be ‘simple’, i.e. one can use for theirs description semiclassical methods.
So, considering VHM events one may assume that the conditions (A) and (B) are hold and one may expect the phase transition phenomena. The paper is organized as follows. In Sec.2 we will offer the qualitative picture of expected phase transition. In sec.3 we will describe quantitatively this phenomena to find the predictions for experiment. In Sec.4 we will give formal derivation of the formalism used in Sec.3.
## 2 Condensation and singularity at $`z=1`$
The $`S`$-matrix interpretation of statistics is based on following definitions. First of all, let us introduce the generating function :
$$T(z,s)=\underset{n}{}z^n\sigma _n(s).$$
(2.1)
Summation is performed over all $`n`$ up to $`n_{max}=\sqrt{s}/m`$ and, at finite CM energy $`\sqrt{s}`$, $`T(z,s)`$ is a polynomial function of $`z`$. Let us assume now that $`z`$ is sufficiently small and by this reason $`T(z,s)`$ depends on upper boundary $`n_{max}`$ weakly. In this case one may formally extend summation up to infinity and in this case $`T(z,s)`$ may be considered as a whole function. This possibility is important being equivalent of thermodynamical limit and it allows to classify the asymptotics over $`n`$ in accordance with position of singularities over $`z`$.
Let us consider $`T(z)`$ (prove will be given below) as the big partition function, where $`z`$ is ‘activity’. It is known that $`T(z)`$ should be regular in the circle of unite radii. The leftist singularity lie at $`z=1`$. This singularity is manifestation of the first order phase transition .
The origin of this singularity was investigated carefully in the paper . It was shown that position of singularities over $`z`$ depends on the number of particles $`n`$ in the system: the two complex conjugated singularity moves to the real $`z`$ axis with rising $`n`$ and in the thermodynamical limit $`n=\mathrm{}`$ they pinch point $`z=1`$. More careful analysis shows that if the system is equilibrium then $`T(z)`$ may be singular at $`z=1`$ and $`z=\mathrm{}`$ only.
The position of singularity over $`z`$ and the asymptotic behavior of $`\sigma _n`$ are related closely. Indeed, for instance, inserting into (2.1) $`\sigma _n\mathrm{exp}\{cn^\gamma \}`$ we find that $`T(z)`$ is singular at $`z=1`$ if $`\gamma <1`$. Generally, using Mellin transformation,
$$\sigma _n=\frac{1}{2\pi i}\frac{dz}{z^{n+1}}T(z)$$
(2.2)
This integral can be calculated expanding it in vicinity of $`z_c`$, where $`z_c`$ is smallest real positive solution of equation:
$$n=z\frac{}{z}\mathrm{ln}T(z).$$
(2.3)
Then integral (2.2) have following estimation:
$$\sigma _ne^{n\mathrm{ln}z_c(n)},z_c>1.$$
(2.4)
Therefore, to have the singularity at $`z=1`$ we should consider $`z_c(n)`$ as a decreasing function of $`n`$. On other hand, at constant temperature, $`\mathrm{ln}z_c(n)\mu _c(n)`$ is the chemical potential, i.e. is a work necessary for one particle creation. So, the singularity at $`z=1`$ means that the system is unstable: the less work is necessary for creation of one more particle if $`\mu (n)`$ is the decreasing function of $`n`$.
The physical explanation of this phenomena is following, see also . Generating function $`T(z)`$ have following expansion:
$$T(z)=\mathrm{exp}\{\underset{l}{}z^lb_l\},$$
(2.5)
where $`b_l`$ are known as the Mayer’s group coefficients. They can be expressed through the inclusive correlation functions and may be used to describe formation of droplets of correlated particles. So, if droplet consist from $`l`$ particles, then
$$b_le^{\beta \xi l^{(d1)/d}}$$
(2.6)
is the mean number of such droplets. Here $`\xi l^{(d1)/d}`$ is the surface energy of $`d`$-dimensional droplet.
Inserting this estimation into (2.5),
$$\mathrm{ln}T(z)\underset{l}{}e^{\beta (l\mu \xi l^{(d1)/d})},\beta \mu =\mathrm{ln}z.$$
(2.7)
First term in the exponent is the volume energy of droplet and being positive it try to enlarge the droplet. The second surface term try to shrink it. Therefore, singularity at $`z=1`$ is the consequence of instability: at $`z>1`$ the volume energy abundance leads to unlimited grow of the droplet.
## 3 Condensation and type of asymptotics over multiplicity
It is important for VHM experiment to have upper restriction on the asymptotics. We wish to show that $`\sigma _n`$ decrease faster than any power of $`1/n`$:
$$\sigma _n<O(1/n).$$
(3.1)
To prove this estimation one should know the type of singularity at $`z=1`$. The detailed derivation of the model used for this purpose we will give in subsequent section.
One can imagine that the points, where the external particles are created, form the gas system. Here we assume that this system is equilibrium, i.e. there is not in this system macroscopical flows of energy, particles, charges and so on.
The lattice gas approximation is used to describe such system. This description is quite general and did not depend on details. Motion of the gas particles leads to necessity sum over all distributions of the particles on cells. For simplicity we will assume that only one particle can occupy the cell.
So, we will introduce the occupation number $`\sigma _i=\pm 1`$ in the $`i`$-th cell: $`\sigma _i=+1`$ we have not particle in the cell and $`\sigma _i=1`$ means that the particle exist. Assuming that the system is equilibrium we may use the ergodic hypothesis and sum over all ‘spin’ configurations of $`\sigma _i`$, with restriction: $`\sigma _i^2=1`$. It is well known this restriction introduce the interactions.
Corresponding partition function in temperature representation
$$\rho (\beta ,H)=D\sigma e^{S\lambda (\sigma )}$$
(3.2)
where integration is performed over $`|\sigma (x)|\mathrm{}`$ and, considering the continuum limit, $`D\sigma =_xd\sigma (x)`$. The action
$$S_\lambda (\sigma )=𝑑x\left\{\frac{1}{2}(\sigma )^2\omega \sigma ^2+g\sigma ^4\lambda \sigma \right\}$$
(3.3)
where
$$\omega \left(1\frac{\beta _{cr}}{\beta }\right),g\frac{\beta _{cr}}{\beta },\lambda \left(\frac{\beta _{cr}}{\beta }\right)^{1/2}\beta H.$$
(3.4)
and $`1/\beta _{cr}`$ is the critical temperature.
### 3.1 Unstable vacuum
We start consideration from the case $`\omega >0`$, i.e. assuming that $`\beta >\beta _{cr}`$. In this case the ground state is degenerate if $`H=0`$. The extra term $`\sigma H`$ in (3.3) can be interpreted as the interaction with external magnetic field $`H`$. This term regulate number of ‘down’ spins with $`\sigma =1`$ and is related to the activity:
$$z^{1/2}=e^{\beta H},$$
(3.5)
i.e. $`H`$ coincide with chemical potential.
The potential
$$v(\sigma )=\omega \sigma ^2+g\sigma ^4,\omega >0,$$
(3.6)
has two minimums at
$$\sigma _\pm =\pm \sqrt{\omega /2g}.$$
If the dimension $`d>1`$ no tunnelling phenomena exist. But choosing $`H<0`$ the system in the right minimum (it correspond to the state without particles) becomes unstable. The system tunneling into the state with absolute minimum of energy.
The partition function $`\rho (\beta ,z)`$ becomes singular at $`H=0`$ because of this instability. The square root branch point gives
$$\mathrm{Im}\rho (b,z)=\frac{a_1(\beta )}{H^4}e^{a_2(\beta )/H^2},a_i>0.$$
(3.7)
Note, $`\mathrm{Im}\rho (b,z)=0`$ at $`H=0`$. Deforming contour of integration in (2.2) on the branch line,
$$\rho _n(\beta )=\frac{1}{\pi }_1^{\mathrm{}}\frac{dz}{z^{n+1}}\frac{8a_1\beta ^4}{\mathrm{ln}^4z}e^{4a_2\beta ^2/\mathrm{ln}^2z}.$$
(3.8)
In this integral
$$z_c\mathrm{exp}\left\{\frac{8a_2\beta ^2}{n}\right\}^{1/3}$$
(3.9)
is essential. This leads to following estimation:
$$\rho _ne^{3(a_2\beta ^2)^{1/3}1/3n^{2/3}}<O(1/n).$$
(3.10)
It is useful to note at the end of this section that
(i) The value of $`\rho _n`$ is defined by $`\mathrm{Im}\rho (b,z)`$ and the metastable states, decay of which gives contribution into $`\mathrm{Re}\rho (b,z)`$, are not important.
(ii) It follows from (3.9) that in the VHM domain
$$HH_c\mathrm{ln}z_c(1/n)^{1/3}0.$$
(3.11)
So, the calculations are performed for the week external field case, when the degeneracy is weekly broken. It is evident that the life time of the unstable (without particles) state is large in this case and by this reason the used semiclassical approximation is rightful. This is important consequence of (1.1).
### 3.2 Stable vacuum
Let us consider now $`\omega <0`$, i.e. $`\beta \beta _{cr}`$. Potential (3.6) have only one minimum at $`\sigma =0`$ in this case. Inclusion of external field shifts the minimum to the point $`\sigma _c=\sigma _c(H)`$. In this case the expansion in vicinity of $`\sigma _c`$ should be useful. In result,
$$\rho (\beta ,z)=\mathrm{exp}\{𝑑x\lambda \sigma _cW(\sigma _c)\},$$
(3.12)
where $`W(\sigma _c)`$ can be expanded over $`\sigma _c`$:
$$W(\sigma _c)=\underset{l}{}\frac{1}{l}\underset{k}{}\{dx_k\sigma _c(x_k;H)\}\stackrel{~}{b}_l(x_1,\mathrm{},x_l).$$
(3.13)
In this expression $`\stackrel{~}{b}_l(x_1,\mathrm{},x_l)`$ is the one-particle irreducible $`l`$-point Green function, i.e. $`\stackrel{~}{b}_l`$ is the virial coefficient. Then $`\sigma _c`$ can be considered as the effective activity of the correlated $`l`$-particle group.
The sum in (3.13) should be convergent and, therefore, $`|s_c|\mathrm{}`$ if $`|H|\mathrm{}`$. But in this case the virial decomposition is equivalent of the expansion over inverse density of particles . In the VHM region it is high and the mean field approximation becomes rightful. In result,
$$\sigma _c\left(\frac{|\lambda |}{4g}\right)^{1/3}:|s_c|\mathrm{}\mathrm{if}|\lambda |\mathrm{},$$
(3.14)
and
$$\rho (\beta ,z)e^{\frac{3|\lambda |^{4/3}}{(4g)^{1/3}}}\left\{12g\left(\frac{|\lambda |}{4g}\right)^{2/3}\right\}^{1/2}.$$
(3.15)
We can use this expression to calculate $`\rho _n`$. In this case
$$z_ce^{4gn^3}\mathrm{}\mathrm{at}n\mathrm{},$$
(3.16)
is essential and in the VHM domain
$$\rho _ne^{gn^4}<O(e^n).$$
(3.17)
This result is evident consequence of vacuum stability. It should be noted once more that the conditions (1.1) considerably simplify calculations.
## 4 Conclusion
In conclusion we wish to formulate once more the main assumptions.
(I). It was assumed first of all that the system under consideration is equilibrium. This condition may be naturally reached in the statistics, where one can wait the arbitrary time till the system becomes equilibrium. Note, in the critical domain the time of relaxation $`t_r(T_c/(TT_c))^\nu `$, $`(TT_c)+0`$, $`\nu >0`$, $`T_c`$ is the critical temperature.
We can not give the guarantee that in the high energy hadron collisions the final state system is equilibrium. The reason of this uncertainty is the finite time the inelastic processes and presence of hidden (confinement) constraints on the dynamics.
But we may formulate the quantitative conditions, when the equilibrium is hold . One should have the Gauss energy spectra of created particles. If this condition is hardly investigated on the experiment, then one should consider the relaxation of ‘long-range’ correlations. This excludes usage of relaxation condition for the ‘short-range’ (i.e. resonance) correlation
(II). The second condition consist in requirement that the system should be in the critical domain, where the (equilibrium) fluctuations of system becomes high. Having no theory of hadron interaction at high energies we can not define where is lie the ‘critical domain’ and even exist it or not.
But anyway, having the VHM ‘cold’ final state we can hope that the critical domain is achieved. Moreover, noting that the entropy reach its maximum at given incident energy, we can hope that the VHM state is equilibrium.
Acknowledgments
We are grateful to V.G.Kadyshevski for interest to discussed in the paper questions. We would like to note with gratitude that the discussions with E.Kuraev was interesting and important. |
warning/0003/hep-ph0003225.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Supersymmetry (SUSY), as one of the most attractive options beyond the standard model (SM), has been studied for the past few decades . From the theoretical point of view it offers a solution to the hierarchy problem. On the other hand, a lot of effort has been devoted to looking at the phenomenological consequences of SUSY both in low-energy processes and at high-energy colliders . One of the candidates for a realistic model is the minimal supersymmetric extension of the SM. In the SM, it is not possible to write down interactions which violate baryon number ($`B`$) or lepton number ($`L`$). In the SUSY version of the SM, particle spectrum is doubled and baryon number and lepton number are assigned to the supermultiplets, leading to $`\mathrm{\Delta }B=1`$ or $`\mathrm{\Delta }L=1`$ interactions in the Lagrangian. In the minimal supersymmetric standard model (MSSM) it is assumed that $`B`$ and $`L`$ are conserved quantum numbers. This is ensured by imposing a discrete multiplicative symmetry called $`R`$ parity which is defined as
$$R=(1)^{L+3B+2S}$$
where $`S`$ is the intrinsic spin of the particle.
It can be checked very easily that $`R`$ equals +1 for standard model particles and -1 for the superpartners. An immediate consequence of $`R`$-parity conservation is that the sparticles appear in pair at each interaction vertex. This leads to the fact that the lightest supersymmetric particle (LSP) is stable. The interactions of the LSP must be of weak strength because they are mediated by virtual sparticles which are known to be quite heavy (of the order of the electroweak scale). The most favorite candidate to become an LSP is the lightest neutralino and the search strategies for supersymmetry guided by the principle of $`R`$-parity conservation are to look for signals with large missing energy and momentum carried by an undetected neutralino . Also the LSP is a good candidate for the cold dark matter of the universe .
The conservation of $`R`$ parity, however, is not prompted by any strong theoretical reason, and theories where $`R`$ parity is violated through nonconservation of either $`B`$ or $`L`$ have been considered. Such scenarios can be studied by generalizing the MSSM superpotential to the following form :
$$W=W_{MSSM}+W_{\mathit{}},$$
(1)
with
$$W_{MSSM}=\mu \widehat{H}_1\widehat{H}_2+h_{ij}^l\widehat{L}_i\widehat{H}_1\widehat{E}_j^c+h_{ij}^d\widehat{Q}_i\widehat{H}_1\widehat{D}_j^c+h_{ij}^u\widehat{Q}_i\widehat{H}_2\widehat{U}_j^c$$
(2)
and
$$W_{\mathit{}}=\lambda _{ijk}\widehat{L}_i\widehat{L}_j\widehat{E}_k^c+\lambda _{ijk}^{}\widehat{L}_i\widehat{Q}_j\widehat{D}_k^c+ϵ_i\widehat{L}_i\widehat{H}_2+\lambda _{ijk}^{\prime \prime }\widehat{U}_i^c\widehat{D}_j^c\widehat{D}_k^c.$$
(3)
Here, $`\widehat{H}_1`$, $`\widehat{H}_2`$ are the $`\mathrm{SU}(2)`$ doublet Higgs superfields which give rise to the masses of down-type and up-type quark superfields, respectively, and $`\widehat{L}`$ $`(\widehat{Q})`$ denote lepton (quark) doublet superfields. $`\widehat{E}^c`$, $`\widehat{D}^c`$, $`\widehat{U}^c`$ are the singlet lepton and quark superfields. $`i,j,k`$ are the generational indices and we have suppressed the $`\mathrm{SU}(2)`$ and $`\mathrm{SU}(3)`$ indices. The $`\lambda _{ijk}`$ are antisymmetric in $`i`$ and $`j`$ while the $`\lambda _{ijk}^{\prime \prime }`$ are antisymmetric in $`j`$ and $`k`$. The first three terms in $`W_{\mathit{}}`$ violate lepton number and the last term violates baryon number. It is obvious that both the $`L`$ and $`B`$ violating terms cannot be present if the proton is stable. In order to get a large proton lifetime ($`10^{40}`$ s) it is sufficient to demand that either $`L`$ or $`B`$ be violated which in turn breaks $`R`$ parity. $`R`$-parity violation leads to considerable changes in the phenomenology. The most important consequence is that the LSP can decay now. Also, the lightest neutralino need not be the LSP because it is no longer a stable particle. The lepton number and baryon number violating terms mentioned above have received a lot of attention and constraints have been derived on these new couplings from present experimental data . Prospects of $`R`$-parity violation have been studied in the context of following present and future colliders: CERN $`e^+e^{}`$ collider LEP, DESY $`ep`$ collider HERA, $`p\overline{p}`$ at Fermilab Tevatron, CERN Large Hadron Collider (LHC), $`e^+e^{}`$ and $`e\gamma `$ Next Linear Collider (NLC) \[ReferencesReferences\]. Here we investigate the signatures of $`R`$-parity breaking at future $`e^{}e^{}`$ linear colliders. Our aim is to study the pair production of right selectrons ($`\stackrel{~}{e}_R`$) which will then decay into an electron and a neutralino. Finally the neutralino will decay into multifermions through different $`R`$-parity-violating couplings.
In this paper we shall discuss the $`R`$ violation in three separate categories for the convenience of the analysis. We will consider, in turn, $`W_{\mathit{}}`$ with either the $`\lambda `$, $`\lambda ^{}`$, or $`\lambda ^{\prime \prime }`$ terms existing in the superpotential at a time. The bilinear term $`ϵ_i\widehat{L}_i\widehat{H}_2`$ is also a viable agent for $`R`$-parity breaking which can induce vacuum expectation values for the sneutrino fields and generates a tree-level mass for one of the neutrinos . This scenario has been studied by several authors in the context of recent results from SuperKamiokande (SK) data on atmospheric neutrinos and attempts have been made to find the correlation between the given pattern of neutrino masses and mixings and collider signatures of supersymmetry . So far, no work has been reported which includes the study of $`R`$-parity violation through the bilinear term in the context of $`e^{}e^{}`$ colliders and we wish to discuss it in our future work which requires a separate analysis altogether .
The paper is organized as follows. In Sec. 2, we describe the physics goals of $`e^{}e^{}`$ colliders and their advantages and disadvantages from the point of view of a supersymmetry search. In Sec. 3 we will discuss the numerical results followed by our conclusions in Sec. 4.
## 2 Search for supersymmetry at $`e^{}e^{}`$ collider
As we know, the current $`e^+e^{}`$ collider at LEP is at the verge of its closing. Apart from putting some lower bounds on different SUSY particles, there has been no sign of new physics beyond the SM from LEP. Perhaps one can hope to see some signals beyond the SM at run II of Tevatron and, of course, the LHC, but the clean environment of the next generation $`e^+e^{}`$ linear collider will definitely complement the signatures from hadron colliders. Even if SUSY is discovered at LHC, NLC can be used as a machine for precision measurements for different SUSY parameters .
Before going on to the discussion of the supersymmetry search, let us first mention in brief the unique features of an $`e^{}e^{}`$ collider which establishes its importance in order to make model independent measurements at future high-energy physics experiments . First of all, it should be emphasized that at linear colliders the replacement of a beam of positrons with a beam of electrons can be achieved in a rather straightforward manner and can lead to the option of colliding electron beams.
At $`e^{}e^{}`$ colliders, the initial energy is well known and both $`e^{}`$ beams can be highly polarized so that the initial states are specified. The backgrounds are, in general, extremely suppressed and they can be reduced further with specific choices of the beam polarizations. However, the total electric charge $`Q`$ and total lepton number $`L`$ of $`e^{}e^{}`$ colliders forbid the pair production of most of the superpartners by virtue of total charge and lepton number conservation. This is one disadvantage of $`e^{}e^{}`$ colliders where only selectrons can be pair produced through the exchange of a Majorana neutralino in the $`\mathrm{t}`$ and $`\mathrm{u}`$ channels as shown in Fig. 1. In contrast, at $`e^+e^{}`$ colliders, selectron pair production occurs through $`\mathrm{s}`$-channel $`\gamma `$ and $`Z`$ exchange as well as through $`\mathrm{t}`$-channel $`\stackrel{~}{\chi }_i^0`$ exchange. The interference between the $`\mathrm{s}`$\- and $`\mathrm{t}`$-channel diagrams is always destructive for $`\sqrt{s}>m_Z`$ . In the $`e^{}e^{}`$ mode, since the $`\mathrm{u}`$-channel diagram is present along with the $`\mathrm{t}`$-channel diagram and the interference between them is constructive, the production cross section is always larger compared to the $`e^+e^{}`$ mode. This cross section can be further increased by choosing the initial electron beam polarization properly. It has been shown that the right selectron pair production cross section is largest for the right-polarized initial electron beam. This can be explained from the fact that in most of the MSSM parameter space the LSP is B-ino dominated, which has a larger coupling with $`e_R\stackrel{~}{e}_R`$ compared to $`e_L\stackrel{~}{e}_L`$. Furthermore, it turns out that the selectron pair production cross section for the unpolarized initial state is smaller than that of right-polarized electron beams.
Another important feature of the $`e^{}e^{}`$ collider is its behavior near threshold which shows a sharp rise in the selectron pair production cross section . This enables one to measure the selectron masses very accurately. In contrast, at an $`e^+e^{}`$ collider the threshold measurement is rather poor, which compels one to determine the $`\stackrel{~}{e}`$ mass (with an error of few GeV) from the measurement of the electron end point energy . The study of slepton flavor violation can also be done very effectively in an $`e^{}e^{}`$ collider.
In Fig. 2, we present contours for the cross section (in fb) for the production of $`\stackrel{~}{e}_R^{}\stackrel{~}{e}_R^{}`$ final states in the ($`\mu ,M_2`$) plane for $`\mathrm{tan}\beta =2,20,40`$ and $`\sqrt{s}=500`$ GeV. The mass of the right selectron is assumed to be 150 GeV for the plots in the left column and 200 GeV for the plots in the right column.
The explanation of the variation of cross section with the parameters which appear in the neutralino mass matrix, as shown in Fig. 2, is as follows. The area ruled out by LEP-2 represents the region which is disallowed by the chargino search at LEP-2 and corresponds to a mass of the lighter chargino ($`\stackrel{~}{\chi }_1^\pm `$) less than 98 GeV . This limit comes purely from kinematic considerations and does not depend on whether $`R`$ parity is conserved or violated. The area which is marked as $`X`$ in the figure is not allowed because here the selectrons become lighter than the LSP and hence the selectron decaying to the lightest neutralino is forbidden <sup>1</sup><sup>1</sup>1In $`R`$-parity-violating models, it is also possible that the selectron, rather than the lightest neutralino, is the LSP, and can decay directly into two leptons (quarks) through the $`\lambda _{ijk}(\lambda _{ijk}^{})`$ couplings, respectively.. Since we are considering right selectron pair production, the contribution to the cross section comes mainly from the lightest neutralino which is dominated by a B-ino over a large part of the parameter space. Here we assume the grand unified theory (GUT) relationship between the $`\mathrm{SU}(2)`$ and $`\mathrm{U}(1)`$ gaugino soft mass parameters $`M_2`$ and $`M_1`$, $`M_1=\frac{5}{3}\mathrm{tan}^2\theta _WM_2`$. As the value of $`M_2`$ increases, the lightest neutralino starts becoming more and more B-ino dominated and hence the strength of the $`e_R\stackrel{~}{e}_R\stackrel{~}{\chi }_1^0`$ coupling increases at the same time. Also, the amplitude in this case requires a $`\mathrm{t}`$-channel neutralino mass insertion. These two effects combined together lead to an increase in the cross section when $`M_2`$ is increased for a fixed value of $`\mu `$ . This feature is evident from Fig. 2. For lower values of $`\mu `$, the B-ino component in the lightest neutralino starts decreasing which means a fall in the cross section and hence in order to get the same cross section the value of $`M_2`$ (consequently the value of $`M_1`$) must be increased. With the increase in selectron mass the available phase space reduces and in order to get the same cross section as in the left column one must go to higher values of $`M_2`$.
The decay of right selectron yields following final state:
$`e^{}e^{}\stackrel{~}{e}_R^{}\stackrel{~}{e}_R^{}e^{}e^{}\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0.`$ (4)
This will give rise to two like-sign electrons and large $`p_T/`$ signature. This kind of a signal as shown in Eq. (4) and the relevant backgrounds have been well studied .
In light of the above discussion, the next question which comes to mind is what could be the potential signatures at an $`e^{}e^{}`$ collider when $`R`$ parity is violated. Recently, the effect of $`R`$ parity violation has been studied for the production process $`e_L^{}e_R^{}\stackrel{~}{e}_L^{}\stackrel{~}{e}_R^{}`$ . In this work we will consider the pair production of right selectrons assuming $`90\%`$ right-polarized electron beams because of the larger cross section in this case. The subsequent analysis will not depend on the choice of initial electron polarization. As we will see in the following section, since the lightest neutralino will decay, it will lead to multilepton final states with missing energy almost free from standard model backgrounds. However, right selectrons can also decay into heavier neutralino states, if it is allowed by kinematics. In that case, the cascade decays of heavier neutralino will produce more complex signals. For the simplicity of our analysis, we will not consider such decay patterns here.
## 3 Decay of $`\stackrel{~}{\chi }_1^0`$ and associated signals
In this section we will discuss the possible signatures arising from the decay of the LSP through different $`R`$-parity-violating interactions, through sfermion (sleptons and squarks) exchange diagrams. These Feynman diagrams and the amplitudes can be found in the literature . Here, we make the assumption that of all the couplings which violate $`R`$ parity, only one is dominant at a time, which is motivated from the fact that in the SM top quark Yukawa coupling is much larger than the others. Furthermore, we assume these couplings to be much smaller than the gauge couplings, though we require them to be large enough to make the LSP decay inside the detector. A generic $`R`$-parity-violating coupling should be larger than $`10^5`$ to satisfy the above requirement . In our subsequent analysis we take these couplings in the range $`10^1`$ \- $`10^2`$. If the $`R`$-parity-violating operator is of the type $`LLE^c`$, the final states will have two charged leptons and a neutrino. The flavor of these leptons are determined by the type of $`\lambda _{ijk}`$ coupling. If the $`R`$-parity-violating operator is of the type $`LQD^c`$, the final sates will have either one charged lepton or a neutrino associated with two quarks. Finally in the presence of baryon number violating coupling $`U^cD^cD^c`$, the final state will have three quarks. Throughout this analysis we assume 250 GeV left-slepton mass \[sneutrino mass is related to left-slepton mass through the $`\mathrm{SU}(2)`$ relation\] and 500 GeV squark mass. All the squarks have been assumed to be degenerate in mass. In our parton level Monte Carlo analysis we treat quarks/partons as jets, and the direction of jets is same as that of the initial quarks/partons. We impose following selection criteria for these leptons and jets:
$`p_T^{\mathrm{}}>5\mathrm{GeV},|\eta _{\mathrm{}}|<3,`$ (5)
$`p_T^j>15\mathrm{GeV},|\eta _j|<3.`$ (6)
We merge two jets into a single jet if their angular separation $`\mathrm{\Delta }R_{jj}<0.7`$, where $`(\mathrm{\Delta }R_{jj})^2(\mathrm{\Delta }\eta )_{jj}^{}{}_{}{}^{2}+(\mathrm{\Delta }\varphi _{jj})^2`$, $`\mathrm{\Delta }\eta _{jj}`$ and $`\mathrm{\Delta }\varphi _{jj}`$ being the difference of pseudorapidities and azimuthal angles, respectively, corresponding to two jets. The lepton is isolated from a jet if $`\mathrm{\Delta }R_{jl}>0.4`$, where $`\mathrm{\Delta }R_{jl}`$ is defined in the same way as above.
### 3.1 Signals from $`\lambda `$-type couplings
Let us now discuss the signals which can be looked for when $`R`$ parity is violated through the terms of the type $`\lambda LLE^c`$. The pair-produced LSPs from the decay of the two right selectrons will lead to the final state consisting of $`e^{}e^{}+4\mathrm{}^\pm +p_T/`$. The flavor of the leptons coming from the neutralino decay will depend on the particular type of coupling involved. For example, $`\lambda _{123}`$ coupling gives
$`\stackrel{~}{\chi }_1^0\nu _e\mu ^{}\tau ^+,e^{}\nu _\mu \tau ^+,\overline{\nu }_e\mu ^+\tau ^{},e^+\overline{\nu }_\mu \tau ^{},`$ (7)
with equal probabilities. Here, for simplicity we have considered a common value for all $`\lambda `$-type couplings taken to be $`0.07(m_{\stackrel{~}{e}}/100\mathrm{GeV})`$, close to the existing indirect bounds relevant for most of those couplings. In order to tag the lepton flavor one must multiply the signal cross section with the efficiency of the corresponding lepton flavor identification.
Since there are two neutrinos in the final states, reconstructing the mass of the LSP in such a case is not possible. However, the kind of final state mentioned above is spectacular in the sense that it is free from standard model background and permits easy detection at a 500 GeV $`e^{}e^{}`$ collider.
In Fig. 3, we have shown the transverse momentum ($`p_T`$) distribution of the charged leptons produced in the final state for $`M_{\stackrel{~}{e}_R}=150`$ GeV, and the following the set of input parameters $`\mu =450`$ GeV, $`M_2=200`$ GeV and $`\mathrm{tan}\beta =2`$. For this set of parameter points $`M_{\stackrel{~}{\chi }_1^0}=103`$ GeV and $`M_{\stackrel{~}{\chi }_2^0}=206`$ GeV. For later studies of the distributions we will use this set of input parameters. It is easy to see from this distribution that all six leptons survive the $`p_T^{\mathrm{}}>5`$ GeV cut. Out of six leptons, two come from the decay of $`\stackrel{~}{e}_R`$; the remaining four leptons come from the decay of $`\stackrel{~}{\chi }_1^0`$.
We display in Table 1, some representative values of the cross sections in order to get an idea about the strength of the signal. In obtaining these numbers, we required six leptons satisfying the criteria given in Eq. (5), and in addition imposed that $`p_T/>15`$ GeV. The $`p_T/`$ requirement ensures that the signal is SM background free. Two values of the right selectron mass, namely, $`m_{\stackrel{~}{e}_R}=150`$ (GeV) and $`m_{\stackrel{~}{e}_R}=200`$ (GeV), have been considered for the calculation of the cross sections. We have considered the actual branching ratios of the decays of right selectrons including the direct decays through $`R`$-parity-violating couplings as well as decays into heavier neutralinos. It has already been mentioned that if the decays into heavier neutralino states are allowed kinematically, they will lead to more complex signals which we have not considered in this work. For a fixed value of $`\mu `$ and $`\mathrm{tan}\beta `$, with increasing $`M_2`$, the LSP mass increases; hence the MSSM decay of $`\stackrel{~}{e}_R`$ decreases because of phase-space suppression, favoring the direct decay of $`\stackrel{~}{e}_R`$ through $`R`$-parity-violating $`\lambda _{231}`$, which in turn reduces our signal. As is evident from this table, large cross sections may be obtained for a considerable region of the parameter space and with a projected integrated luminosity of 50 $`\mathrm{fb}^1`$ at an $`e^{}e^{}`$ collider one could see some thousands of events. It must be noted at this point that if taus are produced in the final state, they would decay mainly into hadrons, but that requires a separate analysis.
### 3.2 Signals from $`\lambda ^{}`$-type couplings
The decay pattern of the LSP changes as we go on to the $`R`$-parity-violating couplings of the type $`\lambda ^{}LQD^c`$. For example, $`\lambda _{123}^{}`$ coupling gives
$`\stackrel{~}{\chi }_1^0\nu _es\overline{b},e^{}cb,\overline{\nu }_e\overline{s}b,e^+\overline{c}b.`$ (8)
As before, we again consider a common value for all $`\lambda ^{}`$-type couplings. To identify the final state flavors one has to take into account the reduction in cross section due to flavor tagging efficiency. It should be mentioned at this point that unlike the $`\lambda `$ case here all final states are not equiprobable. We categorize the signals in the following manner. All these states are assumed to be accompanied by two like-sign dielectrons arising from $`\stackrel{~}{e}_R`$ decay: (1) $`2\mathrm{}^\pm `$ \+ jets; both $`\stackrel{~}{\chi }_1^0\mathrm{}^\pm jj`$; (2) jets + $`p_T/`$ ; both $`\stackrel{~}{\chi }_1^0\nu jj`$; (3) $`\mathrm{}^\pm +\mathrm{jets}+p_T/`$ ; one $`\stackrel{~}{\chi }_1^0\mathrm{}^\pm jj`$, the other $`\stackrel{~}{\chi }_1^0\nu jj`$.
The last channel will be enhanced by a combinatoric factor of 2. We have folded the cross section with the branching fraction of the LSP. The selection cuts (as discussed earlier) are applied to the leptons and jets. After the energy ordering ($`E_{j_1}>E_{j_2}>E_{j_3}>E_{j_4}`$), we study the jet $`p_T`$ distribution as shown in Fig. 4. These jets and charged leptons are also associated with large $`p_T/`$ arising from neutrinos for channels ($`2`$) and ($`3`$) listed above. The $`p_T/`$ distribution is shown in Fig. 5. The distribution $`a`$ corresponds to the case when both the LSP decays into the $`\nu jj`$ channel, where as $`b`$ represents the $`p_T/`$ distribution when one of the LSP decays through the $`\nu jj`$ mode and the other one through the $`\mathrm{}^\pm jj`$ mode.
Finally in Table 2 we give cross sections for signals for two $`\stackrel{~}{e}_R`$ masses 150 GeV and 200 GeV. We required four jets and, respectively, four, three, or two leptons satisfying the criteria given in Eqs. (6) and (5), for channels ($`1`$), ($`2`$), and ($`3`$) listed above. In addition, $`p_T/>15`$ GeV is imposed for channels ($`2`$) and ($`3`$). Cross sections for heavier right-selectron mass ($`=200`$ GeV) are lower than the corresponding quantities for 150 GeV $`\stackrel{~}{e}_R`$ mass, just because of a lack of enough phase space. The difference in the three cross sections in each row can be explained from the branching ratio of $`\stackrel{~}{\chi }_1^0`$ in two different channels $`\mathrm{}^\pm jj`$ and $`\nu jj`$. The inputs remain same as in Table 1. The cross sections for these various channels are fairly large over a wide region of parameter space which is accessible in a 500 GeV $`e^{}e^{}`$ collider. Signals corresponding to $`e^{}e^{}+\mathrm{jets}+p_T/`$ and $`e^{}e^{}+\mathrm{}^\pm +\mathrm{jets}+p_T/`$ final states may have the standard model background coming from $`W^{}W^{}ZZ`$ production. But this cross section is found to be too low ($`<40`$ fb) and does not affect the signal in a significant way.
If the produced LSP is highly relativistic, then its decay products will be confined within a narrow cone around the direction of the LSP. In that case, the lepton (decaying from LSP) in a particular hemisphere is identified and its invariant mass is constructed with all jets in the same hemisphere. A similar thing is done in the opposite hemisphere. Then we demand that these two invariant masses should lie within 10 GeV of each other. If these two invariant masses are equal or nearly equal, we can say that they arise from the same parent particle. In Fig. 6(a) we represent such an invariant mass distribution, which shows a distinct peak at the LSP mass ($`=103`$ GeV). In order to get an estimate of the mass resolution, we have used Gaussian smearing of the energies of jets and leptons to “mimic” the response of a detector:
$`\mathrm{\Delta }E_j/E_j=0.4/\sqrt{E_j}+0.02,\mathrm{\Delta }E_l/E_l=0.15/\sqrt{E_l}+0.01.`$ (9)
In Fig. 6(b), we show the mass distribution after energy smearing. The LSP mass determined in this way has resolution $`\mathrm{\Delta }M/M=4\%`$. It should also be noted that for about 80$`\%`$ of the total events the mass reconstructed in both sides lies within 10 GeV of each other. In principle one can also reconstruct the selectron mass in this way, but a more precise determination can be done by a threshold scan .
### 3.3 Signals from $`\lambda ^{\prime \prime }`$-type couplings
Finally, the presence of $`\lambda ^{\prime \prime }`$ in the superpotential can induce $`B`$ number violating decay of the LSP. In this case, the LSP will simply decay into three hadronic jets:
$`e^{}e^{}\stackrel{~}{e}_R^{}+\stackrel{~}{e}_R^{}e^{}+e^{}+\stackrel{~}{\chi }_1^0+\stackrel{~}{\chi }_1^0e^{}e^{}+6\mathrm{jets},`$ (10)
where the sets of three jets have invariant mass peaking at the neutralino mass (assuming all jets are seen). As before, we impose the selection cuts on leptons and at least four jets. From the $`p_T`$ distribution of six jets in Fig. 7 it is clear that for this value of the LSP mass ($`=103`$ GeV), most of the jets are hard enough to satisfy the jet trigger requirement as discussed previously. From the jet number distributions in Fig. 8, we see that most of the time the cross section prefers to peak at the five-jet channel ($`44.69\%`$ of the events), followed by the six-jet ($`42.86\%`$ of the events) and four-jet ($`11.83\%`$ of the events) channels. The three-jet fraction of the cross section is less than $`0.5\%`$. Imposition of $`p_T^j>15`$ GeV and $`|\eta _j|<3`$ cuts on the jets reduces the jet number. Finally we also merge two jets into a single jet if their angular separation $`\mathrm{\Delta }R_{jj}<0.7`$. The probability of jet merging is highly dependent on the mass of the parent particle from which the jets originate and also on $`\sqrt{s}`$. The larger the boost of the parent particle, the higher the probability of jet merging. In this case, a 103 GeV LSP is produced from the decay of a 150 GeV right selectron. Each of these LSPs then decays into three jets with a reasonable boost, leading six jets to merge into five jets and occasionally into four and three jets.
In Table 3 we give signal cross sections for some representative values of parameters. In this case, we assume the squark mass to be 500 GeV, which enters as a propagator in the decay LSP. One can also reconstruct the LSP mass using the following strategy: selecting the hardest jet in the final state, its invariant mass is then constructed with all other jets in that hemisphere. A similar thing is done in the opposite hemisphere. Then we demand that these two invariant masses should lie within 10 GeV of each other. If these two invariant masses are equal or nearly equal, we can say that they arise from the same parent particle. Though we will not present here the invariant mass distribution, similar kinds of studies have been done by other authors and also by the ALEPH Collaboration in their study of the (now defunct) four-jet anomaly .
Before we conclude, we would like to mention the possible SM backgrounds in this case. We have earlier found that most of the time the signal cross section prefers to peak around five and six jets, free from any SM background. However, there is a small fraction of cross section that goes into four-jet channels, which is less important for our purpose as far as the SM backgrounds are concerned. This particular signal has SM backgrounds from (a) $`e^{}e^{}e^{}e^{}ZZ`$, (b) $`e^{}e^{}e^{}e^{}Z^{}Z`$, and (c) $`e^{}e^{}e^{}e^{}Z^{}Z^{}`$, with hadronic decay of $`Z`$ (assuming all jets are seen).
One can make a rough estimate for this background. After putting the selection cuts and including the relevant branching ratios the cross section for $`e^{}e^{}e^{}e^{}Z`$ is of the order of $`100`$ fb. This cross section will get electroweak suppression if another $`Z`$ boson is radiated; moreover, the $`\mathrm{Br}.(Zq\overline{q})`$ will further reduce this. After all these, if this background is still comparable to the signal, then this can be eliminated by imposing the condition that the pair of dijet invariant mass $`M_{jj}`$ should not peak around $`M_Z`$. However, this may reduce the signal cross section in the region of parameter space where the LSP mass is nearly degenerate with $`M_Z`$. The detailed calculation of the other two backgrounds \[(b) and (c)\] is very cumbersome and we will not perform this here. In this case, our main thrust will be to count the jets in the final state (associated with two electrons) to distinguish it from the background.
## 4 Conclusions
We have discussed the pair production of right selectrons at a 500 GeV $`e^{}e^{}`$ linear collider in the $`R`$-parity-violating supersymmetric model. The decay of right selectrons can yield a final state with an electron and a neutralino, mostly the LSP. Hence, we have two like-sign dielectrons and neutralinos in the final state. We have assumed that $`R`$ parity is weakly violated and thus only the LSP will decay into multifermion states. Different possibilities have been considered and it seems that rather optimistic signals can be seen for this kind of model. The decay of the LSP gives charged leptons, jets, and neutrinos in the final state. The behavior of these leptons, jets, and missing transverse momentum (mainly due to neutrinos) has been analyzed using a parton level Monte Carlo event generator. This also enables us to study the approximate distributions for different kinematic variables of leptons and jets. The decay of the LSP through $`L`$-number-violating coupling ($`\lambda `$) leads to a very distinct signal with hard isolated leptons and large missing transverse momentum. There are no SM processes which can mimic this signal. Similarly, for $`\lambda _{ijk}^{}`$ couplings, the signal basically consists of charged leptons, multiple jets, and/or missing transverse momentum. In addition to this, the Majorana nature of the LSP gives rise to like-sign dilepton signals with practically no SM backgrounds. It has been demonstrated that the reconstruction of the lepton-jet invariant mass can give a rough estimate for the LSP mass. For $`\lambda _{ijk}^{\prime \prime }`$ coupling, the final state will have multiple jets associated with like-sign dielectrons. We have shown that proper jet counting is required to distinguish the signal from the SM backgrounds. In this case also it might be possible to determine the LSP mass from the jet invariant mass reconstruction.
Acknowledgments
The authors are grateful to Biswarup Mukhopadhyaya and Sreerup Raychaudhuri for helpful discussions. |
warning/0003/quant-ph0003036.html | ar5iv | text | # 1 a) The quantized Prisoner’s Dilemma, as described in []. The pair of qubits are prepared in the unentangled state |𝐶𝐶⟩ and then sent through the entangling gate 𝐽̂. Players 𝐴 and 𝐵 then apply their local unitary operations 𝐴̂⊗𝐼̂ and 𝐼̂⊗𝐵̂, respectively. A gate inverse to 𝐽̂ is applied before the final measurement. b) The Prisoner’s Dilemma pay-off table chosen in [].
Comment on “Quantum Games and Quantum Strategies”
In a recent Letter, Eisert et al. presented a quantum mechanical generalization of Prisoner’s Dilemma. In the classical form of this game, rational analysis leads the two players to ‘defect’ against one-another in a mutually destructive fashion . A central result of Eisert et al.’s Letter is the observation that their quantum variant, illustrated in Figure 1, can avoid the ‘dilemma’: the mutually destructive outcome is replaced with an effectively cooperative one. Specifically, it is asserted that the maximally entangled game’s unique Nash equilibrium occurs when both players apply the strategy $`\widehat{Q}=i\sigma _z`$, yielding a pay-off equivalent to cooperative behaviour in the classical game.
In this Comment we show that their observation is incorrect. The mistake follows from the following erroneous assertion:
> It proves to be sufficient to restrict the strategic space to the 2-parameter set of unitary 2 x 2 matrices,
>
> $$\widehat{U}(\theta ,\varphi )=\left(\begin{array}{cc}e^{i\varphi }\mathrm{cos}\theta /2& \mathrm{sin}\theta /2\\ \mathrm{sin}\theta /2& e^{i\varphi }\mathrm{cos}\theta /2\end{array}\right),$$
> (1)
> with $`0\theta \pi `$ and $`0\varphi \pi /2`$.
Here we explicitly consider the complete set of *all* local unitary operations (*i.e.* all of $`SU(2)`$), finding that the properties of the game are wholly different: the strategy $`\widehat{Q}`$ is not an equilibrium; indeed, there is no equilibrium in the space of deterministic quantum strategies. Moreover, it seems unlikely that the restriction to the set $`\widehat{U}(\theta ,\varphi )`$ can reflect any reasonable physical constraint (limited experimental resources, say) because this set is not closed under composition. An ideal counter strategy to $`\widehat{Q}`$, for example, is $`i\sigma _x`$, which is equal to $`\widehat{U}(0,\pi /2)\widehat{U}(\pi ,0)`$. The game of therefore does not constitute a reasonable variant of the general case we consider here.
We will write the operations applied by the players in the form $`\widehat{A}\widehat{B}`$, where $`\widehat{A}`$ is applied to the qubit controlled by $`A`$ and similarly for $`\widehat{B}`$. Suppose that player $`A`$ applies transformation $`\widehat{X}`$ to her qubit, prepared as the first qubit in the maximally entangled state
$$\widehat{J}|CC=\frac{1}{\sqrt{2}}(|CC+i|DD),$$
(2)
where $`\widehat{J}=\mathrm{exp}\{i\pi \widehat{D}\widehat{D}/4\}`$ and $`\widehat{D}=i\sigma _y`$ is the ‘defect’ matrix of . The most general $`\widehat{X}SU(2)`$ is of the form $`\widehat{X}=(x_{ij})`$, where $`x_{11}=x_{22}^{}`$, $`x_{12}=x_{21}^{}`$ and $`det\widehat{X}=1`$. Therefore, $`A`$ produces the state $`(\widehat{X}\widehat{I})\widehat{J}|CC=(I\widehat{Y})J|CC`$ for $`\widehat{Y}=(y_{ij})SU(2)`$, where $`y_{11}=x_{11}`$ and $`y_{12}=ix_{12}`$. In other words, any unitary transformation which $`A`$ applies locally to her qubit is actually equivalent to a unitary transformation applied locally by $`B`$. Consequently, if $`B`$ were to choose $`\widehat{D}\widehat{Y}^{}`$, we would have a final state $`\widehat{J}^{}(\widehat{X}\widehat{D}\widehat{Y}^{})\widehat{J}|CC=\widehat{J}^{}(\widehat{I}\widehat{D}\widehat{Y}^{}\widehat{Y})\widehat{J}|CC=|CD`$, the optimal outcome for $`B`$. Thus, for any given strategy of $`A`$, there is an ideal counter-strategy for $`B`$, and vice-versa; there is no Nash equilibrium of the kind suggested by Eisert et al. .
To obtain such equilibria we must extend the abilities of the players: it suffices to allow them to make probabilistic choices (rather than the full formalism of completely positive maps considered in footnote 14 of ). Suppose that $`A`$ adopts the following strategy: she will choose a move $`\widehat{X}SU(2)`$ at random with respect to the Haar measure on $`SU(2)`$. If $`B`$ responds with $`\widehat{Y}_0SU(2)`$ then the probability that outcome $`i\{CC,CD,DC,DD\}`$ will be measured is
$`P_i(\widehat{Y}_0)`$ $`={\displaystyle _{SU(2)}}\left|i\left|\widehat{J}^{}(\widehat{X}\widehat{Y}_0)\widehat{J}\right|CC\right|^2𝑑\widehat{X}`$
$`={\displaystyle _{SU(2)}}\left|i\left|\widehat{J}^{}(\widehat{X}\widehat{X}_0^{}\widehat{X}_0I)\widehat{J}\right|CC\right|^2𝑑\widehat{X}`$
$`={\displaystyle _{SU(2)}}\left|i\left|\widehat{J}^{}(\widehat{X}\widehat{I})\widehat{J}\right|CC\right|^2𝑑\widehat{X}`$
$`=P_i(\widehat{I}),`$ (3)
where $`\widehat{X}_0SU(2)`$ is chosen such that $`(\widehat{X}_0\widehat{I})\widehat{J}|CC=(\widehat{I}\widehat{Y}_0)\widehat{J}|CC`$ and we have used the right invariance of the Haar measure, assumed to be normalized such that the volume of $`SU(2)`$ is $`1`$. Thus, $`B`$’s choice of strategy does not matter; regardless of his choice, his expected pay-off is simply an unbiased average over the classical pay-offs. Therefore, the situation where both players adopt this random strategy is a Nash equilibrium: neither player can improve his or her payoff by unilaterally altering choice of strategy.
As a final point, we note that one can construct Prisoner’s Dilemma-type pay-off tables for which the quantum equilibrium pay-off we describe above is below the classical equilibrium pay-off, or above it, or even above the classical cooperative pay-off. In this last case the ‘dilemma’ may be said to have been removed . To this extent the behavior of the quantum generalization is qualitatively different from the classical case, although in a way that is perhaps less surprising than originally suggested by Eisert *et al.*
We thank Neil Johnson for helpful conversations. SB is supported by EPSRC. PH acknowledges the support of the Rhodes Trust.
Simon C. Benjamin and Patrick M. Hayden
Centre for Quantum Computation
University of Oxford
Clarendon Laboratory, Parks Road
Oxford OX1 3PU, UK
PACS numbers: 03.67.-a, 02.50.Le, 03.65.Bz |
warning/0003/hep-ph0003221.html | ar5iv | text | # 1 Synopsis
## 1 Synopsis
Essential elements of the fundamental constituents of matter and their interactions have been discovered in the past three decades by operating $`e^+e^{}`$ colliders. A coherent picture of the structure of matter has emerged, that is adequately described by the Standard Model, in many of its facets at a level of very high accuracy. However, the Standard Model does not provide a comprehensive theory of matter. Neither the fundamental parameters, masses and couplings, nor the symmetry pattern can be explained, but they are merely built into the model. Moreover, gravity is not incorporated at the quantum level. First steps to solutions of these problems are associated with the unification of the electroweak and the strong forces, and with the supersymmetric extension of the model which provides a bridge from the presently explored energy scales up to scales close to the Planck mass.
Two strategies can be followed to enter into the area beyond the Standard Model. $`(i)`$ Properties of the particles and forces within the Standard Model will be affected by new energy scales. Precision studies of the top quark, the electroweak gauge bosons and the Higgs boson can thus reveal clues to the physics beyond the Standard Model. $`(ii)`$ Above the mass thresholds, new phenomena can be searched for directly and studied thoroughly so that the underlying basic theories can be reconstructed.
In this dual approach a variety of fundamental problems still remain to be solved within the Standard Model , demanding experiments at energies beyond the range of existing accelerators.
$`(a)`$ The mass of the top quark is much larger than the masses of the electroweak gauge bosons. Understanding the rôle of this particle in Nature is therefore an important goal for the future. In the $`t\overline{t}`$ threshold region of $`e^+e^{}`$ collisions the top quark mass can be measured to an accuracy better than 200 MeV. This is a desirable goal since a future theory of flavor dynamics will lead to relations among the lepton/quark masses and mixing angles in which the heavy top quark is expected to play a key rôle. In addition, stringent tests in the electroweak and Higgs sector of the Standard Model can be carried out when the top mass is known very accurately. Helicity analyses of the $`t\overline{t}`$ production vertex and the $`t`$ decay vertex will determine the magnetic dipole moments of the top quark and the chirality of the ($`tb`$) decay current at the per–cent level. Bounds on the $`𝒞𝒫`$ violating electric dipole moments of the $`t`$ quark can be set to $`10^{18}`$ e cm.
$`(b)`$ Studying the dynamics of the electroweak gauge bosons is another important task at high energy $`e^+e^{}`$ colliders. The form and the strength of the triple and quartic couplings of these particles are uniquely predicted by the non–abelian gauge symmetry of the theory, defining the electroweak charges, the magnetic dipole moments and the electric quadrupole moments of the $`W^\pm `$ bosons in the static limit. Tests of these fundamental symmetry concepts can be performed at an accuracy of $`10^3`$ down to $`10^4`$.
$`(c)`$ A high–luminosity $`e^+e^{}`$ collider with an energy between 300 and 500 GeV will be an ideal instrument to search for Higgs particles throughout the mass range characterized by the scale of electroweak symmetry breaking, and to investigate their properties. The intermediate Higgs mass range below $`200`$ GeV is the theoretically preferred region. In this scenario Higgs particles remain weakly interacting up to the scale of grand unification, thus providing a basis for the renormalization of the electroweak mixing angle $`\mathrm{sin}^2\theta _W`$ from the GUT symmetry value 3/8 down to the experimentally observed value close to 0.2. Once the Higgs particle is found, its properties can be studied thoroughly, the external quantum numbers $`𝒥^{𝒫𝒞}`$ and the Higgs couplings, including the self–couplings of the particle. The measurements of these couplings are the necessary ingredient to establish the Higgs mechanism sui generis experimentally.
Even though many facets of the Standard Model are experimentally supported at a level of very high accuracy, extensions should nevertheless be anticipated as argued before. The next generation of accelerators can shed light on three domains in the area beyond the Standard Model.
The Grand Unification of the gauge symmetries suggests itself quite naturally. This idea can be realized in different scenarios some of which predict new vector bosons and a plethora of new fermions. Mass scales of these novel particles could be as low as a few hundred GeV.
A very important theoretical extension of the Standard Model, which is interrelated with the unification of the gauge symmetries, is Supersymmetry . This novel symmetry concept unifies matter and forces by pairing the associated fermionic and bosonic particles in common multiplets. Several arguments strongly support the hypothesis that this symmetry is realized in Nature indeed. $`(i)`$ Supersymmetry stabilizes light masses of Higgs particles in the context of very high energy scales as demanded by grand unified theories. $`(ii)`$ Supersymmetry may generate the Higgs mechanism itself by inducing the radiative symmetry breaking of $`\mathrm{SU}(2)_\mathrm{L}\times \mathrm{U}(1)_\mathrm{Y}`$ while leaving $`\mathrm{U}(1)_{\mathrm{EM}}`$ and $`\mathrm{SU}(3)_\mathrm{c}`$ unbroken for a top quark mass between 100 and 200 GeV. $`(iii)`$ This symmetry picture is also supported strongly by the successful prediction of the electroweak mixing angle in the minimal version of the theory. The particle spectrum in this theory drives the evolution of the electroweak mixing angle from the GUT value 3/8 down to $`\mathrm{sin}^2\theta _W`$ = 0.2336 $`\pm `$ 0.0017, within a margin $``$ 0.002 to the experimental value $`\mathrm{sin}^2\theta _W^{exp}`$ = 0.2316 $`\pm `$ 0.0002.
A spectrum of several neutral and charged Higgs bosons is predicted in supersymmetric theories. The mass of the lightest Higgs boson is less than $`150`$ GeV in nearly all scenarios while the heavy Higgs particles have masses of the order of the electroweak symmetry breaking scale. Many other novel particles are predicted in supersymmetric theories. The scalar partners of the leptons could have masses in the range of $`200`$ GeV whereas squarks are expected to be considerably heavier. The lightest supersymmetric states are likely to be non–colored gaugino/higgsino states with masses possibly in the 100 GeV range. Searching for these supersymmetric particles will be one of the most important tasks at future $`e^+e^{}`$ colliders. Moreover, the high accuracy which can be achieved when masses and couplings are measured, will allow us to determine the mechanism of supersymmetry breaking and to extrapolate the basic parameters of the theory so that the key elements of the underlying grand unified theories at scales, potentially close to the Planck scale, can be reconstructed.
In the alternative scenario of heavy or no fundamental Higgs bosons, new strong inter- actions between electroweak bosons would be observed, characterized by a scale of order 1 TeV at which the electroweak symmetries would be broken strongly . This scenario could be analyzed by studying the elastic scattering of $`W`$ bosons at high energies. New resonances would be formed, the properties of which would uncover the underlying microscopic interactions.
While new high–mass vector bosons and particles carrying color quantum numbers can be searched for very efficiently at hadron colliders, $`e^+e^{}`$ colliders provide in many ways unique opportunities to discover and explore non–colored particles. This is most obvious in supersymmetric theories. Combining LEP2 analyses with future searches at the Tevatron and the LHC, the light and heavy Higgs bosons can be found individually only in part of the supersymmetry parameter space. Squarks and gluinos can be searched for very efficiently at the LHC. Yet precision studies of their properties are possible only in part of the parameter space. Similarly non-colored supersymmetric particles; a model–independent analysis of gauginos/higgsinos and scalar sleptons can only be carried out at $`e^+e^{}`$ colliders with well–defined kinematics at the level of the subprocesses. The detailed knowledge of all the properties of the colored and non–colored supersymmetric states will reveal the mechanism of supersymmetry breaking and the structure of the underlying theory.
Thus, the physics progamme of $`e^+e^{}`$ linear colliders is in many aspects complementary to the programme of the $`pp`$ collider LHC. The high accuracy which can be achieved at $`e^+e^{}`$ colliders in exploring the properties of the top quark, electroweak gauge bosons, Higgs particles and supersymmetric particles will enable us to cover the energy range above the existing machines up to the TeV region in a conclusive form, eventually providing us with essential clues to the basic structure of matter and the laws of physics.
The discussion will focus on the physics program at an $`e^+e^{}`$ linear collider operating at center–of–mass energies above LEP2 up to about 1 TeV. Primarily high–luminosity runs, collecting integrated luminosities up to 0.5 ab<sup>-1</sup> in one to two years of operation, will be described. Also the results expected from high–luminosity runs at low energies on the $`Z`$ resonance, the GigaZ mode, and near the $`WW`$ threshold will be summarized. Electrons and positrons will in general be assumed polarized to 80% and 60%, respectively. Specific problems which can be solved in $`e^{}e^{}`$, e$`\gamma `$ and $`\gamma \gamma `$ modes of the linear collider will be addressed in the appropriate context.
This summary report is built on the general linear collider review of Ref.. Other material can be found in Refs., experimental aspects particularly in Ref.. For recent summaries of the LHC physics and $`\mu \mu `$ physics programs see Refs. and , respectively.
## 2 Top Quark Physics
Top quarks are the heaviest matter particles in the 3–family Standard Model, introduced to incorporate $`𝒞𝒫`$ violation in the left–handed charged current sector. They may therefore hold the key for aspects of the physics beyond the Standard Model at high–energy scales. Examples in which the large top mass is crucial, are multi–Higgs doublet models, models of dynamical symmetry breaking, compositeness and supersymmetry. Strong indirect evidence for the existence of top quarks, based on the well established gauge symmetry pattern of the Standard Model, had been accumulated quite early. By evaluating the high–precision electroweak data, the value of the top quark mass was estimated to be $`m_t=180\pm 14`$ GeV. Top quarks have recently been observed directly by the two Tevatron experiments , corresponding to a mass of $`m_t=174.3\pm 5.1\text{ GeV}`$ which is in striking agreement with the result of the electroweak analysis.
### 2.1 The Profile of the Top Quark
For a top mass larger than the $`W`$ mass, the channel
$$tb+W^+$$
is the dominant decay mode. For m$`{}_{t}{}^{}`$ 175 GeV the width of the top quark, $`\mathrm{\Gamma }_t`$ 1.4 GeV, is so large compared with the scale $`\mathrm{\Lambda }`$ of the strong interactions that this quark can be treated as a bare quantum which is not dressed by non–perturbative strong interactions .
Chirality of the $`(tb)`$ decay current: The precise determination of the weak isospin quantum numbers does not allow for large deviations of the $`(tb)`$ decay current from the left–handed SM prescription. Nevertheless, since $`V+A`$ admixtures may grow with the masses of the quarks involved \[through mixing with heavy mirror quarks, for instance\], it is necessary to check the chirality of the decay current directly. The $`l^+`$ energy distribution in the semileptonic decay chain $`tW^+l^+`$ depends on the chirality of the current. Any deviation from the standard $`VA`$ current would lead to a stiffening of the spectrum and, in particular, to a non-zero value at the upper end–point of the energy distribution. A sensitivity of about 5% to a possible $`V+A`$ admixture can be reached experimentally .
Non–standard top decays could occur in supersymmetric extensions of the Standard Model: top decays into charged Higgs bosons and/or top decays to stop particles, $`tb+H^+`$ and $`t\stackrel{~}{t}+\stackrel{~}{\chi }_1^o`$. If kinematically allowed, branching ratios of order 10% are expected in both cases so that these decay modes could be observed easily. Decays with signatures as clean as $`tc\gamma ,cZ,cH`$ may be detected for branching ratios of order $`10^4`$ and less.
The main production mechanism for top quarks in $`e^+e^{}`$ collisions is the annihilation channel ,
$$e^+e^{}\stackrel{\gamma ,Z}{}t\overline{t}$$
For $`m_t175`$ GeV, the maximum of the cross section $`\sigma (t\overline{t})800`$ fb is reached about 30 GeV above the threshold, giving rise to a million top quarks in two years of collider operation. If the scale of new areas beyond the Standard Model is much larger than the collider energy, the electroweak production currents can globally be described by form factors which reduce to anomalous $`Z`$ charges, anomalous magnetic dipole moments and electric dipole moments.
Magnetic dipole moments: If the electrons in $`e^+e^{}t\overline{t}`$ are left–handedly polarized, the top quarks are produced preferentially as left–handed particles in the forward direction while only a small fraction is produced as right–handed particles in the backward direction , so that the backward direction is most sensitive to small anomalous magnetic moments of the top quarks. The anomalous magnetic moments can be bounded to the percent level by measuring the angular dependence of the $`t`$ quark cross section in this region.
Electric dipole moments: These moments are generated by $`𝒞𝒫`$–non invariant interactions. Non–zero values of the moments can be detected by means of non–vanishing expectation values of $`𝒞𝒫`$–odd momentum tensors such as $`T_{ij}=(q_+q_{})_i(q_+\times q_{})_j`$ with $`q_\pm `$ being the unit momentum vectors of the $`W`$ decay leptons. Sensitivity limits to $`\gamma ,Z`$ electric dipole moments of $`10^{18}`$ e cm can be reached for an integrated luminosity of $`=100\text{ fb}^1`$ at $`\sqrt{s}=500`$ GeV.
### 2.2 The Top-Quark Mass
Quark–antiquark production near the threshold in $`e^+e^{}`$ collisions is of exceptional interest. The long time which the particles stay close together at low velocities, allows the strong interactions to build up rich structures of bound states and resonances. This picture would have applied to top quarks up to the mass range of $`130`$ GeV. Beyond this value, the picture changes quite dramatically as a result of the rapid top decay: The decay time of the states becomes shorter than the revolution time of the constituents so that toponium resonances cannot be formed any more . For a while, however, remnants of the $`1S`$ state give rise to a peak in the excitation curve, yet it disappears for top masses in excess of 180 GeV. Nevertheless, across this range the resonance remnants induce a steep rise of the cross section near the threshold.
Since the rapid top decay restricts the interaction region to small distances, the excitation curve can be predicted in perturbative QCD , based essentially on the Coulombic interquark potential $`V(R)=4/3\times \alpha _s(R)/R`$. The cross section is built up by the superposition of all $`nS(t\overline{t})`$ states. The form and the height of the excitation curve are very sensitive to the mass of the top quark, cf Fig. 1.
Detailed experimental simulations predict the following sensitivity to the top mass near $`m_t175`$ GeV:
$`\delta m_t`$ $``$ $`200\text{MeV}`$
for an integrated luminosity of $`=50\text{ fb}^1`$, including remnant uncertainties due to higher–order QCD corrections and experimental errors contributing at approximately equal strength. At proton colliders a sensitivity of about 1 to 2 GeV has been predicted for the top mass, based on the reconstruction of top quarks from jet and lepton final states . Thus, $`e^+e^{}`$ colliders will improve the measurement of the top quark mass by at least an order of magnitude.
## 3 Electroweak Gauge Bosons
### 3.1 Standard $`W`$, $`Z`$ Bosons
The fundamental electroweak and strong forces appear to be of gauge theoretical origin. This is one of the outstanding results of theoretical and experimental analyses in the past three decades. However, little direct evidence has been accumulated so far for the non–Abelian nature of the forces in the electroweak $`W^\pm ,Z,\gamma `$ sector. Since deviations from the gauge symmetries manifest themselves in experimental observables with coefficients $`(\beta \gamma )^2`$, high energies will allow stringent direct tests of the self–couplings of the electroweak gauge bosons.
The gauge symmetries of the Standard Model determine the form and the strength of the self–interactions of the electroweak bosons, triple couplings $`WW\gamma ,WWZ`$ and quartic couplings. Deviations from the gauge symmetric form of these vertices could be expected in more general scenarios . In models in which $`W,Z`$ bosons are generated dynamically and interact strongly with each other at high scales $`\mathrm{\Lambda }_{}`$, corrections could alter the vertices to order $`(M_W/\mathrm{\Lambda }_{})^2`$ and induce new types of couplings.
While the experimental analyses of the self–couplings of the electroweak bosons can be carried out at collider energies of 500 GeV with high accuracy, $`WW`$ scattering can only be studied at energies in the TeV range. This is an important process which must be investigated very thoroughly if light Higgs particles do not exist and $`W`$ bosons become strongly interacting particles at high energies.
The properties of the $`Z`$ boson have been studied at LEP and SLC with very high accuracy. By operating TESLA at low energies on the $`Z`$ resonance in the GigaZ mode and near the $`WW`$ threshold, the measurement of fundamental electroweak parameters, the electroweak mixing angle and $`W`$ mass can be improved by yet an order of magnitude, allowing to test electroweak symmetry breaking stringently at the quantum level.
#### The GigaZ Mode:
By building a bypass for the transport of electron and positron bunches, high luminosity can also be reached at low energies in linear colliders. On the $`Z`$ resonance, an ensemble of $`10^9`$ events, 1 GigaZ, can be generated in a year, expanding the LEP sample by two orders of magnitude. With both electron and positron beams longitudinal polarized, the electroweak mixing angle can be determined very accurately by measuring the left–right asymmetry $`A_{LR}=2(14\mathrm{sin}^2\theta _W)/[1+(14\mathrm{sin}^2\theta _W)^2]`$, Ref. :
$`\delta \mathrm{sin}^2\theta _W10^5`$
Similarly the measurement of the $`W`$ mass can be improved to
$`\delta M_W6\text{MeV}`$
by scanning the threshold region, Ref. .
Based on these two measurements, a variety of high–precision tests can be performed in the electroweak sector. Extracting the Higgs mass from the electroweak observables is particularly interesting. From the $`\rho `$–parameter this mass can be predicted to an accuracy of about 6%, improving LHC based predictions by almost an order of magnitude . Comparing this prediction with the direct measurement of the Higgs mass, quantum fluctuations can be tested stringently in a spontaneously broken gauge theory.
#### The Triple Gauge-Boson Couplings:
The couplings $`W^+W^{}\gamma `$ and $`W^+W^{}Z`$ are in general described each by seven parameters. Assuming $`𝒞,𝒫`$ and $`𝒯`$ invariance in the pure electroweak boson sector, the number of parameters can be reduced to three,
$`_k/ig_k`$ $`=`$ $`g_k^1W_{\mu \nu }^{}W_\mu A_\nu +\mathrm{h}.\mathrm{c}.+\kappa _kW_\mu ^{}W_\nu F_{\mu \nu }+{\displaystyle \frac{\lambda _k}{M_W^2}}W_{\rho \mu }^{}W_{\mu \nu }F_{\nu \rho }`$
with $`g_\gamma =e`$ and $`g_Z=e\mathrm{cot}\theta _W`$ for $`k=\gamma ,Z`$. The $`\kappa =1+\mathrm{\Delta }\kappa `$ and the $`\lambda `$ parameters can be identified with the $`\gamma ,Z`$ charges of the $`W`$ bosons and the related magnetic dipole moments and electric quadrupole moments:
$`\mu _\gamma =`$ $`{\displaystyle \frac{e}{2M_W}}\left[2+\mathrm{\Delta }\kappa _\gamma +\lambda _\gamma \right]`$ $`\mathrm{and}\gamma Z`$
$`Q_\gamma =`$ $`{\displaystyle \frac{e}{M_W^2}}\left[1+\mathrm{\Delta }\kappa _\gamma \lambda _\gamma \right]`$ $`\mathrm{and}\gamma Z`$
The gauge symmetries of the SM demand $`\kappa =1`$ and $`\lambda =0`$. The magnetic dipole and the electric quadrupole moments can be measured directly in the production of $`W\gamma `$ and $`WZ`$ pairs at $`p\overline{p}/pp`$ colliders and $`WW`$ pairs at $`e^+e^{}`$ and $`\gamma \gamma `$ colliders.
Detailed experimental analyses have been carried out for the reaction $`e^+e^{}W^+W^{}(l\nu _e)(q\overline{q}^{})`$. The bounds on $`\mathrm{\Delta }\kappa ,\lambda `$ which can be obtained at $`e^+e^{}`$ colliders of 500 GeV are significantly better than the bounds expected from the LHC:
$`\mathrm{\Delta }g_1^Z=2.5\times 10^3`$
$`\mathrm{\Delta }\kappa _\gamma =4.8\times 10^4`$ $`\mathrm{\Delta }\kappa _Z=7.9\times 10^4`$
$`\lambda _\gamma =7.2\times 10^4`$ $`\lambda _Z=6.5\times 10^4`$
Moreover, they improve at 1 TeV by nearly an order of magnitude. The scales $`\mathrm{\Lambda }_{}`$ which can be probed, extend far beyond the energy scales which are accessible directly.
#### Strongly Interacting $`W`$, $`Z`$ Bosons
If the scenario in which $`W`$, $`Z`$, Higgs bosons are weakly interacting up to the GUT scale is not realized in Nature, the alternative scenario is a strongly interacting $`W`$, $`Z`$ sector. Without a light Higgs boson with a mass of less than about 1 TeV, the electroweak bosons must become strongly interacting particles at energies of about 1.2 TeV to comply with the requirements of unitarity for the $`W_LW_L`$ scattering amplitudes. By absorbing the Goldstone particles associated with the spontaneous symmetry breaking of the new strong interactions, the longitudinal degrees of freedom for the massive vector bosons may be built up, as realized in technicolor type theories, for instance. In such scenarios, novel resonances are predicted in the $`𝒪`$(1 TeV) energy range which can be generated in $`W_LW_L`$ collisions.
In scenarios of strongly interacting vector bosons, $`W_LW_L`$ scattering must be studied at energies of order 1 TeV which requires the highest energies possible in $`e^\pm e^{}`$ colliders. (Quasi)elastic $`WW`$ scattering can be investigated by using $`W`$ bosons radiated off the electron and positron beams, $`ee\nu \nu WW`$, or by exploiting final-state interactions in the $`e^+e^{}`$ annihilation to $`W`$ pairs, $`e^+e^{}W^+W^{}`$. All possible (isospin, angular momentum) combinations in $`WW`$ scattering amplitudes $`a_{IJ}`$ can be realized in the first process. The cross sections however are small as long as no resonances are formed.
Building up the electroweak vector boson masses by the interactions of the gauge fields with the Goldstone bosons associated with the spontaneous symmetry breaking of the underlying strong–interaction theory, the longitudinal degrees of freedom of the vector bosons can be identified at high energies with the Goldstone bosons themselves as a result of the equivalence theorem. In analogy to the $`\pi \pi `$ low–energy theorems, the first terms in the energy expansion of the $`WW`$ scattering amplitudes are determined independent of dynamical details:
$`a_{00}=+{\displaystyle \frac{6}{96\pi v^2}}`$ $`a_{20}={\displaystyle \frac{2}{96\pi v^2}}`$
$`a_{11}=+{\displaystyle \frac{1}{96\pi v^2}}`$
These fundamental scattering amplitudes in the threshold region of the strong $`WW`$ interactions can be tested to an accuracy of about 15% at a high–luminosity collider of 1 TeV; new strong-interaction scales extending up to about 3 TeV can be probed in these experiments, thus covering energy scales up to the formation of novel resonances.
The attractive $`I=0`$ and $`I=1`$ channels may form Higgs and $`\rho `$–type resonances at high energies. The formation of resonances would lead to spectacular phenomena in $`WW`$ collisions .
Similar phenomena would also be observed as rescattering effects in the cross section $`\sigma (e^+e^{}W^+W^{})`$ for $`W`$–pair production. $`(I,J)=(1,1)`$ resonance effects would be noticeable at $`\sqrt{s}=1`$ TeV up to resonance masses of about 5 TeV in the angular distributions of the $`W`$ decay final states .
### 3.2 Extended Gauge Theories
The gauge symmetry of the Standard Model, $`\mathrm{SU}(3)\times \mathrm{SU}(2)\times \mathrm{U}(1)`$, is widely believed not to be the ultima ratio. The SM does not unify the electroweak and strong forces since the coupling constants of these interactions are different and appear to be independent. However, one should expect that in a more fundamental theory the three forces are described within a single gauge group and, hence, with only one coupling constant at high energy scales. This grand unified theory will be based on a gauge group containing $`\mathrm{SU}(3)\times \mathrm{SU}(2)\times \mathrm{U}(1)`$ as a subgroup and it will be reduced to this symmetry at low energies.
Two predictions of grand unified theories may have interesting consequences in the energy range of a few hundred GeV :
(a) The unified symmetry group must be spontaneously broken at the unification scale $`\mathrm{\Lambda }_{\mathrm{GUT}}\stackrel{<}{}10^{16}`$ GeV in order to be compatible with the experimental bounds on the proton lifetime. However, the breaking to the SM group may occur in several steps and some subgroups may remain unbroken down to a scale of order 1 TeV. In this case the surviving group factors allow for new gauge bosons with masses not far above the scale of electroweak symmetry breaking. Besides $`\mathrm{SU}(5)`$, two other unification groups have received much attention: In $`\mathrm{SO}(10)`$ three new gauge bosons $`W_R^\pm ,Z_R`$ are predicted, while in E(6) a light neutral $`Z^{}`$ boson may exist in the TeV range.
The virtual effects of a new $`Z_R/Z^{}`$ boson associated with the most general effective theories which arise from breaking E(6) $`\mathrm{SU}(3)\times \mathrm{SU}(2)\times \mathrm{U}(1)\times \mathrm{U}(1)_\mathrm{Y}^{}`$ and $`\mathrm{SO}(10)\mathrm{SU}(2)_\mathrm{L}\times \mathrm{SU}(2)_\mathrm{R}\times \mathrm{U}(1)`$, have been investigated in Ref. . Assuming the $`Z_R/Z^{}`$ bosons to be heavier than the available c.m. energy, the propagator effects on various observables of the process
$`e^+e^{}\stackrel{V}{}f\overline{f}:V=\gamma ,Z\text{and}Z_R/Z^{}`$
have been studied in detail. As shown in Table 1, the effects of new vector bosons can be probed for masses up to 5 TeV at a 500 GeV collider. While they may be produced directly up to about 5 TeV at the LHC, experiments at the $`e^+e^{}`$ collider will measure the couplings of the vector bosons to fermions very precisely, thus identifying the physical nature of the new bosons. Masses up to 10 TeV and 50 TeV can be probed in $`e^+e^{}`$ colliders operating at 800 GeV and 5 TeV, respectively. These two windows extend to much higher scales than the discovery limits anticipated at LHC.
(b) The grand unification groups incorporate extended fermion representations in which a complete generation of SM quarks and leptons can be naturally embedded. These representations accommodate a variety of additional new fermions. It is conceivable that the new fermions acquire masses not much larger than the Fermi scale. This is necessary if the predicted new gauge bosons are relatively light. SO(10) is the simplest group in which the 15 members of each SM generation of fermions can be embedded into a single multiplet. This representation has dimension 16 and contains a right–handed neutrino. The group E(6) contains $`\mathrm{SU}(5)`$ and $`\mathrm{SO}(10)`$ as subgroups. In E(6), each quark–lepton generation belongs to a representation of dimension 27. To complete this representation, twelve new fields are needed in addition to the SM fermion fields. In each family the spectrum includes two additional isodoublets of leptons, two isosinglet neutrinos and an isosinglet quark with charge $`1/3`$.
If the new particles have non–zero electromagnetic and weak charges, and if their masses are smaller than the beam energy of the $`e^+e^{}`$ collider, they can be pair produced. In general, the production processes are built up by a superposition of s–channel $`\gamma `$ and $`Z`$ exchanges, but additional contributions could come from the extra neutral bosons if their masses are not much larger than the c.m. energy . The cross sections are large, of the order of the point–like QED cross section. This leads to samples of several thousands of events. Fermion mixing, if large enough, gives rise to additional production mechanisms for the new fermions: single production in association with the light partners. In this case, masses very close to the total energy of the $`e^+e^{}`$ collider can be reached.
## 4 The Higgs Mechanism
### 4.1 Basis
The Higgs mechanism is the cornerstone in the electroweak sector of the Standard Model. The fundamental SM particles, leptons, quarks and weak gauge bosons, acquire masses by means of the interaction with a scalar field. To accommodate the well–established electromagnetic and weak phenomena, the Higgs mechanism requires the existence of at least one weak iso–doublet scalar field. After absorbing three Goldstone modes to build up the longitudinal polarization states of the $`W^\pm ,Z`$ bosons, one degree of freedom is left–over, corresponding to a real scalar particle.
Three steps are necessary to establish experimentally the Higgs mechanism sui generis as the mechanism for generating the masses of the fundamental SM particles:
The Higgs boson must be discovered – the experimentum crucis;
The couplings of the Higgs particle with gauge bosons and fermions must be proven to increase with their masses;
The Higgs potential generating the non–zero Higgs field in the vacuum and breaking the electroweak symmetry in the scalar sector must be reconstructed by determining the Higgs self–couplings.
The only unknown parameter in the SM Higgs sector is the mass of the Higgs particle. Constraints on the mass can, however, be derived from the upper scale $`\mathrm{\Lambda }_{}`$ of the energy range in which the model is assumed to be valid before the particles become strongly interacting and new dynamical phenomena emerge . Increasing the energy scale, the quartic self–coupling of the Higgs field grows for large values indefinitely. If the Higgs mass is small, the energy cut–off $`\mathrm{\Lambda }_{}`$ is large at which the coupling grows beyond any bound; conversely, if the Higgs mass is large, the cut–off $`\mathrm{\Lambda }_{}`$ is small. The condition $`M_H<\mathrm{\Lambda }_{}`$ sets an upper limit on the Higgs mass in the Standard Model. Detailed analyses lead to an estimate of about 700 GeV for the upper limit on $`M_H`$. If the Higgs mass is less than 180 to 200 GeV, the Standard Model can be extended up to the GUT scale $`\mathrm{\Lambda }_{\mathrm{GUT}}10^{16}`$ GeV, while all particles remain weakly interacting. The hypothesis that the interactions between $`W,Z`$ bosons and Higgs particles remain weak up to the GUT scale, plays a key rôle in deriving the experimental value of the electroweak mixing parameter $`\mathrm{sin}^2\theta _W`$ from grand unified theories. From this hypothesis and the additional requirement of vacuum stability, upper and lower bounds on the Higgs mass can be derived. Based on these arguments, the SM Higgs mass should be expected in the mass window $`130<M_H<180`$ GeV for a top mass value of about 175 GeV.
Several channels can be exploited to search for Higgs particles in the Higgs–strahlung and fusion processes of $`e^+e^{}`$ colliders \[??\]. In the Higgs–strahlung process $`e^+e^{}`$ $`ZH`$, missing–mass techniques can be used in events with leptonic $`Z`$ decays or the Higgs particle may be reconstructed in $`Hb\overline{b},WW`$ directly. The $`WW`$ fusion process $`e^+e^{}`$ $`\overline{\nu }_e\nu _eH`$ requires the reconstruction of the Higgs particle.
Once the Higgs boson is found at LEP, Tevatron or LHC, it will be very important to explore its properties at the $`e^+e^{}`$ linear collider to establish the Higgs mechanism experimentally. This is possible with high precision in the clean environment of $`e^+e^{}`$ colliders in which at high luminosity a large ensemble of order $`10^5`$ Higgs bosons can be generated nearly background–free. The zero–spin of the Higgs particle is reflected in the angular distribution of the Higgs–strahlung process which must approach the $`\mathrm{sin}^2\theta `$ law asymptotically . The strength of the couplings to $`Z`$ and $`W`$ bosons is reflected in the magnitude of the $`e^+e^{}`$ production cross sections. The strength of the couplings to fermions can be measured in the decay branching ratios and the Higgs bremsstrahlung off top quarks. Double Higgs–strahlung can be exploited to measure the trilinear Higgs self–coupling.
From the preceding discussion we conclude that an $`e^+e^{}`$ linear collider with energies in the range of 300 to 500 GeV and high luminosity is the ideal instrument to investigate the Higgs mechanism in the intermediate mass range which, a priori, may be considered the theoretically preferred part in the entire range of possible Higgs mass values.
### 4.2 The Higgs Particle in the Standard Model
The profile of the SM Higgs particle is completely determined if the Higgs mass is fixed. For Higgs particles in the intermediate mass range $`M_ZM_H2M_Z`$ the main decay modes are decays into $`b\overline{b}`$ pairs and $`WW,ZZ`$ pairs with one of the two gauge bosons being virtual below the threshold . Above the $`WW`$ threshold, the Higgs particles decay almost exclusively into these channels, except in the mass range near the $`t\overline{t}`$ decay threshold. Below 140 GeV, the decays $`H\tau ^+\tau ^{},c\overline{c}`$ and $`gg`$ are also important besides the dominating $`b\overline{b}`$ channel. Up to masses of 140 GeV, the Higgs particle is very narrow, $`\mathrm{\Gamma }(H)10`$ MeV. After opening the \[virtual\] gauge boson channels, the state becomes rapidly wider, the width reaching $``$ 1 GeV at the $`ZZ`$ threshold. The width cannot be measured directly in the intermediate mass range. Only above $`M_H200`$ GeV it becomes wide enough to be resolved experimentally.
The main production mechanisms for Higgs particles in $`e^+e^{}`$ collisions are Higgs–strahlung off the $`Z`$ boson line and the $`WW`$ fusion process ,
$`(a)e^+e^{}`$ $`\stackrel{Z}{}`$ $`Z+H`$
$`(b)e^+e^{}`$ $`\stackrel{WW}{}`$ $`\overline{\nu }_e\nu _e+H`$
With rising energy the Higgs–strahlung cross section scales $`\alpha _w^2/s`$ while the fusion cross sections increase logarithmically $`\alpha _w^3M_W^2\text{log }s/M_H^2`$, becoming dominant above 500 GeV:
$`\sigma (e^+e^{}ZH)`$ $``$ $`{\displaystyle \frac{G_F^2M_Z^4}{96\pi s}}[1+(14\mathrm{sin}^2\theta _W)^2]`$
$`\sigma (e^+e^{}\overline{\nu }\nu H)`$ $``$ $`{\displaystyle \frac{G_F^3M_W^4}{4\sqrt{2}\pi ^3}}\mathrm{log}{\displaystyle \frac{s}{M_H^2}}`$
As a general rule, the cross sections and rates \[about $`10^5`$ events\] are sufficiently large to detect Higgs particles with masses up to 70% of the total $`e^+e^{}`$ c.m. energy.
The recoiling $`Z`$ boson in the two–body reaction $`e^+e^{}ZH`$ is mono–energetic and the mass can be derived from the energy of the $`Z`$ boson, $`M_H^2=s2\sqrt{s}E_Z+M_Z^2`$. Initial state bremsstrahlung and beamstrahlung smear out the peak slightly, as shown in Fig. 2. A similarly clear peak can be observed in the fusion process $`e^+e^{}`$ $`\overline{\nu }_e\nu _eH`$ by collecting the decay products of the Higgs boson.
#### Mass and Width:
The mass of the Higgs boson can be measured very accurately by analyzing the recoil $`Z`$ spectrum in Higgs–strahlung events. Experimental simulations have demonstrated that the error in the mass measurement can be reduced to
$`\delta M_H50\text{MeV}`$
in high–luminosity runs.
The width of the SM Higgs boson can be determined in an almost completely model–independent way in the difficult intermediate mass range in which the Breit–Wigner form cannot be reconstructed at an $`e^+e^{}`$ collider. Measuring the branching ratio $`BR_i`$ in the decay and the partial width $`\mathrm{\Gamma }_i`$ in the production process, the total width $`\mathrm{\Gamma }_H`$ can be derived from
$`\mathrm{\Gamma }_H=\mathrm{\Gamma }_i/BR_i`$
The two channels $`i=WW`$ and $`i=\gamma \gamma `$ are useful for this analysis . The partial width $`\mathrm{\Gamma }_{WW}`$ can be extracted from the size of the $`WW`$ fusion cross section while $`\mathrm{\Gamma }_{\gamma \gamma }`$ can be measured in the Compton collider mode . The accuracies of a few percent match the expected accuracy in scanning the Breit–Wigner excitation at a muon–collider.
#### Spin and Parity:
The angular distribution of the $`Z/H`$ bosons in the Higgs–strahlung process is sensitive to the external quantum numbers of the Higgs particle . Since the amplitude is given by $`𝒜(0^+)\epsilon _Z^{}\epsilon _Z^{}`$, the $`Z`$ boson is produced in a state of longitudinal polarization at high energies. As a result, the angular distribution $`\mathrm{d}\sigma /\mathrm{d}\mathrm{cos}\theta \lambda \mathrm{sin}^2\theta +8M_Z^2/s`$ approaches the spin–zero law $`\mathrm{sin}^2\theta `$ asymptotically. This may be contrasted with the distribution $`1+\mathrm{cos}^2\theta `$ for negative parity states, which follows from the transverse polarization amplitude $`𝒜(0^{})\epsilon _Z^{}\times \epsilon _Z^{}k_Z`$. It is also characteristically different from the background process $`e^+e^{}`$ $`ZZ`$ which is strongly peaked in the forward/backward direction.
#### Higgs Couplings:
Since the fundamental SM particles acquire masses by means of the interaction with the Higgs field, the scale of the Higgs couplings to fermions and gauge bosons is set by the masses of the particles:
$`g_{HVV}=2[\sqrt{2}G_F]^{1/2}M_V^2\text{and}g_{Hff}=[\sqrt{2}G_F]^{1/2}M_f`$
It will be a very important task to measure the Higgs couplings to the fundamental particles since they are uniquely predicted by the very nature of the Higgs mechanism. The Higgs couplings to massive gauge bosons can be determined from the measurement of the production cross sections with an accuracy of $`\pm 1\%`$, the $`HZZ`$ coupling in the Higgs–strahlung and the $`HWW`$ coupling in the fusion process. For Higgs couplings to fermions, either loop effects in $`Hgg,\gamma \gamma `$ \[mediated by top quarks\] can be exploited, or the direct measurement of branching ratios $`Hb\overline{b},c\overline{c},\tau ^+\tau ^{},gg`$ in the lower part of the intermediate mass range. This is exemplified in Fig. 3. For $`M_H=120`$ GeV the following accuracy $`\delta BR/BR`$ can be achieved in the determination of the Higgs decay branching ratios:
$`bb:2.4\%`$ $`WW^{}:5.4\%`$
$`cc:8.3\%`$ $`gg:5.5\%`$
$`\tau \tau :6.0\%`$
By measuring the ratio of the $`\tau \tau `$ to the $`bb`$ branching ratios
$`{\displaystyle \frac{BR(H\tau \tau )}{BR(Hbb)}}={\displaystyle \frac{m_\tau ^2}{3m_b^2(M_H)}}`$
the linear dependence of the Yukawa couplings on the fermion masses can be tested very nicely. A direct way to determine the Yukawa coupling of the intermediate mass Higgs boson to the top quark in the range $`m_H120`$ GeV is provided by the bremsstrahlung process $`e^+e^{}t\overline{t}H`$ in high energy $`e^+e^{}`$ colliders . The absolute values of the Yukawa couplings can be reconstructed by combining decay branching ratios with the production cross sections.
#### Higgs Self–couplings:
To generate a non–zero value of the Higgs field in the vacuum, the minimum of the Higgs potential must be shifted away from the origin. Rewriting the potential
$`V`$ $`=`$ $`\lambda [\phi ^2\frac{1}{2}v^2]^2`$
$`=`$ $`{\displaystyle \frac{M_H^2}{2}}H^2+{\displaystyle \frac{M_H^2}{2v}}H^3+{\displaystyle \frac{M_H^2}{8v^2}}H^4`$
in terms of the physical Higgs field $`H`$, the potential can be reconstructed by measuring the trilinear and quartic couplings. At a high–luminosity $`e^+e^{}`$ collider, the trilinear coupling can be tested in the double Higgs–strahlung process:
$`e^+e^{}Z+HH`$
The splitting of a virtual Higgs boson into two real Higgs bosons is determined by the trilinear Higgs coupling: $`e^+e^{}Z+H^{}[HH]`$. Even though the cross section is less than 1 fb , the coupling can nevertheless be measured with an accuracy better than 20%, cf. Ref. . Thus an essential element of the mechanism responsible for the spontaneous symmetry breaking in the scalar sector can be established experimentally at the high–luminosity collider.
## 5 Supersymmetry
Even though there is no direct experimental evidence so far for the realization of supersymmetry in Nature, this concept has so many attractive features that it can be considered as a prime target of present and future experimental particle research . Arguments in favor of supersymmetry are deeply rooted in particle physics. Supersymmetry may play an important r$`\widehat{\mathrm{o}}`$le in a quantum theory of gravity. In relating particles of different spins to each other, i.e. fermions and bosons, low–energy supersymmetry stabilizes the masses of fundamental Higgs scalars in the background of very high energy scales associated with grand unification. Besides solving this hierarchy problem, supersymmetry may even be closely related to the physical origin of the Higgs phenomenon itself: In a supergravity inspired GUT realization with universal scalar masses at the GUT scale, the evolution of one of the scalar masses squared down to the electroweak scale can become negative and can thus give rise to spontaneous symmetry breaking if the top mass has a value between about 100 to 200 GeV while all other squared masses of squarks and sleptons remain positive so that $`\mathrm{U}(1)_{\mathrm{EM}}`$ and $`\mathrm{SU}(3)_\mathrm{C}`$ remain unbroken.
The minimal supersymmetric extension of the Standard Model (MSSM) is based on the SM group $`\mathrm{SU}(3)\times \mathrm{SU}(2)\times \mathrm{U}(1)`$. The MSSM incorporates a spectrum of five Higgs particles (representative for a wide class of models). At tree level, the mass of the lightest Higgs boson $`h^0`$ is smaller than the $`Z`$ mass; the bound is shifted to $`\stackrel{<}{}140`$ GeV by radiative corrections. The masses of the heavy neutral and charged Higgs bosons can be expected in the range of the electroweak symmetry breaking scale. The gauginos are the supersymmetric spin–$`\frac{1}{2}`$ partners of the gauge bosons. The quark and lepton matter particles are associated with scalar supersymmetric particles, squarks and sleptons. To preserve supersymmetry, two Higgs doublets are needed to endow down as well as up–type fermions with masses; the supersymmetric partners of the Higgs bosons are spin–$`\frac{1}{2}`$ higgsinos. \[Charged/neutral higgsinos mix in general with the non–colored gauginos, forming charginos and neutralinos.\] Supersymmetric partners carry a multiplicative quantum number $`R=1`$ ($`R=+1`$ for ordinary particles) which is conserved in this model. Supersymmetric particles are therefore generated in pairs and the lightest supersymmetric particle $`LSP`$ is stable.
Strong support for supersymmetry and the MSSM particle spectrum in the mass range of several hundred GeV follows from the high–precision measurement of the electroweak mixing angle $`\mathrm{sin}^2\theta _W`$ . The value predicted by the MSSM and the value determined by the LEP and other experiments,
$`MSSM:\mathrm{sin}^2\theta _W`$ $`=0.2336\pm 0.0017`$
$`EXP:\mathrm{sin}^2\theta _W`$ $`=0.2316\pm 0.0002`$
match surprisingly well, the difference being less than about 2 per–mille.
A central problem of supersymmetric theories is the breaking mechanism. Several scenarios have been proposed and experimental consequences have been elaborated in a few cases: gravity mediated supersymmetry breaking mSUGRA ; gauge mediated supersymmetry breaking GMSB ; anomaly mediated supersymetry breaking AMSB ; Scherk–Schwarz supersymmetry breaking SSSB . Mass spectra of the supersymmetric particles are quite different in these scenarios so that high–precision measurements of the particle properties will shed light experimentally on this theoretical problem. Moreover, extrapolations can be performed in a stable way which will allow to reconstruct the basic supersymmetric theory eventually at a scale close to the Planck scale. Precision experiments at a high–luminosity $`e^+e^{}`$ linear collider are therefore expected to advance the understanding of supersymmetry in an essential way.
### 5.1 SUSY Higgs Particles
The Higgs spectrum in the minimal supersymmetric extension of the Standard Model consists of five states: $`h^0,H^0,A^0`$ and $`H^\pm `$. Besides the masses, two mixing angles define the properties of the scalar particles and their interactions with gauge bosons, fermions and the self–interactions: the ratio of the two vacuum expectation values $`\mathrm{tg}\beta =v_2/v_1`$ and a mixing angle $`\alpha `$ in the neutral $`𝒞𝒫`$–even sector. Supersymmetry leads to several relations among the masses and mixing angles and, in fact, only two parameters are independent.
#### Neutral Higgs Bosons:
The lightest neutral Higgs boson will decay mainly into fermion pairs since its mass is smaller than $``$ 140 GeV. This is also the dominant decay mode of the pseudoscalar boson $`A^0`$. For values of $`\mathrm{tg}\beta `$ larger than unity and for masses less than $``$ 140 GeV, the main decay modes of the neutral Higgs bosons are decays into $`b\overline{b}`$ and $`\tau ^+\tau ^{}`$ pairs; the branching ratios are of order $`90\%`$ and $`8\%`$, respectively. The decays into $`c\overline{c}`$ pairs and gluons \[proceeding through $`t`$ and $`b`$ quark loops\] are suppressed, for large $`\mathrm{tg}\beta `$ strongly. For large masses, the top decay channels $`H^0,A^0t\overline{t}`$ open up; yet this mode remains suppressed for large $`\mathrm{tg}\beta `$. For large $`\mathrm{tg}\beta `$, the neutral Higgs bosons decay almost universally into $`b\overline{b}`$ and $`\tau ^+\tau ^{}`$ pairs. If the mass is high enough, the heavy $`𝒞𝒫`$–even Higgs boson $`H^0`$ can in principle decay into weak gauge bosons, $`H^0WW,ZZ`$. Since the partial widths are proportional to $`\mathrm{cos}^2(\beta \alpha )`$, they are strongly suppressed and the gold–plated $`ZZ`$ signal of the heavy SM Higgs boson is lost in the supersymmetric extension. The heavy neutral Higgs boson $`H^0`$ can also decay into two lighter Higgs bosons. These modes, however, are restricted to small domains in the parameter space. Other possible channels are decays into supersymmetric particles. While sfermions are likely too heavy to affect Higgs decays in the mass range considered here, Higgs boson decays into charginos and neutralinos could eventually be important since the masses of some of these particles are expected to be of order $`M_Z`$. \[These channels are summarized in Ref..\] The charged Higgs particles decay into fermions but also, if allowed kinematically, into the lightest neutral Higgs boson and a $`W`$ boson. Below the $`tb`$ and $`Wh`$ thresholds, the charged Higgs particles will decay mostly into $`\tau \overline{\nu }_\tau `$ and $`c\overline{s}`$ pairs, the former being dominant for $`\mathrm{tg}\beta >1`$. For large $`M_{H^\pm }`$ values, the top–bottom decay mode $`H^+t\overline{b}`$ becomes dominant.
Adding up the various decay modes, the width of all five Higgs bosons remains very narrow, being of order 10 GeV even for large masses.
The search for the neutral SUSY Higgs bosons at $`e^+e^{}`$ colliders will be a straightforward extension of the search performed at LEP2. This collider is expected to cover the mass range up to $`110`$ GeV for neutral Higgs bosons, depending on $`\mathrm{tg}\beta `$. Higher energies, $`\sqrt{s}250`$ GeV, are required to sweep the entire parameter space of the MSSM . The main production mechanisms of neutral Higgs bosons at $`e^+e^{}`$ colliders are the Higgs–strahlung process and associated pair production,
$`(a)\mathrm{Higgs}\mathrm{strahlung}:e^+e^{}`$ $``$ $`Z+h/H`$
$`(b)\mathrm{Pair}\mathrm{production}:e^+e^{}`$ $``$ $`A+h/H`$
as well as the related fusion processes. The $`𝒞𝒫`$–odd Higgs boson $`A^0`$ cannot be produced in fusion processes because it does not couple to gauge bosons in leading order.
The cross sections for the four Higgs–strahlung and pair production processes can be expressed as
$`\sigma (e^+e^{}Zh)`$ $`=`$ $`\mathrm{sin}^2(\beta \alpha )\sigma _{SM}`$
$`\sigma (e^+e^{}ZH)`$ $`=`$ $`\mathrm{cos}^2(\beta \alpha )\sigma _{SM}`$
and
$`\sigma (e^+e^{}Ah)`$ $`=`$ $`\mathrm{cos}^2(\beta \alpha )\overline{\lambda }\sigma _{SM}`$
$`\sigma (e^+e^{}AH)`$ $`=`$ $`\mathrm{sin}^2(\beta \alpha )\overline{\lambda }\sigma _{SM}`$
where $`\sigma _{SM}`$ is the SM cross section for Higgs–strahlung and the factor $`\overline{\lambda }`$ accounts for the suppression of the $`P`$–wave cross sections near the threshold. The cross sections for the Higgs–strahlung and for the pair production as well as the cross sections for the production of the light and the heavy neutral Higgs bosons $`h^0`$ and $`H^0`$ are mutually complementary to each other, coming either with a coefficient $`\mathrm{sin}^2(\beta \alpha )`$ or $`\mathrm{cos}^2(\beta \alpha )`$. As a result, since $`\sigma _{SM}`$ is large, at least the lightest $`𝒞𝒫`$–even Higgs boson must be detected.
The cross section for $`hZ`$ in the Higgs–strahlung process is large for values of $`M_h`$ near the upper bound. The heavy $`𝒞𝒫`$–even and $`𝒞𝒫`$–odd Higgs bosons $`H`$ and $`A`$, on the other hand, are produced pairwise in this limit: $`e^+e^{}AH`$. The decoupling limit becomes effective for heavy Higgs masses above 250 to 300 GeV. The discovery limit is therefore set by the beam energy independently of the mixing parameter $`\mathrm{tg}\beta `$, in contrast to LHC where the heavy Higgs bosons cannot be detected individually in parts of the parameter space.
#### Charged Higgs Bosons:
The charged Higgs bosons, if lighter than the top quark, can be produced in top decays, $`tb+H^+`$, with a branching ratio varying between $`2\%`$ and $`20\%`$ in the kinematically allowed region. Since for $`\mathrm{tg}\beta `$ larger than unity, the charged Higgs bosons will decay mainly into $`\tau \nu _\tau `$, this results in a surplus of $`\tau `$ final states over $`e,\mu `$ final states in $`t`$ decays, an apparent breaking of lepton universality. For large Higgs masses the dominant decay mode is the top decay $`H^+t\overline{b}`$. In this case the charged Higgs particles must be pair produced in $`e^+e^{}`$ colliders:
$$e^+e^{}H^+H^{}$$
The cross section depends only on the charged Higgs mass. For small Higgs masses the cross section is of order 100 fb at $`\sqrt{s}=500`$ GeV, but it drops very quickly due to the $`P`$–wave suppression $`\beta ^3`$ near the threshold.
SUMMARY: The preceding discussion of the MSSM Higgs sector at $`e^+e^{}`$ linear colliders can be summarized in the following two points:
$`(i)`$ The lightest $`𝒞𝒫`$–even Higgs particle $`h^0`$ can be detected in the entire range of the MSSM parameter space, either in the Higgs–strahlung process $`e^+e^{}hZ`$ or in pair production $`e^+e^{}hA`$ . In fact, this conclusion holds true even at a c.m. energy of 250 GeV, even if invisible neutralino decays are allowed for.
$`(ii)`$ The area in the parameter space where all SUSY Higgs bosons can be discovered individually at $`e^+e^{}`$ colliders, is characterized by $`M_A\stackrel{<}{}\frac{1}{2}\sqrt{s}`$, independently of $`\mathrm{tg}\beta `$. The $`H^0,A^0`$ Higgs bosons can be produced either in Higgs–strahlung or in $`Ah,AH`$ associated production; charged Higgs bosons will be produced in $`H^+H^{}`$ pairs up to the kinematical limit.
### 5.2 Supersymmetric Particles
#### Charginos and Neutralinos
The two charginos $`\stackrel{~}{\chi }_{}^{+}{}_{i}{}^{}`$ and the four neutralinos $`\stackrel{~}{\chi }_{}^{0}{}_{i}{}^{}`$, mixtures of the \[non–colored\] gauginos and higgsinos, are expected to be the lightest supersymmetric particles. In the MSSM with conserved $`R`$–parity, the neutralino $`\stackrel{~}{\chi }_{}^{0}{}_{1}{}^{}`$ with the smallest mass, assumed to be the lightest supersymmetric particle, is stable. The heavier neutralinos and the charginos decay into (possibly virtual) gauge and Higgs bosons plus the $`LSP`$, $`\stackrel{~}{\chi }_{}^{0}{}_{i}{}^{}\stackrel{~}{\chi }_{}^{0}{}_{1}{}^{}+Z`$ and $`\stackrel{~}{\chi }_{}^{0}{}_{1}{}^{}+W`$, or if they are heavy enough, into neutralino/chargino cascades, and leptons plus sleptons .
Neutralinos and charginos are easy to detect at $`e^+e^{}`$ colliders. They are produced in pairs
$`e^+e^{}`$ $``$ $`\stackrel{~}{\chi }_{}^{+}{}_{i}{}^{}\stackrel{~}{\chi }_j^{}[i,j=1,2]`$
$`e^+e^{}`$ $``$ $`\stackrel{~}{\chi }_{}^{0}{}_{i}{}^{}\stackrel{~}{\chi }_{}^{0}{}_{j}{}^{}[i,j=1,..,4]`$
through $`s`$–channel $`\gamma ,Z`$ exchange and $`t`$–channel selectron or sneutrino exchange. Since the cross sections are as large as $`𝒪`$(100 fb), enough events will be produced to discover these particles for masses nearly up to the kinematical limit.
The properties of the charginos and neutralinos can be studied in great detail at $`e^+e^{}`$ colliders. From the fast onset $`\beta `$ of the spin $`\frac{1}{2}`$ excitation curve near the threshold, the masses can be measured very accurately within less than 100 MeV, Fig. 4. Using polarized $`e^\pm `$ beams, the decomposition of the states, $`\stackrel{~}{\chi }_{}^{+}{}_{i}{}^{}=\alpha \stackrel{~}{W}^++\beta \stackrel{~}{H}^+`$ into wino and higgsino components can be determined since left–handed electrons couple to sneutrinos in the $`t`$–channel but right–handed electrons do not, so that the energy and angular dependence of the cross sections is different for the two polarization states . In a similar way the properties of neutralinos can be explored .
#### Sleptons and Squarks:
The superpartners of the right–handed leptons decay into the associated SM partners and neutralinos/charginos. In major parts of the SUSY parameter space the dominant decay mode is $`\stackrel{~}{\mu }_R\mu +\stackrel{~}{\chi }_{}^{0}{}_{1}{}^{}`$ . For the superpartners of the left–handed sleptons, the decay pattern is slightly more complicated since, besides the $`\stackrel{~}{\chi }_1^0`$ channels, decays into leptons and charginos are also possible. In $`e^+e^{}`$ and $`e^{}e^{}`$ collisions, sleptons are produced in pairs:
$`e^+e^{}`$ $``$ $`\stackrel{~}{\mu }_L^+\stackrel{~}{\mu }_L^{},\stackrel{~}{\mu }_R^+\stackrel{~}{\mu }_R^{},\stackrel{~}{\tau }_L^+\stackrel{~}{\tau }_L^{},\stackrel{~}{\tau }_R^+\stackrel{~}{\tau }_R^{}`$
$`e^+e^{}`$ $``$ $`\stackrel{~}{\nu }_L\overline{\stackrel{~}{\nu }}_L`$
$`e^+e^{}`$ $``$ $`\stackrel{~}{e}_L^+\stackrel{~}{e}_L^{},\stackrel{~}{e}_R^+\stackrel{~}{e}_R^{},\stackrel{~}{e}_L^+\stackrel{~}{e}_R^{},\stackrel{~}{e}_R^+\stackrel{~}{e}_L^{}`$
$`e^{}e^{}`$ $``$ $`\stackrel{~}{e}_L^{}\stackrel{~}{e}_L^{},\stackrel{~}{e}_R^{}\stackrel{~}{e}_R^{},\stackrel{~}{e}_L^{}\stackrel{~}{e}_R^{}`$
For charged sleptons, the production proceeds via $`\gamma ,Z`$ exchange in the $`s`$–channel, in the case of selectrons, also by additional $`t`$–channel neutralino exchange. For sneutrinos, the process is mediated by $`s`$–channel $`Z`$–exchange and, in the case of electron–sneutrinos, also by the $`t`$–channel exchange of charginos.
The cross sections for the pair production of sleptons are of the order of $`10^1`$ to $`10^2`$ pb so that their discovery is very easy up to the kinematical limit. From the $`P`$–wave onset $`\beta ^3`$ of the annihilation cross section the masses can in general be determined at a level of 200 to 300 MeV; the sharper onset of selection production in $`e^{}e^{}`$ scattering will reduce this number further. Enough events will be produced to study their detailed properties. The polarization of the $`e^\pm `$ beams will help to identify the couplings of these particles.
The endpoints in the decay spectra of $`\stackrel{~}{\mu }_R\mu +\stackrel{~}{\chi }_1^0`$ can be exploited to determine the LSP mass with an accuracy of 100 MeV, cf. Fig. 4.
If one of the stop states is light enough due to the strong $`LR`$ Yukawa mixing, these particles may be pair produced even at a 500 GeV collider:
$`e^+e^{}\stackrel{~}{t_i}\overline{\stackrel{~}{t_j}}i,j=1,2`$
By measuring the $`LR`$ asymmetry of the production cross sections, the $`\stackrel{~}{t}_L/\stackrel{~}{t}_R`$ mixing angle can be determined to high accuracy $`\delta \mathrm{cos}\theta _{\stackrel{~}{t}}0.01`$ .
### 5.3 Supersymmetry Breaking
The high precision with which masses, couplings and mixing parameters will be determined at $`e^+e^{}`$ colliders, can be exploited to test the mechanism for supersymmetry breaking and the structure of the underlying theory. In minimal supergravity mSUGRA the breaking of supersymmetry is mediated by gravity from a hidden sector to the eigenworld, generating soft SUSY breaking parameters at the grand unification scale<sup>1</sup><sup>1</sup>1 If gravitational interactions would become strong not at a very high scale but near the electroweak scale , collider experiments could probe the additional spatial dimensions through which gravitational fields would propagate . Contact interactions and missing energy events could signal Planck scales in 4 + n dimensions up to about 10 TeV. Thus, the basic space–time structure can be explored in these experiments. . The parameters are generally assumed to be universal at that scale in the gaugino and the scalar sectors. In gauge mediated supersymmetry breaking, gauge interactions connect the mechanism to the eigenworld at a scale possibly between 10 and $`10^3`$ TeV. Mass spectra in mSUGRA and GMSB are characteristically different, the splitting between sleptons and squarks being larger in GMSB. Moreover, since the gravitino in GMSB is very light, $`\stackrel{~}{\chi }_1^0`$ or $`\stackrel{~}{\tau _1}`$ may be long lived giving rise to displaced photons or stable heavy tracks in the decays $`\stackrel{~}{\chi }_1^0\gamma \stackrel{~}{G}`$ or $`\stackrel{~}{\tau _1}\tau \stackrel{~}{G}`$ with the gravitino $`\stackrel{~}{G}`$ escaping undetected .
If minimal supergravity is the underlying theory, the observable properties of the superparticles can be expressed by a small set of parameters defined at the GUT scale. Basically five parameters specify the properties of the particles in the supersymmetric sector. The scalar mass parameter $`m_0`$, the $`\mathrm{SU}(2)`$ gaugino mass $`M_{1/2}`$, the trilinear coupling $`A_0`$ and the sign of the coupling $`\mu `$ of the Higgs doublets in the superpotential, and $`\mathrm{tan}\beta `$, the ratio of the vacuum expectation values $`v_2/v_1`$. Evolving the scalar masses from the GUT scale down to low energies, it turns out that non–colored particles are in general significantly lighter than colored particles. The lightest of the non–colored gauginos/higgsinos and sleptons could have masses in the range of 100 to 200 GeV. Since only a few parameters determine the low energy theory of the evolution from the GUT scale down to the electroweak scale, many relations can be found among the masses of the superparticles which can stringently be tested at $`e^+e^{}`$ colliders . Two examples should illustrate the potential of $`e^+e^{}`$ facilities in this context.
$`(i)`$ The gaugino masses at the scale of $`\mathrm{SU}_2\times \mathrm{U}_1`$ symmetry breaking are related to the common $`\mathrm{SU}_2`$ gaugino mass $`M_{1/2}`$ at the GUT scale by the running gauge couplings:
$`M_i={\displaystyle \frac{\alpha _i}{\alpha _{GUT}}}M_{1/2}i=SU_3,SU_2,U_1`$
with $`\alpha _{GUT}`$ being the gauge coupling at the unification scale. The mass relation in the non–color sector
$`{\displaystyle \frac{M_1}{M_2}}={\displaystyle \frac{5}{3}}\mathrm{tan}^2\theta _W{\displaystyle \frac{1}{2}}`$
can be tested very well by measuring the masses and production cross sections of charginos/neutralinos and sleptons.
$`(ii)`$ In a similar way the slepton masses can be expressed in terms of a common scalar mass parameter $`m_0`$ at the GUT scale, contributions $`M_{1/2}`$ due to the evolution from the GUT scale down to low energies, and the D terms related to the electroweak symmetry breaking. These expressions give rise to relations among the slepton masses:
$`m^2(\stackrel{~}{l}_L)m^2(\stackrel{~}{\nu })`$ $`=`$ $`(12\mathrm{sin}^2\theta _W)\mathrm{cos}2\beta m_Z^2`$
$`m^2(\stackrel{~}{l}_L)m^2(\stackrel{~}{l}_R)`$ $`=`$ $`\kappa M_{1/2}^2{\scriptscriptstyle \frac{1}{2}}(14\mathrm{sin}^2\theta _W)\mathrm{cos}2\beta m_Z^2`$
with $`\kappa =0.37`$. The second relation follows from the hypothesis that the scalar masses are universal at the GUT scale, in particular $`m^2(5^{})=m^2(10)`$ within $`\mathrm{SU}_5`$. This assumption can be tested by relating the mass difference between $`\stackrel{~}{e}_L`$ and $`\stackrel{~}{e}_R`$ to the $`\mathrm{SU}_2`$ gaugino mass.
The typical result of an overall fit to the fundamental mSUGRA parameters at the GUT scale and tg$`\beta `$ is illustrated in Table 2.
In a procedure which reveals more clearly the structure of the underlying theory<sup>2</sup><sup>2</sup>2This procedure is particularly important for non-universal theories., the parameters may not only be fitted by assuming a universal set at the GUT scale from the start, but the set itself may be reconstructed by evolving the mass parameters from the electroweak scale to the unification scale , cf. Fig. 5. For gauginos and sleptons the reconstruction of the universal mass parameters is excellent. Due to mutual cancelations it is much more difficult for the squark and Higgs sectors. Nevertheless, this is the only way to reconstruct operationally the fundamental supersymmetric theory near the Planck scale from experimental observations at the electroweak scale.
Precision tests of supersymmetric particles in $`e^+e^{}`$ collider experiments can thus open a window to energy scales close to the Planck scale where gravity, the fourth of the fundamental forces, becomes an integral part of physics. |
warning/0003/physics0003073.html | ar5iv | text | # Temporary Acceleration of Electrons While Inside an Intense Electromagnetic Pulse
## Abstract
A free electron can temporarily gain a very significant amount of energy if it is overrun by an intense electromagnetic wave. In principle, this process would permit large enhancements in the center-of-mass energy of electron-electron, electron-positron and electron-photon interactions if these take place in the presence of an intense laser beam. Practical considerations severely limit the utility of this concept for contemporary lasers incident on relativistic electrons. A more accessible laboratory phenomenon is electron-positron production via an intense laser beam incident on a gas. Intense electromagnetic pulses of astrophysical origin can lead to very energetic photons via bremsstrahlung of temporarily accelerated electrons.
PACS numbers: 03.65.Sq, 12.15.-y 41.75.Fr, 52.40.Mj, 97.30.-b
The prospect of acceleration of charged particles by intense plane electromagnetic waves has excited interest since the suggestion by Menzel and Salisbury that this mechanism might provide an explanation for the origin of cosmic rays. However, it has generally been recognized that if a wave overtakes a free electron, the latter gains energy from the wave only so long as the electron is still in the wave, and reverts to its initial energy once the wave has past . There is some controversy as to the case of a “short” pulse of radiation, for which modest net energy transfer between a wave and electron appears possible . Acceleration via radiation pressure is negligible . It has been remarked that even in the case of a “long” pulse, some of the energy transferred from the wave to the electron can be extracted if the electron undergoes a scattering process while still inside the wave . This paper is an elaboration of that idea. We do not discuss here the observed phenomenon that an electron ionized from an atom in a strong wave can emerge from the wave with significant energy .
We consider a plane electromagnetic wave (often called the background wave) with dimensionless, invariant field strength
$$\eta =\frac{e\sqrt{A_\mu A^\mu }}{mc^2}=\frac{e_{\mathrm{rms}}}{m\omega _0c}=\frac{e_{\mathrm{rms}}\mathrm{\lambda ̄}_0}{mc^2}.$$
(1)
Here the wave has laboratory frequency $`\omega _0`$, reduced wavelength $`\mathrm{\lambda ̄}_0`$, root-mean-square electric field $`_{\mathrm{rms}}`$, and four-vector potential $`A_\mu `$; $`e`$ and $`m`$ are the charge and mass of the electron, and $`c`$ is the speed of light.
A practical realization of such a wave is a laser been. Laser beams with parameter $`\eta `$ close to one have been used in recent plasma-physics experiments and in high-energy-physics experiments .
When such a wave overtakes a free electron, the latter undergoes transverse oscillation (quiver motion), with relativistic velocities for η
>
1
>
𝜂1\eta\mathrel{\vbox{\kern 0.0pt\hbox{$>$} \kern 0.0pt\hbox{$\sim$} }}1 . The v $`\times `$ B force then couples the transverse oscillation to a longitudinal drift in the direction of the wave. In the nonrelativistic limit, this effect is often said to be due to the “ponderomotive potential” associated with the envelope of the electromagnetic pulse . The resulting temporary energy transfer to the longitudinal motion of the electron can in principle be arbitrarily large.
A semiclassical description of this process exists as well. A quantum-mechanical electron inside a classical plane wave can be described by the Volkov solutions to the Dirac equation . Such electrons are sometimes described as “dressed”, and they can be characterized by a quasimomentum
$$q=p+ϵk_0,$$
(2)
where the invariant $`ϵ`$ is given by
$$ϵ=\frac{m^2\eta ^2}{2(pk_0)},$$
(3)
with $`(pk_0)`$ being the 4-vector product of the 4-momenta $`p`$ of the electron and $`k_0`$ of a photon of the background wave. The factor $`ϵ`$ need not be an integer, and can be thought of as an effective number of wave photons “dragged” along with the electron as a result of a small difference between the large rates of absorption and emission (back into the wave) of wave photons by the electron. (Strictly speaking, the wave used in the Volkov solution is classical and, hence, contains no photons.) As a result, the electron inside the wave has an effective mass, $`\overline{m}`$, that is greater than its free mass $`m`$ :
$$\overline{m}^2=q^2=m^2(1+\eta ^2).$$
(4)
From a classical view, the quasimomentum $`q`$ is the result of averaging over the transverse oscillations (quiver motion) of the electron in the background wave. When discussing conservation of energy and momentum in the classical view, both transverse and longitudinal motion of the electron must be considered; but in a quantum analysis, quasimomentum is conserved and no mention is made of the classical transverse oscillations.
Throughout this paper the background wave propagates in the $`+z`$ direction, and the 4-momentum of a photon of this wave is written
$$k_0=(\omega _0,0,0,\omega _0).$$
(5)
From now on, we use units in which $`c`$ and $`\mathrm{}`$ equal one.
We first consider a relativistic electron moving along the $`+z`$ axis with 4-momentum
$$p=(E,0,0,P)=\gamma m(1,0,0,\beta ),$$
(6)
where $`E`$ and $`P`$ are the energy and the momentum of the electron prior to the arrival of the wave, $`\beta 1`$ is the electron’s velocity and $`\gamma =1/\sqrt{1\beta ^2}1`$. Then
$$(pk_0)=\omega _0(EP)=\frac{m^2\omega _0}{E+P},$$
(7)
so
$$ϵ=\frac{\eta ^2(E+P)}{2\omega _0}\frac{\gamma m\eta ^2}{\omega _0},$$
(8)
where the approximation holds for a relativistic electron. For a wave of optical frequencies (such as a laser), $`ϵ1`$. The quasienergy, $`q_0`$, is then large:
$$q_0=E(1+\eta ^2).$$
(9)
The electron has been accelerated from energy $`E`$ outside the wave to energy $`E(1+\eta ^2)`$ inside the wave. Since $`\eta `$ can in principle be large compared to 1, this acceleration can be very significant.
Can we take advantage of this acceleration in a high-energy-physics experiment? The example of Compton scattering of an electron by one laser beam while in a second laser beam has recently been reported elsewhere . Here, we consider examples of possibly enhanced production of electroweak gauge bosons in high-energy $`ee`$ and $`e\gamma `$ collisions in the presence of an intense laser.
Suppose the electron $`p`$ collides head-on with a positron $`p^{}`$, all inside the background wave. The positron 4-momentum is then
$$p^{}=(E^{},0,0,P^{}),$$
(10)
where $`E^{}m`$ in the relativistic case. Then
$$(p^{}k_0)=\omega _0(E^{}+P^{})2E^{}\omega _0.$$
(11)
The corresponding quasimomentum is
$$q^{}=p^{}+ϵ^{}k_0,$$
(12)
where
$$ϵ^{}=\frac{m^2\eta ^2}{2(p^{}k_0)}\frac{m^2\eta ^2}{2E^{}\omega _0}.$$
(13)
The factor $`ϵ^{}`$ is not large in general, and the energy of a relativistic positron (or electron) moving against an optical wave is almost unchanged.
However, the center-of-mass (cm) energy of the $`e^+e^{}`$ system is increased when the collision occurs inside the background wave. The cm-energy squared is
$$s=(q+q^{})^24EE^{}(1+\eta ^2),$$
(14)
which is enhanced by a factor $`1+\eta ^2`$ compared to the case of no background wave.
For example, the $`Z^0`$ boson could be produced in $`e^+e^{}`$ collisions with 33- rather than 46.6-GeV beams, if the collision took place inside a background wave of strength $`\eta =1`$.
Of course, the background wave Compton scatters off the positron beam at a high rate if η
>
1
>
𝜂1\eta\mathrel{\vbox{\kern 0.0pt\hbox{$>$} \kern 0.0pt\hbox{$\sim$} }}1, which results in substantial smearing of the energy of that beam. In practice, the cm-energy enhancement by a background wave would not be very useful in $`e^+e^{}`$ or $`ee`$ collisions.
Note, however, that Compton scattering is insignificant when the background wave and electron move in the same direction, unless the wave is extraordinarily strong. By an application of the Larmor formula in the (average) rest frame of the electron, we find that the fraction of the electron’s (laboratory) energy radiated in one cycle of its motion in the wave is of order $`\alpha \eta ^2(\omega _0/E)`$, where $`\alpha `$ is the fine-structure constant.
Suppose instead that the electron collides head-on with a high-energy photon of frequency $`\omega `$ and 4-momentum
$$p^{}=k=(\omega ,0,0,\omega ).$$
(15)
Then eq. (14) holds on substituting $`\omega `$ for $`E^{}`$; the cm-energy squared is again enhanced by the factor $`1+\eta ^2`$.
The background wave can, of course, interact directly with the high-energy photon to produce $`e^+e^{}`$ pairs, but if $`4\omega \omega _0<m^2(1+\eta ^2)`$, the pair-production rate is much suppressed . Thus, there is a regime in which $`e`$ \+ photon collisions in a strong background wave are cleaner than $`e^+e^{}`$ or $`ee`$ collisions in the wave.
In practice, we could get the high-energy photon from Compton scattering of the background wave off an electron beam. One might not want to backscatter the wave off a positron beam because of “backgrounds” from $`e^+e^{}Z^0`$.
A physics topic of interest would be the reaction
$$k+e^{}W^{}+\nu ,$$
(16)
which proceeds via the triple-gauge-boson coupling $`\gamma WW`$, and whose angular distribution is sensitive to the magnetic moment of the $`W`$ boson . The “background” process
$$k+e^{}Z^0+e^{}$$
(17)
could be suppressed by suitable choice of polarization of the electron and background wave.
For electron beams of 46.6 GeV as at the Stanford Linear Accelerator Center, green laser light backscatters into photons of energies up to about 30 GeV. Thus if the laser had $`\eta =1`$, the cm energy would extend up to 106 GeV, well above the threshold for reactions (16-17).
However, the enhancement factors $`1+\eta ^2`$ in the electron energy, eq. (9), and in the cm-energy squared, eq. (14), of $`ee`$ or electron-photon collisions are very much dependent on the idealization that the background wave is highly collinear with the electron.
We reconsider the preceding, but now suppose that the electron makes angle $`\theta 1`$ to the $`z`$ axis, The 4-momentum of the electron is
$$p=(E,P\mathrm{sin}\theta ,0,P\mathrm{cos}\theta ),$$
(18)
and
$$(pk_0)=E\omega _0(1\beta \mathrm{cos}\theta )\frac{m\omega _0}{2\gamma }(1+\gamma ^2\theta ^2).$$
(19)
As a consequence, the (quasi)energy of the electron inside the wave is now
$$q_0=p_0+\frac{m^2\eta ^2\omega _0}{2(pk_0)}E\left(1+\frac{\eta ^2}{1+\gamma ^2\theta ^2}\right),$$
(20)
which reduces to eq. (9) as $`\theta `$ goes to zero. However, if $`\theta >\eta /\gamma `$, then the electron is hardly accelerated by the background wave.
Electrons of present interest in high-energy physics typically have energies in the range 1-1000 GeV, corresponding to $`\gamma 10^3`$-$`10^6`$. This places very severe requirements on the alignment of the background wave with the electron beam. Indeed, the angular divergence of an electron beam is often larger than $`1/\gamma `$, so that no alignment of the background wave could impart large energy enhancements to the entire beam.
Furthermore, optical waves with $`\eta 1`$ can only be obtained at present in focused laser beams for which the characteristic angular spread is Δθ
>
0.1
>
Δ𝜃0.1\Delta\theta\mathrel{\vbox{\kern 0.0pt\hbox{$>$} \kern 0.0pt\hbox{$\sim$} }}0.1. So even if the central angle of the beam could be aligned to better than $`1/\gamma `$, only a very small fraction of the beam power would lie within a cone of that angle.
We also note that for the quasimomentum $`q`$ to be meaningful, the electron must have resided inside the strong background field for at least one cycle. A relativistic electron moves distance $`2\gamma ^2(1+\eta ^2)\lambda _0`$ while the background wave advances one wavelength relative to the electron . However, the strong-field region of a focused laser is characterized by its Rayleigh range, which is typically a few hundred wavelengths when $`\eta 1`$. Further, the transverse extent of the (classical) trajectory is of order $`\gamma \eta \lambda _0`$. Hence, in present laser systems, the strong-field region is not extensive enough that the energy transfer (9) could be realized for γ
>
10
>
𝛾10\gamma\mathrel{\vbox{\kern 0.0pt\hbox{$>$} \kern 0.0pt\hbox{$\sim$} }}10.
While physical consequences of the temporary acceleration of relativistic electrons inside an intense laser beam may be difficult to demonstrate, there is also interest in the case where the electron is initially at rest, or nearly so, such as electrons ionized from gas atoms by the passage of the background laser pulse .
An interesting process is so-called trident production,
$$e+Ae^{}+A^{}+e^+e^{},$$
(21)
of an electron-positron pair in the interaction of an ionization electron with a nucleus $`A`$ of a gas atom. For a very heavy nucleus $`A`$, its final state $`A^{}`$ has a different momentum but the same energy. Then the initial electron must provide the energy to create the $`e^+e^{}`$ pair as well as that for the final electron. The least energy required is when all three final-state electrons and positrons are at “rest” (i.e., they have zero net longitudinal momentum; they must always have quiver motion when they are in the wave). Thus, the minimum total quasienergy of the final-state electrons and positrons is $`3\overline{m}`$.
We conclude that the quasienergy $`q_0`$ of the initial electron must be at least $`3\overline{m}`$ for reaction (21) to occur.
If the electron is at rest prior to the arrival of the background wave its 4-momentum is
$$p=(m,0,0,0).$$
(22)
As the electron is overtaken by a wave of strength $`\eta `$ and 4-momentum given by (5), it takes on quasimomentum
$$q=(m(1+\eta ^2/2),0,0,m\eta ^2/2)(\overline{m}\gamma ,0,0,\overline{m}\gamma \beta _z).$$
(23)
Thus, the net longitudinal velocity of the electron inside the wave is $`\beta _z=q_z/q_0=(\eta ^2/2)/(1+\eta ^2/2)`$. As expected, inside a very strong wave the electron can take on relativistic longitudinal motion.
We could have trident production while the electron is still in the wave if the quasienergy $`q_0=m(1+\eta ^2/2)`$ exceeds $`3\overline{m}`$. For an electron initially at rest, this requires $`\eta \sqrt{16+12\sqrt{2}}=5.74`$.
The trident process is still possible within a wave with $`\eta <5.74`$ provided the electron has quasienergy $`q_03\overline{m}`$. This might arise, for example, because of acceleration of the electron by the plasma-wakefield effect .
It is conceivable that the electron creates the pair in a linearly polarized wave at a phase when its (classical) kinetic energy is high, but the final electron and the pair all appear with a lower kinetic energy corresponding to some other phase of the wave. This can’t happen if the interaction takes place at a well-defined point, since the phase of the wave is a unique function space and time. It might occur if the final particles “tunnel” to another space-time point before appearing, and the instantaneous kinetic energy is lower at that point.
However, we will find shortly that such tunneling is not consistent with energy conservation. To be as definite as possible, we consider ordinary energy along the classical trajectories, rather than quasimomentum. The latter is taken into account in the sense that the electron and positron are not created at rest, but with the transverse velocities appropriate to phase of the background wave at the spacetime point at which the pair appears. It is sufficient to consider only those trajectories with zero average momentum (i.e., zero quasi-3-momentum).
For circular polarization of the background wave, the electron trajectory is a circle in the plane perpendicular to the $`z`$ axis, with radius $`a/\omega _0`$, velocity $`\beta =a`$ and Lorentz factor
$$\gamma _{\mathrm{circ}}=\frac{1}{\sqrt{1a^2}}=\sqrt{1+\eta ^2},$$
(24)
where parameter $`a`$ is given by
$$a^2=\frac{\eta ^2}{1+\eta ^2},0a^21.$$
(25)
For a background wave that is linearly polarized in the $`x`$-direction, the trajectory can be parametrized as
$$x=\sqrt{2}\frac{a}{\omega _0}\mathrm{sin}\delta ,z=\frac{a^2}{4\omega _0}\mathrm{sin}2\delta ,$$
(26)
where $`\delta =\omega _0\tau \sqrt{1+\eta ^2}=\omega _0\tau /\sqrt{1a^2}`$, and $`\tau `$ is the proper time. Expression (26) describes the well-known figure-8 trajectory. Now $`dx/d\tau =(dx/dt)(dt/d\tau )=\gamma \beta _x`$, so $`\gamma ^2=1+\gamma ^2\beta ^2=1+(dx/d\tau )^2+(dz/d\tau )^2`$. We find that
$$\gamma _{\mathrm{lin}}=\frac{1+\frac{1}{2}[a^2(\omega _0x)^2]}{\sqrt{1a^2}}.$$
(27)
From expression (26) for the $`x`$-trajectory we see that $`0(\omega _0x)^22a^2`$, so
$$\gamma _{\mathrm{min}}=\frac{1+\eta ^2/2}{\sqrt{1+\eta ^2}},\text{and}\gamma _{\mathrm{max}}=\frac{1+3\eta ^2/2}{\sqrt{1+\eta ^2}}.$$
(28)
These values surround the result that $`\gamma _{\mathrm{circ}}=\sqrt{1+\eta ^2}`$ always for circular polarization. For small $`\eta `$, $`\gamma _{\mathrm{min}}1+\eta ^4/8`$, $`\gamma _{\mathrm{max}}1+\eta ^2`$, and $`\gamma _{\mathrm{circ}}1+\eta ^2/2`$; for large $`\eta `$, $`\gamma _{\mathrm{min}}\eta /2`$, $`\gamma _{\mathrm{max}}3\eta /2`$, and $`\gamma _{\mathrm{circ}}\eta `$.
Suppose an electron interacts with a nucleus at the place where its Lorentz factor is $`\gamma _{\mathrm{max}}`$ and reappears along with an electron-positron pair at a location where $`\gamma _{\mathrm{min}}`$ holds at that moment. The nucleus absorbs the excess momentum of the initial electron. Conservation of (ordinary) energy requires that $`\gamma _{\mathrm{max}}=3\gamma _{\mathrm{min}}`$. But this is not satisfied for any value of $`\eta `$ according to (28). That is, the hypothetical tunneling process is not possible under any circumstances.
In sum, even when in a background wave an electron can produce positrons off nuclei only if the electron has sufficient longitudinal momentum that the corresponding (quasi)energy is three times the (effective) electron mass.
We close by returning to the astrophysical context that began the historical debate on acceleration by intense electromagnetic waves. Gunn and Ostriker have given an extensive discussion the possibility of electron acceleration in the rotating dipole field of a millisecond pulsar, where the field strength $`\eta `$ can be of order $`10^{10}`$. Their argument does not primarily address free electrons overtaken by a wave, but rather electrons “injected” or “dropped at rest’ into the wave. Neutron decay is a candidate process for injection. In very strong fields ($`\eta 1)`$ this decay takes place together with the absorption by the electron (and proton) of a very large number of wave photons, so that the electron is created with (quasi)energy $`m\eta ^2/2`$ (compare eq. (23)) . Because the fields of the pulsar fall off as $`1/r`$ where (coincidentally) $`r_{\mathrm{pulsar}}\lambda _0`$, the wavelength of the rotating dipole radiation, the field region is “short”, and the electron may emerge with some fraction of the large energy it had at the moment of its creation.
An example closer to the theme of the present paper would be an electron that is overtaken by the intense electromagnetic pulse of a supernova (or other transient astrophysical occurrence, perhaps including gamma-ray bursters), and thereby temporarily accelerated to energy $`m\eta ^2/2`$. Such pulses could have significant fields at optical frequencies, where the transverse scale, $`\eta \lambda _0`$, of the motion of accelerated electrons is less than the Chandrasekhar radius for $`\eta <10^{10}`$. In general, the electron has low energy before and after the passage of the pulse. However, high-energy photons can arise via bremsstrahlung of the electron when it interacts with a plasma nucleus while still in the pulse. In this view, the primary astrophysical evidence of temporarily accelerated electrons would be high-energy photons which, of course, could transfer some of their energy to protons and other charged particles in subsequent interactions.
This work was supported in part by DoE grants DE-FG02-91ER40671 and DE-FG05-91ER40627. |
warning/0003/hep-th0003248.html | ar5iv | text | # Introduction
## Introduction
Randall and Sundrum’s recent proposal for an alternative to standard Kaluza-Klein (KK) compactification in Refs. has attracted a lot of attention from many quarters: from a phenomenological point of view, it is a new and fresh proposal to understand the hierarchy between gauge and gravitational interactions, while from a purely gravitational point of view it rises many interesting problems concerning the relation between bulk and brane gravitational phenomena. In any case, these models provide a new arena in which one can study new and old problems of Theoretical Physics.
It is worth trying to extend this framework. Here we will present generalizations of the Randall-Sundrum (RS) scenario which could be used as alternatives to KK compactification. They are solutions of the Einstein equations with arbitrary cosmological constant in $`\widehat{d}`$ dimensions and lead to $`d=(\widehat{d}1)`$-dimensional metrics which solve the Einstein equations with arbitrary $`d`$-dimensional cosmological constant. They have a property which one should require of any framework with extra dimensions: when the $`d`$-dimensional metric is maximally symmetric (and, therefore, is the lower-dimensional vacuum) the corresponding $`\widehat{d}`$-dimensional metric is also maximally symmetric (and, therefore, the upper-dimensional vacuum). This holds in any consistent standard KK compactification: vanishing matter fields and Minkowski metric in lower dimensions (the $`d`$-dimensional vacuum) correspond to Minkowski times a torus metric in upper dimensions (the $`\widehat{d}`$-dimensional vacuum). The same can be said of supersymmetry although there are subtleties that in many cases will make impossible to define lower-dimensional supersymmetry.
In order to exploit these “bulk” solutions for dimensional reduction, we introduce brane sources, find the modified solutions and study the dynamics of gravitons in the new backgrounds. We also find the effective gravity actions and Newton constants in lower dimensions and study supersymmetry on the brane-worlds.
In finding the solutions with branes, we have to allow for cosmological constants that are piecewise constant functions, a fact which comes naturally when dualizing it. Therefore we consider gravity coupled to a volume-form field strength and coupled to a generic $`(\widehat{d}2)`$-brane action.
## 1 Bulk Solutions
We are interested in “warped metrics” of the form<sup>4</sup><sup>4</sup>4We work in arbitrary dimension $`\widehat{d}`$ with mostly minus signature. All $`\widehat{d}`$-dimensional objects carry hats. We choose $`x^{\widehat{d}1}y`$ as the spacelike holographic coordinate and thus, we split the $`\{\widehat{x}^{\widehat{\mu }}\}=\{x^\mu ,y\}`$. Unhatted objects are $`d(=\widehat{d}1)`$-dimensional.
$$d\widehat{s}^2=a^2(y)ds^2dy^2,ds^2=g_{\mu \nu }(x)dx^\mu dx^\nu ,$$
(1.1)
solving the equations
$$\widehat{R}^{\widehat{\mu }\widehat{\nu }}=\widehat{\mathrm{\Lambda }}\widehat{g}^{\widehat{\mu }\widehat{\nu }},R^{\mu \nu }=\mathrm{\Lambda }g^{\mu \nu },$$
(1.2)
where $`\widehat{\mathrm{\Lambda }}`$ and $`\mathrm{\Lambda }`$ are respectively the $`\widehat{d}`$ and $`d`$-dimensional cosmological constants whose signs are, in principle, arbitrary. We define $`\widehat{g}`$ and $`g`$ by
$$\widehat{g}^2=\frac{\widehat{\mathrm{\Lambda }}}{(\widehat{d}1)},g^2=\frac{\mathrm{\Lambda }}{(d1)}.$$
(1.3)
The solutions fall into two classes:
$$\mathrm{𝟏}.\widehat{𝐠}\mathrm{𝟎}a(y)=\frac{1}{2}\sqrt{\pm g^2/\widehat{g}^2}\left(e^{\widehat{g}y}\pm e^{\widehat{g}y}\right),$$
(1.4)
where the sign has to be chosen such as to make $`a(y)`$ real. This is always possible except for the case $`\widehat{g}𝕀,g`$. In the other cases we have<sup>5</sup><sup>5</sup>5The case in which $`a`$ is not real can be fixed by Wick-rotating $`y`$ into a timelike coordinate., with $`g0`$
$$(𝐚)\widehat{𝐠},𝐠a=g/\widehat{g}\mathrm{cosh}\widehat{g}y,$$
(1.5)
$$(𝐛)\widehat{𝐠},𝐠𝕀a=ig/\widehat{g}\mathrm{sinh}\widehat{g}y,$$
(1.6)
$$(𝐜)\widehat{𝐠},𝐠𝕀a=g/\widehat{g}\mathrm{cos}i\widehat{g}y.$$
(1.7)
In this case, the coordinate $`y`$ naturally lives in a circle of length $`\frac{2\pi }{i\widehat{g}}`$.
With $`g=0`$ the only possibility is $`\widehat{𝐠}`$ and
$$a=e^{\pm \widehat{g}y}.$$
(1.8)
$$\mathrm{𝟐}.\widehat{𝐠}=\mathrm{𝟎}a=igy,$$
(1.9)
which means that we must have $`g𝕀`$.
The main property of these solutions is that, if $`g_{\mu \nu }`$ is the maximally symmetric metric in $`d`$ dimensions with curvature given by $`\mathrm{\Lambda }`$, then $`\widehat{g}_{\widehat{\mu }\widehat{\nu }}`$ is the maximally symmetric metric with curvature given by $`\widehat{\mathrm{\Lambda }}`$. The RS solution fits in the $`\widehat{g},g=0`$ case: if $`g_{\mu \nu }=\eta _{\mu \nu }`$, we have upstairs (locally) anti-De Sitter (aDS). Other possibilities that we are introducing here are: to have either aDS or DS both upstairs and downstairs, to have aDS upstairs and DS downstairs and to have Minkowski upstairs and DS downstairs. The most interesting options (at least from the supersymmetry point of view) are the RS solution and the one with Minkowski upstairs and DS downstairs.
In any dimension, in absence of other fields, the gravitino supersymmetry transformation law will take the form<sup>6</sup><sup>6</sup>6Depending on the dimension, we will have one or another kind of minimal spinors associated to representations of the gamma matrices with special properties. This will never be an issue in what follows and our results can be adapted to all the cases of interest.
$$\delta _{\widehat{ϵ}}\widehat{\psi }_{\widehat{\mu }}=\widehat{𝒟}_{\widehat{\mu }}\widehat{ϵ},$$
(1.10)
where $`\widehat{𝒟}_{\widehat{\mu }}`$ is the aDS ($`\widehat{g}`$) or Lorentz ($`\widehat{g}=0`$) covariant derivative<sup>7</sup><sup>7</sup>7Formally we can also consider the DS case ($`\widehat{g}𝕀`$). DS supergravities do exist even though they are inconsistent as quantum theories.
$$\widehat{𝒟}_{\widehat{\mu }}=_{\widehat{\mu }}\frac{1}{4}\widehat{\omega }_{\widehat{\mu }}{}_{}{}^{\widehat{a}\widehat{b}}\widehat{\gamma }_{\widehat{a}\widehat{b}}^{}\frac{i}{2}\widehat{g}\widehat{\gamma }_{\widehat{\mu }}.$$
(1.11)
Then, the Killing-spinor equation $`\delta _{\widehat{ϵ}}\widehat{\psi }_{\widehat{\mu }}=0`$ has the following solutions:
$$\mathrm{𝟏}.\widehat{𝐠}\mathrm{𝟎}\widehat{ϵ}=\frac{1}{2}\left(e^{\widehat{g}y/2}+\phi e^{\widehat{g}y/2}\right)ϵ_++\frac{1}{2}\left(e^{\widehat{g}y/2}\phi e^{\widehat{g}y/2}\right)ϵ_{},$$
(1.12)
where $`\phi =(g/\widehat{g})/|g/\widehat{g}|`$ and where $`ϵ_\pm `$ are two spinors that satisfy
$$\left(𝒟_\mu \frac{i}{2}g\gamma _\mu \right)ϵ_\pm =0,$$
(1.13)
$`𝒟_\mu `$ being the standard Lorentz covariant derivative and $`\gamma _a\widehat{\gamma }_a`$. These equations have maximal number of solutions when the $`d`$-dimensional space is maximally symmetric.
$`\mathrm{𝟐}.\widehat{𝐠}=\mathrm{𝟎}`$. The solution in this case is any $`y`$-independent spinor $`\widehat{ϵ}`$ satisfying
$$\left(𝒟_\mu \frac{i}{2}g\gamma _\mu \right)\widehat{ϵ}=0,$$
(1.14)
where now $`\gamma _a\widehat{\gamma }_a\widehat{\gamma }_y`$. In this case we had to take $`g𝕀`$ and thus this is the $`d`$-dimensional DS covariant derivative. This equation has a maximal number of solutions when the $`d`$-dimensional spacetime is DS.
Observe that, although DS supergravity is inconsistent, any pure gravity solution of that theory can be considered a warped compactification of standard (Poincaré) supergravity in one dimension more.
Although we have managed to reduce the $`\widehat{d}`$-dimensional Killing-spinor equation to a $`d`$-dimensional-looking Killing-spinor equation, this does not mean that we have supersymmetry in the $`d`$-dimensional space. In the $`\widehat{g}0`$ case, we cannot have two different signs for $`g`$. Keeping only one means keeping either $`ϵ_+`$ or $`ϵ_{}`$, but this truncation is only consistent with $`d`$-dimensional Lorentz invariance when $`g=0`$ (the RS case). On the other hand in the $`\widehat{g}=0`$ it seems that there is no problem to have DS supersymmetry. The supersymmetry of the RS solution has also been studied in Refs. and . We will make further comments on their results in the next section.
## 2 Brane-World Solutions
Now, mimicking Randall and Sundrum we consider the gravity plus brane-sources equations
$$\begin{array}{ccc}\hfill \widehat{R}^{\widehat{\mu }\widehat{\nu }}& =& \widehat{\mathrm{\Lambda }}\widehat{g}^{\widehat{\mu }\widehat{\nu }}\widehat{\chi }\left[g^{\rho \sigma }\delta _\rho {}_{}{}^{\widehat{\mu }}\delta _{\sigma }^{}{}_{}{}^{\widehat{\nu }}\frac{1}{\widehat{d}2}\widehat{g}^{\widehat{\mu }\widehat{\nu }}(\widehat{g}^{\rho \sigma }\widehat{g}_{\rho \sigma })\right]_nT_n\delta (yy_n),\hfill \\ & & \\ \hfill R^{\mu \nu }& =& \mathrm{\Lambda }g^{\mu \nu }.\hfill \end{array}$$
(2.1)
Although we write cosmological constants, we will have to allow for piecewise constant functions of $`y`$. Then, by making identifications if necessary we can restrict ourselves to a domain in which they are really constant.
With the same Ansatz for the metric Eq. (1.1) these equations reduce to
$$\{\begin{array}{ccc}\hfill 0& =& a^{\prime \prime }+\frac{\widehat{\mathrm{\Lambda }}}{\widehat{d}1}a+\frac{2\widehat{\chi }}{\widehat{d}2}a_nT_n\delta (yy_n),\hfill \\ & & \\ \hfill 0& =& (a^{})^2+\frac{\widehat{\mathrm{\Lambda }}}{\widehat{d}1}a^2\frac{\mathrm{\Lambda }}{d1}.\hfill \end{array}$$
(2.2)
It is straightforward to see that the solutions take now the form
$$\mathrm{𝟏}.\widehat{𝐠},𝐠\mathrm{𝟎}a(y)=\frac{1}{2}\sqrt{\pm g^2/\widehat{g}^2}\left(e^{_nc_n|yy_n|+C}\pm e^{_nc_n|yy_n|C}\right),$$
(2.3)
where $`\widehat{g}`$ and $`g`$ are defined as before but now $`\widehat{g}`$ takes the value
$$\widehat{g}=\underset{n}{}c_n[2\theta (yy_n)1],$$
(2.4)
and $`g`$ is proportional to $`\widehat{g}`$ with an arbitrary proportionality constant so $`\widehat{g}/g`$ is a true (purely real or imaginary) constant. The simultaneously purely real or imaginary constants $`c_n`$ are given by
$$c_n=\frac{\widehat{\chi }T_n}{2(\widehat{d}2)}\mathrm{tanh}^1\left(\underset{m}{}c_m|yy_m|+C\right)|_{y=y_n}.$$
(2.5)
$$\mathrm{𝟐}.\widehat{𝐠}\mathrm{𝟎},𝐠=\mathrm{𝟎}a(y)=e^{_nc_n|yy_n|},$$
(2.6)
where $`\widehat{g}`$ and the simultaneously purely real constants $`c_n`$ are given by
$$\widehat{g}=\underset{n}{}c_n[2\theta (yy_n)1],c_n=\frac{\widehat{\chi }T_n}{2(\widehat{d}2)},$$
(2.7)
so
$$a(y)=e^{\frac{\widehat{\chi }}{2(\widehat{d}2)}_nT_n|yy_n|}.$$
(2.8)
$$\mathrm{𝟑}.\widehat{𝐠}=\mathrm{𝟎}a=\underset{n}{}c_n|yy_n|+C,$$
(2.9)
with
$$g=\underset{n}{}c_n[2\theta (yy_n)1],c_n=\frac{\widehat{\chi }T_n}{2(\widehat{d}2)}\frac{1}{_mc_m|yy_m|+C}|_{y=y_n}.$$
(2.10)
In general the equations for the constants $`c_n`$ only have solution if all of them (and, therefore, the tensions $`T_n`$) have the same sign. In particular, a system with two branes only has solution if both branes have the same tension. The exception is the $`g=0`$ (RS) case in which one can get solutions for arbitrary tensions (Eq. (2.8)).
The problem of finding the different $`c_n`$’s does not show up if one considers an infinite periodic array of branes and anti-branes with opposite tensions. We can restrict ourselves to a fundamental region bounded by two branes or anti-branes with an anti-brane (resp. brane) in the middle. The system is mirror symmetric with respect to the middle (anti-) brane and we can make a further $`_2`$ identification that leaves us with a piece of spacetime bounded by a brane and an anti-brane in which $`\widehat{\mathrm{\Lambda }}`$ and $`\mathrm{\Lambda }`$ are constant (and in which only one constant $`c_n`$ matters). In these conditions, taking as fundamental region the interval $`y[0,\mathrm{}/2]`$ with an anti-brane placed at $`y=0`$ and a brane at $`y=\mathrm{}/2`$ the warp function $`a(y)`$ takes the same form as if there was only one brane in the whole spacetime:
1. $`\widehat{𝐠},𝐠\mathrm{𝟎}`$
$$(𝐚)\widehat{𝐠},𝐠a=g/\widehat{g}\mathrm{cosh}\left(\widehat{g}|y|+C\right),\widehat{g}=\frac{\widehat{\chi }T\mathrm{coth}(C)}{\widehat{d}2}.$$
(2.11)
$$(𝐛)\widehat{𝐠},𝐠𝕀a=ig/\widehat{g}\mathrm{sinh}\left(\widehat{g}|y|+C\right),\widehat{g}=\frac{\widehat{\chi }T\mathrm{tanh}(C)}{\widehat{d}2}.$$
(2.12)
$$(𝐜)\widehat{𝐠},𝐠𝕀a=g/\widehat{g}\mathrm{cos}\left(i\widehat{g}|y|+C\right),\widehat{g}=i\frac{\widehat{\chi }T\mathrm{coth}(C)}{\widehat{d}2}.$$
(2.13)
In this case, $`\mathrm{}`$ must be an integer fraction of the period of $`y`$ i.e. $`\mathrm{}=\frac{2\pi }{in\widehat{g}}`$.
2. $`\widehat{𝐠}\mathrm{𝟎},𝐠=\mathrm{𝟎}`$. $`\widehat{𝐠}`$
$$a=e^{\widehat{g}|y|},\widehat{g}=\widehat{\chi }T/2.$$
(2.14)
3. $`\widehat{𝐠}=\mathrm{𝟎}`$
$$a=ig|y|+C,g=i\frac{\widehat{\chi }TC}{\widehat{d}2}.$$
(2.15)
Let us now consider the bulk and world-brane supersymmetry of these solutions. We can only have supersymmetry on the brane in the RS case $`\widehat{g}0,g=0`$ and in the Minkowski-DS case $`\widehat{g}=0,g0`$ and imaginary. In these two cases the amount of supersymmetry preserved depends on the $`d`$-dimensional (brane) metric $`g_{\mu \nu }`$. If it is maximally symmetric, then there will be maximal supersymmetry on the brane.
Generic branes generically break $`\widehat{d}`$-dimensional bulk supersymmetry.<sup>8</sup><sup>8</sup>8We are not going to include sources in the supersymmetry transformation rules as in Ref. . We think one really needs proper $`\kappa `$-symmetric brane-sources in order to study in a fully consistent way the supersymmetric source problem. However, in these cases, supersymmetry is not broken locally in the bulk, in between any pair of branes, since there the metric has exactly the same form as in the absence of branes.
One may want to have unbroken supersymmetry globally, an not just in between the branes. First, we need to be able to define the Killing-spinor equation globally. In order to do this, we have to allow for a $`\widehat{g}`$ which is piecewise constant instead of globally constant (the main characteristic of these branes is that the value of $`\widehat{g}`$ is different in both sides). We have implicitly accepted this generalization in this section in order to find the solutions. On the other hand, one can use a dual formulation in which the cosmological constant is replaced by a $`d`$-form potential as in Ref. , an idea which will be investigated in Sec. (4). Once we accept this generalization, the necessary condition to have global unbroken supersymmetry is to be able to match the solutions of the Killing-spinor equation in both sides of a given brane. Let us take, for simplicity, one brane placed at $`y=0`$. Both $`\widehat{g}`$ and $`g`$ change sign across the brane. In the $`y>0`$ side of the brane, the solutions of the Killing-spinor equation are those exhibited in the previous section. In the $`y<0`$ side of the brane we find solutions of the same form, where, now, the spinors appearing in the general solution satisfy
the same equations but with the sign of $`g`$ reversed. We need to set $`g=0`$ which means that in the second case all supersymmetry is broken unless we have a trivial solution.
In the first case, it is not enough to have $`g=0`$ which brings us the the RS case again. It turns out that we also need to impose the condition
$$i\widehat{\gamma }_y\widehat{ϵ}=+\widehat{ϵ},$$
(2.16)
on the Killing-spinor, which reduces supersymmetry to a half. This is the same condition we would impose if we were orbifolding the space between branes.
We would like to stress that our results apply strictly to the cases we are considering: the infinitely thin branes described by the above solutions which make the metric across them discontinuous. Thus, our results do not contradict those of Linde and Kallosh who did not study just pure supergravity but included supersymmetric matter. In that paper the authors tried to find supersymmetric thick domain walls for which the metric is smooth using consistent superpotentials but did not find any.
### 2.1 4-$`d`$ Action and Newton Constant
The action from which the equations of motion Eqs. (2.1) follow has the form
$$\widehat{S}=\frac{1}{2\widehat{\chi }}d^{\widehat{d}}\widehat{x}\sqrt{|\widehat{g}|}\left[\widehat{R}(\widehat{d}2)\widehat{\mathrm{\Lambda }}\right]+\mathrm{branes},$$
(2.17)
and for $`\widehat{d}`$-dimensional metrics of the warped form Eq. (1.1) it reduces to
$$S=\frac{1}{2\widehat{\chi }}𝑑ya^{\widehat{d}3}d^dx\sqrt{|g|}\left[R(d2)\mathrm{\Lambda }\right].$$
(2.18)
Comparing, we find that the $`d`$-dimensional Newton constant $`\chi `$ is related to the $`\widehat{d}`$-dimensional one $`\widehat{\chi }`$ by
$$\widehat{\chi }/\chi =𝑑ya^{\widehat{d}3}.$$
(2.19)
Taking $`\widehat{d}=5`$ for definiteness, we can calculate the proportionality factor in the different cases:
Case 1.a: $`a=g/\widehat{g}\mathrm{cosh}\widehat{g}|y|`$
$$\chi =2\frac{\widehat{g}^3/g^2}{\mathrm{sinh}\left(\widehat{g}\mathrm{}\right)+\widehat{g}\mathrm{}}\widehat{\chi }.$$
(2.20)
Case 1.b: $`a=ig/\widehat{g}\mathrm{sinh}|\widehat{g}||y|`$
$$\chi =2\frac{|\widehat{g}|^3/(ig)^2}{\mathrm{sinh}\left(\widehat{g}\mathrm{}\right)\widehat{g}\mathrm{}}\widehat{\chi }.$$
(2.21)
Case 1.c: $`a=g/\widehat{g}\mathrm{cos}i\widehat{g}y`$
In this case the integration limits are $`0`$ and $`2\pi /i\widehat{g}`$:
$$\chi =\frac{\left(i\widehat{g}\right)^3}{\pi (ig)^2}\widehat{\chi }.$$
(2.22)
Case 3: $`a=ig|y|+C`$
In this case we have:
$$\chi =\frac{3}{2}\frac{ig}{\left(ig\mathrm{}/2+C\right)^3C^3}\widehat{\chi }.$$
(2.23)
## 3 Graviton Dynamics
Expanding the first of Eqs. (2.1) around a background which satisfies the same equation one finds the equation of motion for the perturbation $`\widehat{h}_{\widehat{\mu }\widehat{\nu }}`$ and using the transverse traceless (tt) gauge
$$\widehat{}^{\widehat{\mu }}\widehat{h}_{\widehat{\mu }\widehat{\nu }}=\widehat{h}=0,$$
(3.1)
where $`\widehat{h}=\widehat{g}^{\widehat{\mu }\widehat{\nu }}\widehat{h}_{\widehat{\mu }\widehat{\nu }}`$ we get<sup>9</sup><sup>9</sup>9All indices are raised and lowered with the full $`\widehat{d}`$-dimensional background metric $`\widehat{g}_{\widehat{\mu }\widehat{\nu }}`$.
$$\begin{array}{ccc}\hfill \widehat{}^2\widehat{h}_{\widehat{\mu }\widehat{\nu }}+2\widehat{R}_{\widehat{\rho }(\widehat{\mu }}\widehat{h}^\rho {}_{\widehat{\nu })}{}^{}+2\widehat{R}^{\widehat{\lambda }}{}_{(\widehat{\mu }\widehat{\nu })}{}^{}{}_{}{}^{\widehat{\sigma }}\widehat{h}_{\widehat{\lambda }\widehat{\sigma }}^{}2\widehat{\mathrm{\Lambda }}\widehat{h}_{\widehat{\mu }\widehat{\nu }}& & \\ & & \\ \hfill +2\widehat{\chi }\{\widehat{g}^{\rho \sigma }[2\widehat{h}_{\rho (\widehat{\mu }}\widehat{g}_{\widehat{\nu })\sigma }\frac{1}{\widehat{d}2}(\widehat{h}_{\rho \sigma }\widehat{g}_{\widehat{\mu }\widehat{\nu }}+\widehat{h}_{\widehat{\mu }\widehat{\nu }}\widehat{g}_{\rho \sigma })]& & \\ & & \\ \hfill \widehat{h}^{\rho \sigma }(\widehat{g}_{\rho \widehat{\mu }}\widehat{g}_{\sigma \widehat{\nu }}\frac{1}{\widehat{d}2}\widehat{g}_{\rho \sigma }\widehat{g}_{\widehat{\mu }\widehat{\nu }})\}_nT_n\delta (yy_n)& =& 0.\hfill \end{array}$$
(3.2)
Using now the first of Eqs. (2.1) to eliminate $`\widehat{R}_{\widehat{\rho }\widehat{\mu }}`$ we get
$$\begin{array}{ccc}\hfill \widehat{}^2\widehat{h}_{\widehat{\mu }\widehat{\nu }}+2\widehat{R}^{\widehat{\lambda }}{}_{(\widehat{\mu }\widehat{\nu })}{}^{}{}_{}{}^{\widehat{\sigma }}\widehat{h}_{\widehat{\lambda }\widehat{\sigma }}^{}+2\widehat{\chi }[(\widehat{g}^{\rho \sigma }\widehat{h}_{\rho (\widehat{\mu }}\widehat{h}^{\rho \sigma }\widehat{g}_{\rho (\widehat{\mu }})\widehat{g}_{\widehat{\nu })\sigma }& & \\ & & \\ \hfill +\frac{1}{\widehat{d}2}(\widehat{h}^{\rho \sigma }\widehat{g}_{\rho \sigma }\widehat{g}^{\rho \sigma }\widehat{h}_{\rho \sigma })\widehat{g}_{\widehat{\mu }\widehat{\nu }}]_nT_n\delta (yy_n)& =& 0.\hfill \end{array}$$
(3.3)
Further, using the RS gauge
$$\widehat{h}_{\mu y}=\widehat{h}_{yy}=0,$$
(3.4)
and the fact that the warped general metric Eq. (1.1) is block-diagonal we see that the source terms vanish identically. The equations for $`\widehat{h}_{\mu y},\widehat{h}_{yy}`$ are satisfied identically and do not become constraints. The equation for the remaining piece of the perturbation is
$$\begin{array}{ccc}\hfill a^2\left[^2\widehat{h}_{\mu \nu }+2R^\rho {}_{(\mu \nu )}{}^{}{}_{}{}^{\sigma }\widehat{h}_{\rho \sigma }^{}\right]\widehat{h}_{\mu \nu }^{\prime \prime }& & \\ & & \\ \hfill (\widehat{d}5)a^1a^{}\widehat{h}_{\mu \nu }^{}+2[(\widehat{d}4)a^2(a^{})^2+a^1a^{\prime \prime }]\widehat{h}_{\mu \nu }& =& 0.\hfill \end{array}$$
(3.5)
Now we assume that the perturbation can be expanded in RS modes
$$\widehat{h}_{\mu \nu }(x,y)=\underset{\alpha }{}f_\alpha (y)h_{\mu \nu }^{(\alpha )}(x),$$
(3.6)
of which we only keep the massless one $`h_{\mu \nu }^{(0)}h_{\mu \nu }`$. The sourceless equation of a massless graviton in a maximally symmetric background, in the tt gauge is
$$^2h_{\mu \nu }+2R^\rho {}_{(\mu \nu )}{}^{}{}_{}{}^{\sigma }h_{\rho \sigma }^{}=0,$$
(3.7)
and, thus, we get for $`\widehat{h}_{\mu \nu }=f_0(y)h_{\mu \nu }`$
$$\widehat{h}_{\mu \nu }^{\prime \prime }=(5\widehat{d})a^1a^{}\widehat{h}_{\mu \nu }^{}+2\left[a^1a^{\prime \prime }+(\widehat{d}4)a^2(a^{})^2\right]\widehat{h}_{\mu \nu }.$$
(3.8)
which in $`\widehat{d}=5`$ is solved by
$$f_0(y)=a^2(y).$$
(3.9)
Depending on the specific solution we can have gravity confinement on the brane or not. In general, the inclusion of branes and the orbifolding procedure is necessary to have confinement on just one brane ($`a^2`$ has more than one maximum in the interval of interest). The only exception seems to be the RS case. The DS to DS case ($`\widehat{g},g`$ imaginary) deserves special mention because the holographic coordinate is naturally compact. No branes are needed to make the graviton wave-function normalizable, although we do need them if we want to think in terms of confinement. Some of the general results for a metric of the form (1.1) have also been obtained in Ref. .
## 4 A Brane action for the Randall-Sundrum Scenario
A constant can be understood as the dual of a volume-form field strength. A volume-form (i.e. a $`\widehat{d}`$-form which we will denote by $`\widehat{F}_{(\widehat{d})}`$) is the field strength of a $`d(\widehat{d}1)`$-form potential $`\widehat{A}_{(d)}`$. The equation of motion forces the dual of $`\widehat{F}_{(\widehat{d})}`$ to be constant (or, more generally, piecewise constant). Thus, one can generically substitute a constant in an action, and, in particular the cosmological constant, by a $`d`$-form potential $`\widehat{A}_{(d)}`$. The canonical example is the rewriting of Romans’ massive $`10`$-dimensional type IIA supergravity, which contains a mass parameter $`m`$, by a $`9`$-form Ramond-Ramond potential to which the D8-brane couples .
This implies a generalization of the theory since now one can have solutions in which the value of the cosmological constant is different in different regions of the spacetime: this is precisely what RS-like solutions need. The discontinuities are $`d`$-dimensional topological defects (domain walls) which act as sources for the $`d`$-form potential and can be interpreted as the worldvolumes of $`(\widehat{d}2)`$-branes charged under the $`d`$-form potential (D8-branes in the case of Ref. ). A worldvolume action for these branes should, therefore, contain a Wess-Zumino term: the integral of the pullback of the $`d`$-form potential.
All this seems to work very well in the D8-brane case and, in fact, the rewriting in terms of a 9-form potential proves necessary and even crucial. It is natural to try something similar here. Therefore, we propose an action consisting in a bulk action containing gravity, $`\widehat{g}_{\widehat{\mu }\widehat{\nu }}`$, coupled to a $`d`$-form potential $`\widehat{A}_{(d)}`$ and a bunch of standard worldvolume actions of $`(\widehat{d}2)`$-branes containing the above-mentioned WZ terms with dynamical coordinate fields $`\widehat{X}_n^{\widehat{\mu }}`$, i.e.
$$\begin{array}{ccc}\hfill \widehat{S}& =& \frac{1}{2\chi }d^{\widehat{d}}\widehat{x}\sqrt{|\widehat{g}|}\left[\widehat{R}+\frac{(1)^{\widehat{d}2}}{2\widehat{d}!}\widehat{F}_{(\widehat{d})}^2\right]\hfill \\ & & \\ & & +_n\{\frac{T_n}{2}d^d\xi _n\sqrt{|\gamma _n|}[\gamma _n^{ij}_i\widehat{X}_n^{\widehat{\mu }}_j\widehat{X}_n^{\widehat{\nu }}\widehat{g}_{\widehat{\mu }\widehat{\nu }}(\widehat{X}_n)(\widehat{d}3)]\hfill \\ & & \\ & & +\frac{(1)^d\mu _n}{d!}d^d\xi _n\widehat{A}_{(d)\widehat{\mu }_1\mathrm{}\widehat{\mu }_d}(\widehat{X}_n)_{i_1}\widehat{X}^{\widehat{\mu }_1}\mathrm{}_{i_d}\widehat{X}^{\widehat{\mu }_d}ϵ^{i_1\mathrm{}i_d}\}.\hfill \end{array}$$
(4.1)
The field configurations that minimize this action satisfy the equations of motion for the metric
$$\begin{array}{ccc}\hfill \widehat{G}^{\widehat{\mu }\widehat{\nu }}+\frac{(1)^{(\widehat{d}2)}\mu _n}{2d!}\left[\widehat{F}_{(\widehat{d})}{}_{}{}^{\widehat{\mu }\widehat{\rho }_1\mathrm{}\widehat{\rho }_d}\widehat{F}_{(\widehat{d})}^{}{}_{}{}^{\widehat{\nu }}{}_{\widehat{\rho }_1\mathrm{}\widehat{\rho }_d}{}^{}\frac{1}{2\widehat{d}}\widehat{g}^{\widehat{\mu }\widehat{\nu }}\widehat{F}_{(\widehat{d})}^2\right]+& & \\ & & \\ \hfill +\frac{\chi }{\sqrt{|\widehat{g}|}}_nT_nd^d\xi _n\sqrt{|\gamma _n|}\gamma _n^{ij}_i\widehat{X}_n^{\widehat{\mu }}_j\widehat{X}_n^{\widehat{\nu }}\delta ^{\widehat{d}}(\widehat{x}\widehat{X}_n)& =& 0,\hfill \end{array}$$
(4.2)
the $`d`$-form potential
$$\widehat{}_{\widehat{\mu }}\widehat{F}_{(\widehat{d})}{}_{}{}^{\widehat{\mu }\widehat{\rho }_1\mathrm{}\widehat{\rho }_d}+\frac{2\chi }{\sqrt{|\widehat{g}|}}\underset{n}{}\mu _nd^d\xi _nϵ^{i_1\mathrm{}i_d}_{i_1}\widehat{X}^{\widehat{\mu }_1}\mathrm{}_{i_d}\widehat{X}^{\widehat{\mu }_d}\delta ^{\widehat{d}}(\widehat{x}\widehat{X}_n)=0,$$
(4.3)
the worldvolume metric (after some manipulations)
$$\gamma _{nij}_i\widehat{X}_n^{\widehat{\mu }}_j\widehat{X}_n^{\widehat{\nu }}\widehat{g}_{\widehat{\mu }\widehat{\nu }}(\widehat{X}_n)=\mathrm{\hspace{0.17em}0},$$
(4.4)
and the (coordinate) worldvolume scalars
$$\begin{array}{ccc}\hfill ^2(\gamma )\widehat{X}_n^{\widehat{\mu }}+\widehat{\mathrm{\Gamma }}_{\widehat{\rho }\widehat{\sigma }}{}_{}{}^{\widehat{\mu }}(\widehat{g})_i\widehat{X}_n^{\widehat{\rho }}_j\widehat{X}_n^{\widehat{\sigma }}\gamma _n^{ij}+& & \\ & & \\ \hfill +\frac{(1)^d\mu _n}{T_nd!\sqrt{|\gamma _n|}}\widehat{F}_{(\widehat{d})}{}_{}{}^{\widehat{\mu }}{}_{\widehat{\rho }_1\mathrm{}\widehat{\rho }_d}{}^{}_{i_1}^{}\widehat{X}^{\widehat{\rho }_1}\mathrm{}_{i_d}\widehat{X}^{\widehat{\rho }_d}ϵ^{i_1\mathrm{}i_d}& =& 0.\hfill \end{array}$$
(4.5)
Eq. (4.4) simply states that the worldvolume metrics are those induced on the worldvolumes by the embedding coordinates $`\widehat{X}_n^{\widehat{\mu }}`$. Using worldvolume reparametrization invariance we can set $`d`$ coordinates (static gauge) to the values $`\widehat{X}_n^\mu =\delta _i^\mu \xi _n^i`$. Furthermore, our Ansatz for the remaining coordinate is
$$\widehat{X}_n^dY_n=y_n,$$
(4.6)
where the $`y_n`$’s are constants. We can perform the volume integrals in Eqs. (4.2,4.3) leaving only 1-dimensional delta functions $`\delta (yy_n)`$. Also, this implies for the worldvolume metrics (identifying worldvolume and $`d`$-dimensional spacetime indices)
$$\gamma _{n\mu \nu }=\widehat{g}_{\mu \nu }=a^2(y_n)g_{\mu \nu }.$$
(4.7)
For the potential we have
$$\widehat{A}_{(d)\mu _1\mathrm{}\mu _d}=ca^d\frac{ϵ_{\mu _1\mathrm{}\mu _d}}{\sqrt{|g|}},\widehat{F}_{(\widehat{d})\underset{¯}{y}\mu _1\mathrm{}\mu _d}=cda^d\mathrm{log}^{}a\frac{ϵ_{\mu _1\mathrm{}\mu _d}}{\sqrt{|g|}},$$
(4.8)
where $`c`$ is a constant to be found and $`ϵ`$ is the $`d`$-dimensional Levi-Cività tensor calculated with the $`d`$-dimensional metric.
Let us solve Eqs. (4.5). The equations for the $`\widehat{X}_n^\mu `$’s are automatically solved. $`\widehat{F}_{(\widehat{d})}`$ only contributes to the equations of the $`Y_n`$’s, which are solved for
$$c=T_n/\mu _nT/\mu .$$
(4.9)
This implies that all the quotients $`T_n/\mu _n`$ must have the same value, which is a characteristic of BPS objects. Observe that the $`\mu _n`$’s cannot vanish: had we tried uncharged brane sources we would have never succeeded.
The equation for the potential becomes
$$T/\mu d\mathrm{log}^{\prime \prime }a+2\chi \underset{n}{}\mu _n\delta (yy_n)=0,$$
(4.10)
which is solved by a warp factor of the RS type
$$a=e^{\frac{\chi \mu /T}{d}_n\mu _n|yy_n|}.$$
(4.11)
To solve the Einstein equations we first calculate the energy-momentum tensor of the form potential, i.e.
$$\widehat{F}_{(\widehat{d})}{}_{}{}^{\widehat{\mu }\widehat{\rho }_1\mathrm{}\widehat{\rho }_d}\widehat{F}_{(\widehat{d})}^{}{}_{}{}^{\widehat{\nu }}{}_{\widehat{\rho }_1\mathrm{}\widehat{\rho }_d}{}^{}\frac{1}{2\widehat{d}}\widehat{g}^{\widehat{\mu }\widehat{\nu }}\widehat{F}_{(\widehat{d})}^2=\left(\frac{\chi }{2}\underset{n}{}\mu _n[\theta (yy_n)1]\right)^2\widehat{g}^{\widehat{\mu }\widehat{\nu }},$$
(4.12)
which describes a piecewise cosmological constant. For only two branes with opposite tensions the Einstein equation is exactly the one in Ref. in the intervals in which the cosmological constant is constant and therefore admits the same solutions. In fact, assuming that the $`d`$-dimensional metric is Ricci-flat $`R_{\mu \nu }=0`$, the Einstein equations are solved by the above warp factor if $`(\mu /T)^2=\frac{1}{2}d/(\widehat{d}2)`$ which implies, as in Sec. (2) and Refs. , that
$$a=e^{\frac{\chi }{(\widehat{d}2)}_nT_n|yy_n|}.$$
(4.13)
## 5 Conclusions
We have explored general solutions with warped metrics with and without branes and we have studied their supersymmetry properties and the effective theories on the branes, including supersymmetry.
Two cases are singled out by supersymmetry considerations: the well-known RS case and the case in which the total spacetime is Minkowski and on the brane one has DS spacetime. The brane breaks a half of the available supersymmetry in the RS case, a result also obtained in Ref. , while in the last case a brane seems to break completely the bulk supersymmetry although one can still speak of (DS) supersymmetry on the brane-world with its known problems.
In analogy with the D8-brane effective action, we considered a formulation of the problem in terms of a dynamical $`\widehat{d}2`$ brane with a $`\widehat{d}1`$ form potential.
## Acknowledgments
We would like to thank Bert Janssen and Pedro Silva for many useful conversations. This work was partially supported by the E.U. TMR program FMRX-CT96-0012 and by the Spanish grant AEN96-1655. |
warning/0003/gr-qc0003075.html | ar5iv | text | # A nongravitational wormhole
## I Introduction
There has been an increasing interest in recent years in earth-bound systems that can mimic gravitational and cosmological phenomena. These systems are of a very different nature: superfluid <sup>3</sup>He-A and high-energy superconductors , supersonic acoustic flows , flowing dielectric fluids , and Bose-Einstein condensates . A common feature of some of these systems is that the propagation of disturbances in a fixed flat background is described by the equation of motion of a massless scalar field in a curved “effective” spacetime. The metric of this curved spacetime depends on the physical system under consideration. In this reformulation in terms of Lorentzian geometry it can be shown that, under certain hypothesis, black-hole analogues are possible in some of the abovementioned systems. There are even some chances to probe semiclassical quantum gravity in the laboratory, through the Hawking effect . We would like to study here yet another system in which an effective metric can be introduced to describe the propagation of waves. It is a well-known fact that Maxwell’s Lagrangian acquires nonlinear correction terms due to vacuum polarization. Novello et al showed recently that in any nonlinear electromagnetic (EM) theory photons, in the geometrical optics approximation, do not propagate along null geodesics of the background geometry. They propagate instead along null geodesics of an effective geometry, which depends on the nonlinearities of the dynamics of the theory. This result was derived by Plebańsky in the case of Born-Infeld electrodynamics in 1968 , and was in extended to encompass any nonlinear electromagnetic theory. The propagation of photons in nontrivial QED vacua has also been studied in . Very often these analysis have led to unexpected results. As examples, let us mention the possibility of faster-than-light photons , and the existence of closed timelike curves due to nonlinear EM effects . Morover, it was shown in that the the motion of light through a nontrivial vacuum is equivalent to that of photons propagating according to Maxwell’s theory in a nonlinear medium. The medium induces modifications on the equations of motion which are identical to those produced by the nonlinearities of the field. We would like to address here some general features of this effective geometry, and also to show an example of a specific effective spacetime for photons: an electromagnetic wormhole.
The structure of the paper is the following. In Section II we review the derivation of the effective geometry for photons in nonlinear electrodynamics, and we point out some of its general features. In Section III we show how a wormhole can be constructed out of the effective geometry in the case of Born-Infeld electrodynamics, and we analize its properties. We close with a discussion of the results.
## II Effective geometry for photons
In this section we present the method of the effective geometry . Let us start with a Lagrangian density that is a function of $`F=F_{\mu \nu }F^{\mu \nu }`$ only <sup>1</sup><sup>1</sup>1 The more general case in which $`=(F,F^{})`$, where $`F^{}F_{\mu \nu }F^{\mu \nu }`$ and $`F^{\mu \nu }`$ is the dual of $`F^{\mu \nu }`$, can be found in .. We will follow Hadamard’s technique to study the propagation of the discontinuities of the EM field. Let $`\mathrm{\Sigma }`$ a surface of discontinuity of $`F_{\mu \nu }`$, with equation $`\mathrm{\Sigma }(x^\mu )=\mathrm{const}`$. We assume that the EM tensor is continuous when crossing $`\mathrm{\Sigma }`$, but its derivative has a finite discontinuity. These assumptions can be stated as
$$[F_{\mu \nu }]_\mathrm{\Sigma }=0,[_\lambda F_{\mu \nu }]_\mathrm{\Sigma }=f_{\mu \nu }k_\lambda .$$
(1)
The operation $`[]_\mathrm{\Sigma }`$ on a tensor $`J`$ is defined by
$$[J]_\mathrm{\Sigma }\mathrm{lim}_{\delta 0^+}\left(J|_{\mathrm{\Sigma }+\delta }J|_{\mathrm{\Sigma }\delta }\right),$$
(2)
and represents the discontinuity of the tensor $`J`$ through $`\mathrm{\Sigma }`$. The tensor $`f_{\mu \nu }`$ is the discontinuity of the field, and $`k_\lambda =_\lambda \mathrm{\Sigma }`$ is the propagation vector.
If we apply $`[]_\mathrm{\Sigma }`$ to the nonlinear equations of motion
$$(\sqrt{\gamma }_FF^{\mu \nu })_{,\nu }=0,$$
(3)
($`_F`$ is the derivative of the Lagrangian density w.r.t. $`F`$, $`\gamma _{\mu \nu }`$ is the metric of Minkowski’s spacetime in an arbitrary coordinate system, and $`\gamma `$ is the corresponding determinant) we obtain
$$_Ff^{\mu \nu }k_\nu +2_{FF}\xi F^{\mu \nu }k_\nu =0,$$
(4)
where $`\xi F^{\alpha \beta }f_{\alpha \beta }`$. The same operation on the cyclic identity yields
$$f_{\mu \nu }k_\lambda +f_{\nu \lambda }k_\mu +f_{\lambda \mu }k_\nu =0.$$
(5)
This equation can be rewritten after contraction with $`k_\alpha \gamma ^{\alpha \lambda }F^{\mu \nu }`$. The result is
$$\xi k_\mu k_\nu \gamma ^{\mu \nu }+2F^{\mu \nu }f_\nu ^\lambda k_\lambda k_\mu =0.$$
(6)
We can eliminate the tensor $`f_{\mu \nu }`$ from Eqs.(4) and (6). If we assume that $`\xi `$ and $`_F`$ are nonzero, we get
$$(_F\gamma ^{\mu \nu }4_{FF}F^{\mu \alpha }F_\alpha {}_{}{}^{\nu })k_\mu k_\nu =0.$$
(7)
This expression suggests that the self-interaction of the field $`F^{\mu \nu }`$ can be interpreted as a modification on the spacetime metric $`\gamma _{\mu \nu }`$ which is described by the effective geometry
$$g_{\mathrm{eff}}^{\mu \nu }=_F\gamma ^{\mu \nu }4_{FF}F_{}^{\mu }{}_{\alpha }{}^{}F^{\alpha \nu }.$$
(8)
Note that only in the particular case of linear electrodynamics the discontinuities of the electromagnetic field propagate along null paths in the Minkowski background.
In the derivation of Eq.(8) we assumed that $`\xi `$ and $`_F`$ do not have zeros. It can be proved (see ) that photons such that $`\xi =0`$ propagate along geodesics of the Minkowskian geometry. To analyse the other restriction, note that $`_F`$ is a function of the coordinates $`x^\mu `$ through the EM tensor $`F_{\mu \nu }(x^\mu )`$. If the theory given by $``$ and the field configuration under study are such that $`_F`$ has one or more zeros for certain values $`x_0^\mu `$ of the coordinates, Eq.(4) will in general not be satisfied in those $`x_0^\mu `$. However, the equation is valid for points that are arbitrarily close to $`x_0^\mu `$, and hence the effective metric tensor is everywhere well-defined except for these zeros.
The general expression of the effective geometry can be equivalently written in terms of the vacuum expectation value of the energy-momentum tensor, given by
$$T_{\mu \nu }\frac{2}{\sqrt{\gamma }}\frac{\delta \mathrm{\Gamma }}{\delta \gamma ^{\mu \nu }},$$
(9)
where $`\mathrm{\Gamma }`$ is the effective action
$$\mathrm{\Gamma }d^4x\sqrt{\gamma }.$$
(10)
In the case of one-parameter Lagrangians, $`=(F),`$ we obtain
$$T_{\mu \nu }=4_FF_{\mu }^{}{}_{}{}^{\alpha }F_{\alpha \nu }\eta _{\mu \nu },$$
(11)
where we have chosen a Cartesian coordinate system in which $`\gamma _{\mu \nu }`$ reduces to $`\eta _{\mu \nu }.`$ In terms of this tensor the effective geometry Eq.(8) can be re-written as<sup>2</sup><sup>2</sup>2 For simplicity, we will denote the effective metric as $`g^{\mu \nu }`$ instead of $`g_{\mathrm{eff}}^{\mu \nu }`$ from now on.
$$g^{\mu \nu }=\left(_F+\frac{_{FF}}{_F}\right)\eta ^{\mu \nu }+\frac{_{FF}}{_F}T^{\mu \nu }.$$
(12)
This equation shows that the stress-energy distribution of the field is the true responsible for the departure of the effective geometry from its Minkowskian form. It is also seen that for $`T_{\mu \nu }=0`$, the conformal modification in Eq.(12) clearly leaves the photon paths unchanged.
At this point several remarks are in order:
* Note that there are two metric tensors present in the problem. One is the Minkowskian metric $`\gamma _{\mu \nu }`$. In the absence of forces, ordinary matter moves on geodesics of this metric. The second is the effective metric $`g_{\mu \nu }`$, and only influences the motion of the photons.
* The properties of the effective geometry depend on the specific theory under consideration (through $``$ and its derivatives), and also on the specific field configuration. The latter is determined by solving the nonlinear Eqns.(3) together with $`_{[\mu }F_{\nu \rho ]}=0`$.
* The effective geometry defined by Eq.(8) admits an everywhere nonzero conformal factor. This freedom is due to the fact that the effective geometric structure is valid only for photons. It follows that only concepts that do not depend on the choice of the conformal factor have a definite meaning in this framework.
* The use of the effective geometry is not mandatory. Some results can be also obtained by other methods, although these require an average over the polarization states .
* Although we shall restrict here to flat backgrounds, the effective geometry can and must be used in the presence of curved backgrounds. In this case, the flat spacetime metric $`\eta _{\mu \nu }`$ is replaced by the curved spacetime metric $`g_{\mu \nu }`$. The latter is consistently determined by Einstein’s equations. We must point out that if the change in the path of the photons described by the effective metric is not taken into account, the resultant description of the spacetime under consideration is incomplete. For instance, it was shown in that an apparently “regular” black hole generated by a certain nonlinear electrodynamics is actually singular, the singularities in the effective geometry being associated to the nonlinearities of the theory.
Let us now turn to a far-reaching analogy between nonlinear EM on one side and linear EM in the presence of dielectrics on the other. In was shown that it is possible to describe the propagation of waves governed by Maxwell electrodynamics inside a medium in terms of a modification of the underlying spacetime geometry using the framework developed above. The electromagnetic field inside a material medium can be represented by two antisymmetric tensors, the EM field $`F_{\mu \nu }`$ and the polarization $`P_{\mu \nu }`$, which in the absence of sources obey the equations
$$^\nu P_{\mu \nu }=0,^\nu F_{\mu \nu }^{}=0.$$
(13)
From these equations and Eq.(3) we see that the nonlinear theory in vaccum is equivalent to Maxwell’s theory in an isotropic medium characterized by an electric susceptibility $`ϵ`$ and a magnetic permeability $`\mu `$ given by
$$ϵ=_F\mu =\frac{1}{_F}$$
(14)
(see for details). Therefore, the simple class of effective Lagrangians $`=(F)`$ may be used as a convenient description of Maxwell theory inside isotropic media; conversely, results obtained in the latter context can as well be similarly restated in the former one <sup>3</sup><sup>3</sup>3 Let us remind the reader that $`ϵ`$ and $`\mu `$ are functions of the EM field, and so their derivatives are discontinuous on the wavefront..
## III A wormhole for photons
In this section we show that there exists a field configuration in Born-Infeld nonlinear EM theory that generates an effective geometry which is a wormhole for photons. We shall restrict here to a spherically symmetric and static effective metric. It can be proved (see Appendix) that the effective metric tensor expressed in spherical coordinates of the Minkowskian background takes the form
$$ds_{\mathrm{eff}}^2=\mathrm{\Phi }^1(dt^2dr^2)\mathrm{\Psi }^1r^2d\mathrm{\Omega }^2,$$
(15)
where
$$\mathrm{\Phi }=_F4_{FF}A^2,$$
(16)
$$\mathrm{\Psi }=_F+4_{FF}B^2r^4\mathrm{sin}^2\theta ,$$
(17)
and the only nonzero components of the EM tensor are $`AF^{tr}`$ and $`BF^{\theta \phi }`$. From the nonlinear Eqs.(3) and the cyclic identity, we get
$$A=\frac{\alpha }{_Fr^2},B=\frac{\beta }{r^4sin\theta },$$
(18)
with $`\alpha `$ and $`\beta `$ arbitrary constants. From these equations, it is seen that the electromagnetic tensor $`F_{\mu \nu }`$ must satisfy
$$F=\frac{2}{r^4}\left(\beta ^2\frac{\alpha ^2}{_F^2}\right).$$
(19)
We now adopt as an example Born-Infeld electrodynamics . The Lagrangian density is given by
$$=b^2\left(1\sqrt{1+\frac{F}{2b^2}}\right).$$
(20)
Note that $`_F`$ and $`_{FF}`$ are nonsingular functions of $`r`$ for this Lagrangian. The constraint Eq.(19) gives in this case
$$F=2b^2\frac{\beta ^216\alpha ^2}{b^2r^4+16\alpha ^2}$$
(21)
and the metric functions are
$$\mathrm{\Phi }=\frac{1}{4b^2r^2}\frac{(r^4b^2+16\alpha ^2)^{3/2}}{\sqrt{b^2r^4+\beta ^2}},$$
(22)
$$\mathrm{\Psi }=\frac{1}{4b^2r^4}\frac{\sqrt{b^2r^4+16\alpha ^2}(b^4r^816\alpha ^2\beta ^2)}{(b^2r^4+\beta ^2)^{3/2}}.$$
(23)
We see that the function $`\mathrm{\Psi }`$ has a zero, and as a result, the metric coefficient $`\mathrm{\Psi }^1r^2`$ will be singular. From now on we restrict to the case in which only the electric field is nonzero, i.e. $`\beta =0`$. In this case, the metric functions are
$$\mathrm{\Phi }=\frac{(b^2r^4+16\alpha ^2)^{3/2}}{4b^3r^6},$$
(24)
$$\mathrm{\Psi }=\frac{\sqrt{b^2r^4+16\alpha ^2}}{4br^2}.$$
(25)
Due to the freedom in the conformal factor of the definition of the effective metric (see remarks above) we shall analyze the features of the metric
$$ds^2=dt^2dr^2\mathrm{\Delta }(r)d\mathrm{\Omega }^2,$$
(26)
where $`\mathrm{\Delta }(r)r^2\mathrm{\Phi }/\mathrm{\Psi }`$ is given by
$$\mathrm{\Delta }=\frac{b^2r^4+16\alpha ^2}{b^2r^2}=r^2+\frac{r_{\mathrm{th}}^4}{r^2},$$
(27)
with $`r_{\mathrm{th}}2\sqrt{|\alpha |/b}`$. Note that the metric is singular in $`r=0`$. However, this divergence is an artifact of the coordinate system that comes from the behaviour of the spherical coordinates at $`r=0`$.
From the function $`\mathrm{\Delta }(r)`$ we can calculate the area of the 2-surface $`t=\mathrm{const}`$, $`r=\mathrm{const}`$, which is given by $`A^{(2)}=4\pi \mathrm{\Delta }(r)`$. The function $`A^{(2)}`$ has a minimum at $`r=r_{\mathrm{th}}`$, and does not have zeros. For large values of $`r`$, $`A^{(2)}4\pi r^2`$ and the metric goes into the flat spacetime metric. At $`r=0`$, $`A^{(2)}`$ diverges due to the singularity of the metric. This analysis shows that the effective geometry for the photons is in fact a wormhole , with a throat located at $`r_{\mathrm{th}}`$, and two different asymptotic regions.
The wormhole structure may also be displayed by making a coordinate change to write the metric in the standard form. By setting
$$\rho ^2r^2+\frac{r_{\mathrm{th}}^4}{r^2};$$
(28)
and inverting this relation, we can write the metric Eq.(26) in the form
$$ds^2=dt^2\frac{1}{1\frac{b_\pm (\rho _\pm )}{\rho _\pm }}d\rho _\pm ^2\rho _\pm ^2d\mathrm{\Omega }^2,$$
(29)
with
$$b_\pm (\rho _\pm )=\frac{2\rho _{\mathrm{th}}^4\rho _\pm ^4\pm \rho _\pm ^2\sqrt{\rho _\pm ^4\rho _{\mathrm{th}}^4}}{\rho _\pm ^3\pm \rho _\pm \sqrt{\rho _\pm ^4\rho _{\mathrm{th}}^4}},$$
(30)
and $`\rho _{\mathrm{th}}\rho _\pm <\mathrm{}`$. The “$`+`$” patch covers the region $`r_{\mathrm{th}}r<\mathrm{}`$, and the “$``$” patch, the region $`0<rr_{\mathrm{th}}`$. The functions $`b_\pm (\rho _\pm )`$ satisfy the requirements needed in order to have a wormhole geometry , with a throat located at $`\rho _{\mathrm{th}}\sqrt{2}r_{\mathrm{th}}`$. Because two different shape functions $`b(\rho )`$ are needed, the wormhole is asymmetric under the interchange $`\rho _+\rho _{}`$. Note also that only the region covered by the “+” patch is asymptotically flat.
The same analysis can be carried out for the magnetic case, i.e. $`\alpha =0`$, $`\beta 0`$. It follows that the effective metric also describes a wormhole, with a throat located at $`r_{\mathrm{th}}=\sqrt{|\beta |/b}`$, and asymptotic regions as in the electric case.
To analyze the motion of the photons in the effective geometry, we can use the effective potential. In the case of a static and spherically symmetric spacetime, $`L`$ and $`E`$ represent constants of motion along a geodesic path:
$$g_{\phi \phi }\dot{\phi }=Lg_{tt}\dot{t}=E,$$
(31)
where the dot indicates derivation w.r.t. an affine parameter, and $`\theta =\pi /2`$. If we consider null curves we obtain
$$g_{tt}\dot{t}^2+g_{rr}\dot{r}^2+g_{\phi \phi }\dot{\phi }^2=0,$$
(32)
which can be rewritten as
$$\dot{r}^2=E^2V(r).$$
(33)
The effective potential $`V`$ is given by
$$V\frac{L^2}{g_{rr}g_{\phi \phi }}+E^2\left(1+\frac{1}{g_{rr}g_{tt}}\right).$$
(34)
In the coordinate $`r`$, the metric is given by Eq.(26), and the effective potential is
$$V(r)=\frac{L^2r^2}{r^4+r_{\mathrm{th}}^4}$$
(35)
The maximum of $`V(r)`$ is at $`V_M=L^2/2r_{\mathrm{th}}`$. The following plot gives the effective potential as a function of the $`r`$ coordinate for different values of the angular momentum, and a fixed $`r_{\mathrm{th}}`$.
From this plot we see that photons coming from $`r=\mathrm{}`$ with energy smaller than $`V_M`$ are reflected by the potential barrier. Hence they turn back without reaching the throat. Photons with $`E>V_M`$ will inevitably pass through the throat and will not turn back.
The motion of the photons can also be described from the plot of the effective potential in the $`\rho _\pm `$ coordinates. Let us recall first what happens in an ultrastatic spherically symmetric gravitational wormhole with the following metric:
$$ds^2=dt^2\frac{d\rho ^2}{1\frac{\rho _{\mathrm{th}}^2}{\rho ^2}}\rho ^2d\mathrm{\Omega }^2.$$
(36)
In this case the effective potential is given by
$$V=\left(1\frac{\rho _{\mathrm{th}}^2}{\rho ^2}\right)\left(\frac{L^2}{\rho ^2}E^2\right)+E^2.$$
(37)
From this expression, it is easily seen that only particles for which $`L^2/E^2<\rho _{\mathrm{th}}^2`$ will reach the throat. This feature is displayed in the next plot.
Because the metric given in Eq.(36) describes a symmetric wormhole, after passing through the throat the photon will feel the “mirror image” of this effective potential.
In our case, the expression for $`V`$ in the $`\rho _\pm `$ coordinates is
$$V(\rho _\pm )=2\left(\frac{L^2}{\rho _\pm ^2}E^2\right)\frac{\rho _\pm ^4\rho _{\mathrm{th}}^4}{\rho _\pm ^2(\rho _\pm ^2\sqrt{\rho _\pm ^4\rho _{\mathrm{th}}^4})}+E^2$$
(38)
The following plots give the “$`+`$” and “$``$” parts of the effective potential as a function of the $`\rho _\pm `$ coordinates, for a fixed value of the energy.
Because of the abovementioned asymmetry, here we need both Figs.(3) and (4). All the curves in these two figures intersect at $`\rho =\rho _{\mathrm{th}}`$, which corresponds to $`V=E^2=25`$. In these two plots (as well in Fig.(2)) the curves depend on the energy of the photon. This is illustrated in the subsequent plot.
Comparing Fig.(2) with Figs.(3) and (4) we see that a photon travelling in the effective metric given by Eq.(29) will be under the influence of the same type of effective potential as a photon travelling through a gravitational wormhole. Let us remark that because the path of the photons is invariant under conformal transformations, the wormhole discussed here does not depend on the election of the frame in which the metric is given.
The photons that go from the “$`+`$” patch to the “$``$” patch traverse the throat and inevitably escape to $`\rho _{}\mathrm{}`$ (see Figs.(3) and (4)). This motion takes place in the interval $`(0,r_{\mathrm{th}})`$ of the $`r`$ coordinate (see Fig.(1)).
## IV Conclusion
We have shown how a geometrical structure considered previously in a purely gravitational context can also arise in a nonlinear electromagnetic theory in a flat background. We would like to point out once more that this effective wormhole affects solely the motion of the photons; all other types of matter move according to the flat spacetime dynamical laws. Although we have worked out in detail an example based in Born-Infeld theory, in principle any nonlinear electromagnetic theory could generate analogues of gravitational structures in a given background. This statement also applies to nonabelian gauge theories with nonlinear dynamics.
Let us remark that the gravitational wormhole requires matter that violates the null energy condition as a source of Einstein’s equations. In the case analyzed here, the wormhole is generated through an electromagnetic field that obeys a nonlinear dynamics, and is linked to the metric via Eq.(8).
Finally, it should be stressed that the analogy between nonlinear EM and Maxwell’s EM in the presence of a nonlinear dielectric medium may open the door to an actual realization of this wormhole (and also of any other of these effective gravitational-like structures) in the laboratory. We will tackle this problem in a forthcoming article.
## ACKNOWLEDGMENTS
FB would like to thank the members of the Gravitation and Cosmology group (LAFEX-CBPF) for their hospitality. SEPB would like to acknowledge financial support from CONICET (Argentina).
## appendix
The expression given in Eq.(8)
$`g^{\mu \nu }=_F\gamma ^{\mu \nu }4_{FF}F_{}^{\mu }{}_{\alpha }{}^{}F^{\alpha \nu },`$ (39)
may be viewed as a system of ten equations with the EM field components as unknowns. If we restrict our attention to diagonal $`\gamma _{\mu \nu }`$ and $`g_{\mu \nu }`$, the equations
$$4_{FF}F_{}^{\mu }{}_{\alpha }{}^{}F^{\alpha \nu }=0,\mu \nu ,$$
(40)
impose several constraints on the degrees of freedom of the problem. It turns out that non-trivial solutions are obtained when only two components of $`F^{\mu \nu }`$ are different form zero: one electric and its dual magnetic component <sup>4</sup><sup>4</sup>4This situation is not improved when one considers $`=(F,F^{})`$, due to the property $`F_{}^{\mu }{}_{\lambda }{}^{}F^{\lambda \nu }=\frac{1}{4}F^{\alpha \beta }F_{\alpha \beta }^{}\eta ^{\mu \nu }`$ (see ).. In particular, we can choose a Minkowskian background in spherical coordinates for $`\gamma _{\mu \nu }`$. Setting $`F^{tr}A`$, $`F^{\theta \phi }B`$, and the rest of the components of $`F^{\mu \nu }`$ equal to zero in order to satisfy Eq.(40), we are led to the equations
$`g^{tt}`$ $`=`$ $`_F4_{FF}A^2,`$ (41)
$`g^{rr}`$ $`=`$ $`_F+4_{FF}A^2,`$ (42)
$`g^{\theta \theta }`$ $`=`$ $`{\displaystyle \frac{_F}{r^2}}4_{FF}r^2\mathrm{sin}^2\theta B^2,`$ (43)
$`g^{\phi \phi }`$ $`=`$ $`{\displaystyle \frac{_F}{r^2\mathrm{sin}^2\theta }}4_{FF}r^2B^2,`$ (44)
which is the effective metric given by Eq.(15). |
warning/0003/hep-ph0003146.html | ar5iv | text | # The structure function of semi-inclusive heavy flavour decays in field theory
## 1 Introduction
Nowadays there are many facilities that allow an accurate experimental study of heavy flavour decays. It is therefore becoming more and more important to improve the accuracy and the reliability of the theoretical calculations. In this paper, we study the properties of the decays of heavy flavour hadrons into inclusive hadronic states $`X`$ “tight” in mass, i.e. with an invariant mass $`m_X`$ small with respect to the energy $`E_X`$:
$$m_XE_X.$$
(1)
More specifically, we consider the situation where
$$m_X^2O\left(E_X\mathrm{\Lambda }_{QCD}\right),$$
(2)
so that
$$\frac{m_X^2}{E_X^2}O\left(\frac{\mathrm{\Lambda }_{QCD}}{E_X}\right)1\left(E_X\mathrm{\Lambda }_{QCD}\right).$$
(3)
A formal definition of kinematics (2) is the limit, in the heavy quark rest frame:
$`E_X`$ $``$ $`\mathrm{},`$
$`m_X^2`$ $``$ $`\mathrm{}`$
with
$$\frac{m_X^2}{E_X}=\mathrm{const}.$$
(4)
The divergence of $`m_X^2`$ \- even though slower than the one of $`E_X^2`$ \- implies that the final hadronic state can be replaced with a partonic one, i.e. that the use of perturbation theory is fully justified. Heavy flavour decays are characterized by three mass or energy scales: the mass of the heavy flavour $`m_Q`$, the energy $`E_X`$, and the invariant mass $`m_X`$ of the final hadronic state. Limit (4) implies also the limit of infinite mass for the heavy flavour $`Q`$:
$$m_Q\mathrm{},$$
(5)
since $`m_QE_X`$. Another consequence of (4) is that
$$\frac{m_X^2}{E_X^2}\mathrm{\hspace{0.17em}0},$$
(6)
i.e. we are in the so-called threshold region<sup>1</sup><sup>1</sup>1The converse is not true: limit (6) does not imply limit (4) (we thank G. Veneziano for pointing this out to us)..
The study of these processes has both a theoretical and a phenomenological interest. On the theoretical side, in heavy decays the infrared perturbative structure of gauge theories - the Sudakov form factor \- enters in a rather “pure” form, owing to the absence of initial state mass singularities. On the phenomenological side, the computation of many relevant distributions requires a good theoretical control over the region (1). As examples, let us quote the electron spectrum $`d\mathrm{\Gamma }/dE_e`$ close to the endpoint $`E_em_B/2`$ and the hadron mass distribution $`d\mathrm{\Gamma }/dm_X`$ at small $`m_X`$ in semileptonic $`bu`$ decays, such as
$$BX_u+e+\nu ,$$
(7)
or the photon energy distribution $`d\mathrm{\Gamma }/dE_\gamma `$ close to the endpoint $`E_\gamma m_B/2`$ in $`bs\gamma `$ decays. For the electron or photon spectrum, the region (1) is involved because the requirement of a large energy of the lepton or of the photon pushes down to zero the mass of the recoiling hadronic system. As is well known, the above mentioned distributions in (7) allow an inclusive determination of the CKM-matrix element $`|V_{ub}|`$ , while a large photon energy in the rare decay, $`E_\gamma 2.1`$ GeV is required to cut experimental backgrounds.
In general, the dynamics in region (1) is rather intricate as it involves an interplay of non-perturbative and perturbative contributions. These are related to the Fermi motion of the heavy quark inside the hadron and to the Sudakov suppression in the threshold region (6), respectively. Even though these two effects are physically distinguishable and are treated as independent in various models , they are ultimately both described by the same quantum field theory, QCD. Therefore the problem arises of describing them consistently, i.e. without double countings, inconsistencies, etc. Our idea is to subtract from the hadronic tensor encoding $`all`$ QCD dynamics,
$$W_{\mu \nu }\underset{X}{}H_Q|J_\nu ^+|XX|J_\mu |H_Q\delta ^4(p_Bqp_X),$$
(8)
each of the perturbative components - associated with the Sudakov form factor and with other short-distance corrections - to end up with an explicit representation of the non-perturbative component. In eq. (8), we have defined
$$J_\mu (x)\overline{q}(x)\mathrm{\Gamma }_\mu Q(x),$$
(9)
where $`\mathrm{\Gamma }_\mu `$ is a matrix in Dirac algebra<sup>2</sup><sup>2</sup>2For the left-handed currents of the Standard Model, $`\mathrm{\Gamma }_\mu =\gamma _\mu \left(1\gamma _5\right)`$., $`p_X`$ is the momentum of the final hadronic jet, and $`H_Q`$ is a hadron containing the heavy quark $`Q`$. The non-perturbative component is identified with an ultraviolet (UV) regularized expression for the structure function, or shape function, in the effective theory. The shape function $`f(k_+)`$ has been introduced using the Operator Product Expansion $`(OPE)`$ and can be defined as
$$f(k_+)H_Qh_v^{}\delta (k_+iD_+)h_vH_Q,$$
(10)
where $`h_v`$ is a field in the Heavy Quark Effective Theory (HQET) with 4-velocity $`v`$; $`D_+`$ is the plus component of the covariant derivative, i.e. $`D_+D^0+D^3`$. The shape function represents the probability that the heavy quark has a momentum $`m_Bv+k^{}`$ with a given plus component $`k_+^{}=k_+`$. This function can also be interpreted (see section 4.4) as the probability that $`Q`$ has an effective mass
$$m_{}=m_B+k_+$$
(11)
at disintegration time. The renormalization properties of the shape function have also been analysed . Because of UV divergences affecting its matrix elements, $`f(k_+)`$ needs to be renormalized and it consequently acquires a dependence on the renormalization point $`\mu `$: $`f(k_+;\mu )`$. The non-perturbative information about Fermi motion enters in this framework as the initial value for the $`\mu `$-evolution. The shape function can be extracted from a reference process and used to predict other processes, analogously to the parton distribution functions in usual hard processes such as Deep Inelastic Scattering (DIS) or Drell–Yan . In principle, it can also be computed with a non-perturbative technique, for example lattice QCD .
Our approach aims at a deeper understanding of perturbative and non-perturbative effects with respect to the standard OPE in dimensional regularization (DR). We compare different regularization schemes and find that the factorization procedure is substantially scheme-dependent. By that, we mean a much stronger scheme dependence than the usual one, typically related to the finite part of one-loop amplitudes, corresponding in DR to a replacement of the form $`1/ϵ1/ϵ+`$const. The shape function, in contrast to naive expectations, is not a physical distribution, but it is affected by regularization scheme effects even at the leading, double-logarithmic level. We show, however, that it factorizes most of the non perturbative effects in a class of regularization schemes.
This paper is devoted to a wide audience, i.e. not only to perturbative QCD experts, but also to phenomenologists who are interested in the field theoretic aspects of this area of $`B`$ physics, as well as to lattice-QCD physicists who may wonder about the possibility of simulating the shape function. We have therefore tried to give a plain presentation of our method, together with a self-consistent description of the known results to be found in the literature.
In section 2 we give a simple introduction to the physics of semi-inclusive heavy flavour decays. In section 3 we present our strategy, based on factorization, in order to consistently combine perturbative and non-perturbative contributions and to arrive at a formal definition of the shape function in field theory; we outline the main steps and the relevant issues. In section 4 we review the standard $`OPE`$ derivation of the shape function in the effective theory; this section can be skipped by the experienced reader. In section 5 we return to the strategy outlined in section 4 and apply our factorization procedure in the quantum theory to a specific class of loop corrections. Our technique is completely general, but we believe that it is better illustrated by treating in detail a simple computation, which illustrates most of the general features. In section 6 we discuss factorization in the framework of the effective theory on the light-cone, the so called Large Energy Effective Theory (LEET), which is the relevant effective theory for these processes at low energy. In section 7 we describe the properties of the shape function in the effective theory in various regularizations and its evolution with the UV cutoff or renormalization scale. We also discuss our results on factorization and clarify a controversial factor of 2 in the evolution kernel of the shape function. Section 8 contains the conclusions.
## 2 Physics of semi-inclusive heavy flavour decays
Let us begin by discussing Fermi motion. This phenomenon, originally discovered in nuclear physics, is classically described as a small oscillatory motion of the heavy quark inside the hadron, due to the interaction with the valence quark; in the quantum theory it is also the virtuality of the heavy flavour that matters. Generally, as the mass of the heavy flavour becomes large, i.e. as we take the limit (5), we expect that the heavy particle decouples from the light degrees of freedom and becomes “frozen” with respect to strong interactions. That is indeed true in the “bulk” of the phase space of the decay products, but it is untrue close to kinematical boundaries, as in region (1). This is because a kinematical amplification effect occurs, according to which a small virtuality of the heavy flavour in the initial state produces relatively large variations of the fragmentation mass in the final state. To see how this works in detail, let us begin with a picture of the initial bound state. We assume that the momentum exchanges $`r_\mu `$ between the heavy flavour and the light degrees of freedom are of the order of the hadronic scale,
$$|r_\mu |O\left(\mathrm{\Lambda }_{QCD}\right),$$
(12)
as we take the infinite mass limit (5). In other words, we assume that the momentum transfer does not scale with the heavy mass but remains essentially constant<sup>3</sup><sup>3</sup>3We must specify that we consider an initial hadron containing a single heavy quark: hadrons containing more than one heavy quark, such as for example quarkonium states, need a different theoretical treatment .. This assumption, which is rather reasonable from a physical viewpoint, is at the basis of the application of the HQET . Let us discuss for example the decay (7). The initial meson has momentum
$$p_B=m_Bv,$$
(13)
where $`v`$ is the 4-velocity, which we can take at rest without any loss of generality: $`v^\mu =(1;0,0,0)`$. The final hadronic state $`X`$ has a momentum
$$Q=m_Bvq$$
(14)
and invariant mass
$$m_X^2Q^2.$$
(15)
In eq. (14) $`q_\mu `$ is the momentum of the virtual $`W`$ or, equivalently, of the leptonic pair. We isolate in the decay a hard subprocess consisting in the fragmentation of the heavy quark. If the valence quark - in general the light degrees of freedom in the hadron - have momentum $`k^{}`$, the heavy quark has a momentum<sup>4</sup><sup>4</sup>4For the appearance of $`m_B`$ instead of $`m_b`$, see footnote in section 4.1.
$$p_Q=m_Bv+k^{}$$
(16)
and a virtuality
$$p_Q^2m_B^2=2m_Bvk^{}+k^20\left(\mathrm{in}\mathrm{general}\right).$$
(17)
The final invariant mass of the hard subprocess, i.e. the fragmentation mass, is
$$\widehat{m}_X^2\left(p_Qq\right)^2=\left(Q+k^{}\right)^2=m_X^2+2Qk^{}+k^2m_X^2+2Qk^{}$$
(18)
and this is the mass that controls the kinematic of the hard subprocess, i.e. the Sudakov form factor (the difference between $`m_X^2`$ and $`\widehat{m}_X^2`$ is that we do not include in the latter the momentum of the valence quark). The term $`k^2`$ has been neglected in the last member of eq. (18) because it is small, as gluon exchanges are soft according to the assumption (12). We take the motion of the final $`up`$ quark in the $`z`$ direction, so that the vector $`Q`$ has large zero and third components, both of order $`E_X,`$ and a small square; we have therefore for the average in the meson state:
$$Qk^{}=Qk^{}O\left(E_X\mathrm{\Lambda }_{QCD}\right).$$
(19)
A fluctuation in the heavy quark momentum of order $`\mathrm{\Lambda }_{QCD}`$ in the initial state produces a variation of the final invariant mass of the hard subprocess of order
$$\delta \widehat{m}_X^2O\left(\mathrm{\Lambda }_{QCD}E_X\right).$$
(20)
An amplification by a factor $`E_X`$ has occurred, as anticipated. The fluctuation (20) is of the order of (2) and so it must be taken into account.
We will discuss the shape function at length in sections 4 and 7, but let us introduce now some of its more important properties. If we consider a heavy quark with the given off-shell momentum (16), we find for the shape function<sup>5</sup><sup>5</sup>5The final state consists of a massless on-shell quark at the tree level.
$$f(k_+)^{part}=\delta \left(k_+k_+^{}\right)+O\left(\alpha _S\right),$$
(21)
where
$$k_+\frac{m_X^2}{2E_X}.$$
(22)
Selecting the hadronic final state, i.e. $`k_+`$, we select the light-cone virtuality $`k_+^{}=k_+`$ of the heavy flavour which participates in the decay. After inclusion of the radiative corrections, we find that in general $`k_+^{}k_+`$ <sup>6</sup><sup>6</sup>6To obtain the hadronic shape function, the “elementary” or “partonic” shape function in eq. (21) has to be convoluted with the distribution $`\phi _0\left(k_+^{}\right)`$ of the primordial light-cone virtuality $`k_+^{}`$ of the heavy quark inside the hadron.. Equation (21) is analogous to the relation between the Bjorken variable $`x_Bq^2/\left(2pq\right)`$ ($`p`$ is the momentum of the hadron and $`q`$ that of the space-like photon) and the momentum fraction $`x`$ of partons in the naive parton model, where we have
$$q\left(x\right)^{part}=\delta \left(xx_B\right)+O\left(\alpha _S\right).$$
(23)
In this case, as is well known, by selecting final state kinematics, i.e. $`x_B`$, one selects the momentum fraction $`x=x_B`$ of the partons that participate in the hard scattering. Just as in the heavy flavour decay, radiative corrections lead to a softening of the above condition in $`xx_B`$, due to the emission of collinear partons.
We note that even with the amplification effect (20), Fermi motion effects are irrelevant in most of the phase space, where typical values for the final hadron mass are
$$\widehat{m}_X^2O(E_X^2).$$
(24)
This is in agreement with physical intuition.
As will be proved in section 4.4, the shape function can be interpreted as the distribution of a variable mass. The virtuality of the heavy flavour can be represented by a shift of its mass, $`m_bm_{}`$. In other words, an off-shell particle with a given mass, i.e. with the momentum (16), can be replaced by an on-shell particle with a variable, virtuality dependent, mass, i.e. with a momentum $`m_{}\left(k_+\right)v`$. The physical distribution is obtained by convoluting the distribution of an isolated quark of mass $`m_{}`$ with the probability distribution for such a mass (see eq. (74)): this is the basis of the factorization theory for the semi-inclusive heavy flavour decays.
Fermi motion is a non-perturbative effect in QCD because it involves low momentum transfers to the heavy flavour (cf. eq. (12)), at which the coupling is large; it does however occur also in QED bound states, where it can be treated with perturbation theory <sup>7</sup><sup>7</sup>7Consider for instance an atom composed of a $`\mu `$ and an $`e`$, decaying by $`\mu `$ fragmentation..
The second phenomenon relevant in region (1) is related to soft gluon emission and it is of a perturbative nature - it is a case of the Sudakov form factor in QCD . The $`up`$ quark emitted by the fragmentation of the heavy flavour with a large virtuality - of the order of the final hadronic energy $`E_X`$ \- evolves in the final state, emitting soft and collinear partons, either real or virtual. Since the final state is selected to have a small invariant mass (cf. eq. (6)), real radiation is inhibited with respect to the virtual one. That means that infrared (IR) singularities coming from real and virtual diagrams still cancel, but leave a large residual effect in the form of large logarithms<sup>8</sup><sup>8</sup>8The plus-distribution is defined as
$$\left(\frac{\mathrm{log}x}{x}\right)_+=\theta (x)\frac{\mathrm{log}x}{x}\delta \left(x\right)_0^1\frac{\mathrm{log}y}{y}𝑑y.$$
(25) :
$$\alpha _S\left(\frac{\mathrm{log}m_X^2/E_X^2}{m_X^2/E_X^2}\right)_+.$$
(26)
Schematically, the rate for final states with an invariant mass $`m_X^2`$ has double-logarithmic contributions at order $`\alpha _S`$, of the form:
$$\mathrm{real}=\alpha _S_0^{E_X}_0^1\frac{dϵ}{ϵ}\frac{d\theta ^2}{\theta ^2}\delta \left(ϵ\theta ^2\frac{m_X^2}{E_X}\right)$$
(27)
and
$$\mathrm{virtual}=\alpha _S\delta \left(\frac{m_X^2}{E_X}\right)_0^{E_X}_0^1\frac{dϵ}{ϵ}\frac{d\theta ^2}{\theta ^2},$$
(28)
where $`ϵ`$ is the gluon energy, $`\theta `$ is the angle between the $`up`$ and the gluon, and $`\mathrm{\Theta }=\pi \theta `$ is the polar angle of the gluon 3-momentum. The perturbative corrections of the form (26) blow up at the Born kinematics $`m_X=0`$, which is the threshold of the inelastic channels. For this reason, the above corrections are often called threshold logarithms and need a resummation to any order in $`\alpha _S`$.
## 3 Overview of Factorization
The aim of this paper is a detailed study of factorization in semi-inclusive heavy flavour decays and of the properties of the shape function in field theory. In order to trace all the perturbative and non-perturbative contributions to the process, it is convenient to perform the factorization in two steps. In the first step the heavy flavour is replaced by a static quark. That is accomplished by taking the infinite mass limit (5), keeping $`E_X`$ and $`m_X`$ fixed. With this, the hadronic tensor loses a kinematical scale, namely the heavy flavour mass $`m_Q`$:
$$W_{\mu \nu }(m_Q,E_X,m_X)\stackrel{~}{W}_{\mu \nu }(E_X,m_X),$$
(29)
where the effective hadronic tensor is defined as
$$\stackrel{~}{W}_{\mu \nu }\underset{X}{}H_Q|\stackrel{~}{J}_\nu ^{}|XX|\stackrel{~}{J}_\mu |H_Q\delta ^4(p_Bqp_X),$$
(30)
and it contains the static-to-light currents
$$\stackrel{~}{J}_\mu (x)=\overline{q}(x)\mathrm{\Gamma }_\mu \stackrel{~}{Q}(x).$$
(31)
The difference between the two tensors in eq. (29) is incorporated into a first coefficient function or hard factor. While in full QCD the vector and axial currents are conserved, or partially conserved, so the renormalization constants are UV-finite and anomalous dimensions vanish, this property does not hold anymore in the HQET: the effective current with a static quark is not conserved and it acquires an anomalous dimension $`\stackrel{~}{\gamma _J}`$ <sup>9</sup><sup>9</sup>9 In eq. (32) and (33), we are representing the evolution schematically, without details; f.i., we do not distinguish between the anomalous dimensions of the vector and axial currents.:
$$\left(\frac{d}{d\mathrm{log}\mu }+\stackrel{~}{\gamma _J}\right)\stackrel{~}{J}_\mu =0.$$
(32)
As a consequence also the hadronic tensor acquires an anomalous dimension, which equals twice that of the vector or axial current:
$$\left(\frac{d}{d\mathrm{log}\mu }+2\stackrel{~}{\gamma _J}\right)\stackrel{~}{W}_{\mu \nu }=0.$$
(33)
All this is very easily understood by observing that the original QCD tensor $`W_{\mu \nu }`$ is UV-finite at one loop but it does contain $`\alpha _S\mathrm{log}m_Q`$ terms, and so it is divergent in the infinite-mass limit (5). If this limit is taken ab initio, i.e. before regularization, these terms manifest themselves as new ultraviolet divergences, an heritage of the $`\mathrm{log}m_Q`$ terms of the original tensor. We may say that the dependence on the heavy mass is promoted to UV divergence; in practice
$$\alpha _S\mathrm{log}\frac{m_Q}{E_X}\alpha _S\mathrm{log}\frac{\mathrm{\Lambda }_1}{E_X},$$
(34)
where $`\mathrm{\Lambda }_1`$ is an UV cutoff if we deal with the bare theory, or a renormalization point if we deal with the renormalized theory; in principle $`\mathrm{\Lambda }_1m_Q`$. At the end of the game, the effective hadronic tensor still depends on three scales, just like the original one,
$$\stackrel{~}{W}_{\mu \nu }=\stackrel{~}{W}_{\mu \nu }(E_X,m_X;\mathrm{\Lambda }_1).$$
(35)
The original tensor $`W_{\mu \nu }`$ is parametrized in terms of five independent form factors . For the HQET hadronic tensor (30) there are instead relations between the form factors originating from the spin-symmetry of the HQET. In particular, the structure in $`\mathrm{log}m_Q/E_X`$ of the original QCD tensor can be understood by looking at the UV divergences of $`\stackrel{~}{W}_{\mu \nu }`$<sup>10</sup><sup>10</sup>10In Dimensional Regularization (DR), this means simple poles $`1/ϵ`$..
After the first step $`\stackrel{~}{W}_{\mu \nu }`$ still contains perturbative contributions. The latter are factorized with a second step, which corresponds to the limit (4). Additional UV divergences are introduced also with this second step, which must be regulated with a new cutoff $`\mathrm{\Lambda }_2`$. In principle $`\mathrm{\Lambda }_2E_X`$, since $`E_X\mathrm{}`$. As before with the heavy mass logarithms, soft and collinear logarithms are promoted to ultraviolet logarithms:
$$\alpha _S\left(\frac{\mathrm{log}m_X^2/E_X^2}{m_X^2/E_X^2}\right)_+\alpha _S\left(\frac{\mathrm{log}\left(2k_+/\mathrm{\Lambda }_2\right)}{2k_+/\mathrm{\Lambda }_2}\right)_+.$$
(36)
The second factorization step involves double-logarithmic effects of an infrared nature, in contrast with the single logarithms of the large mass of the first step. In practice, we separate perturbative contributions from non-perturbative ones starting with a cutoff
$$\mathrm{\Lambda }_2E_X$$
(37)
and lowering it to a much smaller value<sup>11</sup><sup>11</sup>11In order to avoid substantial finite cutoff effects, the condition $`\mathrm{\Lambda }_2^{}\mathrm{\Lambda }_{QCD}`$ must hold.
$$\mathrm{\Lambda }_2^{}E_X.$$
(38)
The contributions of the fluctuations with energy between $`\mathrm{\Lambda }_2`$ and $`\mathrm{\Lambda }_2^{}`$ are put into a second coefficient function, while the contributions below $`\mathrm{\Lambda }_2^{}`$ are factorized inside the shape function. The latter is defined in the framework of a low-energy effective theory, with a cutoff given by
$$\mathrm{\Lambda }_{ET}=\mathrm{\Lambda }_2^{}.$$
(39)
Most of the non-perturbative effects in lattice-like regularizations are contained in the shape function, which uniquely determines the final, non-perturbative, hadronic tensor
$$\stackrel{~}{\stackrel{~}{W}}_{\mu \nu }\underset{X}{}H_Q|\stackrel{~}{\stackrel{~}{J}}_\nu ^{}|XX|\stackrel{~}{\stackrel{~}{J}}_\mu |H_Q\delta ^4(p_Bqp_X),$$
(40)
containing the effective-heavy-to-effective-light currents
$$\stackrel{~}{\stackrel{~}{J}}_\mu (x)=\overline{\stackrel{~}{q}}(x)\mathrm{\Gamma }_\mu \stackrel{~}{Q}(x).$$
(41)
It is worth noting that the tensor (40) involves a single form factor, proportional to the shape function itself (see eq. (65)). That is again a consequence of the spin-symmetry of both HQET and LEET , which is more efficient than that one of the HQET alone. The shape function is completely non-perturbative and perturbative factors can no longer be extracted.
The effect of lowering the UV cutoff (eqs. (37) and (38)) is incorporated inside a coefficient function, which, unlike more simple cases such as the light-cone expansion in DIS, is not completely short-distance dominated. Some long-distance effects are left in the coefficient function, but they are expected to be suppressed on physical grounds. Finally, the introduction of ultraviolet divergences with factorization, implies scheme-dependence issues for the shape function, which are rather dramatic because of the double-logarithmic nature of the problem (cf. eqs.(27) and (28)).
In fig. 1, we give a pictorial description of the above procedure.
## 4 OPE
The amplitude for the decay (7), which we take as our example from now on, can be written at the lowest order in the weak coupling as
$$A=\frac{G_F}{\sqrt{2}}l\nu L_\mu 0X|J^\mu |B,$$
(42)
where $`L_\mu `$ is the leptonic current and $`J_\mu `$ is the hadronic one:
$$J_\mu (x)\overline{q}(x)\mathrm{\Gamma }_\mu b(x)$$
(43)
with $`\mathrm{\Gamma }_\mu =\gamma _\mu (1\gamma _5)`$, $`q(x)`$ a light quark field and $`b(x)`$ the beauty quark field. Taking the square of (42) and summing over the final states, we arrive at the hadronic tensor defined in eq. (8). By the optical theorem, we can relate the hadronic tensor $`W_{\mu \nu }`$ to the imaginary part<sup>12</sup><sup>12</sup>12Since $`\mathrm{\Gamma }_\mu `$ is in general complex, we should say, more properly, the absorptive part. of the Green function or forward hadronic tensor $`T_{\mu \nu }`$:
$$W_{\mu \nu }=\frac{1}{\pi }\mathrm{Im}T_{\mu \nu },$$
(44)
where
$$T_{\mu \nu }id^4xe^{iqx}B|T\left(J_\mu ^{}(x)J_\nu (0)\right)|B.$$
(45)
### 4.1 The HQET
We are interested in the evaluation of $`T_{\mu \nu }`$ in the effective theory and we discuss in this section the first factorization step: replacing the beauty quark by a quark in the HQET. As is well known, we can decompose the heavy quark field $`b(x)`$ into two effective quark and antiquark fields $`h_v`$ and $`H_v`$ <sup>13</sup><sup>13</sup>13We prefer to refer to the physical $`B`$-meson mass rather than to the unphysical $`b`$-quark mass. Their difference is of order $`\mathrm{\Lambda }_{QCD}`$, so it is a $`1/m_B`$ correction and can be neglected in our leading-order analysis. Furthermore, in perturbation theory there is no binding energy so that $`m_b=m_B`$.
$$b(x)=e^{im_Bvx}\left[h_v(x)+H_v(x)\right],$$
(46)
satisfying
$$P_+h_v=h_v,P_{}h_v=0,P_{}H_v=H_v,P_+H_v=0,$$
(47)
where $`P_\pm =(1\pm \widehat{v})/2`$ are the projectors over the components with positive and negative energies, respectively. The field $`H_v`$ is neglected (which amounts to neglecting heavy-pair creation), so that
$$b(x)e^{im_Bvx}h_v(x).$$
(48)
By using eq. (48) we obtain
$$\stackrel{~}{T}_{\mu \nu }=id^4xe^{iQx}B(v)|T\overline{h}_v(x)\mathrm{\Gamma }_\mu ^{}q(x)\overline{q}(0)\mathrm{\Gamma }_\nu h_v(0)|B(v).$$
(49)
We now use the Wick theorem and we single out the only contraction that is relevant to $`B`$ decay:
$$\stackrel{~}{T}_{\mu \nu }=d^4xe^{iQx}B|\overline{h}_v(x)\mathrm{\Gamma }_\mu ^{}S(x|0)\mathrm{\Gamma }_\nu h_v(0)|B,$$
(50)
where $`S(x|0)`$ is the light quark propagator. Note that the operator entering the right-hand side of eq. (50) is already normal ordered, since $`h_v`$ has only the component that annihilates heavy quarks, while $`\overline{h}_v`$ only the components that create them. We can express the Fourier transform of the light quark propagator as<sup>14</sup><sup>14</sup>14The notation is very compact. For more explicit representations of the propagator see ref. .
$$S(Q+iD)=\frac{1}{i\widehat{D}+\widehat{Q}+i0}=\frac{i\widehat{D}+\widehat{Q}}{Q^2+2iDQD^2g/2\sigma _{\mu \nu }G^{\mu \nu }+i0}$$
(51)
where $`\sigma _{\mu \nu }i/2[\gamma _\mu ,\gamma _\nu ]`$ is a generator of the Lorentz group, $`G_{\mu \nu }i/g[D_\mu ,D_\nu ]`$ is the field strength and $`D_\mu _\mu igA_\mu `$ is the covariant derivative. In eq. (51), $`0`$ denotes, as usual, an infinitesimal positive number and gives the prescription to deal with pole or branch-cut singularities. There are three different regions according to the value of the jet invariant mass, which are described by three different full or effective theories:
$`i)m_X^2`$ $``$ $`O(E_X^2),`$
$`ii)m_X^2`$ $``$ $`O(\mathrm{\Lambda }_{QCD}^2),`$
$`iii)m_X^2`$ $``$ $`O(\mathrm{\Lambda }_{QCD}E_X).`$ (52)
Since the derivative of the rescaled $`h_v`$ field brings down the residual momentum $`k^{}`$, and it is therefore an operator with matrix elements of order $`O\left(\mathrm{\Lambda }_{QCD}\right)`$, the matrix elements of the operators entering the light quark propagator have a size of the order of
$`i\widehat{D}`$ $``$ $`O(\mathrm{\Lambda }_{QCD}),`$
$`\mathrm{\hspace{0.17em}2}iDQ`$ $``$ $`O(\mathrm{\Lambda }_{QCD}E_X),`$
$`D^2`$ $``$ $`O(\mathrm{\Lambda }_{QCD}^2),`$
$`\sigma _{\mu \nu }G^{\mu \nu }`$ $``$ $`O(\mathrm{\Lambda }_{QCD}^2).`$ (53)
Let us discuss these regions in turn in the next section.
### 4.2 General kinematical regions
This region corresponds to a jet $`X`$ with a large invariant mass, of the order of the energy:
$$m_XO\left(E_X\right).$$
(54)
To a first approximation all the covariant derivative terms can be neglected, so that
$$S(Q+iD)\frac{\widehat{Q}}{Q^2+i0},$$
(55)
i.e. the light quark can be described as a free quark. A higher accuracy is reached when expanding the propagator in powers of the covariant derivative operators up to the required order. We have here an application of the $`1/m_B`$ expansion up to a prescribed (finite) order <sup>15</sup><sup>15</sup>15It is clear that a consistent inclusion of the $`1/m_B`$ corrections involves also the expansion of the heavy quark field $`b(x)`$ into the effective quark field $`h_v(x)`$ up to the required order.. In this region there are no large adimensional ratios of scales, the latter being all of the same order. This implies that in perturbation theory we do not hit large logarithmic corrections to be resummed to all orders in $`\alpha _S.`$ This region is not relevant to the endpoint electron spectrum because the hadronic jets takes away most of the available energy. This region will not be discussed further here.
This region involves a recoiling hadronic system with a mass of the order of the hadronic one: it can be a single hadron or very few hadrons. The dynamics is dominated by the emission, with consequent decay, of few resonances; it is a completely non-perturbative problem. According to the estimates (53), no term can be neglected in the light quark propagator. We are faced with full QCD dynamics as far as the final hadronic state is concerned. This region must be evaluated by an explicit sum over all the kinematically possible hadronic states, and the latter have to be computed with a non-perturbative technique such as a quark model or lattice QCD. This region will not be discussed here either.
This region is intermediate between $`i)`$ and $`ii)`$ and as such it has both perturbative and non-perturbative components. Roughly speaking, we have to take into account non-perturbative effects for the initial state hadron, while we can neglect final state binding effects. This region is characterized by a small ratio of the jet invariant mass to the jet energy, and thus involves the large adimensional ratio in (3). As always is the case, perturbation theory generates logarithms of the above adimensional ratio, eq. (26). The term $`2iDQ`$ at the denominator cannot be brought at the numerator (with a truncated operatorial expansion) because it is of the same order as $`m_X^2.`$ At lowest order, the other covariant derivative terms can be neglected, to give:
$$S(Q+iD)\frac{\widehat{Q}}{m_X^2+2iDQ+i0}.$$
(56)
One can reach a higher level of accurary keeping these latter corrections up to a given order<sup>16</sup><sup>16</sup>16We envisage a relation between the $`1/E_X`$ corrections to the shape function and the power-suppressed perturbative corrections of the form $`\alpha _S/E_X\mathrm{log}^2(E_X/m_X)`$.. The rest of the paper deals with region $`iii)`$ at the lowest order in $`1/m_B`$.
### 4.3 The LEET
In this section we discuss the second factorization step, which involves the description of the final $`up`$ quark in the LEET, according to eq. (56). Let us define the adimensional vector $`n_\mu `$ as:
$$n_\mu =\frac{Q_\mu }{Qv}.$$
(57)
This $`n_\mu `$ has a normalized time component, $`n_0=1`$. In the “semi-inclusive” endpoint region $`iii)`$:
$`n^2`$ $`=`$ $`{\displaystyle \frac{m_X^2}{E_X^2}}`$ (58)
$`=`$ $`O\left({\displaystyle \frac{\mathrm{\Lambda }_{QCD}}{E_X}}\right)1.`$
We will show later that $`n`$ can be replaced by a vector lying exactly on the light-cone, i.e.
$$n\overline{n},$$
(59)
where $`\overline{n}^\mu =(1;0,0,1)`$ ( $`\overline{n}^{\mathrm{\hspace{0.17em}2}}=0`$), representing the direction of the hadronic jet, the $`z`$ axis. We can write
$$S(Q+iD)=\frac{1}{2vQ}\frac{\widehat{Q}}{iD_+k_++i0},$$
(60)
where $`k_+`$ has been defined in eq. (22) and $`D_+\overline{n}D.`$ We can simplify the tensor structure of $`T_{\mu \nu }`$ by using the identity
$$\overline{h}_v\mathrm{\Gamma }_\mu h_v=\frac{1}{2}\mathrm{Tr}(\mathrm{\Gamma }_\mu P_+)\overline{h}_vh_v\frac{1}{2}\mathrm{Tr}(\gamma _\mu \gamma _5P_+\mathrm{\Gamma }_\mu P_+)\overline{h}_v\gamma ^\mu \gamma _5h_v,$$
(61)
which is valid for any $`\mathrm{\Gamma }_\mu `$. The matrix element of the axial vector current between the $`B`$-meson states vanishes by parity invariance, so that <sup>17</sup><sup>17</sup>17A physical argument for the spin factorization is that, in the limit $`m_B\mathrm{}`$, the spin interaction of the $`b`$-quark in the $`B`$ meson vanishes; therefore we can average over the helicity states of the $`b`$ quark .:
$$\stackrel{~}{\stackrel{~}{T}}_{\mu \nu }=s_{\mu \nu }\frac{1}{2vQ}F(k_+),$$
(62)
where we have defined the “light-cone” function
$$F(k_+)B(v)h_v^{}\frac{1}{iD_+k_++i0}h_vB(v),$$
(63)
and
$$s_{\mu \nu }\frac{1}{2}\mathrm{Tr}\left[\mathrm{\Gamma }_\mu ^{}\widehat{Q}\mathrm{\Gamma }_\nu P_+\right]$$
(64)
is the “spin factor”, containg the leading spin effects. The factor $`1/(2vQ)`$ is a jacobian, which appears as we go from the full QCD variable $`Q^2`$ to the effective theory variable $`k_+.`$
Taking the imaginary part of $`T_{\mu \nu }`$, we obtain (see relation (44)):
$$\stackrel{~}{\stackrel{~}{W}}_{\mu \nu }=s_{\mu \nu }\frac{1}{2vQ}f(k_+),$$
(65)
where
$$f(k_+)\frac{1}{\pi }\mathrm{Im}F(k_+)$$
(66)
is the shape function. By using the formula
$$\frac{1}{iD_+k_++i0}=\mathrm{P}\frac{1}{iD_+k_+}i\pi \delta (iD_+k_+),$$
(67)
we recover the definition of the shape function given by eq. (10). Note that it involves the non-local operator $`h_v^{}\delta (k_+iD_+)h_v`$, which results from the resummation of the towers of operators of the form $`\left(QD\right)^n`$.
### 4.4 The Variable Mass
The hadronic tensor can be written in the effective theory in terms of the shape function as:
$$\stackrel{~}{\stackrel{~}{W}}_{\mu \nu }=s_{\mu \nu }_{\mathrm{}}^0𝑑k_+\delta \left(Q^2+2k_+vQ\right)f(k_+);$$
(68)
in the second member, $`k_+`$ is an integration, i.e. dummy, variable. In the free theory, with an on-shell $`b`$-quark (i.e. $`k^{}=0`$ in eq. (16)),
$$f^0(k_+)=\delta (k_+0),$$
(69)
so that
$$W_{\mu \nu }^0=s_{\mu \nu }\delta (Q^20).$$
(70)
The hadronic tensor can be written, up to terms of order $`k_+^2O\left(\mathrm{\Lambda }_{QCD}^2\right),`$ as
$$\stackrel{~}{\stackrel{~}{W}}_{\mu \nu }=s_{\mu \nu }(Q)𝑑m_{}\delta \left(Q_{}^20\right)f(m_{}m_B),$$
(71)
where we have defined
$$Q_{}m_{}vq$$
(72)
and $`m_{}`$ is the “variable” or “fragmentation” mass, defined in eq.(11). Since $`m_{}`$ is just a shift of $`k_+`$, the range is
$$\mathrm{}<m_{}m_B.$$
(73)
Inside $`s_{\mu \nu }`$ we can replace $`Q`$ with $`Q_{},`$ because that amounts only to a correction of order $`k_+=O(\mathrm{\Lambda }_{QCD})`$, so that<sup>18</sup><sup>18</sup>18We replace by 0 the lower limit of integration, because the relevant region is $`m_{}m_BO(\mathrm{\Lambda }_{QCD})`$.
$$\stackrel{~}{\stackrel{~}{W}}_{\mu \nu }(v,Q)=_0^{m_B}𝑑m_{}\phi (m_{})W_{\mu \nu }^0(v,Q_{}),$$
(74)
where
$$W_{\mu \nu }^0(v,Q^{})=s_{\mu \nu }(Q_{})\delta \left(Q_{}^20\right)$$
(75)
is the hadronic tensor in the free theory for a heavy quark of mass $`m_{}`$ and
$$\phi (m_{})f(m_{}m_B)$$
(76)
is the distribution for the effective mass $`m_{}`$ of the $`b`$-quark inside the $`B`$-meson at disintegration time. Equation (74) is the fundamental result of factorization in semi-inclusive heavy flavour decays: it says that the hadronic tensor in the effective theory can be expressed as the convolution of the hadronic tensor in the free theory with a variable mass times a distribution probability for this mass. That offers also the physical interpretation to the shape function anticipated in the introduction: it represents the probability that the $`b`$ quark has an effective mass $`m_{}`$ at the decay time. Since this tensor encodes all the hadron dynamics, $`any`$ distribution can be expressed in a similar factorized form.
## 5 Factorization in the quantum theory
In this section we discuss factorization in the quantum theory, i.e. the separation of short-distance and long-distance contributions, including loop effects.
A shape function $`f(k_+)^{QCD}`$ and a light-cone function $`F(k_+)^{QCD}`$ can also be defined in full QCD by means of the relations :
$$T_{\mu \nu }^{QCD}\left(s_{\mu \nu }+\mathrm{\Delta }s_{\mu \nu }\right)\frac{1}{2vQ}F(k_+)^{QCD}$$
and
$$W_{\mu \nu }^{QCD}\left(s_{\mu \nu }+\mathrm{\Delta }s_{\mu \nu }^{}\right)\frac{1}{2vQ}f(k_+)^{QCD},$$
where $`\mathrm{\Delta }s_{\mu \nu }`$ and $`\mathrm{\Delta }s_{\mu \nu }^{}`$ are defined as the part of the spin structure not proportional to $`s_{\mu \nu }`$. The tensors $`\mathrm{\Delta }s_{\mu \nu }`$ and $`\mathrm{\Delta }s_{\mu \nu }^{}`$ represent residual spin effects not described by the effective theory (ET), which do not contribute to the Double-Logarithmic Approximation (DLA)<sup>19</sup><sup>19</sup>19Note that $`\mathrm{\Delta }s_{\mu \nu }`$ and $`\mathrm{\Delta }s_{\mu \nu }^{}`$ are, in general, different; this was not noted in .. In DLA the forward tensor can therefore be written as
$$T_{\mu \nu }^{QCD}=s_{\mu \nu }\frac{1}{2vQ}F(k_+)^{QCD},$$
(77)
where the “light-cone function” is given by
$$F(k_+)^{QCD}\frac{1}{k_++i0}\left[1+aC\right],$$
(78)
$`a\alpha _S`$
$`C_F/\pi `$ and $`C`$ is the scalar triangle diagram (see fig. 2):
$$CivQ\frac{d^4l}{\pi ^2}\frac{1}{(l+Q)^2+i0}\frac{1}{vl+l^2/2m+i0}\frac{1}{l^2+i0}.$$
(79)
We have set the light quark mass equal to zero .
The hadronic tensor relevant to the decay is obtained by taking the imaginary part according to eq. (44). This transforms the products in convolutions, which are converted again into ordinary products by the well-known Mellin transform .
Infrared singularities (soft & collinear) are regulated by the virtuality $`Q^20`$ of the external $`up`$ quark<sup>20</sup><sup>20</sup>20This is consistent because a virtual massless quark is not degenerate with a quark and a soft and/or collinear gluon.. We may write
$$Q^\mu E_X(1+\frac{n^2}{4};0,0,1+\frac{n^2}{4})=E_X\left(v_{}+\frac{n^2}{4}v_+\right)$$
(80)
and
$$n^\mu v_{}^\mu +\frac{n^2}{4}v_+^\mu ,\overline{n}^\mu =v_{}^\mu $$
(81)
where we defined the light-cone versors
$$v_+(1;0,0,1),v_{}(1;0,0,1).$$
(82)
Let us now consider the properties of the integral $`C`$. First, it is adimensional. Second, it is UV-finite for power counting: the integrand has three ordinary scalar propagators with a total of six powers of momentum at the denominator. This implies that $`C`$ does not depend on an ultraviolet cutoff $`\mathrm{\Lambda }_{UV}`$ as long as it is larger than any physical scale of the process, namely
$$\mathrm{\Lambda }_{UV}m_B.$$
(83)
Third, as already discussed, $`C`$ is also IR-finite as long as $`Q^20.`$ For $`Q^2>0`$ there is an imaginary part, related to the propagation of the real $`up`$ and gluon pair, while for $`Q^2<0`$ the integral is real. Therefore $`C`$ does depend on adimensional ratios of three different scales: $`m_B,E_X`$ and $`m_X.`$ There are only two independent ratios, which we choose as $`m_B/E_X`$ and$`m_X/E_X.`$ We are going to decompose the integral $`C`$ in a sum of various integrals; at the end, one of them will correspond to the double-logarithmic contribution to the shape function $`f(k_+)`$ in the low-energy ET. The other integrals represent additional contributions and they are mostly short-distance dominated in lattice-like regularizations. This decomposition consists of two separate steps, which will be described in the following sections.
### 5.1 From QCD to HQET
In the first step we isolate a hard factor by simply subtracting and adding back the integral with the full beauty quark propagator replaced by a static one (see fig. 3)
$$C=C_s+C_h,$$
(84)
where
$$C_sivQ\frac{d^4l}{\pi ^2}\frac{1}{(l+Q)^2+i0}\frac{1}{vl+i0}\frac{1}{l^2+i0},$$
(85)
and
$$C_hi\frac{vQ}{2m}\frac{d^4l}{\pi ^2}\frac{1}{(l+Q)^2+i0}\frac{1}{vl+i0}\frac{1}{vl+l^2/2m+i0}.$$
(86)
The above decomposition parallels that one performed in section 4 in operatorial language. The light-cone function factorizes at order $`\alpha _S`$ according to
$$F(k_+)^{QCD}=\frac{1}{k_++i0}\left[1+aC\right]=\frac{1}{k_++i0}\left[1+aC_h\right]\left[1+aC_s\right].$$
(87)
We expect that $`C_s`$ has the same infrared behaviour as $`C`$; it will be subjected to a further decomposition in the next sections; $`C_h`$ is the “hard factor”, i.e. the difference between QCD and the static limit for the $`b`$ quark. The latter integral is both UV- and IR-finite. The ultraviolet finiteness stems from power counting: the integrand has two scalar propagators and a static propagator with a total of five powers of momentum in the denominator <sup>21</sup><sup>21</sup>21It is known that ultraviolet power counting may fail in effective theory integrals when there are LEET propagators because of the occurrence of an ultraviolet collinear region . The integral $`C_h`$, however, contains only an HQET propagator.. In the infrared region, all the components of the loop momentum are small
$$\mathrm{IR}:l_\mu \rho l_\mu ,\rho 0,$$
(88)
so we can neglect the terms that are quadratic in $`l^\mu `$ in the propagator denominators:
$$C_{h,IR}d^4l\frac{1}{2lQ+Q^2+i0}\frac{1}{(vl+i0)^2}.$$
(89)
Integrating over $`l_0`$, and closing the contour in the upper half of the $`l_0`$-plane, we see that there are no enclosed poles and the integral vanishes (QED). Inside $`C_h`$ we can therefore make the replacement <sup>22</sup><sup>22</sup>22With this substitution, terms related to higher twist contributions of the forms $`\left(m_X/m_B\right)^{n_1}`$ and $`\left(m_X/E_X\right)^{n_2}`$ are neglected, but this is in agreement with our leading-twist ideology (the indices $`n_1`$ and $`n_2`$ are integers).
$$Q^\mu \overline{Q}^\mu E_Xv_{}(\overline{Q}^{\mathrm{\hspace{0.17em}2}}=0).$$
(90)
It follows that $`C_h`$ depends only on $`m_B`$ and $`E_X`$: $`C_h=C_h(m_B,E_X).`$ Since it is adimensional, it may depend only on the adimensional hadronic energy
$$z\frac{2E_X}{m_B},$$
(91)
i.e. $`C_h=C_h(z).`$ An explicit computation gives
$$C_h=\mathrm{log}\left(1z\right)\mathrm{log}z+Li_2(z)z\mathrm{log}z\left(z1\right),$$
(92)
$`C_h`$ does not contain large logarithmic contributions in the limit $`z0`$, i.e. $`\mathrm{log}z`$ terms <sup>23</sup><sup>23</sup>23The dilogarithm is defined as
$$Li_2(z)_0^z\frac{\mathrm{log}(1x)}{x}𝑑x=\underset{n=1}{\overset{\mathrm{}}{}}\frac{z^n}{n^2}\left(|z|1\right).$$
(93) . This is related to the fact that $`C_h`$ and $`C_s`$ are UV-convergent. In general, single logarithms of the hadronic energy, $`\mathrm{log}z,`$ do appear, representing the difference between the interaction of a full propagating heavy quark, of mass $`m_B,`$ and that one of a static quark. The relevant interaction energies are between the beauty mass $`m_B`$ and the hadronic energy $`E_X`$,
$$\alpha _S_{z^2m_B^2}^{m_B^2}\frac{dk^2}{k^2}=2\alpha _S\mathrm{log}z.$$
(94)
The logarithms (94) are resummed, as usual, by replacing the bare coupling with the running coupling and exponentiating, so that the above formula is corrected into:
$`1+\gamma _0\alpha _S{\displaystyle _{z^2m_B^2}^{m_B^2}}{\displaystyle \frac{dk^2}{k^2}}`$ $``$ $`\mathrm{exp}\left[\gamma _0{\displaystyle _{z^2m_B^2}^{m_B^2}}{\displaystyle \frac{dk^2}{k^2}}\alpha _S(k^2)\right]`$ (95)
$`=`$ $`\mathrm{exp}\left[2\alpha _S\gamma _0\mathrm{log}z+2\gamma _0\beta _0\alpha _S^2\mathrm{log}^2z+\mathrm{}\right],`$
where $`\beta _01/(4\pi )(11/3N_c2/3n_f)`$ and $`\alpha _S\alpha _S(m_B)`$.
Let us summarize the above discussion. A first coefficient function is introduced, which takes into account the fluctuations with energy $`\epsilon `$ in the range
$$E_X<\epsilon <m_B.$$
(96)
In the language of Wilson’s renormalization group, we are lowering the UV cutoff of the effective hamiltonian from $`m_B`$ to $`E_X`$.
### 5.2 Infrared factorization
The second factorization step involves the separation of the various infrared contributions to the process, one of which will ultimately lead to the shape function. This step forces us to introduce explicitly an ultraviolet cutoff $`\mathrm{\Lambda }`$ from which the various factors depend separately. In other words, the decomposition of $`C_s`$ introduces fictitious UV divergences, which cancel in the sum. As anticipated in the overview section, infrared factorization in our scheme involves two different operations:
* the separation of the various pole contributions to the QCD amplitude according to the Cauchy theorem;
* the lowering of the UV cutoff from $`\mathrm{\Lambda }_{UV}E_X`$ to $`\mathrm{\Lambda }_{UV}=\mathrm{\Lambda }_{ET}E_X,`$where $`\mathrm{\Lambda }_{ET}`$ is the cutoff of the final low-energy effective theory, in which the shape function is defined (the latter contains the majority of the non-perturbative effects).
Step 1) will be discussed in this section, while step 2) is treated mostly in the next section.
$`C_s`$ is UV-convergent, as is clear again from power counting, and it does not depend on the beauty mass $`m_B`$, which has been sent to infinity, so that $`C_s=C_s(E_X,m_X).`$ The only adimensional variable that can be constructed out of $`E_X`$ and $`m_X`$ is their ratio or, equivalently, $`n^2`$ (see eq. (58)). Since $`C_s`$ is adimensional it may depend only on $`n^2`$: $`C_s=C_s(n^2).`$ The explicit computation in DLA gives
$$C_s=\frac{1}{2}\mathrm{log}^2(n^2i0)\left(DLA\right).$$
(97)
The infrared factorization is performed by integrating $`C_s`$ over the energy $`l_0`$ using the Cauchy theorem. There are three poles in the lower half of the $`l_0`$-plane related to the propagation of a real static quark, a real gluon and a real $`up`$ quark, located respectively at
$$l_0=i0,l_0=+|\stackrel{}{l}|i0,l_0=Q_0+\sqrt{(Q_3+l_3)^2+l_{}^2}i0.$$
(98)
In the upper half-plane, instead, there are only two poles, related to the gluon and the up propagator:
$$l_0=|\stackrel{}{l}|+i0,l_0=Q_0\sqrt{(Q_3+l_3)^2+l_{}^2}+i0$$
(99)
The poles in (99) are conventionally related to a propagation that goes backward in time; they will therefore be called the antiparticle poles. We close for simplicity the integration contour in the upper half-plane and we have two residue contributions related to the antigluon pole ad the anti-$`up`$ pole respectively (see fig.4):
$$C_s=C_g+C_q.$$
(100)
The antigluon and the anti-$`up`$ contributions are given respectively by
$`C_g`$ $`=`$ $`{\displaystyle \frac{2}{\pi }}vQ{\displaystyle d^3l}{\displaystyle \frac{1}{l_0|\stackrel{}{l}|+i0}}{\displaystyle \frac{1}{vl+i0}}{\displaystyle \frac{1}{Q^2+2lQ+i0}}|_{l_0=|\stackrel{}{l}|+i0}`$ (101)
$`C_q`$ $`=`$ $`{\displaystyle \frac{2}{\pi }}vQ{\displaystyle d^3l\frac{1}{l_0+Q_0|\stackrel{}{l}+\stackrel{}{Q}|+i0}\frac{1}{vl+i0}\frac{1}{l^2+i0}}|_{l_0=Q_0\sqrt{(Q_3+l_3)^2+l_{}^2}+i0}`$
As we see by power counting, $`C_g`$ and $`C_q`$ are UV-divergent and it is therefore necessary to introduce an ultraviolet regularization to treat them separately.
The light-cone function factorizes after this second step as
$$F(k_+)^{QCD}=\frac{1}{k_++i0}\left[1+aC_h\right]\left[1+aC_q\right]\left[1+aC_g\right],$$
(102)
i.e. as a product of three factors.
#### 5.2.1 Wilson line representation
Before explicitly computing these 3-dimensional integrals, let us represent them as 4-dimensional ones, i.e. as one-loop integrals of a properly chosen field theory:
$$C_givQ\frac{d^4l}{\pi ^2}\frac{1}{Q^2+2lQ+i0}\frac{1}{vl+i0}\frac{1}{l^2+i0}$$
(103)
and
$$C_qivQ\frac{d^4l}{\pi ^2}\frac{1}{(l+Q)^2+i0}\frac{1}{vl+i0}\frac{1}{Q^2+2lQ+i0}$$
(104)
The proof of the above equations is just by integration over $`l_0`$: closing the integration contour in the upper half-plane of $`l_0`$, we enclose a single pole, whose residue gives the 3-dimensional integrals in eqs. (101); $`C_g`$ and $`C_q`$ involve one full - i.e. quadratic - propagator and two eikonal - i.e. linear - propagators. It is easy to check that the algebraic sum of $`C_g`$ and $`C_q`$ in the above expressions reproduces the integral $`C_s`$ defined in eq. (85). Introducing the variable $`k_+`$, the integral $`C_g`$ can be written in the “familiar” form
$$C_g(k_+)\frac{i}{2}\frac{d^4l}{\pi ^2}\frac{1}{k_++ln+i0}\frac{1}{vl+i0}\frac{1}{l^2+i0}.$$
(105)
The geometrical interpretation is the following: $`C_g(k_+)`$ is the one-loop correction to a vertex composed of an on-shell Wilson line along the time axis, and a Wilson line along the direction $`n`$ off-shell by $`k_+`$. Note that
$$ln=l_++\frac{n^2}{4}l_{},$$
(106)
so that eq. (105) represents the vertex correction in Feynman gauge to the function
$$F(k_+)_{n^20}B(v)h_v^{}\frac{1}{iD_++in^2/4D_{}k_++i0}h_vB(v).$$
(107)
The imaginary part of the above function equals $`1/\pi `$ times the shape function off the light-cone, $`n^20`$:
$$f(k_+)_{n^20}B(v)h_v^{}\delta \left(k_+iD_+in^2/4D_{}\right)h_vB(v).$$
(108)
In the limit $`n^20`$ (see section 6) we recover the correction to the light-cone function $`F(k_+)`$, which was computed in ref. .
We can give a similar description for $`C_q`$. The 4-dimensional representation for $`C_q`$ involves an on-shell Wilson line along the time direction, a Wilson line along the direction $`n`$ off-shell by $`k_+`$, and a light quark propagator with momentum $`l+Q`$
$$C_q\frac{i}{2}\frac{d^4l}{\pi ^2}\frac{1}{k_++ln+i0}\frac{1}{vl+i0}\frac{1}{(l+Q)^2+i0}.$$
(109)
Note that the expressions for $`C_g`$ and $`C_q`$ are very similar: they differ by an overall sign and by the replacement in the quadratic propagator of $`ll+Q`$. For future reference, let us note that the latter shift involves only the zero and the third components of $`l`$, not the transverse ones.
In fig.5 the decomposition of $`C_s`$ into $`C_g`$ and $`C_q`$ is represented.
#### 5.2.2 Space Momenta Cutoff
We consider the bare theory with the regularization introduced in ref. : a sharp cutoff on the spatial loop momenta $`\mathrm{\Lambda }_S`$ (see next section). Integrating $`C_g`$ over $`l_0`$ by closing the integration contour upward and integrating over the azimuthal angle, we obtain
$$C_g=_0^{\mathrm{\Lambda }_S}𝑑l_1^1d\mathrm{cos}\theta \frac{1}{k_++l\left(1\mathrm{cos}\theta \right)+n^2l/4\left(1+\mathrm{cos}\theta \right)},$$
(110)
where we have used the definition of $`Q^\mu `$ in eq. (80) and $`l|\stackrel{}{l}|`$. Integrating over the polar angle, we obtain
$$C_g=_0^{2\mathrm{\Lambda }_S}\frac{dl}{l}\mathrm{log}\left[\frac{k_++l}{k_++n^2l/4}\right].$$
(111)
Assuming a cutoff much larger than any physical scale in the process, i.e.<sup>24</sup><sup>24</sup>24This is done consistently with the relation (83), in which we have taken a large cutoff for the computation of $`C_s`$.
$$\mathrm{\Lambda }_SE_X,$$
(112)
we obtain in DLA <sup>25</sup><sup>25</sup>25The last member in eq. (113) is an artificial absorptive part that cancels against an opposite one in $`C_q`$ (see eq. (119).
$$C_g=\frac{1}{2}\mathrm{log}^2\frac{E_X}{k_+i0}\mathrm{log}\frac{\mathrm{\Lambda }_S}{E_X}\mathrm{log}\frac{E_X}{k_++i0}.$$
(113)
Three scales enter in $`C_g`$: $`k_+,\mathrm{\Lambda }_S`$ and $`E_X.`$ The appearance of $`k_+`$ and the cutoff $`\mathrm{\Lambda }_S`$ was expected, because these two quantities represent the infrared and the ultraviolet scale, respectively. The noticeable fact is that also the hadronic energy $`E_X`$ makes its appearance. $`C_g`$ contains a double-logarithm of the infrared kinematical scale $`k_+`$ (related to the overlap of the soft and the collinear region, which extends up to $`\mathrm{\Lambda }_S`$); it also contains a single logarithm of the cutoff. The appearance of the hadronic energy $`E_X`$ comes from the necessity of a third mass scale for the function $`C_g`$, which behaves like $`\mathrm{log}^2k_+`$ in $`k_+`$ and like $`\mathrm{log}\mathrm{\Lambda }_S`$ in $`\mathrm{\Lambda }_S.`$
When $`lE_X`$ the integrand behaves as
$$\frac{1}{l}\mathrm{log}\left[\frac{k_++l}{k_++n^2l/4}\right]\frac{1}{l}\mathrm{log}\frac{n^2}{4}$$
(114)
and produces a single-logarithmic ultraviolet divergence. As eq. (114) clearly indicates, $`n^20`$ regulates the collinear or light-cone singularity: up to now we have indeed taken kinematics into account exactly together with a large cutoff. It is interesting to note that if we take a cutoff much smaller than the hadronic energy (as we will do in the “final” low-energy effective theory),
$$\mathrm{\Lambda }_{ET}E_X,$$
(115)
we have
$$n^2l\frac{m_X^2}{E_X^2}\mathrm{\Lambda }_{ET}k_+,$$
(116)
and $`C_g`$ simplifies in
$$C_g\left(\mathrm{\Lambda }_{ET}\right)_0^{2\mathrm{\Lambda }_{ET}}\frac{dl}{l}\mathrm{log}\left(\frac{k_++l}{k_+}\right).$$
(117)
The quantity $`n^2`$ does not enter anymore and the integrand is the same as that with the approximate light-cone kinematics $`n^2=0`$, i.e. with $`n`$ replaced by $`\overline{n}`$. The physical explanation is that soft gluons are not able to distinguish between the two slightly different directions $`n`$ and $`\overline{n}`$. Formally, with the small cutoff (115) we can take the limit (6) inside the integral. In other words, in the low-energy effective theory, we effectively are always in the light-cone limit. For $`lk_+`$, the integrand in eq. (117) has the asymptotic behaviour
$$\frac{1}{l}\mathrm{log}\frac{l}{k_+},$$
(118)
implying a double-logarithmic behaviour with respect to $`\mathrm{\Lambda }_{ET}`$ upon integration over $`l`$, in contrast with the single logarithmic behaviour of the integrand in (114). These properties will be studied systematically in the next section, in which we consider the effective theory on the light-cone.
For the computation of $`C_q`$, it is convenient to first make the shift $`llQ`$ in the expression of $`C_q,`$ eq. (104), and then to compute the residue of the light quark pole at $`l_0=|\stackrel{}{l}|+i0`$: this is legitimate if condition (112) holds. We find
$$C_q=\mathrm{log}\frac{\mathrm{\Lambda }_S}{E_X}\mathrm{log}\frac{E_X}{k_++i0}(\mathrm{\Lambda }_SE_X).$$
(119)
The three scales appearing in $`C_g`$ do appear also in $`C_q`$. We note that $`C_q`$ contains a single logarithm of $`k_+`$, i.e. it is subleading by one logarithm in the infrared counting with respect to $`C_g`$. It has a single-logarithmic UV divergence.
If we compute the integral in eq. (104) with a small cutoff (115), we do not find any infrared logarithm, in contrast with what happens instead with $`C_g`$. Thus the logarithmic contributions to $`C_q`$ come from high-energy gluons and that is an indication that $`C_q`$, unlike $`C_g`$, is short-distance dominated.
One can check that the correct value of $`C_s`$ is reproduced by summing $`C_g`$ and $`C_q`$; in particular, UV divergences cancel.
At the level of logarithmic accuracy, we can replace the strong inequality (112) with a weaker one:
$$\mathrm{\Lambda }_SE_X.$$
(120)
Setting in particular
$$\mathrm{\Lambda }_S=E_X,$$
(121)
the expressions for the gluon and quark-pole residue read
$`C_g\left(\mathrm{\Lambda }_S=E_X\right)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{log}^2{\displaystyle \frac{E_X}{k_+iϵ}}=C_s,`$
$`C_q\left(\mathrm{\Lambda }_S=E_X\right)`$ $`=`$ $`0,`$ (122)
i.e. the gluon-pole term gives the whole contribution while the quark-pole factor vanishes. The term $`C_q`$ therefore has the role of correcting $`C_g`$ when $`\mathrm{\Lambda }_SE_X.`$
Since ultraviolet singularities are single-logarithmic for a large cutoff (eqs. (113) and (119)), other regularizations such as DR give similar results. In other words, the regularization scheme dependence is the usual one: the logarithmic term in the one-loop amplitude is scheme independent while the finite part is scheme dependent.
## 6 The effective theory on the light-cone, the LEET
In full QCD infrared singularities are regulated by the unique quantity $`m_X^20`$. In the effective theory, the light quark propagator is replaced by an eikonal one
$$\frac{1}{\left(l+Q\right)^2+i0}=\frac{1}{Q^2+2lQ+l^2+i0}\frac{1}{Q^2+2lQ+i0}.$$
(123)
In the expression on the right-hand side $`l^2`$ has been neglected and as a consequence $`Q_\mu `$ enters in two distinct and independent ways: its square $`Q^2`$ represents the virtuality of the eikonal line at $`l_\mu =0`$, while its components $`Q^\mu `$ are the coefficients of the linear combination of the loop-momentum components $`l_\mu `$ in the term $`2l_\mu Q^\mu `$. Unlike full QCD, $`Q^2`$ and $`Q^\mu `$ can be considered as independent quantities in the effective theory. We may ask ourselves what happens if we take the limit $`Q^20`$ inside the term $`2lQ2l\overline{Q}`$ while keeping $`Q^2\mathrm{const}0`$, i.e. if we make the replacement
$$\frac{1}{Q^2+2lQ+i0}\frac{1}{Q^2+2l\overline{Q}+i0}.$$
(124)
In the usual notation, the above replacement reads
$$\frac{1}{k_++nl+i0}\frac{1}{k_++\overline{n}l+i0},$$
(125)
corresponding to the limit
$$n^20,k_+\mathrm{const}0.$$
(126)
This means that we are considering an eikonal propagator that lies exactly on the light-cone, with collinear singularities regulated now by $`k_+0`$ only, instead of by $`n^20`$ . As we saw before, the limit $`n^20`$ is “invisible” with a small cutoff, simply because the integrand does not depend on $`n^2`$ in this case. We now want to see what happens in the limit (126) with a large cutoff. To perform IR factorization in the light-cone case, it is convenient to start from the original QCD amplitude $`C_s`$ in which we make the replacement
$$\frac{1}{Q^2+2lQ+l^2+i0}\frac{1}{Q^2+2l\overline{Q}+l^2+i0},$$
(127)
to obtain
$$C_s\overline{C}_sivQ\frac{d^4l}{\pi ^2}\frac{1}{Q^2+2l\overline{Q}+l^2+i0}\frac{1}{vl+i0}\frac{1}{l^2+i0}.$$
(128)
It is straightforward to check that $`C_s`$ and $`\overline{C}_s`$ coincide at the DLA level, i.e. that the approximation (127) preserves the infrared structure,
$$\overline{C}_s=C_s=\frac{1}{2}\mathrm{log}^2(n^2i0)(DLA).$$
(129)
The gluon and the quark pole contributions are given in the light-cone limit by
$`C_{g,n^2=0}`$ $``$ $`ivQ{\displaystyle \frac{d^4l}{\pi ^2}\frac{1}{Q^2+2l\overline{Q}+i0}\frac{1}{vl+i0}\frac{1}{l^2+i0}}`$ (130)
$`C_{q,n^2=0}`$ $``$ $`ivQ{\displaystyle \frac{d^4l}{\pi ^2}\frac{1}{Q^2+2l\overline{Q}+l^2+i0}\frac{1}{vl+i0}\frac{1}{Q^2+2l\overline{Q}+i0}}.`$
The above terms are usually called “soft” and “jet” (or “collinear”) factor respectively, even though we believe that this terminology can be rather misleading, as the redistribution of double logarithmic contributions in $`C_g`$ and $`C_q`$ is substantially dependent on the regularization. We will show later that it is possible, within a specific class of regularization schemes, to confine all the double logarithmic contributions in $`C_g.`$ It is immediate to check that the two above integrands sum up to the integrand of $`\overline{C}_s`$. Making the shift $`ll\overline{Q}`$, the quark factor can also be written as
$$C_{q,n^2=0}ivQ\frac{d^4l}{\pi ^2}\frac{1}{l^2+Q^2+i0}\frac{1}{vlv\overline{Q}+i0}\frac{1}{Q^2+2l\overline{Q}+i0}.$$
(131)
### 6.1 <br>Regularization Effects
The decomposition of $`\overline{C}_s`$ into $`C_{g,n^2=0}`$ and $`C_{q,n^2=0}`$ is strongly dependent on the regularization scheme adopted, as a consequence of the fact that double-infrared logarithms are promoted to double-ultraviolet logarithms with the splitting. We will see that there are substantial regularization scheme effects, even for the leading DLA terms. Two different classes of regularizations are considered. To the first class belongs the regularization considered in ref. : a sharp cutoff on the spatial loop momenta
$$|\stackrel{}{l}|<\mathrm{\Lambda }_S,$$
and a loop energy on the entire real axis,
$$\mathrm{}<l_0<\mathrm{}.$$
(132)
That means, roughly speaking, a discrete space and a continuous time. We believe that this regularization gives the same double-logarithm as the ordinary lattice regularization - the Wilson action . In the latter case all the components of the loop $`4`$-momentum are cutoff, not only the spatial ones
$$|l^\mu |<\mathrm{\Lambda }_4\frac{\pi }{a},$$
(133)
where $`a`$ is the lattice spacing. The physical reason for the equality of the double-logarithmic coefficients in the regularizations (132) and (133) is the following. Soft and collinear logarithms are both related to quasi-real gluon configurations, for which
$$l_0|\stackrel{}{l}|.$$
(134)
Cutting off the spatial momenta therefore should cut off also the relevant energies as far as soft and collinear singularities are concerned.
As a representative of the second class of UV regularizations, consider a sharp cutoff on the transverse momenta (the $`x`$$`y`$ plane):
$$|\stackrel{}{l}_{}|<\mathrm{\Lambda }_{},\mathrm{while}\mathrm{}<l_+,l_{}<\mathrm{}.$$
(135)
This regularization is “effective”, i.e. it is sufficient to cut on the transverse momenta to render the integrals finite. To this class of regularizations belongs the Dimensional Regularization (DR), in which most of the effective field theory computations have been performed. Let us treat the two cases in turn.
#### 6.1.1 Space Momenta Cutoff
An explicit computation of the gluon and the quark pole contributions on the light-cone in the $`\mathrm{\Lambda }_S`$-regularization gives
$`C_{g,n^2=0}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{log}^2{\displaystyle \frac{\mathrm{\Lambda }_S}{k_+i0}}`$
$`C_{q,n^2=0}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{log}^2{\displaystyle \frac{\mathrm{\Lambda }_S}{k_+i0}}{\displaystyle \frac{1}{2}}\mathrm{log}^2{\displaystyle \frac{E_X}{k_+i0}}.`$ (136)
The behaviour with respect to $`k_+`$ is the same as in the case $`n^20.`$ Ultraviolet divergences are now more severe than in the case $`n^20`$, being of double-logarithmic kind. However, the sum is again the correct one
$$C_{g,n^2=0}+C_{q,n^2=0}=C_s.$$
(137)
In other words, the transition to the light-cone theory implies a rearrangement of the ultraviolet structure, but the physical observable, $`C_s`$, is unchanged.
#### 6.1.2 Transverse momenta cutoff
The factor $`C_g`$ is better computed in this case by introducing light-cone coordinates:
$$l_+=l_0+l_3,l_{}=l_0l_3.$$
(138)
Integrating over $`l_{}`$ by closing the integration contour upward and over the transverse momentum, we obtain
$$C_g=2_0^{\mathrm{}}\frac{dyy}{1+y^2}\frac{1}{1+n^2y^2/4}\mathrm{log}\left[1+\frac{\mathrm{\Lambda }_{}}{k_+y}\left(1+\frac{n^2y^2}{4}\right)\right].$$
(139)
Performing the final integration in the case $`n^2=0`$, we obtain
$$C_{g,n^2=0}\left(\mathrm{\Lambda }_{}\right)=\mathrm{log}^2\frac{\mathrm{\Lambda }_{}}{k_+i0}.$$
(140)
For the quark-pole factor $`C_q,`$ the integration over $`l_{}`$ gives
$$C_q=_0^{\mathrm{}}𝑑xx_0^{\mathrm{\Lambda }_{}^2}𝑑l_{}^2\frac{1}{x^2l_{}^2+2E_Xx+1}\frac{1}{k_+x+1i0}.$$
(141)
Integrating over $`l_{}`$ we obtain, in the light-cone limit:
$$C_q=_0^{\mathrm{}}\frac{dx}{x}\frac{1}{1n^2x/4i0}\mathrm{log}\frac{\overline{\mathrm{\Lambda }}_{}^2x^2+x+1}{1+x},$$
(142)
with<sup>26</sup><sup>26</sup>26$`n^2`$ in the above formula has to be interpreted as $`2k_+/E_X`$.
$$\overline{\mathrm{\Lambda }}_{}\frac{\mathrm{\Lambda }_{}}{2E_X}.$$
(143)
The above integral has two double-logarithmic regions for $`\mathrm{\Lambda }_{}E_X,`$
$$(1):\mathrm{\hspace{0.17em}1}x\frac{4}{n^2},(2):\frac{1}{\overline{\mathrm{\Lambda }}_{}}x1.$$
(144)
Performing the integration in the two regions, we find
$$C_{q,n^2=0}\left(\mathrm{\Lambda }_{}\right)=\frac{1}{2}\mathrm{log}^2\frac{\mathrm{\Lambda }_{}^2}{E_X\left(k_+i0\right)}\mathrm{log}^2\frac{\mathrm{\Lambda }_{}}{E_X}(\mathrm{\Lambda }_{}E_X).$$
(145)
The first double-logarithm on the right-hand side is related to region $`(1)`$, the second one to region $`(2)`$. For a smaller UV cutoff, we obtain instead:
$$C_{q,n^2=0}\left(\mathrm{\Lambda }_{}\right)=\frac{1}{2}\mathrm{log}^2\frac{\mathrm{\Lambda }_{}^2}{E_X\left(k_+i0\right)},\left(E_X|k_+|\mathrm{\Lambda }_{}^2E_X^2\right).$$
(146)
Finally, for $`\mathrm{\Lambda }_{QCD}^2\mathrm{\Lambda }_{}^2E_X|k_+|`$, the integral $`C_q`$ vanishes in DLA.
#### 6.1.3 Comments
Let us comment on the results (140) and (145). As with the 3-momentum regularization, $`C_g`$ and $`C_q`$ have double-logarithmic UV divergences, again a consequence of the light-cone limit $`n^2=0`$. The most important point, however, is that $`C_g`$ has an additional factor of 2 with respect to the spatial cutoff case in the coefficient of the double-logarithm of the infrared scale, $`\mathrm{log}^2k_+`$ (cf. eqs. (113) and (140)). With the $`\mathrm{\Lambda }_S`$ regularization, $`C_q`$ has no $`\mathrm{log}^2k_+`$ term, while with the $`\mathrm{\Lambda }_{}`$ regularization it does. The same double logarithm is obtained in the sum $`C_s`$ in both regularizations. In general, the appearance of $`\mathrm{log}^2k_+`$ in $`C_q(\mathrm{\Lambda }_{})`$ implies that, with the $`\mathrm{\Lambda }_{}`$ regularization, $`C_q`$ does not describe only collinear contributions but also soft ones <sup>27</sup><sup>27</sup>27The double logarithm necessarily comes from the overlap.. We interpret this fact by saying that the shape function, in general, does not have any physical meaning, but it just represents the gluon-pole contribution to a physical process: that result is, as far as we know, new. One generally attaches to the shape function a physical meaning - related to the Fermi motion; thus, to understand what is happening, we have to start again from the beginning. The shape function is obtained from the original QCD tensor $`W_{\mu \nu }`$ considering the infrared limit of small momenta compared with the hadronic energy:
$$|l_\mu |E_X.$$
(147)
The tree-level rate in the ET equals the QCD one by construction. However, in loops, the condition (147) is not guaranteed: its validity depends on the regularization scheme adopted. If we cut all the loop-momentum components with a hard cutoff much smaller than the hard scale,
$$|l_\mu |\mathrm{\Lambda }_{UV}E_X,$$
(148)
then the condition (147) is still valid at the loop level. As a consequence, we expect that the leading, double-logarithmic term of the ET shape function will match the QCD one. That is indeed what happens with the spatial momentum regularization, as we have seen explicitly. On the other hand, when one uses a regularization such as DR or $`\mathrm{\Lambda }_{}`$, the equality of the double-logs is no longer guaranteed, and indeed it does not occur in $`\mathrm{\Lambda }_{}`$-regularization, as we have seen explicitly. This is because the longitudinal momentum of the gluon $`l_z`$, or equivalently its energy $`ϵ`$, can become arbitrarily large. For the latter regularizations, even for the double-logarithm, one has to come back to the original QCD loop diagram and perform factorization into a factor $`C_g`$ and a factor $`C_q`$, as we have shown in detail. In ref. it was shown that the factor of 2 in the $`\mathrm{log}^2k_+`$ term in DR is a regularization effect, i.e. it can be removed by going to a non-minimal dimensional scheme. We explicitly see, with the similar $`\mathrm{\Lambda }_{}`$ regularization, that by including $`C_q`$ the scheme-dependence automatically disappears. The origin of the additional factor of 2 in the transverse-momentum regularization is related to the occurrence of a second double-logarithmic region for $`|l_z|\mathrm{\Lambda }_{}`$ (very large rapidity).
Finally, as already noted, let us observe that in the case $`n^20`$ we expect the transverse momentum cutoff to give double-logarithmic results similar to those from the space momentum cutoff. That is because $`n^20`$ cuts the collinear emission at infinite rapidity.
## 7 The shape function in the low-energy effective theory
With the $`\mathrm{\Lambda }_{}`$ regularization, double logarithms are contained in $`C_g`$ as well as in $`C_q`$. Since we want to confine double logarithmic effects inside the shape function only, let us consider from now on the $`\mathrm{\Lambda }_S`$ regularization only. The factor $`C_q`$ is short-distance dominated in the latter regularization, so it is computed once and for all in perturbation theory and “leaves the game”.
Let us therefore return to formula (110) for $`C_g`$. Calling $`ϵ=|\stackrel{}{l}|,`$ and $`t=\theta ^2,`$ $`C_g`$ can be written as
$`C_g(\mathrm{\Lambda }_S,k_+)`$ $``$ $`{\displaystyle _0^{\mathrm{\Lambda }_S}}𝑑ϵ{\displaystyle _0^1}𝑑t{\displaystyle \frac{1}{2k_++ϵt}}`$ (149)
$``$ $`{\displaystyle _0^{\mathrm{\Lambda }_S}}{\displaystyle \frac{dϵ}{ϵ}}{\displaystyle _0^1}{\displaystyle \frac{dt}{t}}\theta \left(ϵtk_+\right)`$
$``$ $`{\displaystyle \frac{1}{2}}\mathrm{log}^2{\displaystyle \frac{\mathrm{\Lambda }_S}{k_+}},`$
where we have assumed $`\mathrm{\Lambda }_SO\left(E_X\right)`$ and we have used the approximation $`1/\left(2k_++ϵt\right)\theta \left(ϵt2k_+\right)/\left(ϵt\right)`$, which is valid in DLA. This form helps visualizing the origin of the double logarithm. We see that contributions come from soft regions, where $`ϵO\left(k_+\right)`$, as well as from hard regions, where $`ϵO\left(\mathrm{\Lambda }_S\right)`$. In order to separate them, the simplest way is to introduce another UV cutoff $`\mathrm{\Lambda }_{ET}`$, this time well below the hadronic energy $`E_X`$ (the hard scale of the process), such as
$$k_+\mathrm{\Lambda }_{ET}\mathrm{\Lambda }_S.$$
(150)
We can write
$$C_g(\mathrm{\Lambda }_S,k_+)=\delta Z(\mathrm{\Lambda }_S,\mathrm{\Lambda }_{ET},k_+)+\delta \overline{F}^{ET}(\mathrm{\Lambda }_{ET},k_+),$$
(151)
where
$`\delta Z(\mathrm{\Lambda }_S,\mathrm{\Lambda }_{ET},k_+)`$ $``$ $`{\displaystyle _{\mathrm{\Lambda }_{ET}}^{\mathrm{\Lambda }_S}}{\displaystyle \frac{dϵ}{ϵ}}{\displaystyle _0^1}{\displaystyle \frac{dt}{t}}\theta \left(ϵtk_+\right)`$ (152)
$`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{log}^2{\displaystyle \frac{\mathrm{\Lambda }_S}{k_+}}+{\displaystyle \frac{1}{2}}\mathrm{log}^2{\displaystyle \frac{\mathrm{\Lambda }_{ET}}{k_+}}`$
is a coefficient function and $`\delta \overline{F}^{ET}(\mathrm{\Lambda }_{ET},k_+)`$ is the one-loop contribution to the light-cone function $`\delta F^{ET}`$, multiplied by the propagator: $`\delta F^{ET}=\delta \overline{F}^{ET}/(k_++i0)`$, as defined in eq. (63),
$`\delta \overline{F}^{ET}(\mathrm{\Lambda }_{ET},k_+)`$ $``$ $`{\displaystyle _0^{\mathrm{\Lambda }_{ET}}}{\displaystyle \frac{dϵ}{ϵ}}{\displaystyle _0^1}{\displaystyle \frac{dt}{t}}\theta \left(ϵtk_+\right)`$ (153)
$`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{log}^2\left({\displaystyle \frac{\mathrm{\Lambda }_{ET}}{k_+}}\right).`$
Note that $`\delta \overline{F}^{ET}`$ depends only on the two scales $`k_+`$ and $`\mathrm{\Lambda }_{ET}`$. This is in line with the idea of a simple low-energy effective theory, which describes infrared phenomena characterized by the scale $`k_+`$, apart from the UV cutoff that enters through loop effects.
We assume that long-distance effects can be traced by the growth of the coupling constant in the proximity of the Landau pole, and that the coupling constant must be evaluated at the transverse momentum squared :
$$\alpha _S\alpha _S\left(k_{}^2\right),$$
(154)
where
$$k_{}^2ϵ^2t.$$
(155)
From the expression of $`\delta Z`$ we see that transverse momenta have a lower bound given by
$$l_{}^2>l_{}^2{}_{\mathrm{min}}{}^{}=\mathrm{\Lambda }_{ET}k_+.$$
(156)
According to our criteria, non-perturbative effects are absent from $`Z`$ as long as
$$l_{}^2{}_{\mathrm{min}}{}^{}\mathrm{\Lambda }_{QCD}^2.$$
(157)
According to eq. (156), this occurs when $`k_+`$ is non-vanishing, as it is for example if
$$k_+O\left(\mathrm{\Lambda }_{QCD}\right),$$
(158)
as expected from Fermi motion (since $`\mathrm{\Lambda }_{ET}\mathrm{\Lambda }_{QCD}`$). However, by taking the imaginary part of $`T_{\mu \nu }`$ to obtain $`W_{\mu \nu }`$, i.e. the rate, the product of factors is converted into a convolution over $`k_+`$ and the point $`k_+=0`$ is included in the integration range. This implies that transverse momenta down to zero contribute to the coefficient function in $`W_{\mu \nu }`$ , i.e. that factorization of short- and long-distance effects breaks down at this point. The breakdown is related to the fact that we are cutting the energies of the gluons, but not the emission angles, which can go down to zero, implying the vanishing of the transverse momenta. That is one of the most important outcomes of our analysis. However, we believe that these long-distance contributions are suppressed. Let us present a qualitative argument. As we can see from inequalities (156) and (157), transverse momenta of the order of the hadronic scale occur in $`Z`$ for a very small slice of values of $`k_+`$,
$$k_+\frac{\mathrm{\Lambda }_{QCD}^2}{\mathrm{\Lambda }_{ET}}\mathrm{\Lambda }_{QCD}.$$
(159)
If the integrand is not singular in this small slice, as it is natural to assume, it gives a reasonally small fraction of the total. Note that the usual factorization of mass singularities is instead “exact”. If we consider for example the moments of DIS cross-section, factorization involves a splitting of the long- and short-distance contributions of the form
$`M_N\left(Q^2\right)`$ $`=`$ $`{\displaystyle _0^1}𝑑x_Bx_B^{N1}\sigma _{DIS}(x_B,Q^2)`$
$`=`$ $`1+\gamma _N\alpha _S{\displaystyle _{m^2}^{Q^2}}{\displaystyle \frac{dl_{}^2}{l_{}^2}}=\left(1+\gamma _N\alpha _S{\displaystyle _{\mathrm{\Lambda }^2}^{Q^2}}{\displaystyle \frac{dl_{}^2}{l_{}^2}}\right)\left(1+\gamma _N\alpha _S{\displaystyle _{m^2}^{\mathrm{\Lambda }^2}}{\displaystyle \frac{dl_{}^2}{l_{}^2}}\right),`$
where $`m`$ is the mass of a light quark.
After the last step (152), the forward hadronic tensor takes the final form
$`T_{\mu \nu }^{QCD}`$ $`=`$ $`{\displaystyle \frac{s_{\mu \nu }}{2vQ}}F(k_+)^{QCD}`$
$`=`$ $`{\displaystyle \frac{s_{\mu \nu }}{2vQ}}{\displaystyle \frac{1}{k_++i0}}\left[1+aC_h\right]\left[1+aC_q\right]\left[1+a\delta Z\right]\left[1+a\delta \overline{F}^{ET}\right],`$
where the various factors have been introduced in eqs. (77), (78), (87) and (102). Taking the imaginary (absorptive) part, according to the optical theorem (44), we have for $`W_{\mu \nu }`$ the multiple convolution
$`W_{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{s_{\mu \nu }}{2vQ}}{\displaystyle 𝑑k_1𝑑k_2𝑑k_3𝑑k_4\delta \left(k_+k_1k_2k_3k_4\right)}`$ (162)
$`\left[\delta \left(k_1\right)+ac_h\left(k_1\right)\right]\left[\delta \left(k_2\right)+ac_q\left(k_2\right)\right]`$
$`\left[\delta \left(k_3\right)+a\delta z\left(k_3\right)\right]\left[\delta \left(k_4\right)+a\delta f^{ET}\left(k_4\right)\right],`$
where
$$f^{ET}(k_+,\mathrm{\Lambda }_{ET})=\delta \left(k_+\right)+a\delta f^{ET}\left(k_+\right)+O(a^2)$$
(163)
is the shape function, defined in eq. (10), for an on-shell quark ($`k_+^{}=0`$); moreover, we have defined
$`c_h\left(k_+\right)`$ $``$ $`{\displaystyle \frac{1}{\pi }}\mathrm{Im}\left[{\displaystyle \frac{1}{k_++i0}}C_h\left(k_+i0\right)\right]`$ (164)
$`=`$ $`\delta \left(k_+\right)C_h\left(k_+\right){\displaystyle \frac{1}{k_+}}\left({\displaystyle \frac{1}{\pi }}\right)\mathrm{Im}C_h\left(k_+i0\right)`$
and analogously for the other factors<sup>28</sup><sup>28</sup>28In ref. , formula (9) should be replaced by $`f^{QCD}(k_+)=𝑑k_1𝑑k_2\delta (k_+k_1k_2)(\delta (k_1)+a\delta z(k_1))f^{ET}(k_2)`$, where $`Z=1+a\delta Z`$.. Typically, by taking the imaginary parts, for double-logarithmic contributions, we have
$`{\displaystyle \frac{1}{\pi }}{\displaystyle \frac{\mathrm{log}^2\left(k_+i0\right)}{k_++i0}}`$ $``$ $`\delta \left(k_+\right)\mathrm{log}^2(k_+)+2\theta \left(k_+\right){\displaystyle \frac{\mathrm{log}\left(k_+\right)}{k_+}}`$ (165)
$`=`$ $`{\displaystyle \frac{d}{dk_+}}\left(\theta \left(k_+\right)\mathrm{log}^2\left(k_+\right)\right)`$
and for single-logarithmic ones
$`{\displaystyle \frac{1}{\pi }}{\displaystyle \frac{\mathrm{log}\left(k_+i0\right)}{k_++i0}}`$ $``$ $`\delta \left(k_+\right)\mathrm{log}(k_+)+\theta \left(k_+\right){\displaystyle \frac{1}{k_+}}`$ (166)
$`=`$ $`{\displaystyle \frac{d}{dk_+}}\left(\theta \left(k_+\right)\mathrm{log}\left(k_+\right)\right).`$
The last members of the above equations have to be interpreted as distributions. In DLA, according to eq. (153), $`f^{ET}`$ up to one loop reads
$`f^{ET}(k_+,\mathrm{\Lambda }_{ET})`$ $`=`$ $`\delta \left(k_+\right)+a\theta \left(k_+\right){\displaystyle \frac{\mathrm{log}\mathrm{\Lambda }_{ET}/\left(k_+\right)}{k_+}}{\displaystyle \frac{a}{2}}\delta \left(k_+\right)\mathrm{log}^2\left({\displaystyle \frac{\mathrm{\Lambda }_{ET}}{k_+}}\right)`$ (167)
$`=`$ $`\delta \left(k_+\right)+{\displaystyle \frac{a}{2}}{\displaystyle \frac{d}{dk_+}}\left(\theta \left(k_+\right)\mathrm{log}^2\left({\displaystyle \frac{\mathrm{\Lambda }_{ET}}{k_+}}\right)\right).`$
### 7.1 Evolution
Taking a derivative with respect to the logarithm of the cutoff, we obtain
$`{\displaystyle \frac{df(k_+,\mathrm{\Lambda }_{ET})}{d\mathrm{log}\mathrm{\Lambda }_{ET}}}`$ $`=`$ $`a\delta \left(k_+\right)\mathrm{log}\left({\displaystyle \frac{\mathrm{\Lambda }_{ET}}{k_+}}\right)+a{\displaystyle \frac{\theta \left(k_+\right)}{k_+}}`$ (168)
$`=`$ $`a{\displaystyle \frac{d}{dk_+}}\left(\theta \left(k_+\right)\mathrm{log}\left({\displaystyle \frac{\mathrm{\Lambda }_{ET}}{k_+}}\right)\right).`$
Comparing the above equation with the evolution equation for the shape function
$$\frac{df(k_+,\mathrm{\Lambda }_{ET})}{d\mathrm{log}\mathrm{\Lambda }_{ET}}=𝑑k_+^{}\mathrm{K}_S(k_+k_+^{};\mathrm{\Lambda }_{ET})f(k_+^{},\mathrm{\Lambda }_{ET}),$$
(169)
and taking into account that, at lowest order in $`\alpha _S`$, $`f(k_+^{},\mathrm{\Lambda }_{ET})=\delta \left(k_+^{}\right)`$ holds, we find for the evolution kernel at one loop
$`\mathrm{K}_S(k_+k_+^{};\mathrm{\Lambda }_{ET})`$ $`=`$ $`a\delta \left(k_+^{}k_+\right)\mathrm{log}\left({\displaystyle \frac{\mathrm{\Lambda }_{ET}}{k_+^{}k_+}}\right)+a{\displaystyle \frac{\theta \left(k_+^{}k_+\right)}{k_+^{}k_+}}`$
$`=`$ $`a{\displaystyle \frac{d}{dk_+}}\left(\theta \left(k_+^{}k_+\right)\mathrm{log}\left({\displaystyle \frac{\mathrm{\Lambda }_{ET}}{k_+^{}k_+}}\right)\right)`$
$`=`$ $`a\left[{\displaystyle \frac{\theta \left(k_+^{}k_+\right)}{k_+^{}k_+}}\delta \left(k_+^{}k_+\right){\displaystyle _0^{\mathrm{\Lambda }_{ET}}}{\displaystyle \frac{d\left(l_+^{}l_+\right)}{l_+^{}l_+}}\right].`$
If we consider the $`\mathrm{\Lambda }_{}`$-regularization, the evolution kernel for the shape function is instead given by (eq. (140)):
$`K_{}(k_+k_+^{};\mathrm{\Lambda }_{})`$ $`=`$ $`2a\delta \left(k_+^{}k_+\right)\mathrm{log}\left({\displaystyle \frac{\mathrm{\Lambda }_{}}{k_+^{}k_+}}\right)+2a{\displaystyle \frac{\theta \left(k_+^{}k_+\right)}{k_+^{}k_+}}`$ (171)
$`=`$ $`2a{\displaystyle \frac{d}{dk_+}}\left(\theta \left(k_+^{}k_+\right)\mathrm{log}\left({\displaystyle \frac{\mathrm{\Lambda }_{}}{k_+^{}k_+}}\right)\right).`$
We notice that there is a factor of 2 between the kernels (7.1) and (171) for the shape function in the two regularizations. The kernel in DR is the same as that in eq. (171), with $`\mathrm{\Lambda }_{}\mu `$.
There is a clear analogy of the evolution of the shape function with the Altarelli–Parisi evolution equation, but with an important difference: the evolution kernel in this case explicitly depends on the cutoff $`\mathrm{\Lambda }_{ET}`$ of the bare theory or on the renormalization point $`\mu `$ if we consider the renormalized theory<sup>29</sup><sup>29</sup>29We thank S. Catani for a discussion on this point.. All this is related to the fact that the Altarelli–Parisi evolution involves a single collinear logarithm for each loop, while our problem is double-logarithmic. Let us discuss this point with a simple analogy. The Altarelli–Parisi evolution, or in general the usual Callan–Symanzik evolution, is analogous to a first-order differential equation, which is autonomous (i.e. time-independent):
$$\frac{dx}{dt}=h(x),$$
(172)
or, in discrete form,
$$x_{n+1}=O\left(x_n\right),$$
(173)
where $`O`$ is a generic operator, such that the formal solution reads
$$x_n=O^n\left(x_0\right).$$
(174)
The evolution in eq. (169) is instead analogous to an evolution equation of the form
$$\frac{dx}{dt}=h(x,t),$$
(175)
or, in discrete form
$$x_{n+1}=O_n\left(x_n\right).$$
(176)
In the latter case there is a different evolution operator at each step<sup>30</sup><sup>30</sup>30In double-logarithmic problems, one can obtain an autonomous differential equation at the price of having a second-order equation, i.e. of the form
$$\frac{d^2x}{dt^2}=h(x).$$
This, anyway, is not an evolution equation..
We clarify at this point a discrepancy of a factor of 2 in the evolution kernel K of the shape function, computed at one loop in DR in both refs. and . We agree with ref. , where the kernel is computed from the Green function in the ET taking a $`\mu `$ derivative, as in eq. (171). We disagree with ref. , where the kernel is computed by taking the difference of the QCD Green function with the ET Green function and then differentiating with respect to $`\mu `$; their kernel is two times smaller than the one in eq. (171). The latter authors give for the QCD amplitude, in our notation, the result
$$F\left(k_+\right)^{QCD}\stackrel{\mathrm{?}}{=}\frac{1}{k_++i0}\left(\frac{1}{2}\right)a\mathrm{log}^2\left(\frac{\mu }{k_+i0}\right).$$
(177)
They find a dependence on the renormalization point $`\mu `$, which we do not find as the QCD diagram is ultraviolet - as well as infrared - finite . If we replace in their renormalization condition, which determines the kernel, our $`\mu `$-independent result for the QCD Green function,
$$F\left(k_+\right)^{QCD}=\frac{1}{k_++i0}\left(\frac{1}{2}\right)a\mathrm{log}^2\left(\frac{m_b}{k_+i0}\right),$$
(178)
we find a vanishing kernel <sup>31</sup><sup>31</sup>31In eq. (178) we have assumed $`E_XO\left(m_b\right)`$.. Since the effective theory is UV-divergent and consequently $`\mu `$-dependent, we believe that there may be a problem with the renormalization conditions. Schematically, the matrix element of a bare operator is of the form
$$p|O_B|p=1+c\frac{\alpha _B^{dim}}{ϵ}\left(p^2\right)^ϵ+(\mathrm{finite}\mathrm{for}ϵ0),$$
(179)
where $`c`$ is a numerical constant, $`p^2`$ refers to an overall momentum scale in the external state, and $`\left(p^2\right)^ϵ`$ comes from the one-loop integral in $`D=42ϵ`$ dimensions; $`\alpha _B^{dim}`$ is the bare coupling of the original $`D`$-dimensional theory: for $`D<4`$ it has a positive mass dimension $`4D=2ϵ`$ , and it must be kept fixed as we vary $`\mu ,`$ which is just an arbitrary mass scale:
$$\frac{d}{d\mu }\alpha _B^{dim}=0.$$
(180)
This implies the well-known condition
$$\frac{d}{d\mu }p|O_B|p=0.$$
(181)
One usually introduces an adimensionalized bare coupling multiplying $`\alpha _B^{dim}`$ by $`\mu ^{2ϵ},`$ where $`\mu `$ is just an arbitrary mass scale as we said before,
$$\alpha _B^{\mathrm{adim}}\mu ^{2ϵ}\alpha _B^{dim},$$
(182)
so that the bare Green function reads
$`p|O_B|p`$ $`=`$ $`1+c{\displaystyle \frac{\alpha _B^{\mathrm{adim}}}{ϵ}}\left({\displaystyle \frac{\mu ^2}{p^2}}\right)^ϵ+(\mathrm{finite}\mathrm{for}ϵ0)`$ (183)
$`=`$ $`1+c{\displaystyle \frac{\alpha _B^{\mathrm{adim}}}{ϵ}}+c\alpha _B^{\mathrm{adim}}\mathrm{log}{\displaystyle \frac{\mu ^2}{p^2}}+\mathrm{}`$
In the minimal-dimensional scheme (MS), we include the pole term in the renormalization constant
$$Z_{MS}=1+c\frac{\alpha _B^{\mathrm{adim}}}{ϵ},$$
(184)
and the remaining terms in the matrix element of the renormalized operator,
$$p|O_{MS}|p=1+c\alpha _B^{\mathrm{adim}}\mathrm{log}\frac{\mu ^2}{p^2}+\mathrm{},$$
(185)
since $`O_B=ZO_R`$. It is only after this splitting that a dependence on $`\mu `$ is introduced separately in the renormalization constant and in the renormalized operator<sup>32</sup><sup>32</sup>32In the notation of ref. , $`\mathrm{log}O_B=\mathrm{ln}\stackrel{~}{f}_B\left(\xi \right)/\mathrm{log}\xi `$, with $`\xi 1/\sqrt{p^2}.`$.
The anomalous dimension is computed from the renormalization constant keeping $`\alpha _B^{dim}`$ fixed:
$$\gamma \frac{d\mathrm{log}Z}{d\mathrm{log}\mu }=\frac{d}{d\mathrm{log}\mu }\left(c\frac{\alpha _B^{dim}\mu ^{2ϵ}}{ϵ}\right)=2c\alpha _B^{\mathrm{adim}}.$$
(186)
It seems to us that a vanishing kernel or anomalous dimension in the effective theory is obtained in ref. because the renormalization constant $`Z`$ has been identified with the whole matrix element (183).
## 8 Conclusions
We have discussed the properties of decays of heavy flavour hadrons into inclusive hadron states $`X`$ with an invariant mass $`m_X`$ small compared with the energy $`E_X`$, $`m_XE_X`$. An explicit factorization procedure has been introduced, which holds on a integral-by-integral basis. It is based on:
* the Cauchy theorem: it is exact and leads to the replacement of ordinary propagators with eikonal propagators in loop integrals;
* the lowering of a hard UV cutoff from $`\mathrm{\Lambda }_{UV}m_b`$ to $`\mathrm{\Lambda }_{UV}=\mathrm{\Lambda }_{ET}`$ $`m_b`$, where $`\mathrm{\Lambda }_{ET}`$ is the UV cutoff of the low-energy effective theory inside which the shape function is defined.
This technique has led us to a clean separation of all the perturbative and non-perturbative contributions. We have found that, while the exact kinematics of the original QCD processes involves a Wilson line off the light-cone for the final light quark, in the low-energy effective theory the light quark is necessarily described by a Wilson line on the light-cone.
We have analyzed the shape function $`f\left(k_+\right)`$ to find out which long-distance, non-perturbative, effects are contained in and which are not, in different regularization schemes. We found that $`f\left(k_+\right)`$, contrary to naive physical expectations, has no direct physical meaning even in double logarithmic approximation, as it represents a partial contribution to the complete physical process. Changing regularization, we have explicitly shown that the leading double-logarithmic contribution to $`f\left(k_+\right)`$ can be changed by a factor of 2, i.e. that the shape function is substantially regularization-scheme dependent. Only after summing the shape function with the other contributions, is a physical, scheme-independent result recovered. We have also shown that in lattice-like regularization the shape function factorizes a large part of the non-perturbative effects: it contains all the double logarithmic contributions of the full QCD process.
Subtracting from the forward hadronic tensor $`T_{\mu \nu }^{QCD}`$, step by step, each of the perturbative components, we end up with an explicit representation of the perturbative and non-perturbative effects. For instance, at one-loop order in DLA we have the result (see eq. (7)):
$$T_{\mu \nu }^{QCD}=\frac{s_{\mu \nu }}{2vQ}\frac{1}{k_++i0}\left[1+aC_h\right]\left[1+aC_q\right]\left[1+a\delta Z\right]\left[1+a\delta \overline{F}^{ET}\right].$$
The coefficient $`C_h`$ is a hard factor that takes into account the fluctuations with energy $`\epsilon `$ in the range $`E_X<\epsilon <m_B.`$ The other two coefficients, $`C_q`$ and $`\delta Z`$, are short-distance-dominated in lattice-like regularization schemes; $`\delta \overline{F}^{ET}`$ is long-distance-dominated in any regularization. The tensor $`W_{\mu \nu }`$, i.e. the rate, is obtained (as usual) by taking the imaginary part of $`T_{\mu \nu }`$.
Another outcome of our analysis is that, contrary to single logarithmic problems, factorization in this (double-logarithmic) problem is not exact, even in lattice-like regularization schemes. Some long-distance effects are present in the coefficient function: they come from gluons with a large energy but with a very small emission angle and consequently with a small transverse momentum. These non-perturbative effects in the coefficient function however are expected to be suppressed on physical grounds, as they occur in a small region of the phase space for a moderately large cutoff of the effective theory.
Finally, we have clarified some discrepancy in the literature about the evolution kernel for the shape-function computed in double logarithmic approximation inside dimensional regularization.
Acknowledgements
We would like to thank G. Martinelli for inspiring discussions. One of us (U.A.) has benefited from many conversations with S. Catani. We also thank D. Anselmi, M. Battaglia, M. Beneke, M. Cacciari, M. Ciafaloni, S. Frixione, M. Grazzini, M. Greco, G. Korchemsky, N. Uraltsev, B. Webber and in particular G. Veneziano. One of us (G.R.) would like to thank the Theory group of CERN, where this work was completed. |
warning/0003/cond-mat0003194.html | ar5iv | text | # Chern-Simons Theory for Quantum Hall Stripes
## Abstract
We develop a Chern-Simons theory to describe a two-dimensional electron gas in intermediate magnetic fields. Within this approach, inhomogeneous states emerge in analogy to the intermediate state of a superconductor. At half filling of the highest Landau level we find unidirectional charge-density-wave (CDW) solutions. With a semiclassical calculation we give an intuitive explanation of the change of CDW orientation in the presence of an in-plane magnetic field. An anisotropy in the electron band mass is suggested as a possible source of the reproducible orientation of the CDW.
In two-dimensional electron systems both interactions and disorder have surprising and striking consequences. A strong perpendicular magnetic field quenches the kinetic energy and allows for the observation of both fractional quantum Hall effect and composite fermions in the lowest Landau level (LLL), where correlations among electrons are especially important. In higher Landau levels on the other hand, interest has mainly been devoted to reentrant localization-delocalization transitions causing the integer quantum Hall effect, with interactions thought of as modifying its non-universal features. The discovery of a pronounced resistivity anisotropy close to half filling of higher LLs with filling factor $`\nu >4`$ has recently attracted attention to interaction effects in this regime. The transport anisotropy is consistent with the formation of a unidirectional CDW state which was predicted theoretically on the basis of Hartree-Fock calculations . Numerical exact diagonalization studies also support the idea of a CDW formation . The orientation of this CDW structure is related to the anisotropy observed in resistivity .
In this letter, we show that a CDW ground state arises naturally in the framework of a Chern-Simons (CS) theory . The CS approach successfully describes the fractional quantum Hall effect and composite fermions in the LLL. In intermediate LLs it is complementary to Hartree-Fock calculations which become exact in the limit of high Landau levels . Our bosonic CS theory has solutions similar to the intermediate state of superconductors and allows for an intuitive interpretation of the striped phase: in the highest, partially filled LL (HLL) domain walls separate incompressible “superconducting” stripes with local filling factor $`\nu _{\mathrm{HLL}}=1`$ from empty stripes ($`\nu _{\mathrm{HLL}}=0`$) penetrated by flux. The CDW wave vector is determined by the competition of the domain-wall energy favoring a long wave length and the Hartree energy which generally increases with CDW wave length but has minima for special wave vectors due to the modulation of the electron charge density by the wave-function form factor. The CDW wave length and structure are calculated numerically within a mean-field approximation. Furthermore, we calculate the wave-function form factor in the presence of an in-plane magnetic field using a semiclassical approximation. We find that for a thin quantum well the in-plane component of the field is parallel to the axis of high resistivity, in agreement with experiments and numerical calculations . The band mass anisotropy is shown to provide a mechanism determining the the experimentally reproducible orientation of the CDW wave vector for perpendicular fields.
Chern-Simons theory. — To employ the CS approach we first map the problem of a partially filled $`N`$th LL onto an effective LLL problem. The scattering of electrons between different LLs can be neglected when the Bohr radius $`a_\mathrm{B}=(4\pi \mathrm{}^2ϵ_0ϵ_\mathrm{r})/(me^2)`$ is much smaller than the magnetic length $`\mathrm{}=\sqrt{\mathrm{}/eB}`$, a condition satisfied in experimentally investigated systems. We describe screening by the completely filled LLs by a wave-vector dependent dielectric constant $`ϵ(q)=ϵ_0ϵ_\mathrm{r}\{1+(2/qa_\mathrm{B})[1\mathrm{J}_0^2(qR_\mathrm{c})]\}`$ with $`\mathrm{J}_0`$ denoting the zeroth order Bessel function and $`R_\mathrm{c}=\sqrt{2N+1}\mathrm{}`$ the cyclotron radius. The intra–LL matrix elements of the Coulomb interaction depend on the LL index $`N`$ only via the Laguerre polynomial $`\mathrm{L}_N`$ in the form factor of the single particle eigenstates . One obtains the same matrix elements for states in the zeroth instead of the $`N`$-th LL but with the effective interaction $`V^{\mathrm{eff}}(𝐪)=\frac{e^2}{2ϵ(q)q}\mathrm{L}_N^2(\mathrm{}^2q^2/2)`$. The Laguerre polynomial can be interpreted as the quotient of the form factor of the $`N`$-th and the zeroth LL as $`\mathrm{L}_N=F_N/F_0`$. Since all LLs span equivalent Hilbert spaces, the new interaction $`V^{\mathrm{eff}}`$ in the LLL gives a faithful representation of the original problem.
The Fermi-Dirac statistics of electrons is mimicked by bosons carrying one flux quantum $`\mathrm{\Phi }_0=h/e`$ with them. The flux quanta are attached by a singular gauge transformation with a statistical gauge field $`𝐚(𝐫)`$ tied to the density of bosons $`\rho (𝐫)=\varphi ^{}(𝐫)\varphi (𝐫)`$ via the condition $`\mathbf{}\times 𝐚(𝐫)=\mathrm{\Phi }_0\rho (𝐫)`$. The bosonic CS Hamiltonian is given by
$$=d^2r\varphi ^{}(𝐫)\frac{1}{2m}\left(i\mathrm{}\mathbf{}+e\delta 𝐚\right)^2\varphi (𝐫)+\frac{1}{2}d^2rd^2r^{}\delta \rho (𝐫)V^{\mathrm{eff}}(𝐫𝐫^{})\delta \rho (𝐫^{}),$$
(1)
with $`\delta 𝐚=𝐀+𝐚`$, $`\delta \rho =\mathrm{\Phi }^{}\mathrm{\Phi }\overline{\nu }_{\mathrm{HLL}}B/\mathrm{\Phi }_0`$, and $`e`$ the charge of the proton.
Mean-field analysis. — We perform a mean-field (MF) analysis of the bosonic CS problem and replace the operators by a complex field. The kinetic energy (of the bosons, not to be confused with the quenched kinetic energy of the electrons) alone would be minimized by the homogeneous saddle points $`\varphi \sqrt{B/\mathrm{\Phi }_0}`$, $`\delta 𝐚0`$ and $`\varphi 0`$, $`\delta 𝐚𝐀`$ corresponding to a completely filled or empty LL, respectively. Without the interaction potential the ground state of a half filled LL would clearly be one domain with filling factor $`\nu _{\mathrm{HLL}}=1`$ and another domain with $`\nu _{\mathrm{HLL}}=0`$. However, due to the long-range nature of the Coulomb interaction, these large domains break up into smaller ones in such a way that the sum of domain-wall and Coulomb energy is minimal, and a charge-density wave is formed in analogy to the intermediate state of a superconductor. The effective field $`\mathbf{}\times \delta 𝐚`$ is expelled from the “superconducting” regions with $`\nu _{\mathrm{HLL}}=1`$. For average fillings $`\overline{\nu }_{\mathrm{HLL}}=1ϵ`$ close to one on the other hand, the solution with lowest energy is given by a flux-tube array in analogy to phases of type-I superconductors in the intermediate state, where each flux tube can carry a single or multiple flux quanta. For the electron problem this is equivalent to the formation of a Wigner crystal.
To make the physics of the mean-field solution transparent, we express the wave function by amplitude and phase, introduce a gauge invariant “velocity” $`𝐐`$, and express the boson current density $`𝐣`$ in these new variables
$`\varphi (𝐫)=\sqrt{{\displaystyle \frac{B}{\mathrm{\Phi }_0}}}f(𝐫)e^{i\theta (𝐫)},𝐐(𝐫)=\mathrm{}^2\left(\mathbf{}\theta +{\displaystyle \frac{2\pi }{\mathrm{\Phi }_0}}\delta 𝐚\right),𝐣(𝐫)={\displaystyle \frac{e\omega _\mathrm{c}}{2\pi \mathrm{}^2}}f^2𝐐(𝐫).`$ (2)
We enforce the constraint on $`𝐚(𝐫)`$ with the help of a Lagrange multiplier $`\lambda `$ (which is the rescaled zero component of the CS field) and find the saddle-point equations
$`f^{\prime \prime }+{\displaystyle \frac{1}{\mathrm{}^4}}𝐐^2f+2(u_\mathrm{H}\lambda )f=0,\mathbf{}\times 𝐐=(v+f^21)\widehat{𝐳},\mathbf{}\lambda ={\displaystyle \frac{1}{\mathrm{}^4}}f^2𝐐\times \widehat{𝐳}`$ (3)
with the vortex density $`𝐯(𝐫)=\mathrm{}^2\mathbf{}\times \mathbf{}\theta (𝐫)`$. The Hartree potential $`u_\mathrm{H}`$ is expressed in units of $`\mathrm{}^2\mathrm{}\omega _\mathrm{c}`$. One immediate consequence of the saddle-point equations is the screening of an electrostatic potential by the CS superconductor in analogy to the screening of a vector potential by a conventional superconductor. A negative potential causes a reduction of $`\nu _{\mathrm{HLL}}=f^2`$ from one and induces a current via the second part of Eq. (3). This argument can be made quantitative by describing the current response to an electric field with the help of the Hall conductivity $`\sigma _\mathrm{H}=e^2/h`$ as $`𝐣=\sigma _\mathrm{H}𝐄\times \widehat{𝐳}`$. From the third part of Eq. (3) one finds $`\mathbf{}\lambda =(e/\mathrm{}\omega _\mathrm{c}\mathrm{}^2)𝐄`$ leading to an exact cancellation of the external potential in the Schrödinger equation for the superfluid density $`f^2`$. Via this mechanism a Bose condensation in a spatially varying external potential is possible. Note that the bosonic wave functions are not restricted to the LLL in the variational approach. This is an valid approximation in the limit where the electron-electron interaction is weak compared to the inter-Landau-level spacing and a mixing of Landau levels is negligible. The same assumption is used also in the Hartree-Fock approaches in order to justify that the ground state does not involve single-particle states above the partially filled highest LL. Since on the scale $`\mathrm{}`$ of the typical particle distance the strength of the bosonic interaction $`V^{\mathrm{eff}}`$ is comparable to that of the original fermionic interaction (for $`N`$ not too large as in experiments), the ground state of the bosonic problem (1) will lie to a good approximation in the LLL, even if we search in an extended space of functions including higher LLs. Dropping the variational constraint in the bosonic problem therefore induces inaccuracies of the same order of magnitude as ignoring the Landau-level mixing in the fermionic problem. As can be seen from the variational result below (Fig. 1), the local filling factor overshoots only by a few percent over one, which indicates that the solution lies to a very good approximation within the LLL although the constraint is not imposed.
1D numerical solution. — In order to confirm the physical picture developed above and to provide a comparison of our CS approach with previous Hartree-Fock approaches , we have numerically determined the MF solution at half filling ($`\overline{\nu }_{\mathrm{HLL}}=\frac{1}{2}`$). This was done approximately by assuming a unidirectional modulation, i.e., that the charge density is modulated with a period $`\mathrm{\Lambda }`$ in $`x`$ direction and is constant in $`y`$ direction. In this approximation, vortices cannot be taken into account as point-like objects. Instead, all vortices within a period have to be concentrated on lines, $`v(𝐫)=_n\frac{\mathrm{\Lambda }}{2}\delta (x(n+1/2)\mathrm{\Lambda })`$. For symmetry reasons, currents then flow only in $`y`$ direction, $`𝐐(𝐫)=Q_y(x)𝐞_y`$. To obtain the MF solution, we use a variational ansatz $`f(x)=_kf_k\mathrm{cos}(2\pi kx/\mathrm{\Lambda })`$ with the constraint $`f((n+\frac{1}{2})\mathrm{\Lambda })=0`$ at the vortex positions. In a first step, we minimize $``$ for given $`\mathrm{\Lambda }`$ by varying the coefficients $`f_k`$ with the constraints of half filling and $`_xQ_y(x)=v+f^21`$. Then the resulting energy $`(\mathrm{\Lambda })`$ is minimized as a function of $`\mathrm{\Lambda }`$ to obtain the CDW period. The choice of $`a_\mathrm{B}=\mathrm{}/\sqrt{\nu }`$ fixes the interaction strength and is representative for experiments . A full 2D numerical solution and the analysis of gauge field fluctuations is beyond the scope of the present analysis.
We present our results for filling factors $`\nu =2N+\frac{1}{2}`$ in the Landau levels $`N=8`$ and $`N=2`$. The structure of the MF solutions is depicted in Fig. 1. For $`N=8`$, $`f^2`$ shows a clear separation between full ($`f^21`$) and empty ($`f^20`$) regions. The width of the transition between the full and empty region is of order $`\mathrm{}`$. Therefore, the separation between full and empty regions is blurred out for smaller $`N`$ (for $`N=2`$, see Fig. 1b). The actual charge and charge-current density, which are obtained after convoluting $`f^2`$ and $`f^2Q_y`$ with the appropriate form factors are much smoother than $`f^2`$ and vary only on the scale of $`R_\mathrm{c}`$. The charge density within this LL shows a relative variation of about 20% for $`N=8`$. This magnitude is roughly consistent with Hartree-Fock calculations for large $`N`$. For $`N=2`$ the charge density shows a stronger modulation of about 60%. Thus, although charge and charge-current density are modulated only on the scale $`\mathrm{\Lambda }`$ for all $`N`$, the relative variation of charge density is found to increase with decreasing $`N`$.
The energy per electron as a function of the periodicity is displayed in Fig. 2. The energy scale in our saddle-point approximation is the cyclotron energy. We expect that gauge-field fluctuations will renormalize it to the Coulomb energy $`e^2/4\pi ϵ_0ϵ_\mathrm{r}\mathrm{}`$. The optimal periodicity $`\mathrm{\Lambda }_N`$, given by the global minimum of this function, is found at $`\mathrm{\Lambda }_83.1R_\mathrm{c}`$ and $`\mathrm{\Lambda }_23.9R_\mathrm{c}`$. Thus, the periodicity is larger than expected from the first zero of the form factor at $`\mathrm{\Lambda }2.6R_\mathrm{c}`$ in agreement with HF calculations . However, we find that $`\mathrm{\Lambda }_N/R_\mathrm{c}`$ increases notably with decreasing $`N`$. In our approach this effect can be understood from the increasing weight of the “kinetic” or domain-wall contribution within the CS energy.
Tilted magnetic field. — So far, we have studied the formation of a CDW by electrons in high Landau levels within the CS approach. In experiments the longitudinal resistivity appears to be systematically higher in $`[1\overline{1}0]`$ directions as compared to $`[110]`$ directions. This raises the question of what mechanism breaks the fourfold crystal symmetry and orienting the CDW wave vector parallel to the $`[1\overline{1}0]`$ directions. For this reason resistivity was studied in tilted magnetic fields and a sufficiently strong in-plane component $`𝐁_{}`$ of the magnetic field was found to orient the directions of large resistivity parallel to $`𝐁_{}`$. We present a mechanism for the coupling between the CDW wave vector and $`𝐁_{}`$ that involves the dynamics of the electrons in the $`z`$ direction perpendicular to the plane. We model the confinement of the electrons in this direction by a harmonic potential with eigenfrequency $`\omega _0`$. While similar models were studied recently in the framework of a Hartree-Fock approach and of Laughlin-like states , we use a simple semiclassical approach to obtain the anisotropic form factor of cyclotron orbits in tilted fields, which then can easily be combined with the theory of untilted fields.
The classical equations of motion for a single electron subject to a tilted magnetic field $`𝐁=(B_{},0,B_z)`$ and a harmonic confining potential are linear and possess two eigenmodes with frequencies $`\omega _\pm ^2=\frac{1}{2}(\omega _0^2+\omega _z^2+\omega _x^2)\pm \frac{1}{2}\sqrt{(\omega _0^2+\omega _z^2+\omega _x^2)^24\omega _0^2\omega _z^2}`$, where $`\omega _x`$ and $`\omega _z`$ are the cyclotron frequencies associated with the in-plane and out-of-plane components of the magnetic field, which typically are small compared to $`\omega _0`$. For vanishing $`B_{}`$, the modes with frequency $`\omega _+^2\omega _0^2[1+\omega _x^2/(\omega _0^2\omega _z^2)]`$ or $`\omega _{}^2\omega _z^2[1\omega _x^2/(\omega _0^2\omega _z^2)]`$ become a pure oscillation in $`z`$ direction or a cyclotron orbit in the $`xy`$ plane, respectively. We determine the semiclassical orbits from the Bohr quantization condition $`𝐩_\sigma 𝐫_\sigma =2\pi \mathrm{}(\frac{1}{2}+n_\sigma )`$, where $`\sigma =\pm `$ and we choose the ground state, $`n_+=0`$, for the oscillator-like mode and the and the $`N`$th Landau level, $`n_{}=N`$, for the cyclotron-like mode. This choice is unique for $`\omega _0>\omega _z`$. Subsequently, we consider only the projection of the orbits onto the $`xy`$ plane, which are ellipses with half axes of lengths
$`R_{x\sigma }^2=(1+2N_\sigma ){\displaystyle \frac{\mathrm{}}{m\omega _\sigma }}{\displaystyle \frac{\omega _z^2}{\omega _z^2+\omega _x^2\omega _\sigma ^4/(\omega _0^2\omega _\sigma ^2)^2}},R_{y\sigma }^2={\displaystyle \frac{\omega _\sigma ^2}{\omega _z^2}}R_{x\sigma }^2.`$ (4)
The cyclotron-like ellipse is larger than the oscillator-like ellipse ($`R_\alpha >R_{\alpha +}`$). In addition, the cyclotron-like ellipse is longer in $`x`$ direction than in $`y`$ direction ($`R_x>R_y`$), whereas for the oscillator-like ellipse the opposite holds true (cf. Fig. 2b). Because of this contrary deformation the resulting anisotropy of the form factor is delicate. The form factor $`F(𝐪)`$ is the Fourier transform of the density $`\rho (𝐫)`$ of finding the electron at a position $`𝐫=𝐫_{}(t)+𝐫_+(t)`$. It factorizes into two contributions from the normal modes, $`F(𝐪)=F_{}(𝐪)F_+(𝐪)`$, which are the Fourier transforms of the density $`\rho _\pm (𝐫)=\frac{1}{T_\pm }_0^{T_\pm }𝑑t\delta (𝐫𝐫_\pm (t))`$ of the normal mode orbits, where $`T_\pm =2\pi /\omega _\pm `$ is the orbit period.
In the presence of the tilted field the form factor is anisotropic and we now discuss the location of its first zero at small momenta, which essentially determines the CDW periodicity and energy. At small momenta, the zeros of $`F`$ are those of $`F_{}`$ since the cyclotron-like orbits have a larger radius. If one estimates the CDW periodicity from the location of the zeros of $`F_{}(𝐪)`$ one finds the periodicity proportional to the radius of the cyclotron ellipse in the direction of the CDW wave vector. The total energy of the CDW is then minimized if its wave vector points parallel to the long direction of the cyclotron ellipse, i.e., parallel to $`𝐁_{}`$. This conclusion is in agreement with experiments .
Finally, we estimate the energy difference between an orientation of the CDW wave vector parallel and perpendicular to the in-plane field in the LL $`N=2`$ (assuming again $`\omega _x\omega _z\omega _0`$). The energy change is related to a change of the CDW period which can be approximated by $`\delta \mathrm{\Lambda }/\mathrm{\Lambda }(R_x/R_\mathrm{c}1)(\omega _x/\omega _0)^2`$. The change of the energy per particle with $`\mathrm{\Lambda }`$ is estimated numerically as $`\delta E0.1\mathrm{}\omega _\mathrm{c}\delta \mathrm{\Lambda }/\mathrm{\Lambda }`$ from the domain-wall energy. The tilt energy per particle is then given by $`E_{\mathrm{tilt}}=\mathrm{tan}^2(\theta )\left(\omega _z/\omega _0\right)^23\mathrm{K}`$ ¿From the threshold angle $`\theta 10^{}`$ where the CDW orientation follows the in-plane component of the field and a typical ratio $`\omega _z0.15\omega _0`$ we find $`E_{\mathrm{tilt}}1`$mK. Consequently, the intrinsic symmetry breaking mechanism in the absence of a tilt must lead to an energy gain of the same magnitude. We now examine two possible mechanisms.
Anisotropy. — In the normal state the electron gas has anisotropic resistivity with “hard” axis parallel to the $`[110]`$ crystallographic direction, i.e., perpendicular to the hard axis observed in the CDW state. This normal anisotropy, which has been reported for some of the samples studied in , can be ascribed to scattering processes caused by an interface roughness with anisotropic correlation lengths $`\xi _{[1\overline{1}0]}>\xi _{[110]}`$ . The observed surface morphology is consistent with such an anisotropic interface roughness. Since the anisotropy in the CDW state is determined by collective pinning, whereas it is determined by single-electron scattering in the normal state, the direction of the hard axis is not necessarily identical and therefore we now examine the collective pinning mechanism. Variations of the interface height $`d`$ result in an electrostatic potential $`\varphi _d(𝐪)=\frac{1}{2}end(𝐪)/ϵ(q)`$. Here $`n`$ is the density of positive donors and the height fluctuations are assumed to be Gaussian correlated, $`[|d(𝐪)|^2]_{\mathrm{av}}=2\pi d_0^2\xi _x\xi _y\mathrm{exp}[(q_x^2\xi _x^2+q_y^2\xi _y^2)/2]`$ with correlation lengths $`\xi _x`$, $`\xi _y\mathrm{\Lambda }`$ and a typical height $`d_02.8`$Å of the order of a GaAs monolayer. Pinning leads to an energy gain via local displacements of the CDW profile which are governed by a smectic elastic response . From the curvature of the CDW energy as a function of the wave length we estimate the compression modulus $`K_x10^{11}`$eVm<sup>-2</sup> and the bending modulus $`K_y(\mathrm{\Lambda }/4\pi )^2K_x`$. For parameters representing the experiments we estimated the pinning energy density in a linear response calculation ,
$`E_{\mathrm{rough}}\left({\displaystyle \frac{n\pi e}{2\nu \mathrm{\Lambda }}}\right)^2{\displaystyle \frac{d^2q}{(2\pi )^2}\frac{\left[|\varphi _d(q_x+2\pi /\mathrm{\Lambda },q_y)|^2\right]_{\mathrm{av}}}{K_xq_x^2+K_yq_y^4}}.`$ (5)
Note that the pinning strength is given by the disorder correlator near the wave vector $`2\pi /\mathrm{\Lambda }`$ and therefore it is exponentially small, $`\left[|\varphi (q_x+2\pi /\mathrm{\Lambda },q_y)|^2\right]_{\mathrm{av}}e^{2(\pi \xi _x/\mathrm{\Lambda })^2}`$. For $`\xi _x\mathrm{\Lambda }`$ this exponential dependence suppresses the pinning energy per particle by orders of magnitude below 1mK per particle as determined above from the threshold tilt angle. Therefore interface roughness is irrelevant for the CDW orientation. In addition it would imply that the CDW wave vector is aligned in the direction of the shorter disorder correlation length, i.e., the $`[110]`$ direction in contradiction to experiments.
As pointed out in Ref. the electric field $`E=en/2ϵ_0ϵ_\mathrm{r}`$ in $`[00\overline{1}]`$ direction between the electron system and the donor layer provides a different symmetry breaking mechanism. We now show that this field generates an anisotropic electronic band mass which induces the observed CDW orientation by an energy gain on the appropriate scale. Since Ga and As carry opposite partial charges they are displaced such that the GaAs bonds in $`[11]`$ direction are stretched whereas the bonds in $`[1\overline{1}]`$ direction are shortened. From pressure experiments it is known that a shortening of bonds leads to an increase of the effective band mass $`m^{}`$. We estimate the effective mass changes $`m_{[1\overline{1}0]/[110]}^{}/m^{}1\pm 10^4`$ by identifying the electric field with an effective pressure . This anisotropic band mass affects the CDW via the form factor. A semiclassical analysis implies elliptic cyclotron orbits with half axes $`R_{[1\overline{1}0]/[110]}/R_c1\pm 10^4`$ which leads to an increased CDW period $`\delta \mathrm{\Lambda }_{[1\overline{1}0]/[110]}/\mathrm{\Lambda }\pm 10^4`$. In analogy to the analysis of the in-plane field this implies that the mass anisotropy prefers the $`[1\overline{1}0]`$ orientation over the $`[110]`$ orientation by about 1mK per electron in agreement with the experimentally observed orientation and the energy scale derived from the tilt experiments.
Conclusion. — We have developed a bosonic Chern-Simons theory to investigate the formation of a CDW state in intermediate Landau levels and the influence of an anisotropic interface potential on the CDW orientation. In a semiclassical calculation we have found that the CDW wave vector is parallel to an in-plane magnetic field for quantum wells with a hard confining potential. The band mass anisotropy is a possible origin for the CDW orientation in perpendicular fields.
Acknowledgments. — We are grateful to B. I. Halperin, M. M. Fogler, M. P. Lilly, R. H. Morf, and F. von Oppen for stimulating discussions and thank A. Kleinschmidt for assistance during the initial stages of the numerical part of this work. B.R. was supported by DFG grant Ro2247/1-1. |
warning/0003/nlin0003011.html | ar5iv | text | # Experimental generation of steering odd dark beams of finite length
## I Introduction
Physically, Dark Spatial Solitons (DSSs) are localized intensity dips existing on stable background beams as a result of an exact counterbalance of diffraction and nonlinearity. A necessary condition for their existence is the presence of a phase dislocation in the wavefront along which the phase is indeterminate and the field amplitude is zero. Besides the intriguing physical picture, the particular interest in the DSSs is motivated by their ability to induce gradient optical waveguides in bulk self-defocusing nonlinear media . The only known truly two-dimensional (2D) DSSs are the Optical Vortex Solitons (OVSs) , whereas in one transverse spatial dimension DSSs manifest themselves as dark stripes . The odd initial conditions required to generate a fundamental 1D DSS correspond to an abrupt $`\pi `$-phase jump centered along the irradiance minimum of the stripe. The OVSs have a more complicated phase profile described by $`\mathrm{exp}(im\phi )`$, where $`\phi `$ is the azimuthal coordinate in a plane perpendicular to the background beam propagation direction and $`m`$ – the so-called topological charge (TC) – is an integer number. This phase function ensures a $`\pi `$-phase jump in each diametrical slice through the vortex core. Fundamental DSSs of these types have the common feature of zero transverse velocity if no perturbations are present. In contrast to that, ring dark solitary waves slowly change their parameters even when born from ideal odd initial conditions .
In the pioneering work of Nye and Berry it is conjectured that mixed edge-screw dislocations cannot exist. Despite that, almost two decades later an indication for their existence was found for two interacting optical vortices of opposite topological charges. In our recent experiments on the generation of quasi-2D DSSs we found that moderate saturation of the nonlinearity can stabilize the snake instability which usually leads to their decay. This made possible the first identification of 1D Odd Dark Beams (ODBs) of finite length with their characteristic edge-screw phase dislocations . The mixed dislocation forces the dark beams to steer in space. This appears to be of practical interest provided that there are effective ways to control the ODBs transverse velocity.
In this article we report the first experimental realization of steering odd dark beams of finite length with mixed phase dislocations under controllable initial conditions. Two approaches to control their transverse velocity are investigated experimentally and compared with numerical simulations.
## II Experimental setup and results
### A Computer-generated holograms
The phase portrait of the mixed edge-screw dislocation (see Fig. 1 in Ref.) consists of a pair of semi-helices with a phase difference of $`\pi `$ to which an effective topological charge of $`\pm 1/2`$ can be ascribed. Their spatial offset $`2b`$ determines the length of the edge part of the dislocation and ensures a phase jump of $`\mathrm{\Delta }\phi `$ in the direction perpendicular to the dark stripe of finite length. The phase function of this mixed phase dislocation can be described by
$$\mathrm{\Phi }_{\alpha ,\beta }(x,y)=\mathrm{\Delta }\phi \left\{\frac{\beta }{\pi }\mathrm{arctan}\left(\frac{\alpha x}{y+b\beta }\right)+\frac{(1\alpha )}{2}\text{sgn}(x)\right\},$$
(1)
where $`x`$ and $`y`$ denote the transverse Cartesian coordinates perpendicular and parallel to the dark beam. $`2b`$ stands for the length of the edge part of the dislocation and
$$\alpha =\{\begin{array}{cc}0\hfill & \text{for }|y|b\hfill \\ 1\text{ and }\beta =1\hfill & \text{for }y>b\hfill \\ 1\text{ and }\beta =1\hfill & \text{for }yb\hfill \end{array}.$$
(2)
The pattern of the Computer-Generated Holograms (CGHs) used to produce this phase distribution consists of parallel lines which become curved at the position of the semi-vortex cores. In the edge part of the dislocation they terminate and reappear shifted, for a $`\pi `$-jump by one half of the pattern period. Holograms with such structures correspond to interference lines shifted along an imaginary line of finite length and to curved lines limiting the dislocations as observed in our previous experiment (see Fig. 8 in ). The binary CGHs used are photolithographically fabricated with a grating period of $`20\mu m`$. Several holograms with various lengths of the edge part of the dislocation are etched on a common substrate. Special attention was paid to align the edge parts of the different dislocations correctly on their common substrate. The simplicity of varying the dislocation length and magnitude of the phase jump is the main advantage of the approach we have chosen. The diffraction efficiency at first orders is measured to be $`9\%`$, close to the theoretical $`10\%`$ limit for binary holograms. The unavoidable quantization inaccuracy of $`\pi /24`$ for holograms of this type is negligible for measurements with phase jumps of $`\mathrm{\Delta }\phi =3\pi /4`$, $`\pi `$, and $`5\pi /4`$ which are presented in this work.
### B Experimental setup
The setup used is similar to that in our previous experiments (see Fig. 1 in Ref.). Briefly, the beam of a single-line $`Ar^+`$ laser ($`\lambda =488nm`$) is used to reconstruct the CGHs. The first-diffraction-order beam with the phase dislocation nested in is filtered through a slit and is focused on the entrance of the $`10cm`$ long Nonlinear Medium (NLM). After the desired propagation path length, the beam is deflected by a prism immersed in the nonlinear liquid and is projected directly on a Charge-Coupled Device (CCD) camera array with a resolution of $`13\mu m`$. The NLM is ethylene glycol dyed with DODCI (Lambdachrome) to reach an absorption coefficient of $`0.107cm^1`$. In a calibration measurement we generated 1D DSSs by using CGHs of the type described in Sec. I. The soliton constant $`Ia^2`$ (i.e. the product of the background-beam intensity $`I`$ and the square of the dark beam width $`a`$ measured at the $`1/e`$-level) was found to reach its asymptotically constant value for input powers of $`P_{\mathrm{sol}}^{1D}33mW`$. It is known that thermally self-defocusing liquids are both nonlocal and saturable. Since the saturation of the nonlinearity is able to modify the ODBs transverse velocity and profile, we needed to estimate it and to account for it in our numerical calculations. In an independent measurement we realized a self-bending scheme similar to that used in . The asymmetry required was introduced by an intentional tilt of the prism immersed in the NLM, which resulted in different nonlinear propagation path-lengths for the different parts of the background beam. The strength of the self-bending effect was measured in the near field. For an absorptive nonlocal medium the choice of a suitable saturation model is not trivial (, see also Sec. IV of ). We found a good fit for the experimental data with the equation $`\mathrm{\Delta }yI/(1+I/I_{sat})^\gamma `$. Using it we estimated $`P_{sat}100mW`$ and $`\gamma =3`$. In addition to the careful alignment of the CGHs on the substrate, the holograms were reproduced to achieve vertical dark beam steering, which is not sensitive to possible undesired weak horizontal self-deflection of the background beam. Changing the ODB parameters (length-to-width ratio and magnitude of the phase jump) was performed by a strict horizontal translation of the substrate. The accuracy of the alignment was tested by checking for equal steering of the ODBs reproduced from two identical holograms placed at opposite ends of the series of aligned CGHs. In this work the ODBs are identified by the corresponding lengths $`2b`$ of the edge portions of the dislocations in CGH pixels ($`1pix.=5\mu m`$) as encoded in the holograms. The deflection $`\mathrm{\Delta }x`$ of the dark beams is measured in units of CCD camera pixels. We estimate that measurements with an encoded dislocation length of $`b/5\mu m`$ CGH pixels corresponds to a dislocation length-to-ODB width ratio $`(2b/a)=1/10b/5\mu m`$ in the numerical simulations (e.g. 14 pix. dislocation length corresponds to $`b/a=1.4`$ in the simulations).
### C ODB steering vs. dislocation length
All data presented in this subsection refer to $`\pi `$ phase jumps across the edge parts of the mixed dislocations. In Fig. 1a we plot the deflection $`\mathrm{\Delta }x`$ of ODBs with different dislocation lengths for input powers of $`1.7mW`$ (circles), $`33mW`$ (triangles) and $`67mW`$ (squares). The data obtained at $`P=1.7mW`$ refer to a linear regime of propagation. The results at $`33mW`$ and $`67mW`$ are extracted from the experimental pictures shown in Fig. 1b and Fig. 1c, respectively, which are recorded for a nonlinear propagation path length of $`z=8.5cm`$. The general tendency of linear increase of the deflection with decreasing the dislocation length is clearly expressed (Fig. 1a, solid line). The shortest mixed phase dislocation encoded was only $`1pix.`$ long. In this case the strong deviation from the linear dependence is caused by the annihilation of the semi-charges due to the shortening of the edge part of the dislocation. It will be shown later that this shortening accelerates for higher input powers (intensities). For that reason, even the ODB with an initially $`10pix.`$ long phase dislocation appears gradually less deflected at $`P=67mW`$ as compared to the case of $`P=33mW`$ (Fig. 1a). In Fig. 1b,c the thick solid lines are intended to denote the positions of the ODBs at the entrance of the NLM. Because ODB steering is present also in the linear regime of propagation (Fig. 1a, dots) these positions (with respect to the dark beam intensity minimum and trailing peak maximum) are identified by numerical simulations for $`b/a=2.5`$. The identification corresponds to an encoded dislocation length of $`25pix.`$, i.e. to the most right frame shown in Fig. 1b,c. Somewhat surprising is the weak sensitivity of the ODB deflection vs. background beam power (intensity) far from TC annihilation. It can be intuitively understood by recalling the known interaction scenario of well-separated OVSs of opposite TCs . In this case the attraction between the OVSs is negligible as compared to their translation as a pair. At a constant input power of $`33mW`$ we measured the ODB deflection vs. nonlinear propagation path length (Fig. 2). As expected, the ODB with $`10pix.`$ long dislocation has higher transverse velocity than that one with $`22pix.`$ long dislocation. The linearity in the dependencies is also well pronounced.
### D Phase control of the ODB steering
As a second possible way to control the dark beam deflection, we considered the variations in the magnitude of the phase jump $`\mathrm{\Delta }\phi `$ across the edge part of the mixed dislocation. In Fig. 3 we compare the experimental dependencies $`\mathrm{\Delta }x(\mathrm{\Delta }\phi )_{b=22pix.}`$ (squares) and $`\mathrm{\Delta }x(b)_{\mathrm{\Delta }\phi =\pi }`$ (dots). The straight lines are the respective linear fits. Because a problem in encoding a larger set of phase jumps in the CGHs was recognized too late, we measured the deflection $`\mathrm{\Delta }x`$ at $`\mathrm{\Delta }\phi =3\pi /4`$, $`\pi `$, and $`5\pi /4`$ only. In view of that the linear fit of the phase dependence in Fig. 3 appears to be assailable. Its linearity, however, is confirmed by numerical simulations (see Sec. III). The dependence of $`\mathrm{\Delta }x`$ on $`\mathrm{\Delta }x(\mathrm{\Delta }\phi )`$ and $`\mathrm{\Delta }x(b)`$ has been plotted in one figure in order to underline the fact that it appears to be easier to deflect the ODB by controlling the phase than by controlling the dislocation length. The measurements are performed at a constant power of $`33mW`$.
### E Power/Intensity dependencies
The ability of the dark spatial solitons (and the dark spatial waves ) to induce gradient all-optical waveguides in bulk self-defocusing NLM originates in the negative nonlinear correction to the linear refractive indices of the media. In view of that, the intensity dependencies remain undoubtfully of interest, despite of the low sensitivity of the ODB steering on the input power (intensity). In Fig. 4 we present experimental data on the power dependence of the length of the edge portion of the mixed dislocation for $`\mathrm{\Delta }\phi =\pi `$ and for two different lengths of $`14pix.`$ and $`22pix.`$ encoded in the holograms. It is easy to understand that the (mixed) phase dislocations do not remain sharp and of an unchanged magnitude when the (odd) dark beam steers . The dislocation lengths are estimated by evaluating the respective longitudinal ODB slices at $`5\%`$ of the background beam intensity (i.e. at the actual noise level in the recorded frames). Generally, the dislocation length decreases monotonically with increasing the input power. Asymptotically, the dislocation flattens and disappears, provided that the ratio $`b/a`$ is less than 2. As mentioned in , at $`b/a4`$ the ODBs should be expected to bend due to the snake instability . In fact we observed such a behavior for ODBs with encoded dislocation lengths ranging from 4 to 6. A vortex-beam creation is recognized by the convergence of two neighboring interference lines in one. However, the vortices formed by this instability remained with highly overlapping cores .
In Fig. 5a,b we plot the measured ODB widths (a) and their lengths (b) at the $`1/e`$-level as a function of the background beam power. The initial lengths of the mixed phase dislocations are denoted in pixels as encoded on the respective CGHs, whereas the widths and lengths at the exit of the NLM are measured in units of CCD camera pixels. The strong decrease in both the ODB width and length up to $`17mW`$ (approximately $`P_{sol}^{1D}/2`$) is followed by an approximate recovering of the width and length at approximately $`P_{sol}^{1D}`$. At higher input powers both transverse quantities decrease but the tendency is slower as compared to the situation below $`P_{sol}^{1D}/2`$ and stabilizes asymptotically at high saturation. As it will be discussed in the next section, the minima in the curves plotted in Fig. 5a,b result from the reshaping of the OD beams at the particular $`1/e`$ intensity-level chosen for evaluation. The estimation there shows that the first rapid decrease in both transverse dimensions of the ODBs does not correspond to a ‘soliton constant’ formation. Generally speaking, the ODBs analyzed are no solitary waves in the widely adopted sense, since they do not survive, for instance, a collision with a second ODB steering in the opposite direction . Nevertheless, the narrowing in both transverse directions for higher powers (intensities) should improve their guiding ability when signal beams or pulses are to be transmitted inside the ODBs and deflected in space. This feature will be addressed elsewhere. It is interesting to note that, independent of the length of the mixed dislocations ($`2b`$), the widths ($`a`$) of the dislocations should be initially equal (in the near field behind the CGHs), but appeared different near the entrance of the NLM. The estimation has shown a CGH-to-NLM distance of approximately 4 Rayleigh diffraction lengths with respect to the initial ODB width. The well pronounced separation between the curves in Fig. 5a should be attributed to the different two-dimensional diffraction at different initial ODB length-to-width ratios.
## III Numerical simulations
In our numerical simulations we tried to model the experimentally obtained dependencies by accounting for the estimated moderate saturation of the nonlinearity ($`P_{sol}^{1D}=33mW`$; $`P_{sat}=100mW`$ at $`\gamma =3`$; see Sec. II B). The $`(2+1)`$-dimensional nonlinear evolution of the steering ODB of finite length in the bulk homogeneous and isotropic NLM is described by the generalized nonlinear Schrödinger equation
$$i\frac{E}{\zeta }+\frac{1}{2}\left(\frac{^2}{\xi ^2}+\frac{^2}{\eta ^2}\right)E\frac{L_{Diff}}{L_{NL}}\frac{|E|^2}{(1+s|E|^2)^\gamma }E=0,$$
(3)
where the transverse spatial coordinates are normalized to the initial dark-beam width $`(\zeta =x/a,\eta =y/a)`$, and the propagation path length is expressed in Rayleigh diffraction lengths $`L_{\mathrm{Diff}}=ka^2`$. Further, $`L_{NL}=(kn_2I_0)^1`$ is the nonlinear length, $`k`$ is the wavenumber inside the NLM, and $`I_0`$ is the background beam intensity. The adopted correction for the nonlinear refractive index is
$$\mathrm{\Delta }n=n_2|E|^2/(1+s|E|^2)^\gamma ,$$
(4)
with $`s=P_{sol}^{1D}/P_{sat}=0.3`$. As it was done in , the slowly varying electric-field amplitude of the ODB was chosen $`tanh`$-shaped
$$E(x,y)=\sqrt{I_0}B(r_{1,0}(x,y))\mathrm{tanh}\left[\frac{r_{\alpha ,\beta }(x,y)}{a}\right]e^{i\mathrm{\Phi }_{\alpha ,\beta }(x,y)},$$
(5)
where $`r_{\alpha ,\beta }(x,y)=\sqrt{x^2+\alpha (y+\beta b)^2}`$ is the effective Cartesian/radial coordinate, $`\mathrm{\Phi }_{\alpha ,\beta }`$ is the phase distribution of the mixed edge-screw dislocation (see Eq. 1), and $`\alpha `$ and $`\beta `$ are given by Eq. 2. The width $`w`$ of the super-Gaussian background beam
$$B(r)=\mathrm{exp}\left\{\left(\sqrt{\frac{x^2+y^2}{w^2}}\right)^{14}\right\}$$
(6)
is chosen to exceed at least 15 times the ODB lengths. The model equation (3) was solved numerically by the beam propagation method on a $`1024\times 1024`$ grid. It should be mentioned, that the initial width $`a(z=0)`$ of the ODB of finite length was chosen to correspond to that of an infinite 1D ODSS ($`a=a_{sol}^{1D}=const.I_0^{1/2}`$). It was proven numerically that the ODB deflection is insensitive to the particular value of $`a`$, and Figs. 6-8 are generated under this assumption. Actually, the nonlinearity causes an appreciable reshaping of the beams, in particular in the first evolution stage when the ODB starts steering, see Fig. 2 in Ref. ). In order to improve the similarity between the experimental data (Figs. 4 and 5) and the numerical results (Figs. 9 and 10), an initial ODB width twice as large as in the experiment is assumed. In Fig. 6 we plot the ODB deflection vs. $`b/a_{z=0}`$ for different input powers (intensities). All data presented in this section refer to a normalized propagation path length of $`z=4L_{NL}`$, which corresponds to that estimated for the experiment. Nevertheless all calculations are carried out up to $`10L_{NL}`$ whereby no qualitative deviation from the tendencies discussed is seen. The linearity in the ODB deflection vs. $`b/a`$ is well obeyed except for $`b/a>2.2`$. The longer ODBs steer slower, bend slightly, and decay into pairs of vortex beams for $`b/a>4`$. In the linear regime of propagation, the ODBs also deflect but the deflection is stronger at higher input powers/intensities. This is more pronounced for shorter dislocations. In the experimental data (Fig. 1a), this behavior is much weaker, rather the deflection remains within the experimental accuracy. Looking for an adequate explanation, in a series of simulations we checked that a $`\pm 30\%`$ inaccuracy in estimating $`P_{sat}`$ results only in $`\pm 5\%`$ deviation in the ODB deflection at $`z=4L_{NL}`$. The observed tendency of a decrease of the ODB steering velocities at increased saturation is well understood but seems insufficient in quantity. We attribute the absence of a well expressed power dependence in Fig. 1a to the NLM nonlocality. In a separate experiment it was estimated, that the nonlocality in this medium is negligible on a spatial scale of several hundred of micrometers only .
In Fig. 7 we plot the ODB deflection vs. the nonlinear propagation path length $`z/L_{NL}`$ for $`b/a=1.0`$ and $`2.2`$. As in the previous figure the magnitude of the edge part of the phase dislocation is set to $`\mathrm{\Delta }\phi =\pi `$. Ones the ODB starts steering, its transverse velocity remains constant (see Fig. 2). The longer ODBs with longer edge dislocations, however, emit dispersive waves in their first evolution stage ($`z<1L_{NL}`$). This causes a ‘delay’ in the steering along the NLM (Fig. 7, lower curve).
In Fig. 8 we present results obtained for the phase-dependent control of the ODB deflection at different input powers (intensities). In qualitative agreement with the experimental observation, the linear increase in the ODB steering with decreasing the magnitude of the phase jump down to $`\mathrm{\Delta }\phi =0.5\pi `$ is evident. The comparison of Figs. 6 and 8 confirms the conclusion that at a fixed nonlinear propagation distance the phase-controlled ODB deflection is more efficient as compared to that by varying the $`b/a`$ ratio.
Fig. 9 is intended to clarify the origin of the non-monotonic power dependencies of the ODB widths (and lengths) observed at input powers below $`P_{sol}^{1D}`$ (see Fig. 5). The solid line represents the ODB full width at half maximum (FWHM), the dashed one the full width at the $`1/e`$-intensity level. In these simulations $`\mathrm{\Delta }\phi =\pi `$ and $`b/a=1.4`$ are assumed. Qualitatively, we obtained the same curves also for $`b/a=2.2`$ by accounting for the initial free-space propagation in the experiment (from the CGH to the entrance of the NLM; $`z4L_{Diff}`$). The minimum in the ODB width evaluated at the $`1/e`$ level originates from the reshaping of the beam profile which is caused by the moderate saturation. A similar reshaping is reported in (see Figs. 2-4 therein). At $`b/a=2.2`$ the data obtained for the ODB length vs. input power (intensity) were found to be even more sensitive to the intensity level of evaluation. Because of the transverse steering of the ODBs of finite length the edge portions of the dislocations shorten monotonically with increasing the background beam power (intensity) (Fig. 10). The numerical results are in a very good qualitative agreement with the experimental ones (Fig. 4).
## IV Conclusion
The results presented show that the inherent steering dynamics of odd dark beams of finite length can be effectively controlled by varying both the magnitude $`\mathrm{\Delta }\phi `$ and the relative length $`b/a`$ of the mixed edge-screw phase dislocation. The background-beam intensity has weak influence on the steering but is important for keeping the optically-induced gradient waveguides steep, which is crucial for all-optical guiding, deflection and switching of signal beams or pulses. Since the mixed phase dislocations shorten and flatten along the nonlinear media (tending asymptotically to washout) the ODBs seem to be promising primarily for future short-range all-optical switching devices.
###### Acknowledgements.
A. D. would like to thank the Alexander von Humboldt Foundation for the award of a fellowship and the opportunity to work in the stimulating atmosphere of the Max-Planck-Institut für Quantenoptik (Garching, Germany). This work was also supported by the National Science Foundation of Bulgaria. |
warning/0003/hep-th0003294.html | ar5iv | text | # Restrictions on Gauge Groups in Noncommutative Gauge Theory
## I Introduction
The Yang-Mills theories naturally arise as low energy limits of the theory of open strings. One can obtain Yang-Mills theories with different gauge groups by studying different D-brane configurations (see e.g. ). For instance, if we place $`N`$ D-branes on top of each other in the flat space, the corresponding open string theory gives rise to the Yang-Mills theory with the gauge group $`U(N)`$. One can also obtain gauge theories with other gauge groups such as $`SO(N)`$ and $`Sp(N)`$ by using the orientifold construction. In more detail, one combines the spatial reflection $`\sigma \pi \sigma `$ on the string world-sheet with the target space reflection, $`X^\mu X^\mu ,\mu =1,\mathrm{},k`$ and $`X^\mu X^\mu ,\mu =(k+1),\mathrm{},10`$. It is the goal of this note to study which gauge groups can be realized in the presence of the background $`B`$-field when the brane world-volume turns into a noncommutative space .
Interaction with the $`B`$-field introduces an extra term into the Polyakov action of the string ,
$$\mathrm{\Delta }S=\frac{i}{2}_\mathrm{\Sigma }B_{\mu \nu }ϵ^{ab}_aX^\mu _bX^\nu .$$
(1)
Here $`a,b=1,2`$ are world-sheet indices, $`\mathrm{\Sigma }`$ is the worldsheet and $`B_{\mu \nu }=B_{\nu \mu }`$ is the $`B`$-field on the target space. One requires the expression (1) to be invariant with respect to the orientifold reflection. This implies the following transformation rules for components of the $`B`$-field,
$`B_{}B_{}`$ $`B_{}B_{}`$ (2)
$`B_{}B_{}`$ $`B_{}B_{}.`$ (3)
Here the symbols $``$ and $``$ stand for the target space indices $`\mu =1,\mathrm{},k`$ and $`\mu =(k+1),\mathrm{},10`$, respectively. In the orientifold construction we finally let the branes lie on the orientifold. The continuity of the $`B`$-field implies that the $`B`$-field on the brane, $`B_{}`$, vanishes. Hence, the brane world-volumes are commutative since it is $`B_{}`$ which is responsible for the noncommutativity . This consideration indicates that one should encounter difficulties in the construction of the gauge theories with gauge groups $`SO(N)`$ and $`Sp(N)`$ on noncommutative spaces. Somewhat surprisingly, we find that the reduction from $`U(N)`$ to $`SU(N)`$ in the framework of noncommutative geometry also fails.
## II The closure of classical Lie algebras under the Moyal commutator
In the flat case the presence of a constant $`B`$-field turns the D-branes into noncommutative spaces, with the ordinary pointwise multiplication of functions replaced by the Moyal product,
$`(XY)(x)`$ $`=`$ $`\mathrm{exp}({\displaystyle \frac{i}{2}}\theta ^{ij}_i^x_j^y)X(x)Y(y)|_{x=y}=`$ (4)
$`=`$ $`XY+{\displaystyle \frac{i}{2}}\theta ^{ij}_iX_jY+\mathrm{}`$ (5)
Here $`X`$ and $`Y`$ are functions on the D-brane world-volume, and $`\theta ^{ij}`$ is a real-valued constant antisymmetric tensor constructed of the metric and $`B`$-field . The Moyal product naturally extends to $`N`$ by $`N`$ matrices, formula (3) still applies. One can also introduce the Moyal commutator by the formula,
$$[X,Y]_{}=XYYX.$$
(6)
In what follows we check whether the matrix Lie algebras of the classical Lie groups $`SO(N),U(N),SU(N)`$ and $`Sp(N)`$ are closed under the Moyal commutator. We choose to work in the fundamental representation of these Lie algebras.
The Lie algebra of $`U(N)`$ consists of anti-Hermitean matrices, $`\overline{X^t}=X`$, where the bar stands for complex conjugation. We first show that this algebra is closed under the Moyal commutator. The key observation is the following property of the Moyal product,
$$\overline{(XY)^t}=\overline{Y^t}\overline{X^t}.$$
(7)
By using the ordinary rules for the transpose of matrices we get,
$`(XY)^t=Y^tX^t+`$ (8)
$`+{\displaystyle \frac{i}{2}}\theta ^{ij}_jY^t_iX^t{\displaystyle \frac{1}{8}}\theta ^{ij}\theta ^{kl}_j_lY^t_i_kX^t+\mathrm{}`$ (9)
The construction for higher order terms is obvious. Now we apply the complex conjugation and rename the indices of $`\theta `$ to obtain,
$`\overline{(XY)^t}=\overline{Y^t}\overline{X^t}+{\displaystyle \frac{i}{2}}\theta ^{ij}_i\overline{Y^t}_j\overline{X^t}`$ (10)
$`{\displaystyle \frac{1}{8}}\theta ^{ij}\theta ^{kl}_i_k\overline{Y^t}_j_l\overline{X^t}+\mathrm{}=\overline{Y^t}\overline{X^t}.`$ (11)
Taking into account $`\overline{X^t}=X`$ and $`\overline{Y^t}=Y`$ yields,
$`\overline{[X,Y]_{}^t}`$ $`=`$ $`\overline{(XY)^t}\overline{YX)^t}=`$ (12)
$`=`$ $`\overline{Y^t}\overline{X^t}\overline{X^t}\overline{Y^t}=`$ (13)
$`=`$ $`YXXY=[X,Y]_{}`$ (14)
which shows that the algebra $`U(N)`$ is closed under the Moyal commutator.
We now turn to the algebras of $`SO(N)`$, $`SU(N)`$ and $`Sp(N)`$. We first show that for $`N=2`$ these algebras are not closed with respect to the Moyal commutator.
The counter examples for both $`SO(2)`$ and $`Sp(2)`$ are given by formulas,
$$X=\left(\begin{array}{cc}0& \alpha \\ \alpha & 0\end{array}\right)Y=\left(\begin{array}{cc}0& \beta \\ \beta & 0\end{array}\right).$$
(15)
and the counterexample for $`SU(2)`$ is
$$X=\left(\begin{array}{cc}i\alpha & 0\\ 0& i\alpha \end{array}\right)Y=\left(\begin{array}{cc}i\beta & 0\\ 0& i\beta \end{array}\right).$$
(16)
Here $`\alpha `$ and $`\beta `$ are coordinates on the manifold chosen so that $`\theta ^{\alpha \beta }0`$. This can always be done unless $`\theta =0`$ and the Moyal product coincides with the ordinary multiplication of matrix-valued functions. With $`X`$ and $`Y`$ as given above one can easily compute the Moyal commutator since all derivatives of order higher than one vanish. The result for both counter examples is
$$[X,Y]_{}=i\theta ^{\alpha \beta }\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$
(17)
Note that this matrix has a nonvanishing trace. Since the Lie algebras of $`SO(2),SU(2)`$ and $`Sp(2)`$ consist of traceless matrices, we conclude that they are not closed under the Moyal commutator. This also applies to $`SO(N),SU(N)`$ and $`Sp(N)`$ for arbitrary $`N`$ because they contain $`SO(2),SU(2)`$ and $`Sp(2)`$ as subgroups.
Acknowledgements First of all I would like to thank Anton Alekseev who has been very helpful during the preparation of this note. I would also like to thank Ulf Danielsson, Martin Kruczenski and Mårten Stenmark for helpful and encouraging discussions. |
warning/0003/nucl-th0003046.html | ar5iv | text | # EVENT-BY-EVENT PHYSICS IN RELATIVISTIC HEAVY-ION COLLISIONS
## 1 Introduction
The importance of event-by-event physics is evident from the following simple analogy: Stick a sheet of paper out of your window on a rainy day. Keeping it there for a long time - corresponding to averaging - the paper will become uniformly wet and one would conclude that rain is a uniform mist. If, however, one keeps the sheet of paper in the rain for a few seconds only, one observes the striking droplet structure of rain<sup>1</sup><sup>1</sup>1Originally, this analogy was given by Prof. A.D. Jackson. Incidentally, one has also demonstrated the liquid-gas phase transition! Analyzing many events gives good statistics and may reveal rare events as snow or hail and thus other phase transitions. The statistics of droplet sizes will also tell something about the fragmentation, surface tension, etc. By varying initial conditions as timing and orienting the paper, one can further determine the speed and direction of the rain drops.
Central ultrarelativistic collisions at RHIC and LHC are expected to produce about $`10^4`$ particles, and thus present one with the remarkable opportunity to analyze, on an event-by-event basis, fluctuations in physical observables such particle multiplicities, transverse momenta, correlations and ratios. Analysis of single events with large statistics can reveal very different physics than studying averages over a large statistical sample of events. The use of Hanbury Brown–Twiss correlations to extract the system geometry is a familiar application of event-by-event physics in nuclear collisions and elsewhere, e.g, in sonoluminoscence . The power of this tool has been strikingly illustrated in study of interference between Bose-Einstein condensates in trapped atomic systems . Fluctuations in the microwave background radiation as recently measured by COBE restrict cosmological parameters for the single Big Bang event of our Universe. Large neutrons stars velocities have been measured recently which indicate that the supernova collapse is very asymmetrical and leads to large event-by-event fluctuations in “kick” velocities during formation of neutron stars.
The tools applied to study these phenomena do, however, vary in order to optimize the analysis and due to limited statistics. The COBE and the interference in Bose-Einstein condensates require study of fluctuations within a single event. The HBT studies in heavy-ion collisions and sonoluminosence requires further averaging over many events in order to obtain sufficient statistics; one has not yet studied fluctuations in source radii event-by-event. Anisotropic flow requires an event-plane reconstruction in each event but again averaging over many events is necessary to obtain a statistically relevant measurement of the flow. The event-by-event fluctuations in heavy-ion collisions (and neutron star kick velocities) go a step further by studying variations from event to event.
Studying event-by-event fluctuations in ultrarelativistic heavy ion collisions to extract new physics was proposed in a series of papers , in which the analysis of transverse energy fluctuations in central collisions was used to extract evidence within the binary collision picture for color, or cross-section, fluctuations. More recent theoretical papers have focussed on different aspects of these fluctuations, such as searching for evidence for thermalization , correlations between transverse momentum and multiplicity , critical fluctuations at the QCD phase transition and other correlations between collective quantities .
Intermittency studies of factorial moments of multiplicities are related to event-by-event fluctuations. One of the motivations for intermittency studies was the idea of self-similarity on small scales, an idea borrowed from chaos theories. The factorial moments of particle multiplicities did find approximate power law behavior when the intervals of rapidity and angles were made increasingly smaller, at least until a certain small scale. The power law scaling in nucleus-nucleus collisions was, however, weaker than in proton-proton collisions. This indicated that the stronger correlations in proton-proton collisions were mainly due to resonances, minijets and other short range correlations, but that they were averaged out in nuclear collisions by summing over the many individual participating nucleons. The scaling was not a collective phenomenon and indications of new physics were not found . In more recent event-by-event fluctuation studies the self-similar scaling idea is abandoned. They concentrate on the mean and the variance of the particle multiplicities per event and correlations between different particle species, transverse momentum, azimuthal angle, etc. One directly compares to expectations from proton-proton collisions scaled up by the number of participants. One follows these fluctuations and correlations for heavy-ion collisions as function of centrality and system size searching for anomalous behavior as compared to proton-proton collisions.
Recently NA49 has presented a prototypical event-by-event analysis of fluctuations in central Pb+Pb collisions at 158 GeV per nucleon at the SPS, which produce more than a thousand particles per event . The analysis has been carried out on $``$100.000 such events measuring fluctuations in multiplicities, particle ratios, transverse momentum, etc.
Results from the RHIC collider are eagerly awaited . The hope is to observe the phase transition to quark-gluon plasma, the chirally restored hadronic matter and/or deconfinement. This may be by distinct signals of enhanced rapidity and multiplicity fluctuations in conjunction with J/$`\mathrm{\Psi }`$ suppression, strangeness enhancement, $`\eta ^{}`$ enhancement, constant (critical) temperatures vs. transverse enery or rapidity density , transverse flow or other collective quantities as function of centrality, transverse energy or multiplicity as will be discussed in detail below.
The purpose of this review is to understand these and other possible fluctuations. We find that the physical observables, including multiplicity, kaon to pion ratios, and transverse momenta agree well with recent NA49 data at the SPS, and indicate that such studies do not yet reveal the presence of new physics. Predictions for RHIC and LHC energies are given. The centrality dependence with and without a phase transition to a quark-gluon plasma is discussed - in particular, how the physical quantities are expected to display a qualitative different behavior in case of a phase transition, and how a first order phase transition could be signaled by very large fluctuations.
## 2 Phase Transitions and Fluctuations
Lattice QCD calculations find a phase transition in strongly interacting matter which is accompanied by a strong increase of the number of effective degrees of freedom . The Early Universe underwent this transition at a time $`t=0.30.4(T_c/MeV)^2`$ seconds. For a hadronic gas melting temperature of $`T_c=150`$ MeV this occurred around $`15`$ microseconds after the Big Bang. By colliding heavy nuclei we expect to reproduce this transition at sufficiently high collisions energies.
### 2.1 Order of the QCD phase transition
The nature and order of the transition is not known very well. Lattice calculations can be performed for zero quark and baryon chemical potential only, $`\mu _B=0`$, where they suggest that QCD has a weak first order transition provided that the strange quark is sufficiently light , that is for 3 or more massless quark flavors. The transition is due to chiral symmetry restoration and occur at a critical temperature $`T_C150`$ MeV. In pure SU(3) gauge theory (that is no quarks, $`N_f=0`$) the transition is a deconfinement transition which is of first order and occurs at a higher temperature $`T_c260`$ MeV.
However, when the strange or the up and down quark masses become massive, the QCD transition changes to a smooth cross over. The phase diagram is then like the liquid-gas phase diagram with a critical point above which the transition goes continuously through the vapor phase. For reasonable values for the strange quark mass, $`m_s150`$ MeV and small up and down quark masses, lattice calculations find either a weak first order transition or a smooth soft cross-over . In case of a weak first order transition, the latent heat and density discontinuities and the signals, that depend on these quantities, will be small.
For exactly two massless flavors, $`m_{u,d}=0`$ and $`m_s=\mathrm{}`$, the transition is second order at small baryon chemical potential. Random matrix theory finds a 2nd order phase transition at high temperatures which, however, change into a 1st order transition above a certain baryon density - the tricritical point. For small up and down quark masses the transition changes to a continuous cross-over at zero baryon chemical potential but remains a first order at large baryon chemical potential. A critical point must therefore exist at small but finite baryon chemical potential which may be searched for in relativistic heavy-ion collisions .
### 2.2 Density, rapidity, temperature and other fluctuations
Fluctuations are very sensitive to the nature of the transition. In case of a second order phase transition the specific heat diverges, and this has been argued to reduce the fluctuations drastically if the matter freezes out at $`T_c`$ . For example, the temperature fluctuations have a probability distribution
$`w\mathrm{exp}(C_V({\displaystyle \frac{\mathrm{\Delta }T}{T}})^2);`$ (1)
a diverging specific heat near a 2nd order phase transition would then remove fluctuations if matter is in global thermal equilibrium. The implications of such critical phenomena near second order phase transitions and critical points are discussed in detail in . It is found that the expansion of the systems slows the growth of the correlation lengths associated with the critical phenomena and the systems “slows out of equilibrium”, which affects the experimental signatures related to transverse momenta and temperatures.
First order phase transitions are contrarily expected to lead to large fluctuations due to droplet formation or more generally density or temperature fluctuations. These hot droplets will expand and hadronize in contrast to cold static quark matter droplets that may exist in cores of neutron stars . In case of a first order phase transition relativistic heavy ion collisions lead to interesting scenarios in which matter is compressed, heated and undergoes chiral restoration. If the subsequent expansion is sufficiently rapid, matter will pass the phase coexistence curve with little effect and supercool . This suggests the possible formation of “droplets” of supercooled chiral symmetric matter with relatively high baryon and energy densities in a background of low density broken symmetry matter. These droplets can persist until the system reachs the spinodal line and then return rapidly to the now-unique broken symmetry minimum. A large mismatch in density and energy density seems to be a robust prediction for a first order transition at large baryon densities. At high temperatures, which is more relevant for relativistic heavy-ion collisions (see Fig. (1)), the transition is probably at most weakly first order as discussed above.
Density fluctuations may appear both for a first order phase transition and for a smooth cross-over. If the transitions is first order, matter may supercool and subsequently create fluctuations in a number of quantities. Density fluctuations in the form of hot spots or droplets of dense matter with hadronic gas in between is a likely outcome (see Fig. (2)). Even if the transition is a smooth cross-over, the resulting soft equation of state has a small sound speed, $`c_s^2=P/ϵ`$. The equation of state $`P(ϵ)`$ has in both cases a flat region that may be hard to distinguish in a finite systems existing for a short time only. We do not know the early non-equilibrium stages of relativistic nuclear collisions and the resulting initial density fluctuations, hot spots, etc. If the system becomes thermalized at some stage, then a smaller $`c_s^2`$ is likely to allow for larger density fluctuations since the pressure difference is smaller. Furthermore, in the subsequent expansion the density fluctuations are not equilibrated as fast when $`c_s^2`$ is small because the pressure differences, that drive the differential expansion, are small. The dissipation of an initial density fluctuations can be estimated by a stability analysis . Linearizing the hydrodynamic equations in small fluctuations around the Bjorken scaling solution, an entropy fluctuation is typically damped by a factor (see Appendix A for details)
$`{\displaystyle \frac{\delta S_{final}}{\delta S_{initial}}}\left({\displaystyle \frac{\tau _0}{\tau _f}}\right)^{|Re[\lambda _\pm ]|}.`$ (2)
Here, a typical formation time is $`\tau _01`$ fm/c and freezeout time $`\tau _f8`$ fm/c as extracted from HBT studies . The eigenvalues $`\lambda _\pm `$ depend on the sound speed and the wave length of the rapidity fluctuations. As described in more detail in Appendix A, one of the eigenvalues are small and vanish for $`c_s=0`$. The resulting suppression of an initial density fluctuations during expansion is typically smaller than a factor 0.5. If density fluctuations are enhanced initially due to a softening of the equation of state due to smooth cross over, then they will largely be preserved later on. Yet, such fluctuations will be smaller than for a true first order transition forming supercooled droplets.
Let us assume that hadrons emerge from a collection of density fluctuations or droplets with a Boltzmann distribution with temperature $`T`$ and from a more or less continous background obeying approximate Bjorken scaling. The resulting particle distribution is
$$\frac{dN}{dyd^2p_t}\underset{i}{}f_ie^{m_t\mathrm{cosh}(y\eta _i)/T}+\mathrm{background}.$$
(3)
Here, $`y`$ is the particle rapidity and $`p_t`$ its transverse momentum, $`f_i`$ is the number of particles hadronizing from each droplet i, and
$$\eta _i=\frac{1}{2}\mathrm{log}\frac{t_i+z_i}{t_iz_i}=\frac{1}{2}\mathrm{log}\frac{1+v_i}{1v_i}$$
(4)
is the rapidity of droplet i. The size, number and separation between droplets or density fluctuations will depend on the violent initial conditions. Between droplets a relatively continuous background of hadrons is expected in coexistence. In (3) the droplet is assumed not to expand internally neither longitudinally nor transversely. If it does expand, the emerging hadrons will have a wider distribution of rapidities which will be harder to distinguish from the background.
When $`m_t/T1`$, we can approximate $`\mathrm{cosh}(y\eta _i)1+\frac{1}{2}(y\eta _i)^2`$ in Eq.(3). The Boltzmann factor determines the width of the droplet rapidity distribution as $`\sqrt{T/m_t}`$. The rapidity distribution will display fluctuations in rapidity event by event when the droplets are separated by rapidities larger than $`|\eta _i\eta _j|\stackrel{>}{}\sqrt{T/m_t}`$. If they are evenly distributed by smaller rapidity differences, the resulting rapidity distribution (3) will appear flat.
The droplets are separated in rapidity by $`|\eta _i\eta _j|\mathrm{\Delta }z/\tau _0`$, where $`\mathrm{\Delta }z`$ is the correlation length in the dense and hot mixed phase and $`\tau _0`$ is the invariant time after collision at which the droplets form. Assuming that $`\mathrm{\Delta }z1`$fm — a typical hadronic scale — and that the droplets form very early $`\tau _0\stackrel{<}{}1`$fm/c, we find that indeed $`|\eta _i\eta _j|\stackrel{>}{}\sqrt{T/m_t}`$ even for the light pions. If strong transverse flow is present in the source, the droplets may also move in a transverse direction. In that case the distribution in $`p_t`$ may be non-thermal and azimuthally asymmetric.
Even if the transition is not first order, fluctuations may still occur in the matter that undergoes a transition. The fluctuations may be in density, chiral symmetry , strangeness, or other quantities and show up in the associated particle multiplicities. The “anomalous” fluctuations depend not only on the type and order of the transition, but also on the speed by which the collision zone goes through the transition, the degree of equilibrium, the subsequent hadronization process, the amount of rescatterings between hadronization and freezeout, etc. It may be that any sign of the transition is smeared out and erased before freezeout. That no anomalous event-by-event fluctuations have been found at CERN within experimental accuracy indicate that no transition took place or that the signals were erased before freezeout. Whether they remain at RHIC is yet to be discovered and we shall provide some tools for the analysis in the following sections.
## 3 Multiplicity Fluctuations in Relativistic Heavy-Ion Collisions
In order to be able to extract new physics associated with fluctuations, it is necessary to understand the role of expected statistical fluctuations. Our aim here is to study the sources of these fluctuations in collisions. As we shall see, the current NA49 data (see Fig. (6)) can be essentially understood on the basis of straightforward statistical arguments. Expected sources of fluctuations include impact parameter fluctuations, fluctuations in the number of primary collisions and in the results of such collisions, fluctuations in the relative orientation during the collision of deformed nuclei , effects of rescattering of secondaries, and QCD color fluctuations. Since fluctuations in collisions are sensitive to the amount of rescattering of secondaries taking place, we discuss in detail two limiting cases, the participant or “wounded nucleon model” (WNM) , in which one assumes that particle production occurs in the individual participant nucleons and rescattering of secondaries is ignored, and the thermal limit in which scatterings bring the system into local thermal equilibrium.
Data at AGS, SPS and RHIC energies show that multiplicities are enhanced by $``$30% in central collisions between heavy ($`A200`$) nuclei as compared to the WNM prediction (see Fig. (3)). Whether rescatterings increase relative fluctuations through greater production of multiplicity, transverse momenta, etc., or decrease fluctuations by involving a greater number of degrees of freedom, is not immediately obvious . Rescatterings probably increase both the average multiplicity and its variance but whether the relative amount of fluctuations are increased is model dependent. It has even been found in relativistic heavy ion collisions that the multiplicity fluctuations increase in the first few rescattering but then decrease again as the thermal limit is approached. VENUS simulations showed that rescattering had negligible effects on transverse energy fluctuations.
We first review known multiplicities and fluctuations in the basic $`pp`$ collisions, go on to study nucleus-nucleus collisions, and finally show in a simple model how phase transitions are capable of producing very significant fluctuations in particle multiplicities.
### 3.1 Charged particle production in $`pp`$ and $`p\overline{p}`$ reactions
Participant models or WNMs are basically a superposition of NN collisions. Such models have been studied extensively at these energies within the last decades at several particle accelerators and we here give a brief compilation of relevant results.
The average number of charged particles produced in high energy $`pp`$ and ultrarelativistic $`p\overline{p}`$ collisions can be parametrized by
$`N_{ch}4.2+\mathrm{\hspace{0.17em}4.69}\left({\displaystyle \frac{s}{\mathrm{GeV}^2}}\right)^{0.155},`$ (5)
for cms energies $`\sqrt{s}\stackrel{>}{}2`$ GeV. At ultrarelativistic energies the charged particle production is very similar in $`pp`$, $`pn`$ and $`p\overline{p}`$ collisions and the parametrization of Eq. (5) applies in a wide range of cms energies 2 GeV$`\stackrel{<}{}s^{1/2}\stackrel{<}{}`$ 2 TeV as shown in Fig. (4). At SPS, RHIC and LHC energies, $`\sqrt{s}20`$, 200, 5000 GeV, we find $`N_{ch}7.3,20,60`$, respectively.
At high energies KNO scaling is a good approximation. KNO scaling implies that multiplicity distributions are invariant when scaled with the average multiplicity. Thus all moments scale like
$`N_{ch}^qc_qN_{ch}^q,`$ (6)
at high energies where $`c_q`$ are constants independent of collision energy. The fluctuations,
$`\omega _N(N^2N^2)/N,`$ (7)
therefore scale with average multiplicity, $`N`$, and therefore increase with collision energy as in Eq. (5). The fluctuations in the charged particle multiplicity can be parametrized rather accurately by
$`\omega _{N_{ch}}0.35{\displaystyle \frac{(N_{ch}1)^2}{N_{ch}}}`$ (8)
as shown in Fig. (5) for $`pp`$ and $`p\overline{p}`$ collisions in the same wide range of energies. At the very high energies breakdown of KNO scaling has been observed in the direction that the fluctuations are slightly larger. At SPS, RHIC and LHC energies we find $`\omega _{N_{ch}}2.0,6.2,20`$, respectively in $`pp`$ and $`p\overline{p}`$ collisions.
In nuclear (AA) collisions the number of participating nucleons $`N_p`$ grow with centrality and nuclear mass number $`A`$. Therefore the average charged particle multiplicity and variance grows with $`N_p`$, whereas the ratio and therefore the fluctuation $`\omega _{N_{ch}}`$ is independent of $`N_p`$, and equal to the fluctuations in pp collisions. (Other fluctuations such as impact parameter will be included below.) Higher moments of the multiplicity distributions are large in high energy pp and $`p\overline{p}`$ collisions due to KNO scaling but in nuclear collisions such higher moments are suppressed by factors of $`1/N_p`$ and are therefore less interesting than the second moment. This justifies our detailed analyzes of the variance (or rms width) of the fluctuations.
### 3.2 Fluctuations in the participant model
In the participant or wounded nucleon models nucleus-nucleus collisions at high energies are just a superposition of nucleon-nucleon (NN) interactions. In peripheral collisions there are only few NN collisions, the collision zone is small, rescatterings few and the WNM should therefore apply. For central nuclear collisions, however, multiple NN scatterings, energy degradation, rescatterings between produced particles and other effects complicate the particle production and do enhance the multiplicities by ca. 30% as seen in experiment (see Fig. (3)). Thermal models may better describe central collisions as will be investigated afterwards. Yet, the WNM provides a simple baseline to compare to, when going from peripheral towards central collisions.
Let us first calculate fluctuations in the participant model. Although the multiplicities are somewhat underestimated, the measured multiplicity and transverse energy in nuclear collisions at AGS and SPS energies are known to scale approximately linearly with the number of participants . In this picture
$`N={\displaystyle \underset{i}{\overset{N_p}{}}}n_i,`$ (9)
where $`N_p`$ is the number of participants and $`n_i`$ is the number of particles produced in the acceptance by participant $`i`$. In the absence of correlations between $`N_p`$ and $`n`$, the average multiplicity is $`N=N_pn`$. For example, NA49 measures charged particles in the rapidity region $`4<y<5.5`$ and finds $`N270`$ for central Pb+Pb collisions. Finite impact parameters $`(b\stackrel{<}{}3.5`$ fm) as well as surface diffuseness reduce the number of participants from the total number of nucleons $`2A`$ to $`N_p350`$ estimated from Glauber theory; thus $`n0.77`$. Squaring Eq. (9) and again assuming no correlations between different wounded nucleon emission, $`n_in_j=n_in_j`$ for $`ij`$, we find the multiplicity fluctuations (see Appendix B)
$`\omega _N=\omega _n+n\omega _{N_p},`$ (10)
where $`\omega _N`$, $`\omega _n`$ and $`\omega _{N_p}`$ are the multiplicity fluctuations in the total number of particles (within the acceptance), in each source, and in the number of sources respectively.
A major source of multiplicity fluctuations per participant, $`\omega _n`$, is the limited acceptance. While each participant produces $`\nu `$ charged particles, only a smaller fraction $`f=n/\nu `$ are accepted. Without carrying out a detailed analysis of the acceptance, one can make a simple statistical estimate assuming that the particles are accepted randomly, in which case $`n`$ is binomially distributed with $`\sigma (n)=\nu f(1f)`$ for fixed $`\nu `$. Including fluctuations in $`\nu `$ we obtain, similarly to Eq. (10),
$`\omega _n=1f+f\omega _\nu .`$ (11)
In nucleon-nucleon collisions at SPS energies, the charged particle multiplicity is $`7.3`$ and $`\omega _\nu 1.9`$ ; as the multiplicity should be divided between the two colliding nucleons, we obtain $`\nu 3.7`$ and thus $`f=n/\nu =0.21`$ for the NA49 acceptance. Consequently, we find from Eq. (11) that $`\omega _n1.2`$. The random acceptance assumption can be improved by correcting for known rapidity correlations in charged particle production in $`pp`$ collisions .
Multiplicities generally increase with centrality of the collision. We will use the term centrality as impact parameter $`b`$ in the collision. It is not a directly measurable quantity but is closely correlated to the transverse energy produced $`E_T`$, the measured energy in the zero degree calorimeter and the total particle multiplicity $`N`$ measured in some large rapidity interval. The latter is within the WNM approximately proportional to the number of participating nucleons
$$N_p(𝐛)=_{overlap}\left[\rho (𝐫+\frac{𝐛}{2})+\rho (𝐫\frac{𝐛}{2})\right]d^3r.$$
(12)
For sharp sphere nuclei the number of participants drops from $`N_p(b=0)=2A`$ in central collisions to $`N_p(b=2R)=0`$ in grazing collisions. For realistic nuclei with diffuse surface and with collision probabilities given by Glauber theory, the number of participants are 5-10% smaller in central collisions but slightly larger in peripheral collisions.
As a consequence of nuclear correlations, which strongly reduce density fluctuations in the colliding nuclei, the fluctuations $`\omega _{N_p(b)}`$ in $`N_p`$ are very small for fixed impact parameter $`b`$ . Almost all nucleons in the nuclear overlap volume collide and participate. \[By contrast, the fluctuations in the number of binary collisions are non-negligible.\] Cross section fluctuations play a small role in the WNM . Fluctuations in the number of participants can arise when the target nucleus is deformed, since the orientations of the deformation axes vary from event to event . The fluctuations, $`\omega _{N_p}`$, in the number of participants are dominated by the varying impact parameters selected by the experiment. In the NA49 experiment, for example, the zero degree calorimeter selects the 5% most central collisions, corresponding to impact parameters smaller than a centrality cut on impact parameter, $`b_c3.5`$ fm. We have
$`\omega _{N_p}N_p={\displaystyle \frac{1}{\pi b_c^2}}{\displaystyle _0^{b_c}}d^2bN_p(b)^2N_p^2,`$ (13)
where $`N_p=(1/\pi b_c^2)_0^{b_c}d^2bN_p(b)`$. The number of participants for a given centrality, calculated in , can be approximated by $`N_p(b)N_p(0)(1b/2R)`$ for $`0b\stackrel{<}{}3.5`$ fm; thus
$`\omega _{N_p}={\displaystyle \frac{N_p(0)}{18}}\left({\displaystyle \frac{b_c}{2R}}\right)^2.`$ (14)
For NA49 Pb+Pb collisions with $`N_p(0)400`$ and $`(b_c/2R)^25\%`$ we find $`\omega _{N_p}1.1`$. Impact parameter fluctuations are thus important even for the centrality trigger of NA49. Varying the centrality cut or $`b_c`$ to control such impact parameter fluctuations (14) should enable one to extract better any more interesting intrinsic fluctuations. Recent WA98 analyzes confirm that fluctuations in photons and pions grow approximately linearly with the centrality cut $`(b_c/2R)^2`$ as predicted by Eq. (14). The impact parameter fluctuations associated with the range of the centrality cut, such at total transverse energy or multiplicity, can therefore be removed. However, fluctuations in impact parameter may still remain for a given centrality. The Gaussian multiplicity distribution found in central collisions changes for minimum bias to a plateau-like distribution .
Calculating $`\omega _N`$ for the NA49 parameters, we find from Eq. (10), $`\omega _N1.2+(0.77)(1.1)=2.0`$, in good agreement with experiment, which measures a multiplicity distribution $`\mathrm{exp}[(NN)^2/2N\omega _N^{exp}]`$, where $`\omega _N^{exp}`$ is $`2.01`$ (see Fig. (6)).
### 3.3 Fluctuations in the thermal model
Let us now consider, in the opposite limit of strong rescattering, fluctuations in thermal models. In a gas in equilibrium, the mean number of particles per bosonic mode $`n_a`$ is given by
$`n_a=\left(\mathrm{exp}(E_a/T)1\right)^1,`$ (15)
with fluctuations
$`\omega _{n_a}=1+n_a.`$ (16)
The total fluctuation in the multiplicity, $`N=_an_a`$, is
$`\omega _N^{BE}=1+{\displaystyle \underset{a}{}}n_a^2/{\displaystyle \underset{a}{}}n_a.`$ (17)
If the modes are taken to be momentum states, bosons/fermions have thermal fluctuations, $`\omega _N=1\pm n_p^2/n_p`$ where $`n_p=(\mathrm{exp}(ϵ_p/T)1)^1`$ is the boson/fermion distribution function, which are slightly larger/smaller than those of Poisson statistics for a Boltzmann distribution, $`\omega _N=1`$. The resulting fluctuations are $`\omega _N^{BE}=\zeta (2)/\zeta (3)=1.37`$ for massless bosons as, e.g., gluons. Massive bosons have smaller fluctuations with, for example, $`\omega _\pi =1.11`$ and $`\omega _\rho =1.01`$ when $`T=m_\pi `$. Massless fermions, e.g. quarks, have $`\omega _F=2\zeta (2)/3\zeta (3)0.91`$ independent of temperature.
Resonances are implicitly included in the WNM fluctuations. In the thermal limit resonances are found to increase total multiplicity fluctuations but decrease, e.g., net charged particle fluctuations . In high energy nuclear collisions, resonance decays such as $`\rho 2\pi `$, $`\omega 3\pi `$, etc., lead to half or more of the pion multiplicity. Only a small fraction $`r2030`$% produce two charged particles in a thermal hadron gas or in RQMD (see also ). Not all of the decay particles from the same resonance always fall into the NA49 acceptance, $`4<y<5.5`$, and the fraction of pairs will be smaller; we estimate $`r0.1`$. Including such resonance fluctuations in the BE fluctuations gives, similarly to Eq. (10),
$`\omega _N^{BE+R}=r{\displaystyle \frac{1r}{1+r}}+(1+r)\omega _N^{BE}.`$ (18)
With $`r0.1`$ we obtain $`\omega _N^{BE+R}1.3`$. In the estimated effect of resonances is about twice ours: $`\omega _N1.5`$, not including impact parameter fluctuations.
Fluctuations in the effective collision volume add a further term $`N\sigma (V)/V^2`$ to $`\omega _N^{BE+R}`$. Assuming that the volume scales with the number of participants, $`\omega _V/V\omega _{N_p}/N_p`$, we find from Eq. (10) that $`\omega _N=\omega _N^{BE+R}+n\omega _{N_p}2.1`$, again consistent with the NA49 data. Because of the similarity between the magnitudes of the thermal and WNM multiplicity fluctuations, the present measurements cannot distinguish between these two limiting pictures.
### 3.4 Centrality dependence and degree of thermalization
It is very unfortunate that the WNM and thermal models predict the same multiplicity fluctuations in the NA49 acceptance - and that they agree with the experiment. If the numbers from the two models had been different and the experimental number in between these two, then one would have had quantified the degree of thermalization in relativistic heavy ion collisions.
The similarity of the fluctuation in the thermal and WNM is, however, a coincidence at SPS energies. As seen from Fig. (5) the fluctuations in $`pp`$ collisions increase with collision energy and just happen to cross the thermal fluctuations, $`\omega _{thermal}2.2`$, at SPS energies.<sup>2</sup><sup>2</sup>2As discussed below the thermal fluctuations in positive or negative particles are $`\omega _\pm 1.1`$ in a thermal hadron gas. The fluctuation in total charge is twice that due to overall charge neutrality which relates the number of positive to negative particles.
At RHIC or LHC energies the situation will be much clearer. Here the charged particle fluctuations in $`pp`$ collisions are much larger as seen in Fig. (5), namely $`\omega _{N_{ch}}^{pp}=6.5`$, 20 at RHIC and LHC energies respectively. The thermal fluctuations remain as $`\omega _{thermal}2.2`$. Therefore a dramatic reduction in event-by-event fluctuations are expected at higher energies at the nuclear collisions become more central as shown in Fig. (7).
This can be exploited to define a “Degree of thermalization” as the measured fluctuations at a given centrality relative to those in the thermal and $`pp`$ limits
$`\mathrm{Degree}\mathrm{of}\mathrm{thermalization}{\displaystyle \frac{\omega _N^{WNM}\omega _N^{exp}}{\omega _N^{WNM}\omega _N^{thermal}}},`$ (19)
which ranges from unity in the thermal limit to zero in the WNM. Whereas both $`\omega _N^{WNM}`$ and $`\omega _N^{exp}`$ may depend on the acceptance the degree of thermalization Eq. (19) should not. Contributions from volume or impact parameter fluctuations may, however, be centrality dependent and should therefore be subtracted. Alternatively, the fluctuations in a ratio, e.g. $`N_{}/N_+`$, should be taken for limited acceptances.
At RHIC and LHC it should be straight forward to measure the degree of thermalization as function of centrality. This is interesting on its own and a necessary requirement for studies of anomalous fluctuations from a phase transition.
### 3.5 Enhanced fluctuations in first order phase transitions
First order phase transitions can lead to rather large fluctuations in physical quantities. Thus, detection of enhanced fluctuations, beyond the elementary statistical ones considered to this point, could signal the presence of such a transition. For example, before it became clear that the chiral symmetry restoring phase transition in hot QCD is not a strong first order phase transition, it was suggested that matter undergoing a transition from chirally symmetric to broken chiral symmetry could, when expanding, supercool and form droplets, resulting in large multiplicity versus rapidity fluctuations . Let us imagine that $`N_D`$ droplets fall into the acceptance, each producing $`n`$ particles, i.e., $`N=N_Dn`$. The corresponding multiplicity fluctuation is (see Appendix B)
$`\omega _N=\omega _n+n\omega _{N_D}.`$ (20)
As in Eq. (11), we expect $`\omega _n1`$. However, unlike the case of participant fluctuations, the second term in (20) can lead to huge multiplicity fluctuations when only a few droplets fall into the acceptance; in such a case, $`n`$ is large and $`\omega _{N_D}`$ of order unity. The fluctuations from droplets depends on the total number of droplets, the spread in rapidity of particles from a droplet, $`\delta y\sqrt{T/m_t}`$, as well as the experimental acceptance in rapidity, $`\mathrm{\Delta }y`$. When $`\delta y\mathrm{\Delta }y`$ and the droplets are binomially distributed in rapidity, $`\omega _{N_D}1\mathrm{\Delta }y/y_{\mathrm{tot}}`$, which can be a significant fraction of unity.
In the extreme case where none or only one droplet falls into the acceptance with equal probability, we have $`\omega _{N_D}=1/2`$ and $`n=2N`$. The resulting fluctuation is $`\omega _NN`$, which is more than two orders of magnitude larger than the expected value of order unity as currently measured in NA49. It should be said immediately that a much smaller enhancement is realistic as the transition probably is at most weakly first order and many effects will smear the signal. Yet, this simple example clearly demonstrates the importance of event-by-event fluctuations accompanying phase transitions, and illustrates how monitoring such fluctuations versus centrality becomes a promising signal, in the upcoming RHIC experiments, for the onset of a transition. It is the hope and expectation that the higher RHIC energies probe deeper into the QGP phase by creating higher temperatures and energy densities whereby larger regions of QGP are produced. The larger event multiplicities should make it possible to improve on statistics and thereby also the ability to detect anomalous fluctuations. The potential for large fluctuations (orders of magnitude) from a transition makes it worth looking for at RHIC considering the relative simplicity and accuracy (percents) of multiplicity measurements.
Let us subsequently consider a less extreme model in which a transition leads to enhanced fluctuations of some kind. Assume that the total multiplicity within the acceptance arises from a normal hadronic background component ($`N_{HG}`$) and from a second component ($`N_{QGP}`$) that has undergone a transition:
$`N=N_{HG}+N_{QGP}.`$ (21)
Its average is $`N=N_{HG}+N_{QGP}`$. Assuming that the multiplicity of each of these components is statistically independent, the multiplicity fluctuation becomes
$`\omega _N=\omega _{HG}+(\omega _{QGP}\omega _{HG}){\displaystyle \frac{N_{QGP}}{N}}.`$ (22)
Here, $`\omega _{HG}`$ is the standard fluctuation in hadronic matter $`\omega _{HG}1`$. The fluctuations due to the component that had experienced a phase transition, $`\omega _{QGP}`$, depend on the type and order of the transition, the speed with which the collision zone goes through the transition, the degree of equilibrium, the subsequent hadronization process, the number of rescatterings between hadronization and freezeout, etc. If thermal and chemical equilibration eliminate all signs of the transition, then $`\omega _{QGP}\omega _{HG}`$.
The amount of QM and thus $`N_{QM}`$ depends on centrality, energy and nuclear masses in the collision. For a given centrality the densities vary from zero at the periphery of the collision zone to a maximum value at the center. Furthermore, the more central the collision the higher energy densities are created. The transverse energy, $`E_T`$, the total multiplicity and/or the energy in the zero-degree calorimeter, $`E_{ZDC}`$, have been found to be good measures of the centrality of the collision at SPS energies. Therefore, it would be very interesting to study fluctuations vs. centrality which are proportional to energy density. By varying the binning size for centrality one can also remove impact parameter fluctuations as discussed above.
If the energy density in the center of the collision zone exceeds the critical energy density for forming QM at a certain centrality, $`E_1`$, then a mixed phase of QM and HM is formed. At a higher energy density, where the critical energy density plus the latent heat for the transition is exceeded, which we shall assume occur at a centrality $`E_2`$, then a pure QM phase is produced in the center. These quantities will depend on the amount of stopping at a given centrality, the geometry, $`T_c`$, etc. In the mixed phase $`E_1E_T<E_2`$, the relative amount of QM, $`N_{QM}/N`$, is proportional to both the volume of the mixed phase. and the fraction of the volume that is in the QM phase. The latter varies in the volume such that it vanishes at HM/QM boundary.
In Fig. (8) a schematic plot of the fluctuations of Eq. (22) is shown as function of centrality for various $`\omega _{QM}`$. Up to centrality $`E_1`$ the fluctuations are unchanged. Above the central overlap zone undergoes the transition to the QM/HM mixed phase and fluctuations start to grow when $`\omega _{QM}>\omega _{HM}`$. At the higher centrality, $`E_2`$ the central overlap zone is in the pure QM phase but the maximum fluctuations $`\omega _{QM}`$ are not reached because the surface regions of the collision zone is still in the HM phase. On the other hand, if the hadronization of the QM state is smooth and does not lead to enhanced fluctuations (i.e. if $`\omega _{QGP}=\omega _{HG}`$), it cannot be observed in such a study.
The multiplicity fluctuations can be studied for any kind of particles, total or ratios. Total multiplicities describe total multiplicities whereas, e.g. the ratio $`\pi ^0/(\pi ^++\pi ^{})`$ can reveal fluctuations in chiral symmetry. The onset and magnitude of such fluctuations would reveal the symmetry and other properties of the new phase.
## 4 Correlations between total and net charge, baryon number or strangeness
By a combined analysis of fluctuations in, e.g., positive, negative, total and net charge as well as ratios, the intrinsic and other fluctuations as well as correlations can be extracted and exploited to reveal interesting physics as will be demonstrated in the following.
### 4.1 General analysis of fluctuations and correlations
Multiplicity fluctuations between various kinds of particles can be strongly correlated. As a first example, consider the multiplicities of positive and negative pions, $`N_+`$ and $`N_{}`$, in a rapidity interval $`\mathrm{\Delta }y`$ for any relativistic heavy-ion experiment. Similar analyzes can be performed for any two kinds of distinguishable particles.
The net positive charge from the protons in the colliding nuclei is much smaller than the total charge produced in an ultrarelativistic heavy-ion collision. For example, $`N_+`$ exceeds $`N_{}`$ by only $`15`$% at in Pb$`+`$Pb collisions at SPS energies. The fluctuations in the number of positive and negative (or neutral) pions are also very similar, $`\omega _{N_+}\omega _N_{}`$. Charged particle fluctuations have been estimated in thermal as well as participant nucleon models including effects of resonances, acceptance, and impact parameter fluctuations. By varying the acceptance and centrality, the degree of thermalization can actually be determined empirically. Detailed analysis indicates that the fluctuations in central Pb+Pb collisions at the SPS are thermal whereas peripheral collisions are a superposition of pp fluctuations .
The fluctuations in the total ($`N_{ch}=N_++N_{}`$) and net ($`Q=N_+N_{}`$) charge are defined as
$`{\displaystyle \frac{(N_+\pm N_{})^2N_+\pm N_{}^2}{N_++N_{}}}=`$
$`{\displaystyle \frac{N_+}{N_{ch}}}\omega _{N_+}+{\displaystyle \frac{N_{}}{N_{ch}}}\omega _N_{}\pm C,`$ (23)
where the correlation is given by
$`C={\displaystyle \frac{N_+N_{}N_+N_{}}{N_{ch}/2}}.`$ (24)
Fluctuations in positive, negative, total and net charge can be combined to yield both the intrinsic fluctuations in the numbers of $`N_\pm `$ and the correlations in their production as well as a consistency check. These quantities can change as a consequence of thermalization and a possible phase transition.
In practice, $`\omega _{N_+}\omega _N_{}`$, so that the fluctuation in total charge simplifies to
$`\omega _{N_{ch}}`$ $``$ $`{\displaystyle \frac{N_{ch}^2N_{ch}^2}{N_{ch}}}=\omega _{N_+}+C,`$ (25)
and that for the net charge becomes
$`\omega _Q`$ $``$ $`{\displaystyle \frac{Q^2Q^2}{N_{ch}}}=\omega _{N_+}C.`$ (26)
The fluctuation in net charge can be related to the fluctuation in the ratio of positive to negative particles
$`\omega _QN_+/N_{}N_{ch}\omega _{N_{}/N_+}/4,`$ (27)
plus volume (or impact parameter) fluctuations . The virtue of this expression is that volume fluctuations can in principle be extracted empirically. Alternatively one can vary the centrality bin size or the acceptance. Furthermore, the volume fluctuations for net and total charge are proportional to the net ($`N_+N_{}`$) and total ($`N_++N_{}`$) charge respectively with the same prefactor. In the following we shall assume that such “trivial” volume fluctuations have been removed.
The analysis has so far been general and Eqs. (23-26) apply to any kind of distinguishable particles, e.g. positive and negative particles, pions, kaons, baryons, etc. - irrespective of what phase the system may be in, or whether it is thermal or not. In the following, we shall consider thermal equilibrium, which seems to apply to central collisions between relativistic nuclei, in order to reveal possible effects on fluctuations of the presence of a quark-gluon plasma.
### 4.2 Charge fluctuations in a thermal hadron gas
In a thermal hadron gas (HG) as created in relativistic in nuclear collisions, pions can be produced either directly or through the decay of heavier resonances, $`\rho ,\omega ,\mathrm{}`$. The resulting fluctuation in the measured number of pions is
$`\omega _{N_+}=\omega _N_{}=f_\pi \omega _\pi +f_\rho \omega _\rho +f_\omega \omega _\omega +\mathrm{}.,`$ (28)
where $`f_r`$ is the fraction of measured pions produced from the decay of resonance $`r`$, and $`_rf_r=1`$. These mechanisms are assumed to be independent, which is valid in a thermal system.
The heavier resonances such as $`\rho ^0,\omega ,\mathrm{}`$ decay into pairs of $`\pi ^+\pi ^{}`$ and thus lead to a correlation
$`C^{HG}={\displaystyle \frac{1}{3}}f_\rho +f_\omega +\mathrm{}..`$ (29)
Resonances reduce the fluctuations in net charge in a HG in chemical equilibrium at temperature $`T=170`$ MeV and baryon chemical potential $`\mu _b=270`$ MeV and strangeness chemical potential $`\mu _s=74`$ MeV to $`\omega _Q=0.70`$ . In the value $`\omega _Q=0.70`$ is found.
In addition, overall charge conservation reduces fluctuations in net charge when the acceptance is large and thus increases correlations as will be discussed below.
### 4.3 Charge fluctuations in a quark-gluon plasma
A phase transition to the QGP can alter both fluctuations and correlations in the production of charged pions. To the extent that these effects are not eliminated by subsequent thermalization of the HG, they may remain as observable remnants of the QGP phase. As shown in Refs. , net charge fluctuations in a plasma of u, d quarks and gluons are reduced partly due to the intrinsically smaller quark charge and partly due to correlations from gluons
$`\omega _Q={\displaystyle \frac{N_q}{N_{ch}}}\omega _F{\displaystyle \frac{1}{N_f}}{\displaystyle \underset{f=u,d,\mathrm{}}{\overset{N_f}{}}}q_f^2,`$ (30)
where $`N_f`$ is the number of quark flavors, $`q_f`$ their charges, and $`N_q`$ the number of quarks. The total number of charged particles (but not the net charge) can be altered by the ultimate hadronization of the QGP. Assuming a pion gas as the final state, this effect can be estimated by equating the entropy of all pions to the entropy of the quarks and gluons. Since 2/3 of all pions are charged and since the entropy per fermion is 7/6 times the entropy per boson in a two-flavor QGP
$`N_{ch}{\displaystyle \frac{2}{3}}(N_g+{\displaystyle \frac{7}{6}}N_q),`$ (31)
where the number of gluons is $`N_g=(16/9N_f)N_q`$. Inserting this result in (30), we see that the resulting fluctuations are $`\omega _Q=0.18`$ in a two-flavor QGP (and $`\omega _Q=0.12`$ for three flavors). As pointed out in , lattice results give $`\omega _Q0.25`$. <sup>3</sup><sup>3</sup>3It is amusing to note that this number gives a very bad (i.e., negative) estimate for $`N_g/N_q`$ in Eq. (31). However, according to a substantial fraction of the pions are decay products from the HG, and the entropy of the HG exceed that of a pion gas by a factor $`1.751.8`$. As described in the net charge fluctuations should be increased by this factor in the QGP, i.e. $`\omega _Q0.33`$ in a two-flavor QGP, whereas it remains similar in the HG, $`\omega _Q0.6`$.
The above models are all grand canonical ones, i.e. no net charge conservation, as opposed to microcanonical models that now will be discussed. If the high density phase is initially dominated by gluons with quarks produced only at a later stage of the expansion by gluon fusion, the production of positively and negatively charged quarks will be strongly correlated on sufficiently small rapidity scales. An increased entropy density in the collisions volume will lead to enhanced multiplicity as compared to a standard hadronic scenario if total entropy is conserved. The associated particle production must conserve net charge on large rapidity scales ($`\mathrm{\Delta }y\stackrel{>}{}1`$) due to causality because fields cannot communicate over large distances and therefore must conserve charge within the “event horizon”. Therefore the net charge, $`N_{ch}`$, is approximately conserved whereas the total charge, $`Q`$, increase by an amount proportional to the additional entropy produced. If the entropy density increases from $`s_{HG}`$ to $`s_{QGP}`$ going from a HG to QGP without additional net charge production, fluctuations in net charge will be reduced correspondingly,
$`\omega _Q^{QGP}{\displaystyle \frac{s_{HG}}{s_{QGP}}}\omega _Q^{HG}.`$ (32)
The resulting fluctuation in net charge is necessarily smaller than that from thermal quark production as given by Eq. (30). A similar phenomenon occurs in string models where particle production by string breaking and $`q\overline{q}`$ pair production results in flavor and charge correlations on a small rapidity scale .
If droplets or density fluctuations appear, they are expected not to produce additional net charge. Consequently, the net charge fluctuations should still vanish $`\omega _Q0`$ whereas $`\omega _{ch}2\omega _{N_+}2\omega _{QGP}`$.
The strangeness fluctuation in kaons $`K^\pm `$ might seem less interesting at first sight since strangeness is not suppressed in the QGP: The strangeness per kaon is unity, and the total number of kaons is equal to the number of strange quarks. However, if strange quarks are produced at a late stage in the expansion of a fluid initially dominated by gluons, the net strangeness will again be greatly reduced on sufficiently small rapidity scale. Consequently, fluctuations in net/total strangeness would be reduced/enhanced.
The baryon number fluctuations have been estimated in a thermal model in a grand canonical model. It is, however, not known how possible variations in baryon stopping event-by-event and subsequent diffusion and annihilation of the baryons and antibaryons in the hadronic phase affect these results. If only charged particles are detected, but not $`K^0`$, $`\overline{K}^0`$, neutrons and antineutrons, the fluctuations have smaller correlations as compared to the total and net strangeness or baryon number.
### 4.4 Total charge conservation
Total charge conservation is important when the acceptance $`\mathrm{\Delta }y`$ is a non-negligible fraction of the total rapidity. It reduces the fluctuations in the net charge as calculated within the canonical ensemble, Eqs. (28-31). If the total positive charge (which is exactly equal to the total negative charge plus the incoming nuclear charges) is randomly distributed, the resulting fluctuations are smaller than the intrinsic ones by a factor $`(1f_{acc})`$, where
$`f_{acc}=(N_{ch}^{tot})^1{\displaystyle _{\mathrm{\Delta }y}}{\displaystyle \frac{dN_{ch}}{dy}}𝑑y`$ (33)
is the acceptance fraction or the probability that a charged particle falls into the acceptance $`\mathrm{\Delta }y`$ assuming full $`𝐩_t`$ coverage. Since charged particle rapidity distributions are peaked near midrapidity, charge conservation effectively kills fluctuations in the net charge even when $`\mathrm{\Delta }y`$ is substantially smaller than the laboratory rapidity, $`y_{lab}6`$ (11) at SPS (RHIC) energies. Total charge conservation also has the effect of increasing $`\omega _{ch}`$ towards $`2\omega _{N_+}`$ according to Eq. (25). Similar effects can be seen in photon fluctuations when photons are produced in pairs through $`\rho ^02\gamma `$. In the WA98 experiment, $`\omega _\gamma 2`$ is found after the elimination of volume fluctuations .
On the other hand, if the acceptance $`\mathrm{\Delta }y`$ is too small, particles can diffuse in and out of the acceptance during hadronization and freezeout . Furthermore, correlations due to resonance production will disappear when the average separation in rapidity between decay products exceeds the acceptance. Each of these effects tends to increase all fluctuations towards Poisson statistics when $`\mathrm{\Delta }y\stackrel{<}{}\delta y`$, where $`\delta y`$ denotes the rapidity interval that particles diffuse during hadronization, freezeout and decay. We find approximately
$`\omega _Q^{exp}\left({\displaystyle \frac{\mathrm{\Delta }y}{\mathrm{\Delta }y+2\delta y}}\omega _Q+{\displaystyle \frac{2\delta y}{\mathrm{\Delta }y+2\delta y}}\right)(1f_{acc}),`$ (34)
where $`\omega _Q`$ is the canonical thermal fluctuation of Eqs. (29,30) and $`\omega _Q^{exp}`$ is the fluctuation corrected for both $`\delta y`$ and total charge conservation.
The resulting fluctuations in total and net charge are shown in Fig. (9) assuming $`\omega _{N_+}=\omega _\pi 1.1`$ and $`\delta y=0.5`$. As mentioned above, $`f_{acc}`$ and $`\mathrm{\Delta }y`$ are related by the measured charge particle rapidity distributions . The total charge fluctuations in a HG ($`C=0.4`$) from Eq. (31) agree well with NA49 data after subtraction of residual impact parameter fluctuations. Data on charge particle ratios, which do not contain impact parameter fluctuations, will be able to test the net charge fluctuations of Eq. (34) to higher accuracy. Predictions from UrQMD are also shown for comparison . The sensitivity to diffusion is small as seen in Fig. (9) where for the fluctuations are also shown for $`\delta y=0.8`$ as recently used in . The curves in Fig. (9) apply to RHIC energies as well after scaling $`\delta y`$ with $`\mathrm{\Delta }y`$.
## 5 Fluctuations in particle ratios
By taking ratios of particles, e.g. $`K/\pi `$, $`\pi ^+/\pi ^{}`$, $`\pi ^0/\pi ^\pm `$, …, one conveniently removes volume and impact parameter fluctuations to first approximation. Simply increasing/decreasing the volume or centrality, the average number of particles of both species scales up/down by the same amount and thus cancel in the ratio.
### 5.1 $`\pi ^+/\pi ^{}`$ ratio and entropy production
Most particles produced in relativistic nuclear collisions are pions and they therefore constitute most of the number of positive and negatively charged particles. The fluctuations in the $`\pi ^+/\pi ^{}`$ ratio and thus the ratio of positive and negative particles are intimately related to the fluctuations in net charge
$`\omega _{N_{}/N_+}={\displaystyle \frac{4}{N_{ch}}}N_+/N_{}\omega _Q+\omega _{ipf},`$ (35)
where $`\omega _{ipf}`$ is the impact parameter or volume fluctuations and $`\omega _Q`$ are the net charge fluctuations as given by Eq. (27).
The $`\pi ^+/\pi ^{}`$ ratio has been studied in detail in . Resonances such as $`\rho ,\omega ,\mathrm{}`$ decaying into two or three pions correlate the $`\pi ^+`$ and $`\pi ^{}`$ production as for positively and negatively charged particles discussed above. Consequently, the fluctuation in the $`\pi ^+/\pi ^{}`$ ratio is reduced by $`30`$% in agreement with NA49 data .
### 5.2 $`K/\pi `$ ratio and strangeness enhancement
To second order in the fluctuations of the numbers of K and $`\pi `$, we have
$`K/\pi ={\displaystyle \frac{K}{\pi }}\left(1+{\displaystyle \frac{\omega _\pi }{\pi }}{\displaystyle \frac{K\pi K\pi }{K\pi }}\right).`$ (36)
The corresponding fluctuations in $`K/\pi `$ are given by
$`D^2{\displaystyle \frac{\omega _{K/\pi }}{K/\pi }}={\displaystyle \frac{\omega _K}{K}}+{\displaystyle \frac{\omega _\pi }{\pi }}2{\displaystyle \frac{K\pi K\pi }{K\pi }}.`$ (37)
The fluctuations in the kaon to pion ratio is dominated by the fluctuations in the number of kaons alone. The third term in Eq. (37) includes correlations between the number of pions and kaons. It contains a negative part from volume fluctuations, which removes the volume fluctuations in $`\omega _K`$ and $`\omega _\pi `$ since such fluctuations cancel in any ratio. In the NA49 data shown in Fig. (10) the average ratio of charged kaons to charged pions is $`K/\pi =0.18`$ and $`\pi 200`$. Excluding volume fluctuations, we take $`\omega _K\omega _\pi 1.21.3`$ as discussed above. The first two terms in Eq. (37) then yield $`D0.200.21`$ in good agreement with preliminary measurements $`D=0.23`$ . Thus at this stage the data gives no evidence for correlated production of K and $`\pi `$, as described by the final term in Eq. (37), besides volume fluctuations. The similar fluctuations in mixed event analyses $`D_{mixed}=0.208`$ confirm this conclusion.
Strangeness enhancement has been observed in relativistic nuclear collisions at the SPS. For example, the number of kaons and therefore also $`K/\pi `$ is increased by a factor of 2-3 in central Pb+Pb collisions. It would be interesting to study the fluctuations in strangeness as well. By varying the acceptance one might be able to gauge the degree of thermalization as discussed above. The fluctuations in the $`K/\pi `$ ratio as function of centrality would in that case reveal whether strangeness enhancement is associated with thermalization or other mechanisms lie behind. In a plasma of deconfined quarks strangeness is increased rapidly by $`ggs\overline{s}`$ and $`q\overline{q}s\overline{s}`$ processes and lead to enhancement of total strangeness $`s+\overline{s}`$ whereas the net strangeness $`s\overline{s}`$ remains zero. The fluctuations in net and total strangeness will qualitatively behave like net and total charge, however, with unit strangeness quantum numbers as compared to the fractional charges.
### 5.3 $`\pi ^0/\pi ^\pm `$ ratio and chiral symmetry restoration.
Fluctuations in neutral relative to charged pions would be a characteristic signal of chiral symmetry restoration in heavy ion collisions. If, during expansion and cooling, domains of chiral condensates gets “disoriented” (DCC) , anomalous fluctuations in $`\pi ^0/\pi ^\pm `$ ratios could results if DCC domains are large. For single DCC domain the probability distribution of ratios $`d=\pi ^0/(\pi ^0+\pi ^++\pi ^{})`$ is $`P(d)=1/\sqrt{2d}`$ with mean $`d=1/3`$ and fluctuation $`\omega _d=4/15`$, i.e. much larger than ordinary fluctuations in such ratios (see Eq. (37)) which decrease inversely with the number of pions.
Neutral pions are much harder to measure than charged pions but with respect to fluctuations, it suffices to measure the charged pions only. The anomalous fluctuations in $`\pi ^0`$ due to a DCC are anti-correlated to $`\pi ^\pm `$, i.e. they are of same magnitude but opposite sign. A DCC can equally well be searched for in total charge fluctuations as in the $`\pi ^0/\pi ^\pm `$ ratio, except for the troublesome impact parameter fluctuations.
### 5.4 $`J/\mathrm{\Psi }`$ multiplicity correlations and absorption mechanisms
$`J/\mathrm{\Psi }`$ suppression has been found in relativistic nuclear collisions and it is yet unclear how much is due to absorption on participant nucleons and produced particles (comovers). Whether “anomalous suppression” is present in the data is one of the most discussed signals from a hot and dense phase at early times . It was originally suggested that the formation of a quark-gluon plasma would destroy the $`c\overline{c}`$ bound states .
In relativistic heavy ion collisions very few $`J/\mathrm{\Psi }`$’s are produced in each collision. Of these only 6.9% branch into dimuons that can be detected and so the chance to detect two dimuon pairs in the same event is very small. Therefore, it will be correspondingly difficult to measure fluctuations and other higher moments of the number of $`J/\mathrm{\Psi }`$.
Another more promising observable is the correlation between the multiplicities in, e.g., a rapidity interval $`\mathrm{\Delta }y`$ of a charmonium state $`\psi =J/\mathrm{\Psi },\psi ^{},..`$ ($`N_\psi `$) and all particles ($`N`$) . The correlator $`NN_\psi NN_\psi `$ also enters the in the ratio $`\psi /N`$ (see Eq. (36)). The correlator has as good statistics as the total number of $`\psi `$ and it may contain some very interesting anti-correlations, namely that $`\psi `$ absorption grows with multiplicity $`N`$. The physics behind can be comover absorption, which grows with comover density, or formation of quark-gluon plasma, which may lead to both anomalous suppression of $`\psi `$ and large multiplicity in $`\mathrm{\Delta }y`$. Contrarily, direct Glauber absorption should not depend on the multiplicity of produced particles $`N`$ since it is caused by collisions with participating nucleons.
To quantify this anti-correlation we model the absorption/destruction of $`\psi `$’s by simple Glauber theory
$`{\displaystyle \frac{N_\psi }{N_\psi ^0}}=e^{\sigma _{c\psi }\rho _cl}e^{\gamma N/N},`$ (38)
where $`N_\psi ^0`$ is the number of $`J/\mathrm{\Psi }`$’s before comover or anomalous absorption sets in but after direct Glauber absorption on participant nucleons. In Glauber theory the exponent is the absorption cross section times the absorber density and path length traversed in matter. The density and therefore also the exponent is proportional to the multiplicity $`N`$ with coefficient
$`\gamma ={\displaystyle \frac{d\mathrm{log}N_\psi }{d\mathrm{log}N}}_{|N=N}.`$ (39)
In a simple comover absorption model for a system with longitudinal Bjorken scaling, it can be calculated approximately
$`\gamma {\displaystyle \underset{c}{}}{\displaystyle \frac{dN_c}{dy}}{\displaystyle \frac{v_{c\psi }\sigma _{c\psi }}{4\pi R^2}}\mathrm{log}R/\tau _0,`$ (40)
where $`dN_c/dy`$, $`\sigma _{c\psi }`$, $`v_{c\psi }`$ and $`\tau _0`$ are the comover rapidity density, absorption cross section, relative velocity and formation time respectively.
On average comover or anomalous absorption is responsible for a suppression factor $`e^\gamma `$. It is difficult to determine because only the total $`\psi `$ suppression including direct Glauber absorption on participants is measured.
The anti-correlation is straight forward to calculate when the fluctuations in the exponent are small (i.e. $`\gamma \sqrt{\omega _N/N}1`$). It is
$`NN_\psi NN_\psi =\gamma \omega _NN_\psi .`$ (41)
It is negative and proportional to the amount of comover and anomalous absorption and obviously vanishes when the absorption is independent of multiplicity ($`\gamma =0`$). The anti-correlation can be accurately determined as the current accuracy in determining $`N_\psi `$ is a few percent (NA50 minimum bias ) in each $`E_T`$ bin.
The anti-correlations in Eq. (41) may seem independent of the rapidity interval. However, if it is less than the typical relative rapidities between comovers and the $`\psi `$, the correlations disappear. Preferrably, the rapidity interval should be of the order of the typical rapidity fluctuations due to density fluctuations.
The anticorrelations of Eq. (41) quantify the amount of comover or anomalous absorption and can therefore be exploited to distinguish between these and direct Glauber absorption mechanisms. In that respect it is similar to the elliptic flow parameter for $`\psi `$ for the comover absorption part but differs for the anomalous absorption.
### 5.5 Photon fluctuations: thermal emission vs. $`\pi ^02\gamma `$
WA98 have measured photon and charged particle multiplicities and their fluctuations versus centrality and $`E_T`$ binning size. As mentioned above impact parameter fluctuations are proportional to the $`E_T`$ binning size; the WA98 analysis nicely confirms this, and can subsequently remove impact parameter fluctuations. The resulting charged particle multiplicity fluctuations with impact parameter fluctuations subtracted, $`\omega _Nn\omega _{N_p}1.11.2`$ were shown in Fig. (7).
The fluctuations in photon multiplicities were found to be almost twice as large as for charged particles $`\omega _\gamma n\omega _{N_p}2.0`$. This has the simple explanation that photons mainly are produced in $`\pi ^02\gamma `$ decays. The fluctuations are then the double of the fluctuations in $`\pi ^0`$ to first approximation as seen from Eq. (18).
If the photons were directly produced from a “shining” thermal fireball one would expect that they would exhibit Bose-Einstein fluctuations, $`\omega _\gamma =\omega _N^{BE}=1.37`$ for massless particles. In addition the $`\pi ^0`$’s in the hadronic background will produce photons with $`\omega _\gamma =\omega _N^{BE}=2.0`$. The measured fluctuation in the number of photons will therefore lie between these two numbers and can be exploited to quantify the amount of thermal photon emission vs. $`\pi ^02\gamma `$ decay from a hadronic gas
$`{\displaystyle \frac{N_\gamma ^{thermal}}{N_\gamma ^{thermal}+N_\gamma ^{\pi ^0}}}={\displaystyle \frac{2.0\omega _\gamma ^{exp}}{2.01.37}}.`$ (42)
The impact parameter fluctuations must be subtracted from the measured photon fluctuations $`\omega _\gamma ^{exp}`$ by, e.g., taking the ratio of photons to some other particle with known behavior.
## 6 Transverse momentum fluctuations
Fluctuations in average transverse momentum were among the first event-by-event analyses studied. In a series of papers Mrówczyński et al. have studied transverse momentum fluctuations in heavy-ion collisions with the purpose of studying thermalization and other effects. Fluctuations in temperature and thus average transverse momentum event-by-event were studied by a number of people in connection with critical phenoma relevant if the transition is close to a critical point. Experimental analyses by NA49 reveal that a careful evaluation of systematic effects are required before substantial equilibration can be claimed in central heavy-ion collisions from transverse momentum fluctuations. They also have found strong correlations between multiplicity and transverse momentum.
The total transverse momentum per event
$`P_t={\displaystyle \underset{i=1}{\overset{N}{}}}p_{t,i},`$ (43)
is very similar to the transverse energy, for which fluctuations have been studied extensively . The mean transverse momentum and inverse slopes of distributions generally increase with centrality or multiplicity. Assuming that $`\alpha d\mathrm{log}(p_t_N)/d\mathrm{log}N`$ is small, as is the case for pions , the average transverse momentum per particle for given multiplicity $`N`$ is to leading order
$`p_t_N=p_t(1+\alpha (NN)/N).`$ (44)
where $`p_t`$ is the average over all events of the single particle transverse momentum. With this parametrization, the average total transverse momentum per particle in an event obeys $`P_t/N=p_t`$. When the transverse momentum is approximately exponentially distributed with inverse slope $`T`$ in a given event, $`p_{t,i}=2T`$, and $`\sigma (p_{t,i})=2T^2=p_t^2/2`$.
The total transverse momentum and also the transverse energy contains both fluctuations in multiplicity and fluctuations in the individual particle transverse momenta and energy (see Appendix C). An interesting quantity is therefore the total transverse momentum per particle, $`P_t/N`$, where the multiplicity fluctuations are removed to first order although important correlations remain.
The total transverse momentum per particle in an event has fluctuations
$`N\sigma (P_t/N)`$ $`=`$ $`\sigma (p_{t,i})+\alpha ^2p_t^2\omega _N+{\displaystyle \frac{1}{N}}{\displaystyle \underset{ij}{}}(p_{t,i}p_{t,j}p_t^2).`$ (45)
The three terms on the right are respectively:
i) The individual fluctuations $`\sigma (p_{t,i})=p_{t,i}^2p_t^2`$, the main term. In the NA49 data, $`p_t=377`$ MeV and $`N=270`$. From Eq. (45) we thus obtain $`(\sigma (P_t/N)^{1/2}/p_t1/\sqrt{2N}=4.3`$%, which accounts for most of the experimentally measured fluctuation 4.65% . The data contains no indication of intrinsic temperature fluctuations in the collisions.
ii) Effects of correlations between $`p_t`$ and $`N`$, which are suppressed with respect to the first term by a factor $`\alpha ^2`$. In NA49 the multiplicity of charged particles is mainly that of pions for which $`Tp_t/2`$ increases little compared with pp collisions, and $`\alpha 0.050.1`$. Thus, these correlations are small for the NA49 data. However, for kaons and protons, $`\alpha `$ can be an order of magnitude larger as their distributions are strongly affected by the flow observed in central collisions .
iii) Correlations between transverse momenta of different particles in the same event. In the WNM the momenta of particles originating from the same participant are correlated. In Lund string fragmentation, for example, a quark-antiquark pair is produced with the same $`p_t`$ but in opposite direction. The average number of pairs of hadrons from the same participant is $`n(n1)`$, where $`n`$ is the number of particles emitted from the same participant nucleon, and therefore the latter term in Eq. (45) becomes $`(n(n1)/n)(p_{t,i}p_{t,ji}p_t^2)`$. To a good approximation, $`n`$ is Poisson distributed, i.e., $`n(n1)/n=n`$, equal to 0.77 for the NA49 acceptance, so that this latter term becomes $`(p_{t,i}p_{t,ji}p_t^2)`$. The momentum correlation between two particles from the same participant is expected to be a small fraction of $`\sigma (p_{t,i})`$.
To quantify the effect of rescatterings, the difference between $`N\sigma (P_t/N)`$ and $`\sigma (p_t)`$ has been studied in detail via the quantity
$`\mathrm{\Phi }(p_t)\sqrt{N\sigma (P_t/N)}\sqrt{\sigma (p_{t,i})}.`$ (46)
As we see from Eq. (45), in the applicable limit that the second and third terms are small,
$`\mathrm{\Phi }(p_t)`$ $``$ $`{\displaystyle \frac{1}{\sqrt{\sigma (p_{t,i})}}}\left(\alpha ^2p_t^2\omega _N+(p_{t,i}p_{t,ji}p_t^2)\right).`$
In the Fritiof model, based on the WNM with no rescatterings between secondaries, one finds $`\mathrm{\Phi }(p_t)4.5`$ MeV. In the thermal limit the correlations in Eq. (46) should vanish for classical particles but the interference of identical particles (HBT correlations) contributes to these correlations $`6.5`$ MeV ; they are again slightly reduced by resonances. The NA49 experimental value, $`\mathrm{\Phi }(p_t)=5`$ MeV (corrected for two-track resolution) seems to favor the thermal limit . Note however that with $`\alpha 0.050.1`$, the second term on the right side of Eq. (LABEL:phii) alone leads to $`\mathrm{\Phi }14`$ MeV, i.e., the same order of magnitude. If $`(p_{t,i}p_{t,ji}p_t^2)`$ is not positive, then one cannot a priori rule out that the smallness of $`\mathrm{\Phi }(p_t)`$ does not arise from a cancellation of this term with $`\alpha ^2p_t^2\omega _N`$, rather than from thermalization.
A comparison of the transverse momentum fluctuations of charged particles to those in mixed events, where correlations thus are removed, showed a small enhancement of only $`0.002\pm 0.002`$ . It was estimated that Bose effects should enhance this ratio by 1-2% but that total energy conservation introduces an anticorrelation that partially cancels the Bose enhancement . Experimental problems with two-track resolution have also been estimated to lead to a ratio that is 1-2% lower. Consequently, the numbers seem to be compatible.
The covariance matrix between multiplicity and transverse momentum has been analyzed by NA49 . Strong but trivial correlations is found due to the fact that higher multiplicity gives larger total transverse momentum event-by-event. This correlation is removed in the quantity $`P_t/N`$ and its covariance matrix with multiplicity appears diagonal.
## 7 Event-by-Event Fluctuations at RHIC
The theoretical analysis above leads to a qualitative understanding of event-by-event fluctuations and speculations on how phase transitions may show up. It gives a quantitative description of AGS and SPS data without the need to invoke new physics. We shall here look ahead towards RHIC experiments and attempt to describe how fluctuations may be searched for.
General correlators between all particle species should be measured event-by-event, e.g., the ratios
$`{\displaystyle \frac{N_i/N_j}{N_i/N_j}}1+{\displaystyle \frac{\omega _{N_j}}{N_j}}{\displaystyle \frac{N_iN_jN_iN_j}{N_iN_j}},`$ (48)
where $`N_{i,j}`$ are the multiplities in acceptances $`i`$ and $`j`$ of any particle. Volume fluctuations are automatically removed in such ratios, their fluctuations and correlations. If the energy deposition, transverse energy or momentum are measured, these latter will have additional fluctuation due to the multiplicity fluctuations as explained in Appendix C.
More generally we define the multiplicity correlations between any two bins
$`\omega _{ij}={\displaystyle \frac{N_iN_jN_iN_j}{\sqrt{N_iN_j}}},`$ (49)
also referred to as the covariance. When $`i,j`$ refer to two rapidity bins the covariance is also proportional to the rapidity (auto-)correlation function $`C(y_iy_j)`$.
It is instructive to consider first completely random (uncorrelated or statistical) particle emission. For a fixed total multiplicity $`N_{Tot}`$, the probability for a particle to end up in bin $`i`$ is $`p_i=N_i/N_{Tot}E_i/E_{Tot}`$. The distribution is a simple multinomial distribution for which
$`\omega _{ij}=\left\{\begin{array}{ccc}1p_i\hfill & \hfill ,& i=j\hfill \\ \sqrt{p_ip_j}\hfill & \hfill ,& ij\hfill \end{array}\right\}.`$ (52)
The $`i=j`$ result is the well known one for a binomial distribution. The $`ij`$ result is negative because particles in different bins are anti-correlated: more (less) particles in one bin leads to less (more) in other bins on average due to a fixed total number of particles.
As shown above there are nonstatistical fluctuations due to various sources: Bose-Einstein fluctuations, resonances, etc., and — in particular — density fluctuations. As in Eq. (21) we assume that the multiplicity consist of particles from a HM and a QM phase. The covariances in Eq. (52) are derived analogously to Eq. (22)
$`\omega _{ij}=\omega _{ij,HM}+(\omega _{ij,QM}\omega _{ij,HM}){\displaystyle \frac{N_{i,QM}}{N_i}},`$ (53)
when $`N_i=N_j`$; when different the general formula is a little more complicated. Now, the hadronic fluctuations $`\omega _{ij,HM}`$ is of order unity for $`i=j`$, smaller for adjacent bins and vanishes or even becomes slightly negative according to (52) for bins very different in pseudorapidity or azimuthal angle $`\varphi `$. The QM fluctuations can be much larger: $`\omega _{i,QM}N_{i,QM}`$ (see the discussion after Eq. (20)). To discriminate the QM fluctuations from the hadronic ones, Eq. (53) requires
$`N_{i,QM}\stackrel{>}{}\sqrt{(\omega _{ij}\omega _{ij,HM})N_i}.`$ (54)
The charged particle multiplicity in central $`Au+Au`$ collisions at RHIC is $`dN_{ch}/d\eta 500600`$ per unit pseudo-rapidity . To see a clear increase in fluctuations, say $`\mathrm{\Delta }\omega \omega _{ij}\omega _{ij,HM}1`$, a density fluctuation of only $`N_{i,QM}\stackrel{>}{}\sqrt{N_i}25`$ particles are required per unit rapidity corresponding to a few percent of the average. By analyzing many events (of the same total multiplicity) the accuracy by which fluctuations are measured experimentally is greatly improved. Generally, $`\mathrm{\Delta }\omega 1/\sqrt{N_{events}}`$, and so fluctuations can in principle be determined with immense accuracy.
It may be advantageous to correlate bins with the same pseudorapidity but different azimuthal angles since the hadronic correlations between these are small whereas QM fluctuations remain.
No experimental determination of the purely statistical uncertainties associated with any one-body distribution — such as multiplicity as a function of rapidity — can be performed without measuring and diagonalizing the correlation matrix $`C_{ij}=N_iN_jN_iN_j`$. While it is conventional to assign uncertainties according to the diagonal elements $`M_{ii}`$, the correlations in the covariance matrix are required for a correct error analysis and can also reveal physical important results.
## 8 Summary
A phase transitions in high energy nuclear collisions, whether it is first order or a soft cross-over, density fluctuations may be expected that show up in rapidity and multiplicity fluctuations event-by-event. The fluctuations can be enhanced significantly in case of droplet formation as compared to that from an ordinary hadronic scenario. A combined analyses of, e.g., positive, negative, total and net charge, allows one to extract the various fluctuations and correlations uniquely. Likewise a number of other observables as charged and neutral pions, kaons, photons, $`J/\mathrm{\Psi }`$, etc., and their ratios can show anomalous correlations and enhancement or suppression of fluctuations. This clearly demonstrates the importance of event-by-event fluctuations accompanying phase transitions, and illustrates how monitoring such fluctuations versus centrality becomes a promising signal, in the upcoming RHIC experiments, for the onset of a transition. The potential for enhanced or suppressed fluctuations (orders of magnitude) from a transition makes it worth looking for at RHIC considering the relative simplicity and accuracy of multiplicity fluctuation measurements.
An analysis of fluctuations in central Pb+Pb collisions as currently measured in NA49 does, however, not show any sign of anomalous fluctuations. Fluctuations in multiplicity, transverse momentum, $`K/\pi `$ and other ratios can be explained by standard statistical fluctuation and additional impact parameter fluctuations, acceptance cuts, resonances, thermal fluctuations, etc. This understanding by “standard” physics should be taken as a baseline for future studies at RHIC and LHC and searches for anomalous fluctuations and correlations from phase transitions that may show up in a number of observables.
By varying the centrality one should be able to determine quantitatively the amount of thermalization in relativistic heavy ion collisions as defined in Eq. (19) . For peripheral collisions, where only few rescatterings occur, we expect the participant model (WNM) to be approximately valid and the degree of thermalization to be small. For central collisions, where many rescatterings occur among produced particles, we expect to approach the thermal limit and the degree of thermalization should be close to 100%. At RHIC and LHC energies the fluctuations in the number of charged particles consequently decrease drastically with centrality whereas at SPS energies the two limits are accidentally very close.
Event-by-event physics is an important tool to study thermalization and phase transitions through anomalous fluctuations and correlations — as in rain.
## Acknowledgements
Thanks to G. Baym and A.D. Jackson for inspiration and collaboration on some of the work described in this report. Discussion with S. Voloshin and G. Roland (NA49), J.J. Gårdhøje and collaborators in NA44 and BRAHMS, T. Nayak(WA98), J. Bondorf, S. Jeon, V. Koch, and many suggestions for improvement from an anonymous referee are gratefully acknowledged.
## Appendix A: Damping of initial density fluctuations
Hydrodynamic flow with Bjorken scaling is stable according to a stability analysis carried out in . Bt linearing the hydrodynamic equations in small perturbations in entropy density $`\delta s`$ and rapidity $`\delta y`$ around the Bjorken scaling solution and looking for solutions in the form of harmonic perturbations, $`e^{ik\eta }`$, the hydrodynamic equations could be written in matrix form (Eq. A.13 in )
$`\tau {\displaystyle \frac{\delta }{\delta \tau }}\left(\begin{array}{c}\delta s/s\hfill \\ \delta y\hfill \end{array}\right)=\left(\begin{array}{ccc}0\hfill & & ik\hfill \\ ikc_s^2\hfill & & (1c_s^2)\hfill \end{array}\right)\left(\begin{array}{c}\delta s/s\hfill \\ \delta y\hfill \end{array}\right).`$ (61)
The eigenvalues of the above matrix
$`\lambda _\pm ={\displaystyle \frac{1}{2}}(1c_s^2)\pm \sqrt{{\displaystyle \frac{1}{4}}(1c_s^2)^2c_s^2k^2},`$ (62)
always have real negative part for $`c_sk0`$ and fluctuations are therefore damped. For long wave length fluctuations in rapidity and not too soft equations of state, $`c_sk>1c_s^2`$, the solution is a damped oscillator. Note that the long wave length solution $`k=0`$ reproduces the Bjorken scaling.
The exact solution for the entropy density fluctuation
$`{\displaystyle \frac{\delta s}{s}}=c_+e^{\lambda _+\mathrm{ln}(\tau /\tau _0)}+c_{}e^{\lambda _{}\mathrm{ln}(\tau /\tau _0)},`$ (63)
is sensitive to the equation of state through $`c_s`$, the initial conditions for the rapidity density fluctuations (the constants $`c_\pm `$), and their wave length $`k^1`$.
At large times the eigenvalue with the largest real part dominates and
$`{\displaystyle \frac{\delta s}{s}}\left({\displaystyle \frac{\tau _0}{\tau _f}}\right)^{|Re[\lambda _\pm ]|}.`$ (64)
Here the oscillating factor has been ignored, leaving the power law fall-off of fluctuations with exponent
$`Min|Re[\lambda _\pm ]|={\displaystyle \frac{1}{2}}(1c_s^2)Re[\sqrt{{\displaystyle \frac{1}{4}}(1c_s^2)^2c_s^2k^2}]`$ (65)
One notes that density fluctuations are undamped for soft equation of states ($`c_s=0`$). They are also undamped if their wave length is long ($`k0`$).
To estimate the resulting damping we take a typical rapidity fluctuation for a droplet $`\delta y\sqrt{T/m_t}1`$ discussed above, which corresponds to a wave-number $`k1`$. For an ideal equation of state with sound speed $`c_s=1/\sqrt{3}`$ the last term in Eq. (62) is then either imaginary or small and real, and the real part of the eigenvalue is dominated by the first term of Eq. (62), $`Re[\lambda _\pm ]1/3`$. If we take a typical formation time $`\tau _01`$ fm/c and a freezeout time $`\tau _f8`$ fm/c as extracted from HBT studies , the resulting suppression of a density fluctuation during expansion is a factor $`8^{1/3}=0.5`$ according to Eq. (64).
## Appendix B: Fluctuations in source models
As fluctuations for a source model appears again and again (see Eqs. 10,11,18,22) we shall derive this simple equation in detail.
We define the fluctuations for any stochastic variable $`x`$ as
$`\omega _x={\displaystyle \frac{x^2x^2}{x}}.`$ (66)
It is usually of order unity and therefore more convenient than variances. For a Poisson distribution, $`P_N=e^\alpha \alpha ^N/N!`$, the fluctuation is $`\omega _N=1`$. For a binomial distribution with tossing probability $`p`$ the fluctuation is $`\omega _N=1p`$, independent of the number of tosses. In heavy ion collisions several processes add to fluctuations so that typically $`\omega _N^{exp}12`$. Correlations can in some cases double the fluctuations as, for example, $`\pi ^02\gamma `$ doubles the fluctuations in photon multiplicity and net charge conservation doubles the fluctuation in total charge. Impact parameter fluctuations further increases the total charge fluctuations to $`\omega _{N_{ch}}=35`$ in peripheral nuclear collisions .
Generally, when the multiplicity ($`N`$) arise from independent sources $`(N_p)`$ such as participants, resonances, droplets or whatever,
$`N={\displaystyle \underset{i=1}{\overset{i=N_p}{}}}n_i,`$ (67)
where $`n_i`$ is the number of particles produced in source $`i`$. In the absence of correlations between $`N_p`$ and $`n`$, the average multiplicity is $`N=N_pn`$. Here, $`..`$ refer to averaging over each individual (independent) source as well as the number of sources. The number of sources vary from event to event and average is performed over typically $`N_{events}100.000`$ events as in NA49 or $`N_{events}10^6`$ in WA98.
Squaring Eq. (67) assuming that the source emit particles independently, i.e. $`n_in_j=n_in_j`$ for $`ij`$, the square consists of the diagonal and off-diagonal elements:
$`N^2=N_pn_i^2+N_p(N_p1)n_i^2.`$ (68)
With (66) we obtain the multiplicity fluctuations
$`\omega _N={\displaystyle \frac{N^2N^2}{N}}=\omega _n+n\omega _{N_p},`$
as in Eq. (10).
## Appendix C: Fluctuations in the energy deposited
Many experiments do not measure individual particle tracks or multiplicities but instead the energy deposited in arrays of detector segments, $`E_i`$, in a given event. One could also project the energy transversely by weighting with the sine of the scattering angle to study fluctuations in transverse energy . Since particles mostly have relativistic speeds in relativistic heavy-ion collisions, the transverse energy is almost the same as the total transverse momentum in an event.
The total energy deposited in the event is
$`E_{Tot}={\displaystyle \underset{i}{\overset{D}{}}}E_i,`$ (69)
and can be used as a measure of the centrality of the collision. The energy deposited in each element (or group of elements) is the sum over the number of particle tracks ($`N_i`$) hitting detector $`i`$ of the individual ionization energy of each particle ($`ϵ_i`$)
$`E_i={\displaystyle \underset{n}{\overset{N_i}{}}}ϵ_n.`$ (70)
The average is: $`E_i=N_iϵ`$. The energy will approximately be gaussian distributed, $`d\sigma /dE_i\mathrm{exp}((E_iE_i)^2/2\omega _{E_i}E_i)`$, with fluctuations (see Appendix B)
$`\omega _{E_i}{\displaystyle \frac{E_i^2E_i^2}{E_i}}=\omega _ϵ+ϵ\omega _{N_i}.`$ (71)
Here, the fluctuation in ionization energy per particle
$`{\displaystyle \frac{\omega _ϵ}{ϵ}}={\displaystyle \frac{ϵ^2}{ϵ^2}}1,`$ (72)
depends on the typical particle energies in the detector and the corresponding ionization energies for the detector type and thickness. For the BRAHMS detectors we estimate $`\omega _ϵ/ϵ0.3`$ . This number will, however, depend on rapidity since the longitudinal velocity enters the ionization power. As these are “trivial” detector parameter, we shall exclude the fluctuations $`\omega _ϵ`$ in most analyses and concentrate on the second term in Eq. (71) which is the fluctuations in the number of particles as examined in detail above. |
warning/0003/math0003035.html | ar5iv | text | # Branched cyclic covers and finite type invariants
## 1. Introduction
Let K be the free abelian group generated by (oriented) equivalence classes of pairs $`(M,K)`$, where $`M`$ is an integral homology three-sphere and $`K`$ is an oriented knot in it. Let $`𝖬`$ be the free abelian group generated by homeomorphism classes of three-manifolds. Let the mapping
$$\mathrm{\Sigma }^p:\text{K}𝖬$$
be the linear extension of the operation defined on a pair $`(M,K)`$ as the $`p`$-fold branched cyclic covering of $`M`$ branched over $`K`$. We are chiefly interested in knots in the three-sphere, for which the corresponding space will be denoted $`\text{K}(S^3)`$, when required.
The subject of this investigation is the set of knot invariants obtained by composing the branched cyclic cover construction, for a fixed degree, with finite-type invariants of three-manifolds (of various descriptions).
This paper is in two parts. The main theorem of the first part is Theorem 2.0.3, and its enhancement, Theorem 2.0.6. This theorem describes the effect of surgery on complete graph Y-links on the $`p`$-fold branched cyclic covers of a knot, with respect to the Goussarov-Habiro filtration. A calculus is developed in Lemmas 2.0.9 and 2.0.10 which describes this action.
In the second part of the paper, we illustrate this calculus by exploring some very general questions about the invariants obtained by composing projections of the LMO invariant with the operation of taking a $`p`$-fold branched cyclic cover. The main theorem of this part is the realisation theorem, Theorem 4.0.2, which has the following, possibly unsurprising corollary:
###### Corollary 1.0.1.
Take a positive integer $`p`$. If $`v`$ is a rational valued three-manifold invariant factoring non-trivially through the LMO invariant on integral homology spheres, then $`v\mathrm{\Sigma }^p`$ is not a finite-type invariant of knots.
It is interesting that we can use a finite-type property (in one sense) to show that an invariant is not finite-type (in some different sense).
The generality of the above theorem does not seem to follow from the more descriptive investigations of specific examples and projections of the LMO invariant that have been recorded: Hoste gave a formula for the Casson invariant of the $`p`$-fold branched cyclic cover over an untwisted double of a knot \[Hos\]; Davidow gave formulae for the case of iterated torus knots \[Dav1\], and for some $`1`$-twisted doubles \[Dav2\]; Mullins calculated the composition of the Casson-Walker invariant with the 2-fold branched cyclic cover in terms of the Jones polynomial, when the left hand side was well-defined \[Mul1, Mul2\]; Ishibe used this work to give formulae for the general case of $`m`$-twisted doubles \[Ish\]; Garoufalidis showed that Mullins formulae was valid in general, using the Casson-Walker-Lescop invariant \[Gar\]; and also in that paper Garoufalidis initiated the investigation into the LMO invariant on branched cyclic covers, describing the form of the result for covers over twisted doubles.
The main motivation for this work was to see whether the techniques we employed in \[K\] (an application of clasper theory) had wider application. The other main motivation was to explore a possible approach to Lev Rozansky’s program for a “finite type theory of knots’ complements” \[Roz\]. We expect the chief interest in this article to be the analysis of a suggestive family of operations which will be important in the development of a more structural relationship between the loop expansion of the Kontsevich integral, and the degree expansion of $`Z^{LMO}\mathrm{\Sigma }^p`$. (That is the informal conjecture we are seeking to motivate with this paper.)
A review of some of the theory of Goussarov-Habiro filtration is included as an appendix.
###### Acknowledgements 1.0.2.
The author is supported by a Japan Society for the Promotion of Science Postdoctoral Fellowship, and thanks Tomotada Ohtsuki and the Department of Mathematical and Computing Sciences at the Tokyo Institute of Technology for their support. The author would also like to thank Stavros Garoufalidis, Kazuo Habiro, Hitoshi Murakami and Theodore Stanford (and many others) for interesting comments.
## 2. Part 1: Complete clasper moves and branched cyclic covers
The Goussarov-Habiro filtration \[G3, Hab\] is a descending filtration of $`𝖬`$ (see Definition A.0.2):
$$𝖬𝖥_1^Y𝖬𝖥_2^Y𝖬\mathrm{}$$
On the other hand we have clasper calculus, introduced by Habiro \[HT, Hab\], certain aspects of which were independently developed by Goussarov \[G2\]. This generalised Matveev’s “Borromeo” move \[Mat\] and Murakami-Nakanishi’s “delta-unknotting” operation \[MN\].
In this work, we will identify a subclass of clasper operations which the branched cyclic covering construction will translate to $`n`$-equivalences of three-manifolds. For technical reasons, this work will use the language of Y-links; note that our language is slightly different to that adopted in \[GL2\].
In Appendix A we will fix the term graph Y-link to mean a Y-link decoration of a knot $`K`$ such that each leaf of every Y-component links either another leaf, in a ball (see Appendix A for diagrammatic conventions):
or links the knot thus:
Note that we also consider graph Y-links which may be without legs. Note that surgery on a graph Y-link without legs will change the homeomorphism class of the ambient manifold; surgery on a connected graph Y-link with at least one leg will not.
To every graph Y-link there is canonically associated (up to sign) a uni-trivalent diagram (which here refers to the familiar uni-trivalent graphs, which may have univalent vertices located on a skeleton; when we take linear combinations, we introduce STU, IHX and AS relations). This association is made by replacing Y-components by trivalent vertices, replacing leaves linking the knot by legs, and by joining tines of trivalent vertices by an edge when the corresponding leaves were linked. A graph Y-link is called connected when the dashed graph of this associated diagram is connected. The degree of a graph Y-link is half the number of vertices of the dashed graph of the associated uni-trivalent diagram.
To introduce the moves in question, some definitions are required.
###### Definition 2.0.1.
A graph Y-link decorating a knot $`K`$ is complete if it is connected and if every leg of the associated uni-trivalent diagram leads to a seperate trivalent vertex.
In particular, we are ruling out the chord:
and graph Y-links which end in a fork:
It is important to point out that this is quite a strong restriction: for example, a move of such a kind on a link will not affect its Milnor invariants.
Now we need an alternative measure on graph Y-links.
###### Definition 2.0.2.
The surplus of a uni-trivalent diagram is defined to be the number of trivalent vertices minus the number of univalent vertices of its dashed graph. The surplus of a graph Y-link is defined to be the surplus of its associated uni-trivalent diagram.
For example, the surplus of a complete graph Y-link is simply the grade of the underlying uni-trivalent diagram, after all the legs have been removed. The main theorem is expressed in these terms.
###### Theorem 2.0.3 (Main theorem).
If a knot $`K^L`$ is obtained from a knot $`K`$ by surgery on a complete graph Y-link $`L`$ of surplus $`s`$, where $`s2`$, then
$$\mathrm{\Sigma }_{K^L}^p\mathrm{\Sigma }_K^p𝖥_s^Y𝖬.$$
There is another important invariant for which a similar property holds. The following theorem \[KGr\], and the theory of which it is an application, may appear in a future manuscript.
###### Theorem 2.0.4.
If a knot $`K^L`$ is obtained from a knot $`K`$ by surgery on a complete graph Y-link $`L`$ of surplus $`s`$, where $`s2`$, then
$$Z^{LMO}(K^L)Z^{LMO}(K)=\{\text{series of diagrams of surplus }s\},$$
where $`Z^{LMO}`$ is the Kontsevich-LMO (Le-Murakami-Ohtsuki) invariant of knots in integral homology three-spheres.
The point being that it is true for every grade. We leave the reader to make their own conjectures.
Returning to the main theorem of this paper, the next natural question is how these moves operate on the leading term in the space of associated graded quotients. Specifically, what is
$$\mathrm{\Sigma }_{K^L}^p\mathrm{\Sigma }_K^p𝖦_s^Y𝖬\mathrm{?}$$
Lemma 2.0.9 and Lemma 2.0.10 calculate this term. To report this calculation, let the ambient integral-homology three-sphere be presented by some surgery link, and isotope $`K`$ and the decorating graph Y-link in that three-manifold to appear in the complement of that surgery link in the three-sphere.
###### Definition 2.0.5.
Let $`\lambda _L`$ be a graph Y-link in $`S^3`$ obtained by sawing off the legs of $`L`$
$$\text{}\text{}$$
and then forgetting $`K`$ and the surgery link.
Let $`\lambda _L^{}`$ be the corresponding graph Y-link in $`\mathrm{\Sigma }_K^p`$ induced by the connect-sum of that three-sphere into $`\mathrm{\Sigma }_K^p`$. Let $`(\mathrm{\Sigma }_K^p)^{\lambda _L^{}}`$ denote the result of surgery on that graph Y-link.
Lemma 3.2.1 asserts that $`(\mathrm{\Sigma }_K^p)^{\lambda _L^{}}\mathrm{\Sigma }_K^p`$ is well-defined in $`𝖦_s^Y𝖬`$. Lemma 2.0.9 and Lemma 2.0.10 show that:
###### Theorem 2.0.6 (Main theorem, part 2).
$$\mathrm{\Sigma }_{K^L}^p\mathrm{\Sigma }_K^p(\mathrm{\Sigma }_K^p)^{\lambda _L^{}}\mathrm{\Sigma }_K^p𝖦_s^Y𝖬.$$
###### Remark 2.0.7.
The definition of $`(\mathrm{\Sigma }_K^p)^{\lambda _L^{}}\mathrm{\Sigma }_K^p`$ is perhaps more easily understood in terms of the surjective map from the space of marked uni-trivalent graphs of degree $`s`$ to $`𝖦_s^Y𝖬([M])`$. In this sense, the first non-vanishing term is proportional to the image under that map of the diagram obtained by removing all the legs from a diagram associated to the initial decoration. The procedure we employ in this paper is used so that we can be precise, in a natural way, about the sign of the element referred to (see Section 3.1).
The following two lemmas can be used to immediately calculate this term. Calculation proceeds in two steps. The first lemma is used to reduce the number of legs on $`K`$, so that eventually $`\mathrm{\Sigma }_{K^L}^p\mathrm{\Sigma }_K^p`$ is expressed as a linear combination of differences $`\mathrm{\Sigma }_{K_i}^p\mathrm{\Sigma }_K^p`$ where each $`K_i`$ has been obtained from $`K`$ by surgery on a connected graph Y-link without legs. The second lemma calculates an expression for the terms in that sum.
###### Remark 2.0.8.
This procedure is the reason it is more natural in this context to work in the generality of integral homology three-spheres. It is interesting that the calculation necessitates this extension, not unlike the results of \[GL2\].
###### Lemma 2.0.9.
Let $`K`$ be a knot. Let $`L`$, $`L^{}`$ and $`L^{\prime \prime }`$ denote three decorations of $`K`$ by complete graph Y-links of surplus $`s`$ that differ in a ball in the following way:
$$\text{},$$
$$\text{}\text{and}\text{}.$$
Then
$$(\mathrm{\Sigma }_{K^L}^p\mathrm{\Sigma }_K^p)=(\mathrm{\Sigma }_{K^L^{}}^p\mathrm{\Sigma }_K^p)(\mathrm{\Sigma }_{K^{L^{\prime \prime }}}^p\mathrm{\Sigma }_K^p)𝖦_s^Y𝖬.$$
The statement of the next lemma requires an embedding of a graph in $`MK`$ that will be associated to a decoration of $`K`$ by a graph Y-link $`L`$ in the complement of $`K`$. This embedding is constructed by contracting the clasper graph presenting $`L`$ onto a spine: call this embedded graph $`E_L`$. In the statement of the lemma, a loop on $`L`$ will mean precisely a path in $`MK`$ that corresponds to a loop on $`E_L`$.
###### Lemma 2.0.10.
Let a knot $`K^L`$ be obtained from a knot $`K`$ by surgery on some connected graph Y-link $`L`$ without legs of surplus $`s`$. Then
$$\mathrm{\Sigma }_{K^L}^p\mathrm{\Sigma }_K^p=\{\begin{array}{ccc}p((\mathrm{\Sigma }_K^p)^{\lambda _L^{}}\mathrm{\Sigma }_K^p)\hfill & \text{if}\hfill & \text{every loop on }L\text{ has zero linking}\hfill \\ & & \text{number with }K\text{ modulo }p,\hfill \\ & & \\ 0\hfill & & \text{otherwise}.\hfill \end{array}𝖦_s^Y𝖬.$$
Intuitively speaking, this is what one would expect. The implementation of that intuition in this setting, however, is not immediate (Section 3.4).
Let us illustrate this calculus with some examples based on the most familiar Goussarov-Habiro finite type invariant, the Casson-Walker-Lescop invariant, denoted $`\lambda _{CWL}`$. The connection will be made with the following lemma. Extend the symbol $`|H_1(M)|`$ to be zero in the case that $`M`$ has non-vanishing first Betti number.
###### Lemma 2.0.11.
1. $`\lambda _{CWL}(𝖦_3^Y𝖬)=0.`$
2. Let some three-manifold $`M`$ have some graph Y-link $`\mathrm{\Theta }`$ embedded in it, whose underlying diagram is the “theta” graph. Then
$$\lambda _{CWL}(M^\mathrm{\Theta })=\lambda _{CWL}(M)\pm 2|H_1(M)|,$$
where the sign can be determined in a given situation.
Perhaps the easiest way to see this is to observe the identification of the co-efficient of the theta graph in the LMO invariant with half the Casson-Walker-Lescop invariant, \[LMMO\] (see also \[GH, HB, L\]), and then observe that the difference of the LMO invariant on a manifold and on that manifold surgered along a $`\mathrm{\Theta }`$ is precisely $`\pm |H_1(M)|\theta `$ plus higher order terms, where the sign can be determined in a given situation (see \[LeGr\] for such first non-vanishing term calculations).
###### Example 2.0.12.
Fix an integer $`p`$ and take a knot $`K`$. Consider some graph $`Y`$-link decoration $`\mathrm{\Theta }`$ of $`K`$ with underlying diagram the “theta” graph, and let $`K^\mathrm{\Theta }`$ be the knot (in some integral homology three-sphere) obtained from surgery on $`\mathrm{\Theta }`$. Let $`\theta `$ denote the spatial graph in $`S^3K`$ associated to this decoration. Let $`l_1`$ denote a choice of a knot in $`S^3K`$ obtained by forgetting some edge of $`\theta `$, and let $`l_2`$ denote a knot obtained by forgetting some other choice of edge. Then:
$$\lambda _{CWL}(\mathrm{\Sigma }_{K^\mathrm{\Theta }}^p)=\lambda _{CWL}(\mathrm{\Sigma }_K^p)\pm \{\begin{array}{cc}2p|H_1(\mathrm{\Sigma }_K^p)|\hfill & \text{if lk}(l_i,K)=0\text{mod p, }i=1,2\text{ ,}\hfill \\ 0\hfill & \text{otherwise,}\hfill \end{array}$$
where the sign can be determined in a given situation. This follows from Lemma 2.0.10.
How might we organise such a calculation in generality? That is, where we have a decoration of a knot by a theta graph with some legs attached. Such an example follows.
###### Example 2.0.13.
We indicate a basis $`\{a,b\}`$ for the first homology of the associated trivalent graph, and some labels $`\{ϵ_1,ϵ_2\}`$ have been affixed to the legs. Let $`U`$ indicate the unknot, and let the decorating graph Y-link be denoted by $`\kappa `$.
$$\text{}U$$
Theorem 2.0.3 indicates that under $`\mathrm{\Sigma }^p`$ this difference lies in $`𝖦_s^Y𝖬`$. To calculate the term we proceed by applying Lemma 2.0.9 to each of the legs. This produces four terms, which will be recorded by setting each of the $`ϵ_i`$ to either 0 or 1. For example the term that corresponds to $`(ϵ_1,ϵ_2)=(0,1)`$ will be:
$$\text{}U$$
Now, according to Lemma 2.0.10, precisely those terms for which the cycles $`a`$ and $`b`$ have mod-p linking number zero with $`U`$ will contribute. For a given $`(ϵ_1,ϵ_2)`$, the linking numbers will be (giving $`U`$ the clockwise orientation):
$`lk(a,U)`$ $`=`$ $`ϵ_1,`$
$`lk(b,U)`$ $`=`$ $`ϵ_2+1.`$
If we collect this information into the monomial $`F_{ϵ_1,ϵ_2}(t_a,t_b)=t_a^{lk(a,U)}t_b^{lk(b,U)}`$, then for a given $`p`$, for some determinable sign $`\alpha `$:
$$\lambda _{CWL}=\alpha 2p\underset{(ϵ_1,ϵ_2)=(0,0)}{\overset{(1,1)}{}}\left(\frac{1}{p}\underset{r=0}{\overset{p1}{}}\frac{1}{p}\underset{s=0}{\overset{p1}{}}F_{ϵ_1,ϵ_2}(e^{r\frac{2\pi i}{p}},e^{s\frac{2\pi i}{p}})\right).$$
\[ To see this, note the following identity:
$$\underset{s=0}{\overset{p1}{}}(e^{\frac{2\pi i}{p}s})^t=\{\begin{array}{cc}p\hfill & \text{if}p|t,\hfill \\ 0\hfill & \text{otherwise.]}\hfill \end{array}$$
We can seperate out the dependence on $`p`$ as follows. Defining the polynomial:
$$F^{}(t_a,t_b)=\underset{(ϵ_1,ϵ_2)=(0,0)}{\overset{(1,1)}{}}F_{ϵ_1,ϵ_2}(t_a,t_b),$$
then the result is:
$$\lambda _{CWL}(\mathrm{\Sigma }_{U^\kappa }^p)=\alpha \frac{2}{p}\underset{r=0}{\overset{p1}{}}\underset{s=0}{\overset{p1}{}}F^{}(e^{r\frac{2\pi i}{p}},e^{s\frac{2\pi i}{p}}).$$
###### Remark 2.0.14.
Observe that such a procedure can be implemented for any decoration of any knot by a complete graph Y-link whose underlying diagram is a theta graph with some legs attached. Thus we have a flexible new construction of knots for which the Casson-Walker-Lescop invariants of their $`p`$-fold branched cyclic covers may be calculated, for all choices of $`p`$. It is interesting that this sequence is determined by the values of an associated polynomial at the roots of unity, evoking the familiar formula for the order of the first homology in terms of the Alexander-Conway polynomial \[F, HK\] of the branching knot.
Observe also that similar constructions and arguments can be applied to any projection of the LMO invariant (see Example 4.1.2).
Section 2.1 will prove part 1 of the main theorem. Section 3 will prove the Lemmas 2.0.9 and 2.0.10 which facilitate the extension.
### 2.1. Branched cyclic covers
Excellent expositions of this construction can be found in the literature, for example in \[Rol\] and in \[Lic\].
Let $`N(K)`$ be a regular neighbourhood of a knot $`K`$ in an integral homology three-sphere $`M`$, and let $`X_K`$ be the closure of its exterior. Let $`X_K^p`$ be the $`p`$-fold cyclic cover of $`X_K`$.
In this work, it will be convenient to construct this space directly with a Seifert surface $`F`$ for $`K`$. This direct construction uses the space $`W_F`$, which is obtained by “splitting $`X_K`$ open along $`F`$”. In practice, $`W_F`$ is obtained by removing the intersection of an open collar neighbourhood of the surface $`F`$ with $`X_K`$: that is, if $`N:F\times [1,1]M`$ is a bicollar for $`F`$ then $`W`$ is obtained as the removal of $`X_KN(F\times (ϵ,ϵ))`$ from $`X_K`$. The subspaces $`X_KN(F\times \{ϵ\})`$ and $`X_KN(F\times \{ϵ\})`$ are homeomorphic images of $`F`$, call them $`F^+`$ and $`F^{}`$, and there will exist a homeomorphism between them $`\varphi :F^{}F^+`$, a gluing along which recovers the exterior of the knot, $`X_K`$.
The unique $`p`$-fold cyclic cover of the exterior of the knot, $`X_K^p`$, is constructed from $`p`$ copies of $`W_F`$: denote them $`W_F^i`$, where the index $`i`$ is an element of $`_p`$. Glue $`F^{}`$ in $`W_F^i`$ to $`F^+`$ in $`W_F^{i+1}`$ using $`\varphi `$ composed with the identification $`W_F^iW_F^{i+1}`$.
The $`p`$-fold cyclic cover of $`M`$ branched over $`K`$ is obtained as a completion of the space $`X_K^p`$. Specifically, a solid torus is glued into the torus boundary of $`X_K^p`$ so that its meridian projects to a meridian of $`K`$ traversed $`p`$ times.
In this work, we will encounter a situation where a knot $`K^L`$ is recovered from another knot $`K`$ by surgery on some framed link $`L`$ decorating $`K`$, and where there is given some Seifert surface $`F`$ for $`K`$ in the complement of $`L`$. In such a situation, a Seifert surface $`F^L`$ for $`K^L`$ may be canonically chosen: namely, just choose the position of $`F`$ in the surgered manifold. The pair $`(W_{F^L},\varphi ^L)`$ can then be obtained from the pair $`(W_F,\varphi )`$ by surgery on the link $`L`$ in $`W_F`$. Thus, by construction, $`\mathrm{\Sigma }_{K^L}^p`$ is recovered by surgery on the framed link in $`\mathrm{\Sigma }_K^p`$ comprising of one copy of the original framed link in each of the $`p`$ copies of $`W_F`$ included in $`\mathrm{\Sigma }_K^p`$.
### 2.2. Complete graph Y-links
Before we describe this proof, we introduce a term.
###### Definition 2.2.1.
A mixed Y-link in a three-manifold is a an embedding of a set of Y-components and possibly some claspers, into that three-manifold.
Let us turn to the scenario described by the main theorem of this section. We have two knots: $`K`$ and a knot $`K^L`$ obtained from $`K`$ by surgery on a complete graph Y-link $`L`$ of surplus $`s`$. The strategy is to replace $`L`$ with a mixed Y-link also presenting the knot $`K^L`$, but which is in the complement of a Seifert surface $`\stackrel{ˇ}{F}`$ for $`K`$.
To begin, choose some Seifert surface $`F`$. We will now isotope the graph Y-link $`L`$ so that it is in a special position with respect to $`F`$. We require two conditions of this position. The first requirement is that the band of a clasper associated to a leg does not intersect $`F`$ (although one of the leaves certainly will). Given such a position, which can always be obtained, then around every leaf we can find a ball so that the arrangement is obtained by some homeomorphism of that ball with the following ball in $`S^3`$:
The second requirement is that the only other intersections of $`L`$ with $`F`$ are transversal intersections of bands with $`F`$.
A graph Y-link in such a position will now be modified in two steps. The first step is make an adjustment of $`L`$ around each leg, where it appears as in the ball above. The adjustment is made with the move described in Lemma A.0.5:
There is one such adjustment to be made for every (if there are any, that is) leg of the graph Y-link, leaving the decoration with one transverse intersection of a band with $`F`$ for every leg, and with possible further transverse intersections of internal bands with $`F`$.
The second step is to adjust the resulting mixed graph Y-link and Seifert surface at each such intersection (so one for each leg with possible further adjustments). Using the move described in Lemma A.0.4 the band can be broken into a clasp at that intersection, then $`F`$ can be tubed around the introduced annulus (this is shown in the diagram following the proof).
After these adjustments, we have a mixed graph Y-link in the complement of a Seifert surface for $`K`$. In this mixed graph Y-link, one can identify $`s`$ Y-components (one for each surplus trivalent vertex) and a possible number of other claspers that have emerged from the modifications. By an abuse of notation call this mixed Y-link $`L`$. This has been constructed to be in the complement of some Seifert surface for $`K`$, call it $`\stackrel{ˇ}{F}`$.
Use this Seifert surface to construct $`\mathrm{\Sigma }_K^p`$. Denote by $`\stackrel{~}{L}`$ the Y-link comprised of a copy of $`L`$ for each of the $`p`$ copies of $`W_{\stackrel{ˇ}{F}}`$ included into $`\mathrm{\Sigma }_K^p`$. By construction, surgery on this mixed Y-link yields $`\mathrm{\Sigma }_{K^L}^p`$.
$$\text{}\text{}$$
### 2.3. Proof of Theorem 2.0.3
Our task is to show that $`\mathrm{\Sigma }_K^p`$ and $`\mathrm{\Sigma }_{K^L}^p`$ are $`s1`$-equivalent. It is sufficient, by Lemma A.1.3, to exhibit an $`s`$-scheme relating $`\mathrm{\Sigma }_{K^L}^p`$ to $`\mathrm{\Sigma }_K^p`$. This $`s`$-scheme will be based upon the mixed Y-link $`\stackrel{~}{L}`$ in $`\mathrm{\Sigma }_K^p`$.
The appropriate selection of Y-sublinks is immediate. Order the Y-components of $`L`$, denoting them Y<sub>1</sub> through to Y<sub>s</sub>. Let $`\stackrel{~}{L}_i`$ denote the Y-sublink of $`\stackrel{~}{L}`$ comprised of the set of $`p`$ copies of Y<sub>i</sub>. These $`s`$ disjoint Y-sublinks form the required relating $`s`$-scheme: $`\{\stackrel{~}{L};\stackrel{~}{L}_1,\mathrm{},\stackrel{~}{L}_s\}`$. Recall that $`\stackrel{~}{L}_{i_1,\mathrm{},i_s}`$ denotes the link obtained by forgetting those sublinks $`\stackrel{~}{L}_j`$ for which $`i_j=1`$. By construction, the result of surgery on $`\stackrel{~}{L}_{i_1,\mathrm{},i_s}`$ is the $`p`$-fold branched cyclic cover branched over the knot obtained from $`K`$ by doing surgery on $`L_{i_1,\mathrm{},i_s}`$. If this multiplet is not $`\{0,\mathrm{},0\}`$ then this knot will be $`K`$.
### 2.4. A convenient expression
Lemma A.1.4 can be employed to give an expression for $`\mathrm{\Sigma }_{K^L}^p\mathrm{\Sigma }_K^p`$ that will be employed frequently in the sequel. Let $`\stackrel{~}{L}^{(a_1,\mathrm{},a_s)}`$ denote the sublink of $`\stackrel{~}{L}`$ that forgets, for each $`i`$ between $`1`$ and $`s`$, every copy of Y<sub>i</sub> except the copy in the included subspace $`W_{\stackrel{ˇ}{F}}^{a_i}`$.
###### Lemma 2.4.1.
$$\mathrm{\Sigma }_{K^L}^p\mathrm{\Sigma }_K^p=\underset{(a_1,\mathrm{},a_s)=(1,\mathrm{},1)}{\overset{(p,\mathrm{},p)}{}}[\mathrm{\Sigma }_K^p,\stackrel{~}{L}^{(a_1,\mathrm{},a_s)}].$$
## 3. Proofs of the lemmas
Before we prove these lemmas, we must first recollect how to compare two graph Y-links in a given three-manifold.
### 3.1. Comparing graph Y-links in three-manifolds
To an n-component graph Y-link $`L`$ in a three-manifold $`M`$ is canonically associated a trivalent graph $`D_L`$: take a trivalent vertex for every Y-component, and join tines with an edge when the associated leaves are linked in $`M`$. Note that, at this point, we are not associating any extra structure to this diagram (like cyclic orderings of edges at trivalent vertices, or orientations of edges). Let $`L^{}`$ be another $`n`$-component graph Y-link and let there exist a graph isomorphism between their associated diagrams $`D_L`$ and $`D_L^{}`$.
###### Lemma 3.1.1.
$$[M,L]=\pm [M,L^{}]𝖦_n^Y𝖬.$$
Proof. Take some surgery link presenting $`M`$. The clasper graphs presenting the Y-links $`L`$ and $`L^{}`$ can be be isotoped in $`M`$ to appear in the complement of the surgery presentation in $`S^3`$. We now “move” the position of $`L^{}`$ to that of $`L`$. That is, first isotope the vertices (discs of the clasper graph) of $`L^{}`$ (in the complement of the surgery link) to occupy the same position as the corresponding vertices of $`L`$, choosing the identification of two discs which aligns the appropriate edges. Now $`L^{}`$ can be adjusted to be in the position of $`L`$ by crossing changes between its internal bands, and by crossing changes between its internal bands and surgery components (using Lemma A.0.7). A number of half-twists will possibly be introduced as the last step (using Lemma A.0.1).
In this work, we will employ the following procedure to determine the sign of the above comparison. Begin by choosing an orientation for each Y-component of $`L`$ (that is, orient the associated thrice-punctured disc). This induces an orientation on the associated boundary components, which will appear as follows (or its mirror image):
A graph Y-link with such a choice at every Y-component will be referred to as an oriented graph Y-link, and the choice will be called an orientation for it. This choice is equivalent to a choice of cyclic ordering of the leaves at every Y-component. To fix this correspondence, for example, say that the edges of the above Y-component have a counter-clockwise cyclic ordering. Record this information on the associated graph (that is, cyclically order the edges at every trivalent vertex of it). Call this enhancement an orientation for the graph.
To each edge of an oriented graph, associated to an oriented graph Y-link, a sign will be associated: the twist of that edge. The twist is defined as the linking number of the (oriented) inner circles of the two leaves whose linking that edge represents. Let $`T(e)`$ denote the twist of the edge $`e`$.
This information can now be used to compare two graph Y-links, $`L`$ and $`L^{}`$, equipped with an isomorphism $`\varphi :D_LD_L^{}`$. Namely, choose an orientation for $`L`$, inducing an orientation on $`D_L`$. The graph isomorphism induces an orientation on $`D_L^{}`$, and hence on $`L^{}`$.
###### Lemma 3.1.2.
$$[M,L]=\left(\underset{e\text{ an edge}}{}T(e)T^{}(\varphi (e))\right)[M,L^{}]𝖦_n^Y𝖬$$
Proof. We enhance the proof of Lemma 3.1.1. Notice that the crossing changes that are made as $`L^{}`$ is “moved” to $`L`$ do not affect the twist of any edge. Introducing a half twist, however, affects the twist of the associated edge by multiplying it by minus one.
It will also prove necessary, in the sequel, to have an expression for the twist of an edge whose associated pair of leaves are linked indirectly by a chain of claspers. That is, where a leaf $`L_i`$ of a Y-component $`Y_i`$ is linked in a ball with a leaf of a clasper $`C_1`$, whose other leaf links in a ball with a leaf of another clasper $`C_2`$, and so on, until the clasper $`C_\mu `$ whose other leaf links a leaf $`L_f`$ of the other Y-component $`Y_f`$.
To calculate the twist in such a situation (that is, the twist that results from surgeries on all those claspers), begin by choosing an orientation for each clasper in the chain (that is, an orientation for the twice-punctured disc). It will appear as below (or its mirror image):
Let $`l_0`$ be the linking number of the inner circles of the annuli corresponding to the appropriate pair of leaves, one from $`L_i`$ and one from $`C_1`$; let $`l_1`$ be the linking number of the inner circles of the appropriate pair of leaves, one from $`C_1`$ and one from $`C_2`$; and so on. The following fact is straightforward:
###### Lemma 3.1.3.
The following product is well-defined (independent of the choices made in orienting the claspers) and
$$T(e)=l_0\underset{i=0}{\overset{\mu }{}}(l_i).$$
###### Remark 3.1.4.
Actually, this is more information than we will use. All we need to know is that the resulting twist is a function of the linking numbers along the chain. One can prove this fact by considering the effect on the left-most linking number of a surgery on the left-most clasper; and proceed inductively.
### 3.2. Proof of Lemma 3.2.1
###### Lemma 3.2.1.
$`(\mathrm{\Sigma }_K^p)^{\lambda _L^{}}\mathrm{\Sigma }_K^p`$ is well-defined in $`𝖦_s^Y𝖬`$.
Proof. If the Y-components of $`L`$ are oriented (according to the Section 3.1), then any graph Y-links in $`\mathrm{\Sigma }_K^p`$ resulting from the definition of $`\lambda _L^{}`$ will have isomorphic underlying vertex-oriented graphs with identical twists along the edges, so that the well-definedness follows from Lemma 3.1.2.
### 3.3. Proof of Lemma 2.0.9
The graph Y-link $`L`$ is modified in the ball, following the manoeuvres of Section 2.2, to appear in that ball in the following way:
Surgery on this mixed Y-link recovers the knot $`K^L`$. The knot $`K^L^{}`$ is recovered by surgery on the sublink that is obtained by forgetting the claspers B and C, and the knot $`K^{L^{\prime \prime }}`$ is obtained by surgery on the subgraph that is obtained by forgetting the clasper A and smoothing the half-twist. By an abuse of notation, we will refer to these mixed Y-links as $`L`$, $`L^{}`$ and $`L^{\prime \prime }`$.
The branched cyclic covering $`\mathrm{\Sigma }_{K^L}^p`$ (resp. $`\mathrm{\Sigma }_{K^L^{}}^p`$, resp. $`\mathrm{\Sigma }_{K^{L^{\prime \prime }}}^p`$) is obtained from $`\mathrm{\Sigma }_K^p`$ by performing surgery on the mixed Y-link that is comprised of a copy of $`L`$ (resp. $`L^{}`$, resp. $`L^{\prime \prime }`$) for each of the $`p`$ subspaces $`W_{\stackrel{ˇ}{F}}`$ included in $`\mathrm{\Sigma }_K^p`$. Call these mixed Y-links $`\stackrel{~}{L}`$, $`\stackrel{~}{L^{}}`$ and $`\stackrel{~}{L^{\prime \prime }}`$. The mixed Y-link $`\stackrel{~}{L}`$ will have $`p`$ copies of each of the components illustrated above. This set of components will be denoted by affixing a $`_p`$-valued superscript to the labels Y,A,B and C. We will use the same symbols to denote the corresponding components in the sublinks $`\stackrel{~}{L^{}}`$ and $`\stackrel{~}{L^{\prime \prime }}`$.
To fix this notation fully, it remains to specify which boundary component after the removal of $`\stackrel{ˇ}{F}`$ corresponds to $`\stackrel{ˇ}{F}^+`$. This choice can be specified if we declare that we are making the choice that results in $`C^i`$ linking $`B^{i+1}`$ in a ball, as follows:
Order the Y-components of $`L`$ so that the one pictured above is Y<sub>1</sub>. Let $`\stackrel{~}{L}^{(a_1,\mathrm{},a_s)}`$ denote the sublink of $`\stackrel{~}{L}`$ that is obtained by forgetting every Y-component, except the copy of Y<sub>i</sub> in $`W_{\stackrel{ˇ}{F}}^{a_i}`$, for each $`i`$. Similarly use the symbol $`\stackrel{~}{L^{}}^{(a_1,\mathrm{},a_s)}`$ (resp. $`\stackrel{~}{L^{\prime \prime }}^{(a_1,\mathrm{},a_s)}`$) for the sublink obtained from $`\stackrel{~}{L}^{(a_1,\mathrm{},a_s)}`$ by forgetting every B<sup>i</sup> and C<sup>j</sup> (resp. forgetting every A<sup>k</sup> and smoothing the half-twist in every B<sup>l</sup>).
Now, Lemma 2.4.1 indicates that
$$\mathrm{\Sigma }_{K^L}^p\mathrm{\Sigma }_K^p=\underset{(a_1,\mathrm{},a_s)=(1,\mathrm{},1)}{\overset{(p,\mathrm{},p)}{}}[\mathrm{\Sigma }_K^p,\stackrel{~}{L}^{(a_1,\mathrm{},a_s)}],$$
with identical expressions holding when $`L`$ is replaced by $`L^{}`$ or $`L^{\prime \prime }`$. The lemma in question then follows from the following equation, which we will subsequently demonstrate:
(3.3.1)
$$[\mathrm{\Sigma }_K^p,\stackrel{~}{L}^{(a_1,\mathrm{},a_s)}]=[\mathrm{\Sigma }_K^p,\stackrel{~}{L^{}}^{(a_1,\mathrm{},a_s)}][\mathrm{\Sigma }_K^p,\stackrel{~}{L^{\prime \prime }}^{(a_1,\mathrm{},a_s)}].$$
Consider the mixed Y-link $`\stackrel{~}{L}^{(a_1,\mathrm{},a_s)}`$ corresponding to some $`s`$-tuplet $`(a_1,\mathrm{},a_s)`$. Let the mixed Y-links $`R^{(a_1,\mathrm{},a_s)}`$ and $`S^{(a_1,\mathrm{},a_s)}`$ be obtained from it by modifying it in the included subspace $`W_{\stackrel{ˇ}{F}}^{a_1}`$ as follows:
,
and .
Lemma A.0.6, indicates that
$$[\mathrm{\Sigma }_K^p,\stackrel{~}{L}^{(a_1,\mathrm{},a_s)}]=[\mathrm{\Sigma }_K^p,R^{(a_1,\mathrm{},a_s)}]+[\mathrm{\Sigma }_K^p,S^{(a_1,\mathrm{},a_s)}].$$
Equation 3.3.1 follows from the identifications
(3.3.2) $`[\mathrm{\Sigma }_K^p,R^{(a_1,\mathrm{},a_s)}]`$ $`=`$ $`[\mathrm{\Sigma }_K^p,\stackrel{~}{L^{\prime \prime }}^{(a_1,\mathrm{},a_s)}],`$
(3.3.3) $`[\mathrm{\Sigma }_K^p,S^{(a_1,\mathrm{},a_s)}]`$ $`=`$ $`[\mathrm{\Sigma }_K^p,\stackrel{~}{L^{}}^{(a_1,\mathrm{},a_s)}].`$
We leave the reader to check these relations. The identifications proceed by removing claspers when they have a leaf that bounds a disk in the complement of the rest of the graph. Note that the half-twist in B$`^{a_1}`$ is smoothed at the expense of the minus one (using Lemma A.0.1).
### 3.4. Proof of Lemma 2.0.10
Choose a Seifert surface $`F`$ for $`K`$, and an orientation for $`K`$, and hence an orientation for the Seifert surface. Represent this orientation by a normal vector field. Denote the embedded trivalent graph that is associated to a position of $`L`$ that only meets this Seifert surface in transversal intersections of internal edges, by $`E_L`$.
Let $`L`$ be used again (abusing the notation) to denote the mixed Y-link that results from the modifications of this embedding with respect to this Seifert surface that are described in Section 2.2. In this case, there are no legs to worry about, so the only modifications occur when edges of the associated clasper graph intersect the Seifert surface.
$$\text{}\text{}$$
The mixed Y-link $`L`$ comprises of $`s`$ Y-components (order them Y$`{}_{1}{}^{},\mathrm{},`$Y<sub>s</sub>) and a number of other claspers, in the complement of some Seifert surface $`\stackrel{ˇ}{F}`$.
Recall that Lemma 2.4.1 indicates that
$$\mathrm{\Sigma }_{K^L}^p\mathrm{\Sigma }_K^p=\underset{(a_1,\mathrm{},a_s)=(1,\mathrm{},1)}{\overset{(p,\mathrm{},p)}{}}[\mathrm{\Sigma }_K^p,\stackrel{~}{L}^{(a_1,\mathrm{},a_s)}],$$
using notation explained there.
Let us consider one of these terms, corresponding to some $`p`$-tuplet $`(a_1,\mathrm{},a_s)`$. We will show that corresponding to each edge of $`E_L`$, to whose endpoints the two Y-components Y<sub>i</sub> and Y<sub>j</sub> are associated, there will be an equation in $`_p`$ relating $`a_i`$ and $`a_j`$ which must be satisfied in order that the corresponding term $`[\mathrm{\Sigma }_K^p,\stackrel{~}{L}^{(a_1,\mathrm{},a_s)}]`$ be non-zero.
Orient the edges of $`E_L`$. Define functions $`i`$ and $`f`$ from the edges of $`E_L`$ to the labels of the Y-components of $`L`$ so that $`Y_{i(e)}`$ corresponds to the origin of the edge $`e`$, and the Y-component $`Y_{f(e)}`$ corresponds to its end. Let $`<e,F>`$ denote the signed sum of intersections of the edge $`e`$ (oriented from $`i(e)`$ to $`f(e)`$) with $`F`$ (where a plus is an intersection in the direction of the orienting normal vector field).
###### Lemma 3.4.1.
If for the $`s`$-tuplet $`(a_1,\mathrm{},a_s)`$,
(3.4.1)
$$a_{f(e)}a_{i(e)}+<e,F>_p,$$
for some edge $`e`$ of $`E_L`$, then
$$[\mathrm{\Sigma }_K^p,\stackrel{~}{L}^{(a_1,\mathrm{},a_s)}]=0.$$
Proof
Take such a pair of Y-components in $`L`$, Y<sub>i(e)</sub> and Y<sub>f(e)</sub>. Every time the corresponding edge of $`E_L`$ intersected the Seifert surface $`F`$ it was broken into a clasp. Thus, in general, Y<sub>i(e)</sub> and Y<sub>f(e)</sub> will be linked in $`L`$ via some chain of claspers. Denote them $`C_1`$ up to $`C_\mu `$ where a leaf of $`C_1`$ links a leaf of Y<sub>i(e)</sub>, with its other leaf linking a leaf of the clasper $`C_2`$ etc. Denote the copies of each of these claspers that occur in the graph $`\stackrel{~}{L}^{(a_1,\mathrm{},a_s)}`$ by affixing a $`_p`$-valued superscript.
To show this equation, we proceed to remove copies of the $`C_j`$ when they have a leaf that bounds a disc (in the complement of the rest of the mixed Y-link is to be understood when we use this phrase). If, finally, we are left with a mixed Y-link where one of the Y-components has a leaf that bounds a disk, then we know that the contribution $`[\mathrm{\Sigma }_K^p,\stackrel{~}{L}^{(a_1,\mathrm{},a_s)}]`$ will be zero.
If the edge $`e`$ has no (resp. 1 positive, resp. 1 negative) intersection with the Seifert surface $`F`$, then the term $`[\mathrm{\Sigma }_K^p,\stackrel{~}{L}^{(a_1,\mathrm{},a_s)}]`$ will be zero unless $`a_{f(e)}=a_{i(e)}`$ (resp. $`a_{f(e)}=a_{i(e)}+1`$, resp. $`a_{f(e)}=a_{i(e)}1`$), for otherwise the Y-components will have leafs bounding discs.
Consider, then, the situation where the edge $`e`$ does intersect $`F`$ at at least two points, so that some extra claspers arise. If the first intersection is a positive (resp. negative) intersection, then we can remove all copies of $`C_1`$ except the copy $`C_1^{a_{i(e)}+1}`$ (resp. $`C_1^{a_{i(e)}1}`$). Proceeding, we remove all copies of $`C_2`$ except the one that intersects the surviving copy of $`C_1`$. And so on. One can check that as a result of this process, the surviving copy of $`C_\mu `$ will intersect precisely the Y-component $`Y_j^{a_{i(e)}+<e,F>}`$.
###### Lemma 3.4.2.
If the $`s`$-tuplet $`(a_1,\mathrm{},a_s)`$ satisfies the equations
(3.4.2)
$$a_{f(e)}=a_{i(e)}+<e,F>_p,$$
for every edge $`e`$ of $`E_L`$, then
$$[\mathrm{\Sigma }_K^p,\stackrel{~}{L}^{(a_1,\mathrm{},a_s)}]=(\mathrm{\Sigma }_K^p)^{\lambda _L^{}}\mathrm{\Sigma }_K^p𝖦_s^Y𝖬.$$
Proof.
Take the mixed Y-link $`L`$ that is obtained following Section 2.2, where every intersection with a Seifert surface is broken into a clasp. Orient each Y-component and clasper of $`L`$. This introduces an orientation on each Y-component and clasper of $`\stackrel{~}{L}^{(a_1,\mathrm{},a_s)}`$. Forget every clasper with a leaf that bounds a disc (as in the previous proof).
Note that the linking numbers between pairs of bounding circles in $`L`$ and the corresponding pairs in $`\stackrel{~}{L}^{(a_1,\mathrm{},a_s)}`$ are the same. Thus, the theorem follows from Lemma 3.1.2 and Lemma 3.1.3.
It is clear that the Equations 3.4.2 are precisely the requirement that every loop on $`E_L`$ have zero mod-$`p`$ linking number with $`K`$. Note that when at least one solution exists, there are precisely $`p`$ solutions, corresponding to the possible values of $`a_1`$.
## 4. Part 2: The relation to Vassiliev theory
In Part 2 we consider some very general questions about the relationship with Vassiliev theory of knot invariants obtained by composing finite-type three-manifold invariants with the branched cyclic covering construction.
The Vassiliev theory of finite type invariants of oriented knots in the three-sphere centers on a filtration of $`\text{K}(S^3)`$, the free abelian group generated by ambient isotopy classes of oriented knots:
$$\text{K}(S^3)𝖥_1\text{K}𝖥_2\text{K}\mathrm{}$$
Our first speculation might be that any finite-type three-manifold invariant $`v`$, composed with $`\mathrm{\Sigma }^p`$, is a finite-type invariant of knots. To begin, then, we will see that this hypothesis can be ruled out using simple arguments (for any choice of $`p`$). The method of proof also provides some other interesting information: that the function $`|H_1(\mathrm{\Sigma }_K^p)|`$ can be unbounded on a sequence of knots of strictly increasing $`n`$-triviality.
###### Theorem 4.0.1.
Fix a positive integer $`p`$. For every integer $`n>0`$, there exists a non-trivial element $`E_n𝖥_n\text{K}`$ such that
(4.0.1)
$$\mathrm{\Sigma }^p(E_n)𝖥_1^Y𝖬.$$
Proof.
Let $`\mathrm{\Omega }_2`$ be the knot obtained from surgery on the following decoration of the unkot by a 0-framed surgery link (that is, the link obtained by replacing twice-punctured discs by surgery pairs, as in Appendix A). Similarly let $`\mathrm{\Omega }_n`$ denote the obvious extension to a knot obtained from surgery on a wheel with $`n`$ spokes.
(4.0.2)
$$\text{}.$$
$`\mathrm{\Omega }_n`$ is $`(n1)`$-trivial, being obtained from surgery on a graph Y-link of degree $`n`$ (alternatively we can identify it as a Kanenobu-Ng wheel \[Kan, Ng\]).
Our theorem will follow if we can identify an $`n^{}n`$ so that $`\mathrm{\Sigma }^p(\mathrm{\Omega }_n^{})`$ is not an integral homology three sphere. Then the proof would be completed by setting $`E_n=\mathrm{\Omega }_n^{}U`$, where $`U`$ is the unknot (see Matveev’s theorem, Theorem A.0.1).
To this end we can exploit a formula for the order of the first homology of the $`p`$-fold branched cyclic covering of a knot $`K`$ in terms of its Alexander polynomial $`A_K(t)`$ \[G, F, HK\]. (The symbol $`|H_1(\mathrm{\Sigma }^p(K))|`$ below is extended to be zero if $`\mathrm{\Sigma }^p(K)`$ has non-vanishing first Betti number).
$$|H_1(\mathrm{\Sigma }^p(K))|=\underset{q=0}{\overset{p1}{}}|A_K(e^{2\pi i\frac{q}{p}})|.$$
In \[K\] we calculated (alternatively, (do the work to) identify this as one of the knots considered in \[Kan\])
$$A_{\mathrm{\Omega }_n}(t)=(1(1t)^n)(1(1t^1)^n),$$
so that
$$|H_1(\mathrm{\Sigma }^p(\mathrm{\Omega }_n))|=\underset{q=0}{\overset{p1}{}}|(1(1e^{2\pi i\frac{q}{p}})^n)|^2.$$
Define a function $`f(p,n)`$ to be $`|H_1(\mathrm{\Sigma }^p(\mathrm{\Omega }_n))|`$. We are interested in the behaviour of $`f(p,n)`$ as $`n`$ goes to infinity whilst $`p`$ is fixed.
Consider a factor $`(1\alpha (p,q)^n)`$ in the above, where $`\alpha (p,q)=1e^{2\pi i\frac{q}{p}}`$. If $`|\alpha (p,q)|<1`$, then $`lim_n\mathrm{}|1\alpha (p,q)^n|=1`$. Alternatively, if $`|\alpha (p,q)|>1`$, then $`lim_n\mathrm{}|1\alpha (p,q)^n|=\mathrm{}`$. Note, however, that there are precisely two possible choices of $`\frac{q}{p}`$ such that $`|\alpha (p,q)|=1`$; namely $`\frac{q}{p}=\pm \frac{1}{6}`$ (which will occur in $`f(p,q)`$ whenever $`6|p`$). In these cases $`\alpha (6,\pm 1)=e^{2\pi i\frac{1}{6}}`$, and $`lim_n\mathrm{}(1\alpha (6,\pm 1)^n)`$ does not exist. We will get around this by choosing a subsequence; namely $`lim_n\mathrm{}(1\alpha (6,\pm 1)^{6n+3})=2.`$
Finally, observe that there is always at least one factor $`|1\alpha (p,q)|`$ such that $`|\alpha (p,q)|>1`$. Then
$$\underset{n\mathrm{}}{lim}f(p,6n+3)=\mathrm{}.$$
See \[Gor\] for similar techniques.
We next consider a measure of the independence of these theories. The following theorem is modelled on Stanford–Trapp’s definition of an invariant which is “independent of finite-type invariants”, although note that the stated property (appears to be) substantially weaker.
It says that we may realise (some vector proportional to) any combination of primitive diagrams as the first non-vanishing term of the LMO invariant on the $`p`$-fold branched cyclic covering branched over a knot which is $`n`$-trivial, where the $`n`$-triviality is a free parameter.
According to the notation of \[LeGr\], $`\widehat{Z}^{LMO}`$ denotes the LMO invariant normalised by powers of the order of the first homology (the normalisation that satisfies the simple connect-sum formula).
###### Theorem 4.0.2.
Take positive integers $`n,p`$ and $`s`$, and a knot $`K`$ in $`S^3`$ such that the $`p`$-fold branched cyclic covering of $`S^3`$ branched over $`K`$ is a rational homology three-sphere. Then, there exists a non-zero integer $`\alpha (s,n,p)`$ such that for every $``$-linear combination of primitive degree $`s`$ uni-trivalent diagrams on an empty skeleton, $`D`$, there exists a knot $`K(s,n,p,D)`$ satisfying:
* $`K(s,n,p,D)`$ is $`n`$-equivalent to $`K`$,
* $`\widehat{Z}^{LMO}(\mathrm{\Sigma }_{K(s,n,p,D)}^p)\widehat{Z}^{LMO}(\mathrm{\Sigma }_K^p)=p\alpha (s,n,p)D+`$ terms of higher order.
See Example 4.1.2 for an illustration of the simple idea behind this theorem.
###### Remark 4.0.3.
In some sense, leaving the $`n`$-triviality of the realising knot a free parameter is the hard part of the above theorem. It is likely, for example, that such a realisation theorem (with $`n`$ fixed at $`s`$) is an extension of Garoufalidis’s investigations of the LMO invariant on branched cyclic covers over doubled knots.
Note that if the $`n`$-triviality is not required to be free above, and $`p>2`$, then a realisation with $`\alpha =1`$ exists. Can we find an improvement of the above theorem in the stated generality where $`\alpha (s,n,p)=1`$?
A corollary of Theorem 4.0.2 is:
###### Corollary 4.0.4.
Take a positive integer $`p`$. If $`v`$ is a rational valued three-manifold invariant factoring non-trivially through the LMO invariant on integral homology spheres, then $`v\mathrm{\Sigma }^p`$ is not a finite-type invariant of knots.
The proof of this corollary is indicated in Section 4.3.
### 4.1. Graph Y-links and the LMO invariant
Graph Y-links in rational homology three-spheres realise particular first non-vanishing terms of the LMO invariant. It is convenient here to cite a direct calculation \[LeGr\]. Below, let $`\widehat{Z}^{LMO}`$ denote the version of the LMO invariant normalised by powers of the order of the first homology, following the notation of \[LeGr\].
###### Theorem 4.1.1.
Let $`M`$ be a rational homology three-sphere and let $`L`$ be a connected graph Y-link in $`M`$ with some associated uni-trivalent diagram $`D_L`$ (i.e. represented by the underlying graph of $`L`$ with some choice of orientation at each vertex). Then,
(4.1.1)
$$\widehat{Z}^{LMO}(M^L)\widehat{Z}^{LMO}(M)=\pm D_L+\text{terms of higher order}.$$
Using Lemma A.0.1, one can always realise the opposite sign by introducing a half-twist.
###### Example 4.1.2.
Before we consider the details of the realisation theorem, let us consider an example to highlight the simple idea involved. Consider the following decoration of the unknot $`U`$ by a graph Y-link $`\kappa _n`$:
According to Theorem 2.0.6, $`\mathrm{\Sigma }_{U^{\kappa _n}}^p\mathrm{\Sigma }_U^p`$ lies in $`𝖥_4^Y𝖬`$, and moreover is proportional to the difference:
$$\text{}S^3𝖦_4^Y𝖬$$
To calculate the proportionality constant, we follow the procedure described in Example 2.0.13. Start by introducing a variable for each leg $`(ϵ_1,\mathrm{},ϵ_n)`$. If $`ϵ_j`$ is zero, then that leg is to be sawn off; if it is 1 then that leg is sawn off with an extra wrap around the knot (see Example 2.0.13).
The underlying trivalent graph has three cycles. Two of the cycles are unlinked from the knot for every $`n`$-tuplet $`(ϵ_1,\mathrm{},ϵ_n)`$. The linking number of the third cycle with the knot will be $`ϵ_1+\mathrm{}+ϵ_n`$.
Following Example 2.0.13, letting
$$F(n,t)=\underset{(ϵ_1,\mathrm{},ϵ_n)=(0,\mathrm{},0)}{\overset{(1,\mathrm{},1)}{}}(1)^{ϵ_1+\mathrm{}+ϵ_n}t^{ϵ_1+\mathrm{}+ϵ_n}=(1t)^n,$$
then
$$\mathrm{\Sigma }_{U^{\kappa _n}}^pS^3=\left(\underset{q=0}{\overset{p1}{}}F(n,e^{\frac{2\pi iq}{p}})\right)((S^3)^{\lambda _{\kappa _n}}S^3)𝖦_4^Y𝖬.$$
Thus
$$\widehat{Z}^{LMO}(\mathrm{\Sigma }_{U^{\kappa _n}}^p)=1+ϵ(\underset{q=0}{\overset{p1}{}}(1e^{\frac{2\pi iq}{p}})^n)\kappa +\text{terms of higher degree},$$
for some $`ϵ=\pm 1`$.
Given a knot $`K`$, an oriented trivalent diagram $`D`$, and a positive integer $`n`$, we choose a knot $`K_{(D,n)}`$ as follows. Select a graph Y-link $`L`$ in the three-sphere so that
$$\widehat{Z}^{LMO}((S^3)^LS^3)=D+\text{terms of higher degree}.$$
Then, locate $`L`$ in a ball disjoint from $`K`$ in the three-sphere. Select some edge of $`L`$, give it an orientation, and add $`n`$ legs joining that edge to the knot $`K`$ in that ball, so that each edge is added according to the following orientation (but otherwise, in any fashion):
There are many choices in this definition. Nevertheless:
###### Lemma 4.1.3.
$$\widehat{Z}^{LMO}(\mathrm{\Sigma }_{K_{(D,n)}}^p\mathrm{\Sigma }_K^p)=\left(\underset{q=0}{\overset{p1}{}}F(n,e^{\frac{2\pi iq}{p}})\right)D+\text{terms of higher degree}.$$
This follows from the same logic as in the above example. It is clear that if $`D`$ has surplus $`s`$, then $`K_{(D,n)}`$ is obtained from $`K`$ by surgery on a graph Y-link of degree $`\frac{s}{2}+n`$, and so is $`(\frac{s}{2}+n)`$-equivalent to $`K`$ \[Hab\].
### 4.2. Proof of Theorem 4.0.2
Without loss of generality, let the combination to be realised be written $`S=_{i=1}^tD_i`$. Take some positive integer $`l`$. Let $`K_{(S,l)}`$ denote a knot:
$$(\mathrm{}((K_{(D_1,l)})_{(D_2,l)})\mathrm{})_{(D_t,l)}.$$
It is clear that
$$\widehat{Z}^{LMO}(\mathrm{\Sigma }^p(K_{(S,l)}K))=\left(\underset{q=0}{\overset{p1}{}}F(l,e^{\frac{2\pi iq}{p}})\right)S+\text{terms of higher order}.$$
The number of legs to be added, $`l`$, is a free parameter, so that one can make the knot $`K_{(S,l)}`$ $`n`$-equivalent to $`K`$ for arbitrarily large $`n`$.
The remaining difficulty is covered by the following lemma.
###### Lemma 4.2.1.
$$\text{lim}_l\mathrm{}\left(\underset{q=0}{\overset{p1}{}}F(l,e^{\frac{2\pi iq}{p}})\right)0.$$
Proof
Let $`\omega =e^{\frac{2\pi i}{p}}.`$ Define a function $`f(l)=_{q=0}^{p1}F(l,\omega ^q)`$. Take some positive integer $`l`$. We will show that $`f(l^{})`$ is non-zero for at least one $`l^{}`$ such that $`ll^{}<l+p`$.
We assume that $`f(l+i)=0`$ for $`0i<p`$, and derive a contradiction. Now, we have assumed that $`f(l)=0`$, so that, letting $`T_q=(1\omega ^q)`$,
$$f(l+1)=\underset{q=0}{\overset{p1}{}}(1\omega ^q)T_q^l=\underset{q=0}{\overset{p1}{}}\omega ^qT_q^l=0.$$
Proceeding with this assumption, one finds that
$$f(l+j)=(1)^j\underset{q=0}{\overset{p1}{}}\omega ^{jq}T_q^l=0.$$
In other words, under this assumption, we have a matrix equation:
$$\left(\begin{array}{c}\hfill f(l)\\ \hfill f(l+1)\\ \hfill f(l+2)\\ \hfill .\\ \hfill .\\ \hfill .\\ \hfill (1)^{p1}f(l+p1)\end{array}\right)=\left(\begin{array}{ccccccc}1\hfill & 1\hfill & 1\hfill & .\hfill & .\hfill & .\hfill & 1\hfill \\ 1\hfill & \omega \hfill & \omega ^2\hfill & .\hfill & .\hfill & .\hfill & \omega ^{p1}\hfill \\ 1\hfill & \omega ^2\hfill & \omega ^4\hfill & .\hfill & .\hfill & .\hfill & \omega ^{2(p1)}\hfill \\ .\hfill & .\hfill & .\hfill & .\hfill & .\hfill & .\hfill & .\hfill \\ .\hfill & .\hfill & .\hfill & .\hfill & .\hfill & .\hfill & .\hfill \\ .\hfill & .\hfill & .\hfill & .\hfill & .\hfill & .\hfill & .\hfill \\ 1\hfill & \omega ^{(p1)}\hfill & \omega ^{2(p1)}\hfill & .\hfill & .\hfill & .\hfill & \omega ^{(p1)(p1)}\hfill \end{array}\right)\left(\begin{array}{c}T_0^l\hfill \\ T_1^l\hfill \\ T_2^l\hfill \\ .\hfill \\ .\hfill \\ .\hfill \\ T_{p1}^l\hfill \end{array}\right)$$
The vector on the right in clearly non-zero, and the matrix is a Vandermonde matrix, and is non-singular, so the vector on the left cannot be zero. This is our contradiction.
### 4.3. Proof of Corollary 1.0.1
This is equivalent to finding for any vector $`DA_s`$ an element $`E_{(D,n)}𝖥_n\text{K}`$ such that
$$Z^{LMO}(\mathrm{\Sigma }^p(E_{(D,n)}))=D+\text{higher order terms}.$$
This follows easily from the Theorem 4.0.2 (freely multiplying by rationals) and the observation that for knots $`K`$ and $`L`$:
(4.3.1)
$$\mathrm{\Sigma }^p(K\mathrm{\#}L)=\mathrm{\Sigma }^p(K)\mathrm{\#}\mathrm{\Sigma }^p(L).$$
## Appendix A The Goussarov-Habiro theory of finite-type invariants of three-manifolds
The theory to be recalled in this section is due, independently, to Goussarov \[G3\] and to Habiro \[Hab\]. The original finite type theory (in the setting of $`HS^3`$s) is due to Ohtsuki \[Oht\]; in that setting there are also definitions due to Garoufalidis and Levine \[GL\]; and the first definition to consider the set of all three-manifolds was due to Cochran and Melvin \[CM\]. The Goussarov-Habiro theory is defined in terms of a move on three-manifolds due to Matveev \[Mat\], closely related to a move of Murakami and Nakanishi’s, on links \[MN\]. Our summary will be concise: we expect expositions of this theory to appear in the future (\[GGP\]).
A clasper is an embedding of a standard band-summed pair of annuli into a three-manifold. The annuli are called the leaves of the clasper. This gives placement information for a two-component framed link, as follows, and a move on that clasper replaces that manifold with the manifold obtained by doing surgery on that link. The tubes in the diagram below can contain part of the surgery link presenting the three-manifold, and any other objects that may be embedded in the three-manifold:
$$\text{}\text{}$$
In this work, claspers in $`S^3`$ will be depicted via blackboard framed diagrams. That is, circles should be thickened to annuli, and the other arcs should be thickened to bands, in the plane of the diagram. The following symbols will be used to indicate half-twists of bands:
$$\text{}\text{represents}\text{}\text{and}\text{}\text{represents}\text{}.$$
A Y-component is an embedding of a standard triple of annuli band-summed into a disk, into a three-manifold. This gives placement information for a six-component surgery link. A move on that Y-component replaces that manifold with the manifold obtained by doing surgery on that link. This is Matveev’s Borromeo move and Murakami-Nakanishi’s Delta-unknotting move. It is convenient, here, to build a standard Y-component from three claspers, as follows:
$$\text{}\text{}$$
where the following convention has been used:
$$\text{}\text{is indicated by}\text{}.$$
It can be shown that this move generates an equivalence relation on the set of closed oriented three-manfolds.
###### Theorem A.0.1 (\[Mat\]).
For two closed oriented three-manifolds $`M`$ and $`N`$, there is an isomorphism $`H_1(M)H_1(N)`$ preserving the linking form on the torsion if and only if $`N`$ is obtained from $`M`$ via a finite sequence of Y-moves.
The Goussarov-Habiro theory is based on a filtration of $`𝖬`$, the free abelian group generated by homeomophism classes of closed oriented three-manifolds.
The definition uses a Y-link in a three-manifold, which is a collection of disjointly embedded Y-components. Take a $`\mu `$-component Y-link $`L`$, and let $`M_L`$ denote the three-manifold obtained by performing the moves on the Y-components of $`L`$. Define a vector $`[M,L]𝖬`$ corresponding to a pair of some $`M`$ and some Y-link $`L`$ in $`M`$ by
$$[M,L]=\underset{L^{}L}{}(1)^{\mathrm{\#}L\mathrm{\#}L^{}}M_L^{},$$
where the sum is over all Y-sublinks of $`L`$ (so that there will be $`2^\mu `$ such).
###### Definition A.0.2.
Define the subspace $`𝖥_n^Y𝖬𝖬`$ to be the subspace spanned by all vectors $`[M,L]`$ with $`L`$ an $`n`$-component Y-link.
It is clear that this filters $`𝖬`$,
$$𝖬𝖥_1^Y𝖬𝖥_2^Y𝖬𝖥_3^Y𝖬\mathrm{},$$
and that a theory of finite-type invariants of three-manifolds can be constructed by declaring an invariant to be finite-type of order $`n`$ if it vanishes on the subspace $`𝖥_{n+1}^Y𝖬`$.
The associated graded quotients $`\frac{𝖥_n^Y𝖬}{𝖥_{n+1}^Y𝖬}`$, denoting them $`𝖦_n^Y𝖬`$, are finite dimensional, with a finite spanning set of generators of the form $`[M,L]`$ where $`L`$ is a graph Y-link. These are Y-links where every leaf of every Y-component links precisely one other leaf in a ball, as follows:
or is some standard link fixed to represent some element of homology. The full class will not occur in this work, so when we refer to a graph Y-link, we will mean specifically a Y-link where every leaf meets another in a ball.
To a graph Y-link is associated a uni-trivalent graph by taking a trivalent vertex for every Y-component, and joining edges when the associated leaves link in a ball. If the graph associated to some graph Y-link is connected, call the associated graph Y-link connected.
###### Lemma A.0.3.
If $`L`$ is a connected graph Y-link, and $`L^{}`$ is a Y-sublink of $`L`$ not $`L`$, then $`M_L^{}M`$. A consequence is that
$$[M,L^{}]=0.$$
This follows from the fact that any clasper or Y-component with a leaf that bounds a disc in the complement of the rest of the Y-link, may be removed without affecting the result:
$$\text{}\text{}.$$
In the remainder of this section, we describe some aspects of the manipulation of clasper graphs that will be useful in ensuing calculations.
In this work, we will encounter situations where it is convenient to present a Y-link by decorating some other Y-link with a collection of claspers: the desired Y-link being recovered by surgery on those claspers. So, on occasions when precision is required, we will call such a link a mixed Y-link. The definition of the vector $`[M,L]`$ is extended to mixed Y-links, taking the alternating sum over the Y-components only. A sublink of a mixed link, is obtained by forgetting a number of Y-components or claspers.
Now we collect some moves that will be useful in this paper.
###### Lemma A.0.4.
The result of surgery on two mixed Y-links which differ in a ball as follows is the same.
###### Lemma A.0.5.
The result of surgery on two mixed Y-links which differ in a ball as follows is the same.
###### Lemma A.0.6.
Let $`L_A`$, $`L_B`$ and $`L_C`$ be $`n`$-component Y-links that differ in a ball as follows. The dashed part of the leaf indicates that that part follows some path in the three-manifold before returning to the ball in question.
$$[M,L_A]=[M,L_B]+[M,L_C]𝖦_n^Y𝖬.$$
###### Lemma A.0.7.
Let the Y-links $`L_A`$ and $`L_B`$ differ as follows. The tube can contain parts of surgery components or part of the rest of the graph:
$$[M,L_A]=[M,L_B]𝖦_n^Y𝖬.$$
Finally, let $`L_A`$ and $`L_B`$ differ in a ball as follows:
###### Lemma A.0.8.
(A.0.1)
$$[M,L_A]=[M,L_B]𝖦_n^Y𝖬.$$
### A.1. n-equivalence
The property of $`n`$-equivalence was introduced by Ohyama in the setting of knots \[Ohy\]; its importance for Vassiliev theory was observed by Goussarov \[G1\].
If two three-manifolds are $`n`$-equivalent then their difference lies in $`𝖥_{n+1}^Y𝖬`$ and on this pair all finite-type invariants of order less than or equal to $`n`$ agree.
###### Definition A.1.1.
A $`n+1`$-scheme for $`M`$ is a mixed Y-link $`L`$ in $`M`$ together with a set of $`n+1`$ disjoint Y-sublinks of $`L`$, $`L_1`$ up to $`L_{n+1}`$.
For some $`n+1`$-tuplet $`\{i_1,\mathrm{},i_{n+1}\}`$, where $`i_k`$ is either 0 or 1, the notation $`L_{i_1,\mathrm{},i_j}`$ is used to denote the Y-link that is obtained by forgetting those sublinks whose associated index is $`1`$.
###### Definition A.1.2.
A three-manifold $`N`$ is $`n`$-equivalent to a three-manifold $`M`$ if $`M`$ has an $`n+1`$-scheme $`\{L;L_1,\mathrm{}L_{n+1}\}`$ such that
* $`M_{L_{0,\mathrm{},0}}N`$,
* $`M_{L_{i_1,\mathrm{},i_{n+1}}}M`$ for any other multiplet.
In such a situation we will say that $`\{L;L_1,\mathrm{}L_{n+1}\}`$ is an $`n+1`$-scheme relating $`N`$ to $`M`$.
###### Lemma A.1.3.
If $`M`$ is $`n`$-equivalent to $`N`$, then
$$MN𝖥_{n+1}^Y𝖬.$$
In such a situation we also have a nice expression for $`MN`$ in the graded space $`\frac{𝖥_{n+1}^Y𝖬}{𝖥_{n+2}^Y𝖬}`$.
To introduce this expression we need to be more specific with some notation. Denote the relating $`n+1`$-scheme in $`M`$ by $`\{L;L_1,\mathrm{},L_{n+1}\}`$. Let $`\sigma (i)`$ be the function giving the number of Y-components of the Y-sublink $`L_i`$. Order the Y-components of each Y-sublink $`L_i`$. For an $`n+1`$-tuple $`(a_1,\mathrm{},a_{n+1})`$, where $`1a_i\sigma (i)`$, let $`L^{(a_1,\mathrm{},a_{n+1})}`$ be the mixed Y-link in $`M`$ obtained by forgetting all Y-components of each of these Y-sublinks, except precisely one Y-component from each $`L_i`$: that is, from $`L_i`$ choose $`a_i`$.
###### Lemma A.1.4.
$$NM=\underset{(a_1,\mathrm{},a_{n+1})=(1,\mathrm{},1)}{\overset{(\sigma (1),\mathrm{},\sigma (n+1))}{}}[M,L^{(a_1,\mathrm{},a_{n+1})}]𝖦_{n+1}^Y𝖬.$$ |
warning/0003/math-ph0003040.html | ar5iv | text | # References
THOMAS-FERMI THEORY – Sometimes called the ‘statistical theory’, it was invented by L. H. Thomas\[TH\] and E. Fermi\[EF\], shortly after Schrödinger invented his quantum-mechanical wave equation, in order to approximately describe the *electron density*, $`\rho (x)`$, $`x𝐑^3`$, and the *ground state energy*, $`E(N)`$ for a large atom or molecule with a large number, $`N`$, of electrons. Schrödinger’s equation, which would give the exact density and energy, cannot be easily handled when $`N`$ is large.
A starting point for the theory is the *TF energy functional*. For a molecule with $`K`$ nuclei of charges $`Z_i>0`$ and locations $`R_i𝐑^3(i=1,\mathrm{},K)`$, it is
$`(\rho )`$ $`:=`$ $`{\displaystyle \frac{3}{5}}\gamma {\displaystyle _{𝐑^3}}\rho (x)^{5/3}dx{\displaystyle _{𝐑^3}}V(x)\rho (x)dx`$ (1)
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{𝐑^3}}{\displaystyle _{𝐑^3}}{\displaystyle \frac{\rho (x)\rho (y)}{|xy|}}dxdy+U`$
in suitable units. Here,
$`V(x)`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{K}{}}}Z_j|xr_j|^1,`$
$`U`$ $`=`$ $`{\displaystyle \underset{1i<jK}{}}Z_iZ_j|R_iR_j|^1,`$
and $`\gamma =(3\pi ^2)^{2/3}`$. The constraint on $`\rho `$ is $`\rho (x)0`$ and $`_{𝐑^3}\rho =N`$. The functional $`\rho (\rho )`$ is convex.
The justification for this functional is this:
$``$ The first term is roughly the minimum quantum-mechanical kinetic energy of $`N`$ electrons needed to produce an electron density $`\rho `$.
$``$ The second term is the attractive interaction of the $`N`$ electrons with the $`K`$ nuclei, via the Coulomb potential $`V`$.
$``$ The third is approximately the electron-electron repulsive energy.
$``$ $`U`$ is the nuclear-nuclear repulsion and is an important constant.
The TF energy is defined to be
$$E^{\mathrm{TF}}(N)=inf\{(\rho ):\rho L^{5/3},\rho =N,\rho 0\},$$
i.e., the TF energy and density is obtained by minimizing $`(\rho )`$ with $`\rho L^{5/3}(𝐑^3)`$ and $`\rho =N`$. The Euler-Lagrange equation, called the Thomas-Fermi equation, is
$$\gamma \rho (x)^{2/3}=\left[\mathrm{\Phi }(x)\mu \right]_+,$$
(2)
where $`[a]_+`$ = $`\mathrm{max}\{0,a\}`$, $`\mu `$ is some constant (Lagrange multiplier) and $`\mathrm{\Phi }`$ is the TF potential:
$$\mathrm{\Phi }(x)=V(x)_{𝐑^3}|xy|^1\rho (y)dy.$$
(3)
The following essential mathematical facts about the TF equation were established by E.H. Lieb and B. Simon \[LS\] (cf. the review article \[EL\]).
1. There is a density $`\rho _N^{\mathrm{TF}}`$ that minimizes $`(\rho )`$ if and only if $`NZ:=_{j=1}^KZ_j`$. This $`\rho _N^{\mathrm{TF}}`$ is unique and it satisfies the TF equation (2) for some $`\mu 0`$. Every positive solution, $`\rho `$, of (2) is a minimizer of (1) for $`N=\rho `$. If $`N>Z`$ then $`E^{\mathrm{TF}}(N)=E^{\mathrm{TF}}(Z)`$ and any minimizing sequence converges weakly in $`L^{5/3}(𝐑^3)`$ to $`\rho _Z^{\mathrm{TF}}`$.
2. $`\mathrm{\Phi }(x)0`$ for all $`x`$. (This need not be so for the real Schrödinger $`\rho `$.)
3. $`\mu =\mu (N)`$ is a strictly monotonically decreasing function of $`N`$ and $`\mu (Z)=0`$ (the neutral case). $`\mu `$ is the chemical potential, namely
$$\mu (N)=\frac{E^{\mathrm{TF}}(N)}{N}.$$
$`E^{\mathrm{TF}}(N)`$ is a strictly convex, decreasing function of $`N`$ for $`NZ`$ and $`E^{\mathrm{TF}}(N)=E^{\mathrm{TF}}(Z)`$ for $`NZ`$. If $`N<Z`$, $`\rho _N^{\mathrm{TF}}`$ has compact support.
When $`N=Z`$, (2) becomes $`\gamma \rho ^{2/3}=\mathrm{\Phi }`$. By applying the Laplacian $`\mathrm{\Delta }`$ to both sides we obtain
$$\mathrm{\Delta }\mathrm{\Phi }(x)+4\pi \gamma ^{3/2}\mathrm{\Phi }(x)^{3/2}=4\pi \underset{j=1}{\overset{K}{}}Z_j\delta (xR_j),$$
which is the form in which the TF equation is usually stated (but it is valid only for $`N=Z`$).
An important property of the solution is Teller’s theorem \[ET\] (proved rigorously in \[LS\]) which implies that the TF molecule is always unstable, i.e., for each $`NZ`$ there are $`K`$ numbers $`N_j(0,Z_j)`$ with $`_jN_j=N`$ such that
$$E^{\mathrm{TF}}(N)>\underset{j=1}{\overset{K}{}}E_{\mathrm{atom}}^{\mathrm{TF}}(N_j,Z_j),$$
(4)
where $`E_{\mathrm{atom}}^{\mathrm{TF}}(N_j,Z_j)`$ is the TF energy with $`K=1,Z=Z_j`$ and $`N=N_j`$. The presence of $`U`$ in (1) is crucial for this result. The inequality is strict. Not only does $`E^{\mathrm{TF}}`$ decrease when the nuclei are pulled infinitely far apart (which is what (4) says) but any dilation of the nuclear coordinates $`(R_j\mathrm{}R_j,\mathrm{}>1)`$ will decrease $`E^{\mathrm{TF}}`$ in the neutral case (positivity of the pressure\[EL\]\[BL\]. This theorem plays an important role in the stability of matter.
An important question concerns the connection between $`E^{\mathrm{TF}}(N)`$ and $`E^\mathrm{Q}(N)`$, the ground state energy (= infimum of the spectrum) of the Schrödinger operator, $`H`$, it was meant to approximate.
$$H=\underset{i=1}{\overset{N}{}}\left[\mathrm{\Delta }_i+V(x_i)\right]+\underset{1i<jN}{}|x_ix_j|^1+U,$$
which acts on the antisymmetric functions $`^NL^2(𝐑^3;𝐂^2)`$ (i.e., functions of space and spin). It used to be believed that $`E^{\mathrm{TF}}`$ is asymptotically exact as $`N\mathrm{}`$ but this is not quite right; $`Z\mathrm{}`$ is also needed. Lieb and Simon \[LS\] proved that if we fix $`K`$ and $`Z_j/Z`$ and we set $`R_j=Z^{1/3}R_j^0`$, with fixed $`R_j^0𝐑^3`$, and set $`N=\lambda Z`$, with $`0\lambda <1`$ then
$$\underset{Z\mathrm{}}{lim}E^{\mathrm{TF}}(\lambda Z)/E^\mathrm{Q}(\lambda Z)=1.$$
(5)
In particular, a simple change of variables shows that $`E_{\mathrm{atom}}^{\mathrm{TF}}(\lambda ,Z)=Z^{7/3}E_{\mathrm{atom}}^{\mathrm{TF}}(\lambda ,1)`$ and hence the true energy of a large atom is asymptotically proportional to $`Z^{7/3}`$. Likewise, there is a well-defined sense in which the quantum mechanical density converges to $`\rho _N^{\mathrm{TF}}`$ (cf. \[LS\]).
The TF density for an atom located at $`R=0`$, which is spherically symmetric, scales as
$`\rho _{\mathrm{atom}}^{\mathrm{TF}}(x;N=\lambda Z,Z)`$ $`=`$
$`Z^2\rho _{\mathrm{atom}}^{\mathrm{TF}}(`$ $`Z^{1/3}x;`$ $`N=\lambda ,Z=1).`$
Thus, a large atom (i.e., large $`Z`$) is smaller than a $`Z=1`$ atom by a factor $`Z^{1/3}`$ in radius. Despite this seeming paradox, TF theory gives the correct electron density in a real atom — so far as the bulk of the electrons is concerned — as $`Z\mathrm{}`$
Another important fact is the large $`|x|`$ asymptotics of $`\rho _{\mathrm{atom}}^{\mathrm{TF}}`$ for a neutral atom. As $`|x|\mathrm{}`$,
$$\rho _{\mathrm{atom}}^{\mathrm{TF}}(x,N=Z,Z)\gamma ^3(3/\pi )^3|x|^6,$$
independent of $`Z`$. Again, this behavior agrees with quantum mechanics — on a length scale $`Z^{1/3}`$, which is where the bulk of the electrons are to be found.
In light of the limit theorem (5), Teller’s theorem can be understood as saying that as $`Z\mathrm{}`$ the quantum mechanical binding energy of a molecule is of lower order in $`Z`$ than the total ground state energy. Thus, Teller’s theorem is not a defect of TF theory (although it is sometimes interpreted that way) but an important statement about the true quantum mechanical situation.
For finite $`Z`$ one can show, using the Lieb-Thirring inequality \[LT\] and the Lieb-Oxford inequality \[LO\], that $`E^{\mathrm{TF}}(N)`$, with a modified $`\gamma `$, gives a lower bound to $`E^\mathrm{Q}(N)`$.
Several ‘improvements’ to Thomas-Fermi theory have been proposed, but none have a fundamental significance in the sense of being ‘exact’ in the $`Z\mathrm{}`$ limit. The von Weizsäcker correction consists in adding a term
$$(\mathrm{const}.)_{𝐑^3}|\sqrt{\rho (\mathrm{x})}|^2\mathrm{dx}$$
to $`(\rho )`$. This preserves the convexity of $`(\rho )`$ and adds (const.)$`Z^2`$ to $`E^{\mathrm{TF}}(N)`$ when $`Z`$ is large. It also has the effect that the range of $`N`$ for which there is a minimizing $`\rho `$ is extend from \[0,Z\] to \[0,Z + (const.) K\].
Another correction, the Dirac exchange energy, is to add
$$(\mathrm{const}.)_{𝐑^3}\rho (x)^{4/3}\mathrm{d}x$$
to $`(\rho )`$. This spoils the convexity but not the range \[0,Z\] for which a minimizing $`\rho `$ exists cf. \[LS\] for both of these corrections.
When a uniform external magnetic field $`B`$ is present, the operator $`\mathrm{\Delta }`$ in $`H`$ is replaced by
$$|i+A(x)|^2+\sigma B(x),$$
with curl $`A=B`$ and $`\sigma `$ denoting the Pauli spin matrices. This leads to a modified TF theory that is asymptotically exact as $`Z\mathrm{}`$, but the theory depends on the manner in which $`B`$ varies with $`Z`$. There are five distinct regimes and theories: $`BZ^{4/3},BZ^{4/3},Z^{4/3}BZ^3,BZ^3,Z^3`$. These theories \[LSY1\]\[LSY2\] are relevant for neutron stars. Another class of TF theories with magnetic fields is relevant for electrons confined to two-dimensional geometries (quantum dots) \[LSY3\]. In this case there are three regimes. A convenient review is \[LSY4\].
Still another modification of TF theory is its extension from a theory of the ground states of atoms and molecules (which corresponds to zero temperature) to a theory of positive temperature states of large systems such as stars (cf. \[JM\]\[WT\]).
Elliott H. Lieb
Departments of Mathematics and Physics
Princeton University
©1998 by Elliott H. Lieb |
warning/0003/cond-mat0003179.html | ar5iv | text | # Self organization in a minority game: the rôle of memory and a probabilistic approach
## Abstract
A minority game whose strategies are given by probabilities $`p`$, is replaced by a ‘simplified’ version that makes no use of memories at all. Numerical results show that the corresponding distribution functions are indistinguishable. A related approach, using a random walk formulation, allows us to identify the origin of correlations and self organization in the model, and to understand their disappearance for a different strategy’s update rule, as pointed out in a previous work.
Keywords: Minority game, Organization, Evolution
The minority game (MG), introduced by Challet and Zhang , addresses the problem of self organization of a population without direct interactions between its members, but with a feedback mechanism related with its collective behavior. Each person (dubbed ‘agent’, because of the use of the model in economic problems) has to choose from a simple alternative, without knowing what the other agents will do. At the end of the game, there are two groups, one for each alternative: agents belonging to the smaller group (the ‘minority’) will be the winners. Feedback is established by a reward system for winners and losers.
Different methods to choose one or the other alternative give rise to different versions of the model. It is usual to refer to these methods as ‘strategies’. In the following, we will use the model proposed by Johnson et.al. . In this version, each agent knows beforehand the previous $`m`$ outcomes (a ‘history’), of the game, as well as the ‘next move’ of the most recent occurences of all $`2^m`$ possibilities. What is distinctive of Johnson’s et.al. formulation is the assignement to each of the $`N`$ agents of a single number as its strategy, $`p_j`$ $`(0p_j1)`$: given a history, each agent will either choose the same outcome as that stored in the memory, with probability $`p_j`$, or will choose the opposite with probability $`q_j=1p_j`$ $`(1jN).`$
Winners, i.e. those in the minority group, will gain $`G`$ points ($`G1`$ in Ref. ); those in the majority group lose a point. Strategies can be modified, following the evolution of the game: if the number of points of an agent is below a threshold value $`d0`$, his ‘account’ is reset to zero, and he/she gets a new strategy, whose value $`p^{}`$ is chosen with an uniform probability from the interval $`(pR/2,p+R/2)`$, where $`|R|2;`$ in what follows we will use the simpler notation $`pp^{}=p\pm \mathrm{\Delta }p.`$ Whenever necessary, we used reflective boundary conditions.
Johnson et.al. have shown that, as a result of correlations, the system self organizes mainly in two distinct groups of agents, with extreme values of their strategies: $`p0`$ or $`p1.`$ The frequency distribution $`P(p)`$ is shown in Fig.1. Moreover, they also found that their results did not change if memories are not updated, or even if they are randomly chosen; a similar result was obtained by Cavagna in relation with the work of Challet and Zhang.
In this work, we show that a ‘simplified’ version of the model, making no use of the memories, also displays self organization, and the resulting distribution $`P(p)`$ is indistinguishable from the original model, for all odd $`N3`$. Another approach (called ‘probabilistic’ in what follows) serves to present a rather detailed interpretation of this MG, explaining the mechanisms establishing correlations between the agents, and their relations with the rules of the game. This also allows us to explain the smallness of correlations found in a previous publication, where a different update rule of the strategies was used .
In our ‘simplified’ version of the model, agent $`j`$ chooses one of the two options (option “1”, say) with probability $`p_j,`$ or the other (option “0”) with probability $`q_j=1p_j`$, without making recourse to any history at all. All other rules, like the determination of the minority, the reward system, the upgrade of strategies $`p`$, $`etc.`$ are the same as before. In Fig. 1 we compare results obtained with this version and with the original formulation, for $`N=101`$, $`d=4`$, $`R=0.2`$ . A single realization of the game involved $`n_t=10^5`$ time steps, and the distribution was averaged over $`n_s=10^4`$ samples. As it was already mentioned, both results are indistinguishable.
Even if one lets the memories totally outside of the game, there still are many parameters in the model (the variables $`G,d,N,R)`$, and it is worth to see how it depends on these variables. All our simulations were made in such a way that we can assure that the expected numerical fluctuations for the initial, uniform distribution, were small. To this end, we requested that the standard deviation $`1/\sqrt{Nn_s/c}`$ be no greater than 0.02 . We always used c=100 channels in $`p`$ (each of width 0.01). Moreover, to be able to compare results for different values of $`N`$, we have normalized all our data so that $`P(p)𝑑p=1.`$
We collected data for $`N=3,5,7,9,11,21,51`$ and $`101`$, $`d=0,2,4`$ and $`5`$, for fixed values of $`R=0.2`$ , and $`G=1`$. Figure 2 shows the density values for both extrema, $`P(0),P(1)`$ $`vs`$ $`1/N`$. It is rather clear that $`N=101`$ is already near the ‘thermodynamic’ limit $`(N\mathrm{})`$. Also, there is only a moderated dependence on the value of the threshold $`d.`$
Results for the ‘simplified’ version of the model show that we need not to consider memories, and constitute the starting point for the following development.
Our second, probabilistic, approach originated in the observation that the change of strategies, $`pp^{},`$ can be thought of as the movement of points in $`p`$ space, suggesting the formulation of the model as some kind of generalized random walk (RW). It is of some interest to mention that there is a formal similarity of $`P(p)`$ with a property of a RW; specifically, this is the case of the expresion for the probability of the last visit to the origin, $`x=0,`$ for such a system . In the simplest version of a RW, a point moves regularly with a constant step $`S`$, randomness being only related with the sign of $`S.`$ Applied to the case of $`N`$ points randomly distributed in a one dimensional (1D) box, it produces a uniform distribution in space, $`P(x)=P_0`$. In a more general case, the probability $`m(x)`$ to move a point will be, in general, a function of position $`x`$. This would be the case, for instance, of a system with absorbing walls, or a gas with a thermal gradient. Moreover, there are situations where the movement needs not to be regular (in time). The stationary state would be established when $`m(x)P(x)=constant`$.
We will consider the movement of $`N`$ agents in a 1D space of probabilities, $`p,`$ with a variable step $`\mathrm{\Delta }p`$, as it was already made in Johnson’s formulation of the model. On the other hand, we will use the set $`\{p_i,G,d\}`$ to decide if an agent moves or not. It is possible to write $`\mu _i,`$ the probability of agent $`i`$ to be in the majority, in terms of $`\{p_j\}`$. If $`G=d=0`$, then it is simple to see that $`\mu _i`$ is equal to the probability $`m(p_i)`$ to move the agent, $`i.e.`$ $`p_ip_i^{}.`$ In the following we will refer to $`m`$ as the mobility. Note that, in general, the actual value of $`m`$ will depend on all the $`p_j;`$ we choose to write $`m(p_i)`$ for agent $`i`$, to emphasize that its value changes with the position of the agent in $`p`$-space. We cannot find a closed expresion for the mobility if $`G`$ and $`d`$ are $`0`$; nevertheless, as the mobility must follow in general the behavior of $`\mu `$, we still can use it to describe the system’s behavior.
Let us consider the case $`N=3`$, for simplicity. The system is characterized by the strategies $`p_1,p_2,p_3`$ (and $`q_i=1p_i`$). The probabilities for every agent to be in the minority are
$$\lambda _1=p_1q_2q_3+q_1p_2p_3,\lambda _2=q_1p_2q_3+p_1q_2p_3,\lambda _3=q_1q_2p_3+p_1p_2q_3$$
(1)
and the corresponding probabilities of being in the majority are $`\mu _i=1\lambda _i`$.
These expressions can be easily generalized to the case of $`N=2n+1`$ agents.
Introducing $`x_i=p_i/q_i`$, $`y_i=1/x_i`$, $`U=_{j=1}^{j=N}p_i`$, $`Q=_{j=1}^{j=N}q_i`$ it is
$`\lambda _i`$ $`=`$ $`x_iQ(1+{\displaystyle \underset{k_2}{}}x_{k_2}+{\displaystyle \underset{k_2,k_3}{}}x_{k_2}x_{k_3}+\mathrm{}+{\displaystyle \underset{k_2,k_3,..,k_n}{}}x_{k_2}x_{k_3}\mathrm{}x_{k_n})`$ (3)
$`+y_iU(1+{\displaystyle \underset{k_2}{}}y_{k_2}+{\displaystyle \underset{k_2,k_3}{}}y_{k_2}y_{k_3}+\mathrm{}+{\displaystyle \underset{k_2,k_3,..,k_n}{}}y_{k_2}y_{k_3}\mathrm{}y_{k_n})`$
Equations(1)-(3) show that $`\mu _i`$ depends on all the strategies, $`\{p_j\}`$. In other words, it illustrates that the origin of correlations in the distribution $`P(p)`$ can be traced back to the rules defining the minority game.
We have made numerical simulations based on Eqs.(1)-(3), for $`N=3,5,7,9`$ and $`11`$: at every step of the game, an agent gains (loses) one point with probability $`\lambda `$ $`(\mu )`$. This procedure makes explicit how agents relate their behaviors through the strategies. Figure 3 shows results for $`N=`$ $`11`$, together with the corresponding results from our ‘simplified’ version. The similarity between the results obtained with both methods it is remarkable; the small differences seen in the very narrow region near both extrema $`(p0,p1),`$ are probably due to the noise attributable to our simulations. This is not to say that this approach is identical to the original model. In fact, we can only expect the probabilistic approach to be equivalent to the ‘simplified’ version in a statistical sense, but it is not possible to compare both methods at each time step. To illustrate this difference, notice that in this formulation one will accept some outcomes which are not allowed in the original MG; thus, for instance, as agents win a point with probability $`\lambda _i,`$ there exists a finite probability, $`W_v(N)0,`$ that the majority of the agents can win a point, in an apparent violation of the basic rule of the game (indeed, it can even happen that all agents are simultaneously winners). If $`N=3,`$
$$W_v(3)=\lambda _1\lambda _2\lambda _3+\mu _1\lambda _2\lambda _3+\lambda _1\mu _2\lambda _3+\lambda _1\lambda _2\mu _3$$
(4)
and similar relations for all $`N.`$
Using Eq.(3), and the generalization of Eq.(4), we calculate for a uniform distribution that $`W_v(3)0.14`$, while it can be estimated that $`W_v(\mathrm{})0.25`$. Similar values are obtained for non uniform distributions with a shape analogous to that of Fig.3.
Figure 4 has results for $`N=3,`$ $`d=0,1`$ and $`G=0,1.`$ In this case, results obtained with the ‘simplified’ version (not shown here) are indistinguishable from those coming from Eq.(1). If $`G=0,`$ the ensuing self-organization is small; we have verified the same type of behavior for all $`N11`$. On the contrary, if $`G=1`$ self-organization is very important. In both cases, $`d`$ has a smaller influence on the behavior of the system. We have included results for $`N=3,d=G=0`$ in Fig. 4, because in this case $`m(p)=\mu (p),`$ and it is possible to get a clear picture of the resulting (small) organization. Assume agents are numbered so that $`p_1<p_2<p_3`$, and consider the case $`p_1<1/2,`$ $`p_3>1/2`$. It follows from Eq.(1) that $`m_2>m_1,m_3`$. This describes a situation where agents near $`p0`$ and $`p1`$ have a tendency to remain in their positions, while the agent in between moves more frequently. Eventually this will change if, as a result of the movement, either $`p_2<p_1`$ or $`p_2>p_3,`$ increasing the accumulation of agents near $`p=0`$ and $`p=1`$ . Incidentally, we can use this picture to understand why the use of a different updating rule can destroy self-organization, as previously reported by one of us . Assume the same situation as before, $`i.e.`$ an agent with $`p_21/2`$ and high mobility, and two agents near $`p_10`$ and $`p_31`$, with smaller mobility. If now strategies are updated using $`p^{}=1p\pm \mathrm{\Delta }p,`$ as the agent with $`p_10`$ moves, he will go near $`p_11.`$ It is easy to see from Eq.(1) that the corresponding mobilities are $`m_1m_31>m_2.`$ In words, correlations are broken in a single step, so that only a very tiny indication of self-organization remains. It is worthwhile to mention that this type of evolution can no longer be sensibly described as a random walk.
The probabilistic approach provides a natural starting point for an analytic formulation of this model. A sample is described by a point in an $`N`$ dimensional space of probabilities $`p_j`$; the set of all the $`n_s`$ points can be thougth of as a non interacting gas that evolves towards a stationary distribution, driven by the rules of the game. We are presently working on the implementation of these ideas .
In summary, we have made mainly two contributions to the knowledge of this version of the MG: $`(i)`$ we have shown the irrelevance of memory in the resulting self organization of the system; in this respect, therefore, this version of the MG behaves differently than that of , where it has been claimed that all agents need to receive the same information, whether it be true or false, to be able to self organize; $`(ii)`$ our probabilistic formulation proved to be a very good approximation to the model and, equally important, allows us to understand in detail how the game’s rules establish correlations between the agents.
We thank F. Parisi for useful discussions and comments. E.B. was partially supported by CONICET of Argentina, PICT-PMT0051; H.C. was partially supported by EC Grant ARG/B7-3011/94/27, Contract 931005 AR. |
warning/0003/hep-th0003126.html | ar5iv | text | # Relativistic wave equations with fractional derivatives and pseudo-differential operators
## 1 Introduction
The relativistic covariant wave equations represent an intersection of ideas of the theory of relativity and quantum mechanics. The first and best known relativistic equations, the Klein-Gordon and particularly Dirac equation, belong to the essentials, which our present understanding of the microworld is based on. In this sense it is quite natural, that the searching for and the study of the further types of such equations represent a field of stable interest. For a review see e.g. and citations therein. In fact, the attention has been paid first of all to the study of equations corresponding to the higher spins $`(s1)`$ and to the attempts to solve the problems, which have been revealed in the connection with these equations, e.g. the acausality due to external fields introduced by the minimal way.
In this paper we study the class of equations obtained by the ’factorization’ of the D’Alambertian operator, i.e. by a generalization of the procedure, by which the Dirac equation is obtained. As the result, from each degree of extraction $`n`$ we get a multi-component equation, hereat the case $`n=2`$ corresponds to the Dirac equation. However the equations for $`n>2`$ differ substantially from the cases $`n=1,2`$ since they contain fractional derivatives (or pseudo-differential operators), so in the effect their nature is non-local.
In the first part (Sec. 2), the generalized algebras of the Pauli and Dirac matrices are considered and their properties are discussed, in particular their relation to the algebra of the $`SU(n)`$ group. The second, main part (Sec. 3) deals with the covariant wave equations generated by the roots of the D’Alambertian operator, these roots are defined with the use of the generalized Dirac matrices. In this section we show the explicit form of the equations, their symmetries and the corresponding transformation laws. We also define the scalar product and construct the corresponding Green functions. The last section (Sec. 4) is devoted to the summary and concluding remarks.
Let us remark, the application of the pseudo-differential operators in the relativistic equations is nothing new. The very interesting aspects of the scalar relativistic equations based on the square root of the Klein-Gordon equation are pointed out e.g. in the papers -. Recently, an interesting approach for the scalar relativistic equations based on the pseudo-differential operators of the type $`f(\mathrm{})`$ has been proposed in the paper . One can mention also the papers , in which the square and cubic roots of the Dirac equation were studied in the context of supersymmetry. The cubic roots of the Klein-Gordon equation were discussed in the recent papers , .
It should be observed, that our considerations concerning the generalized Pauli and Dirac matrices (Sec. 2) have much common with the earlier studies related to the generalized Clifford algebras (see e.g. - and citation therein) and with the paper , even if our starting motivation is rather different.
## 2 Generalized algebras of Pauli and Dirac matrices
Anywhere in the next by the term matrix we mean the square matrix $`n\times n,`$ if not stated otherwise. Considerations of this section are based on the matrix pair introduced as follows.
###### Definition 1
For any $`n2`$ we define the matrices
$$S=\left(\begin{array}{cccccc}0& & & & & 1\\ 1& & & & & \\ & 1& & & & \\ & & & & & \\ & & & & & \\ & & & & 1& 0\end{array}\right),$$
(2.1)
$$T=\left(\begin{array}{cccccc}1& & & & & \\ & \alpha & & & & \\ & & \alpha ^2& & & \\ & & & & & \\ & & & & & \\ & & & & & \alpha ^{n1}\end{array}\right)$$
(2.2)
where $`\alpha =\mathrm{exp}(2\pi i/n)`$ and in the remaining empty positions are zeros.
###### Lemma 2
Matrices $`X=S,T`$ satisfy the following relations
$$\alpha ST=TS,$$
(2.3)
$$X^n=I,$$
(2.4)
$$XX^{}=X^{}X=I,$$
(2.5)
$$detX=(1)^{n1},$$
(2.6)
$$\text{Tr }X^k=0,k=1,2\mathrm{}n1,$$
(2.7)
where $`I`$ denotes the unit matrix.
*Proof:*
All the relations easily follow from the Definition 1.
###### Definition 3
Let $`𝒜`$ be some algebra on the field of complex numbers, $`(p,m)`$ be a pair of natural numbers, $`X_1,X_2,\mathrm{},X_m𝒜`$ and $`a_1,a_2,\mathrm{},a_m`$ $`C`$. The $`pth`$ power of the linear combination can be expanded:
$$\left(\underset{k=1}{\overset{m}{}}a_kX_k\right)^p=\underset{p_j}{}a_1^{p_1}a_2^{p_2}\mathrm{}a_m^{p_m}\{X_1^{p_1},X_2^{p_2},\mathrm{},X_m^{p_m}\};p_1+\mathrm{}+p_m=p,$$
where the symbol $`\{X_1^{p_1},X_2^{p_2},\mathrm{},X_m^{p_m}\}`$ represents the sum of the all possible products created from elements $`X_k`$ in such a way that each product contains element $`X_k`$ just $`p_ktimes.`$ This symbol we shall call combinator.
###### Example 4
$$\{X,Y\}=XY+YX,$$
(2.8)
$$\{X,Y^2\}=XY^2+YXY+Y^2X,$$
(2.9)
$$\{X,Y,Z\}=XYZ+XZY+YXZ+YZX+ZXY+ZYX.$$
(2.10)
Now, we shall prove some useful identities.
###### Lemma 5
Let us assume $`z`$ is a complex variable, $`p,r0`$ and denote
$$q_p(z)=(1z)(1z^2)\mathrm{}(1z^p),q_0(z)=1,$$
(2.11)
$$F_{rp}(z)=\underset{k_p=0}{\overset{r}{}}\mathrm{}\underset{k_2=0}{\overset{k_3}{}}\underset{k_1=0}{\overset{k_2}{}}z^{k_1}z^{k_2}\mathrm{}z^{k_p},$$
(2.12)
$$G_p(z)=\underset{k=0}{\overset{p}{}}\frac{z^k}{q_{pk}(z^1)q_k(z)},$$
(2.13)
$$H_p(z)=\underset{k=0}{\overset{p}{}}\frac{1}{q_{pk}(z^1)q_k(z)}.$$
(2.14)
Then the following identities hold for $`z0,z^j1;j=1,2,\mathrm{},p`$:
$$q_p(z)=(1)^pz^{\frac{p(p+1)}{2}}q_p(z^1),$$
(2.15)
$$G_p(z)=0,$$
(2.16)
$$H_p(z)=1,$$
(2.17)
$$F_{rp}(z)=\underset{k=0}{\overset{p}{}}\frac{z^{kr}}{q_{pk}(z)q_k(z^1)}$$
(2.18)
and in particular for $`z^{p+r}=1`$
$$F_{rp}(z)=0.$$
(2.19)
*Proof:*
1) Relation (2.15) follows immediately from definition (2.11):
$$q_r(z)=(1z)(1z^2)\mathrm{}(1z^r)=zz^2\mathrm{}z^r(z^11)\mathrm{}(z^r1)$$
$$=(1)^rz^{\frac{r(r+1)}{2}}q_r(z^1).$$
2) Relations (2.16), (2.17):
First, if we invert the order of adding in the relations (2.13), (2.14) making substitution $`j=pk,`$ then
$$G_p(z)=\underset{k=0}{\overset{p}{}}\frac{z^k}{q_{pk}(z^1)q_k(z)}=z^p\underset{j=0}{\overset{p}{}}\frac{z^j}{q_j(z^1)q_{pj}(z)}=z^pG_p(z^1),$$
(2.20)
$$H_p(z)=\underset{k=0}{\overset{p}{}}\frac{1}{q_{pk}(z^1)q_k(z)}=\underset{j=0}{\overset{p}{}}\frac{1}{q_j(z^1)q_{pj}(z)}=H_p(z^1).$$
(2.21)
Now, let us calculate
$$H_p(z)H_{p1}(z)=\underset{k=0}{\overset{p}{}}\frac{1}{q_{pk}(z^1)q_k(z)}\underset{k=0}{\overset{p1}{}}\frac{1}{q_{p1k}(z^1)q_k(z)}$$
(2.22)
$$=\frac{1}{q_p(z)}+\underset{k=0}{\overset{p1}{}}\frac{1}{q_{pk}(z^1)q_k(z)}\underset{k=0}{\overset{p1}{}}\frac{1}{q_{pk1}(z^1)q_k(z)}$$
$$=\frac{1}{q_p(z)}+\underset{k=0}{\overset{p1}{}}\frac{1(1z^{kp})}{q_{pk}(z^1)q_k(z)}=\underset{k=0}{\overset{p}{}}\frac{z^{kp}}{q_{pk}(z^1)q_k(z)}=G_p(z^1).$$
The last relation combined with Eq. (2.21) implies
$$G_p(z^1)=G_p(z),$$
(2.23)
which compared with Eq. (2.20) gives
$$G_p(z^1)=0;z0,z^j1,j=1,2,\mathrm{}p.$$
(2.24)
So the identity (2.16) is proved. Further, relations (2.24), (2.22) imply
$$H_p(z)H_{p1}(z)=0,$$
(2.25)
therefore
$$H_p(z)=H_{p1}(z)=\mathrm{}=H_0(z)=1$$
(2.26)
and the identity (2.17) is proved as well.
3) The relation (2.18) can be proved by the induction, therefore first let us assume $`p=1`$, then its l.h.s. reads
$$\underset{k_1=0}{\overset{k_2}{}}z^{k_1}=\frac{1z^{k_2+1}}{1z}$$
and r.h.s. gives
$$\frac{1}{q_1(z)}+\frac{z^{k_2}}{q_1(z^1)}=\frac{1}{1z}+\frac{z^{k_2}}{1z^1}=\frac{1z^{k_2+1}}{1z},$$
so for $`p=1`$ the relation is valid. Now let us suppose the relation holds for $`p`$ and calculate the case $`p+1`$
$$\underset{k_{p+1}=0}{\overset{k_{p+2}}{}}\mathrm{}\underset{k_2=0}{\overset{k_3}{}}\underset{k_1=0}{\overset{k_2}{}}z^{k_1}z^{k_2}\mathrm{}z^{k_{p+1}}=\underset{k_{p+1}=0}{\overset{k_{p+2}}{}}z^{k_{p+1}}\mathrm{}\underset{k_2=0}{\overset{k_3}{}}\underset{k_1=0}{\overset{k_2}{}}z^{k_1}z^{k_2}\mathrm{}z^{k_p}$$
$$=\underset{k_{p+1}=0}{\overset{k_{p+2}}{}}z^{k_{p+1}}\underset{k=0}{\overset{p}{}}\frac{z^{kk_{p+1}}}{q_{pk}(z)q_k(z^1)}=\underset{k=0}{\overset{p}{}}\frac{1}{q_{pk}(z)q_k(z^1)}\underset{k_{p+1}=0}{\overset{k_{p+2}}{}}z^{(k+1)k_{p+1}}$$
$$=\underset{k=0}{\overset{p}{}}\frac{1}{q_{pk}(z)q_k(z^1)}\frac{1z^{(k+1)(k_{p+2}+1)}}{1z^{k+1}}=\underset{k=0}{\overset{p}{}}\frac{z^{k1}z^{(k+1)k_{p+2}}}{q_{pk}(z)q_k(z^1)(z^{k1}1)}$$
$$=\underset{k=0}{\overset{p}{}}\frac{z^{(k+1)k_{p+2}}z^{k1}}{q_{pk}(z)q_{k+1}(z^1)}=\underset{k=1}{\overset{p+1}{}}\frac{z^{kk_{p+2}}z^k}{q_{p+1k}(z)q_k(z^1)}=\underset{k=0}{\overset{p+1}{}}\frac{z^{kk_{p+2}}z^k}{q_{p+1k}(z)q_k(z^1)}$$
$$=\underset{k=0}{\overset{p+1}{}}\frac{z^{kk_{p+2}}}{q_{p+1k}(z)q_k(z^1)}\underset{k=0}{\overset{p+1}{}}\frac{z^k}{q_{p+1k}(z)q_k(z^1)}.$$
The last sum equals $`G_{p+1}(z^1),`$ which is zero according to Eq. (2.16), so we have proven relation (2.18) for $`p+1.`$ Therefore the relation is valid for any $`p`$.
4) The relation (2.19) is a special case of Eq. (2.18). The denominators in the sum (2.18) can be with the use of the identity (2.15) expressed
$$q_{pk}(z)q_k(z^1)=(1)^pz^sq_{pk}(z^1)q_k(z),s=\left(\frac{p}{2}k\right)(p+1)$$
and since $`z^{rk}=z^{pk}`$, the sum can be rewritten
$$\underset{k=0}{\overset{p}{}}\frac{z^{kr}}{q_{pk}(z)q_k(z^1)}=(1)^p\underset{k=0}{\overset{p}{}}\frac{z^sz^{pk}}{q_{pk}(z^1)q_k(z)}$$
$$=(1)^pz^{\frac{p(p+1)}{2}}\underset{k=0}{\overset{p}{}}\frac{z^k}{q_{pk}(z^1)q_k(z)}.$$
Obviously, the last sum coincides with $`G_p(z)`$, which is zero according to already proven identity (2.16).
Let us remark, last lemma implies also the known formula
$$x^ny^n=(xy)(x\alpha y)(x\alpha ^2y)\mathrm{}(x\alpha ^{n1}y),\alpha =\mathrm{exp}(2\pi i/n).$$
(2.27)
The product can be expanded
$$x^ny^n=\underset{j=0}{\overset{n}{}}c_jx^{nj}(y)^j$$
and one can easily check that
$$c_0=1,c_n=\alpha \alpha ^2\alpha ^3\mathrm{}\alpha ^{n1}=(1)^{n1}.$$
For the remaining $`j,`$ $`0<j<n`$ we get
$$c_j=\underset{k_j=j1}{\overset{n1}{}}\mathrm{}\underset{k_2=1}{\overset{k_31}{}}\underset{k_1=0}{\overset{k_21}{}}\alpha ^{k_1}\alpha ^{k_2}\mathrm{}\alpha ^{k_j}$$
and after the shift of summing limits we obtain
$$c_j=\alpha \alpha ^2\alpha ^3\mathrm{}\alpha ^{j1}\underset{k_j=0}{\overset{nj}{}}\mathrm{}\underset{k_2=0}{\overset{k_3}{}}\underset{k_1=0}{\overset{k_2}{}}\alpha ^{k_1}\alpha ^{k_2}\mathrm{}\alpha ^{k_j}.$$
This multiple sum is a special case of the formula (2.12) and since $`\alpha ^n=1`$, the identity (2.19) is satisfied. Therefore for $`0<j<n`$ we get $`c_j=0`$ and formula (2.27) is proved.
###### Definition 6
Let us have a matrix product created from some string of matrices $`X,Y`$ in such a way that matrix $`X`$ is in total involved $`ptimes`$ and $`Yrtimes`$. By the symbol $`P_j^+`$ ($`P_j^{}`$) we denote permutation, which shifts the leftmost (rightmost) matrix to right (left) on the position in which the shifted matrix has $`j`$ matrices of different kind left (right). (Range of $`j`$ is restricted by $`p`$ or $`r`$ if the shifted matrix is $`Y`$ or $`X`$).
###### Example 7
$$P_3^+XYXYYXY=YXYYXXY$$
(2.28)
Now, we can prove the following theorem.
###### Theorem 8
Let $`p,r>0`$ and $`p+r=n`$ (i.e. $`\alpha ^{p+r}=1)`$. Then the matrices $`S,T`$ fulfill
$$\{S^p,T^r\}=0.$$
(2.29)
Proof:
Obviously, all the terms in the combinator $`\{S^p,T^r\}`$ can be generated e.g. from the string
$$\underset{𝑝}{\underset{}{SS\mathrm{}.S}}\underset{𝑟}{\underset{}{TT\mathrm{}.T}}=S^pT^r$$
by means of the permutations $`P_j^+`$
$$\{S^p,T^r\}=\underset{k_p=0}{\overset{r}{}}\mathrm{}\underset{k_2=0}{\overset{k_3}{}}\underset{k_1=0}{\overset{k_2}{}}P_{k_1}^+P_{k_2}^+\mathrm{}P_{k_p}^+S^pT^r.$$
(2.30)
Now the relation (2.3) implies
$$P_j^+S^pT^r=\alpha ^jS^pT^r$$
and Eq. (2.30) can be modified
$$\{S^p,T^r\}=\left(\underset{k_p=0}{\overset{r}{}}\mathrm{}\underset{k_2=0}{\overset{k_3}{}}\underset{k_1=0}{\overset{k_2}{}}\alpha ^{k_1}\alpha ^{k_2}\mathrm{}\alpha ^{k_p}\right)S^pT^r.$$
(2.31)
Apparently the multiple sum in this equation coincides with r.h.s. of Eq. (2.12) and satisfies the condition for Eq. (2.19), thereby the theorem is proved.
Let us remark, that alternative use of permutations $`P_j^{}`$ instead of $`P_j^+`$ would lead to the equation
$$\{S^p,T^r\}=\left(\underset{k_r=0}{\overset{p}{}}\mathrm{}\underset{k_2=0}{\overset{k_3}{}}\underset{k_1=0}{\overset{k_2}{}}\alpha ^{k_1}\alpha ^{k_2}\mathrm{}\alpha ^{k_r}\right)S^pT^r.$$
(2.32)
Comparison of Eqs. (2.31), (2.32) with the relation for $`F_{pr}`$ defined by Eq. (2.12) implies
$$F_{pr}(\alpha )=F_{rp}(\alpha ).$$
(2.33)
Obviously this equation is valid irrespective of the assumption $`\alpha ^{p+r}=1`$, i.e. it holds for any $`n`$ and $`\alpha =\mathrm{exp}(2\pi i/n)`$. It follows, that Eq. (2.33) is satisfied for any $`\alpha `$.
###### Definition 9
By the symbols $`Q_{pr}`$ we denote $`n^2`$ matrices
$$Q_{pr}=S^pT^r,p,r=1,2,\mathrm{},n$$
(2.34)
###### Lemma 10
$$Q_{rs}Q_{pq}=\alpha ^{sp}Q_{kl};k=\text{mod}(r+p1,n)+1,l=\text{mod}(s+q1,n)+1,$$
(2.35)
$$Q_{rs}Q_{pq}=\alpha ^{sprq}Q_{pq}Q_{rs},$$
(2.36)
$$\left(Q_{rs}\right)^n=(1)^{(n1)rs}I,$$
(2.37)
$$Q_{rs}^{}Q_{rs}=Q_{rs}Q_{rs}^{}=I,$$
(2.38)
$$Q_{rs}^{}=\alpha ^{rs}Q_{kl};k=nr,l=ns,$$
(2.39)
$$detQ_{rs}=(1)^{(n1)(r+s)}$$
(2.40)
and for $`rn`$ or $`sn`$
$$\text{Tr }Q_{rs}=0.$$
(2.41)
Proof:
The relations follow from the definition of $`Q_{pr}`$ and relations (2.3)-(2.7).
###### Theorem 11
The matrices $`Q_{pr}`$ are linearly independent and any matrix $`A`$ (of the same dimension) can be expressed as their linear combination
$$A=\underset{k,l=1}{\overset{n}{}}a_{kl}Q_{kl},a_{kl}=\frac{1}{n}\text{Tr}(Q_{kl}^{}A).$$
(2.42)
Proof:
Let us assume matrices $`Q_{kl}`$ are linearly dependent, i.e. there exists some $`a_{rs}0`$ and simultaneously
$$\underset{k,l=1}{\overset{n}{}}a_{kl}Q_{kl}=0,$$
which with the use of the previous lemma implies
$$\text{Tr}\underset{k,l=1}{\overset{n}{}}a_{kl}Q_{rs}^{}Q_{kl}=a_{rs}n=0.$$
This equation contradicts our assumption, therefore the matrices are independent and obviously represent a base in the linear space of matrices $`n\times n`$, which with the use of the previous lemma implies the relations (2.42).
###### Theorem 12
For any $`n2`$, among $`n^2`$ matrices (2.34) there exists the triad $`Q_\lambda ,Q_\mu ,Q_\nu `$ for which
$$\{Q_\lambda ^p,Q_\mu ^r\}=\{Q_\mu ^p,Q_\nu ^r\}=\{Q_\nu ^p,Q_\lambda ^r\}=0;0<p,r,p+r=n$$
(2.43)
and moreover if $`n3`$, then also
$$\{Q_\lambda ^p,Q_\mu ^r,Q_\nu ^s\}=0;0<p,r,s,p+r+s=n.$$
(2.44)
Proof:
We shall show the relations hold e.g. for indices $`\lambda =1n,\mu =11,\nu =n1.`$ Let us denote
$$X=Q_{1n}=S,Y=Q_{11},Z=Q_{n1}=T,$$
(2.45)
then the relation (2.36) implies
$$YX=\alpha XY,ZX=\alpha XZ,ZY=\alpha YZ.$$
(2.46)
Actually the relation $`\{X^p,Z^r\}=0`$ is already proven in the Theorem 8, obviously the remaining relations (2.43) can be proved exactly in the same way.
The combinator (2.44) can be similarly as in the proof of Theorem 8 expressed
$$\{X^p,Y^r,Z^s\}$$
(2.47)
$$=\underset{j_p=0}{\overset{r+s}{}}\mathrm{}\underset{j_2=0}{\overset{j_3}{}}\underset{j_1=0}{\overset{j_2}{}}P_{j_1}^+P_{j_2}^+\mathrm{}P_{j_p}^+X^p\underset{k_p=0}{\overset{s}{}}\mathrm{}\underset{k_2=0}{\overset{k_3}{}}\underset{k_1=0}{\overset{k_2}{}}P_{k_1}^+P_{k_2}^+\mathrm{}P_{k_r}^+Y^rZ^s,$$
which for matrices obeying relations (2.46) give
$$\{X^p,Y^r,Z^s\}$$
$$=\left(\underset{j_p=0}{\overset{r+s}{}}\mathrm{}\underset{j_2=0}{\overset{j_3}{}}\underset{j_1=0}{\overset{j_2}{}}\alpha ^{j_1}\alpha ^{j_2}\mathrm{}\alpha ^{j_p}\right)\left(\underset{k_p=0}{\overset{s}{}}\mathrm{}\underset{k_2=0}{\overset{k_3}{}}\underset{k_1=0}{\overset{k_2}{}}\alpha ^{k_1}\alpha ^{k_2}\mathrm{}\alpha ^{k_r}\right)X^pY^rZ^s.$$
Since the first multiple sum (with indices $`j`$) coincides with Eq. (2.12) and satisfy the condition for Eq. (2.19), r.h.s. is zero and the theorem is proved.
Now let us make few remarks to illuminate content of the last theorem and meaning of the matrices $`Q_\lambda `$. Obviously, the relations (2.43), (2.44) are equivalent to the statement: any three complex numbers $`a,b,c`$ satisfy
$$(aQ_\lambda +bQ_\mu +cQ_\nu )^n=(a^n+b^n+c^n)I.$$
(2.48)
Further, the theorem speaks about existence of the triad but not about their number. Generally for $`n>2`$ there is more than one triad defined by the theorem, but on the other hand not any three various matrices from the set $`Q_{rs}`$ comply with the theorem. Simple example are some $`X,Y,Z`$ where e.g. $`XY=YX,`$ which happens for $`YX^p,2p<n`$. Obviously in this case at least the relation (2.43) surely is not satisfied. Computer check of the relation (2.47) which has been done with all possible triads from $`Q_{rs}`$ for $`2n20`$ suggests, that a triad $`X,Y,Z`$ for which there exist the numbers $`p,r,s1`$ and $`p+r+sn`$ so that $`X^pY^rZ^sI`$ also does not comply with the theorem. Further, the result on r.h.s. of Eq. (2.47) generally depends on the factors $`\beta _k`$ in relations
$$XY=\beta _3YXYZ=\beta _1ZYZX=\beta _2XZ$$
(2.49)
and computer check suggests the sets in which for some $`\beta _k`$ and $`p<n`$ there is $`\beta _k^p=1`$ also contradict the theorem. In this way the number of different triads obeying the relations (2.43), (2.44) is rather complicated function of $`n`$ \- as shown in the table
$$\begin{array}{cccccccccccccccccccc}n:& \text{2}& \text{3}& \text{4}& \text{5}& \text{6}& \text{7}& \text{8}& \text{9}& \text{10}& \text{11}& \text{12}& \text{13}& \text{14}& \text{15}& \text{16}& \text{17}& \text{18}& \text{19}& \text{20}\\ \mathrm{\#}3:& \text{1}& \text{1}& \text{1}& \text{4}& \text{1}& \text{9}& \text{4}& \text{9}& \text{4}& \text{25}& \text{4}& \text{36}& \text{9}& \text{16}& \text{16}& \text{64}& \text{9}& \text{81}& \text{16}\end{array}$$
Here the statement triad $`X,Y,Z`$ is different from $`X^{},Y^{},Z^{}`$ means that after any rearrangement of the symbols $`X,Y,Z`$ for marking of matrices in the given set, there is always at least one pair $`\beta _k\beta _k^{}.`$
Naturally, one can ask if there exists also the set of four or generally $`N`$ matrices, which satisfy a relation similar to Eq. (2.48)
$$\left(\underset{\lambda =0}{\overset{N1}{}}a_\lambda Q_\lambda \right)^n=\underset{\lambda =0}{\overset{N1}{}}a_\lambda ^n.$$
(2.50)
For $`2n10`$ and $`N=4`$ the computer suggests the negative answer - in the case of matrices generated according to the Definition 9. However, one can verify: if $`U_l,l=1,2,3`$ is the triad complying with the theorem (or equivalently with the relation (2.48)), then the matrices $`n^2\times n^2`$
$$Q_0=IT=\left(\begin{array}{cccccc}I& & & & & \\ & \alpha I& & & & \\ & & \alpha ^2I& & & \\ & & & & & \\ & & & & & \\ & & & & & \alpha ^{n1}I\end{array}\right),$$
(2.51)
$$Q_l=U_lS=\left(\begin{array}{cccccc}0& & & & & U_l\\ U_l& & & & & \\ & U_l& & & & \\ & & & & & \\ & & & & & \\ & & & & U_l& 0\end{array}\right)$$
(2.52)
satisfy relation (2.50) for $`N=4`$. Generally, if $`U_\lambda `$ are matrices complying with Eq. (2.50) for some $`N3`$, then the matrices created from them according to the rule (2.51), (2.52) will satisfy Eq. (2.50) for $`N+1`$. The last statement follows from the following equalities. Let us assume
$$\underset{k=0}{\overset{N}{}}p_k=n,$$
then
$$\{Q_0^{p_0},Q_1^{p_1},\mathrm{}Q_N^{p_N}\}=\underset{j_{p_N}=0}{\overset{np_N}{}}\mathrm{}\underset{j_1=0}{\overset{j_2}{}}\underset{j_0=0}{\overset{j_1}{}}P_{j_0}^{}P_{j_1}^{}\mathrm{}P_{j_{p_N}}^{}\{Q_0^{p_0},\mathrm{}Q_{N1}^{p_{N1}}\}Q_N^{p_N}$$
$$=\underset{j_{p_N}=0}{\overset{np_N}{}}\mathrm{}\underset{j_1=0}{\overset{j_2}{}}\underset{j_0=0}{\overset{j_1}{}}P_{j_0}^{}P_{j_1}^{}\mathrm{}P_{j_{p_N}}^{}\{(U_0S)^{p_0},\mathrm{}(U_{N1}S)^{p_{N1}}\}(IT)^{p_N}$$
$$=\underset{j_{p_N}=0}{\overset{np_N}{}}\mathrm{}\underset{j_1=0}{\overset{j_2}{}}\underset{j_0=0}{\overset{j_1}{}}\alpha ^{j_0}\alpha ^{j_1}\mathrm{}\alpha ^{j_{p_N}}\{(U_0S)^{p_0},\mathrm{}(U_{N1}S)^{p_N1}\}(IT)^{p_N}$$
$$=\left(\underset{j_{p_N}=0}{\overset{np_N}{}}\mathrm{}\underset{j_1=0}{\overset{j_2}{}}\underset{j_0=0}{\overset{j_1}{}}\alpha ^{j_0}\alpha ^{j_1}\mathrm{}\alpha ^{j_{p_N}}\right)\{U_1^{p_1},\mathrm{}U_{N1}^{p_{N1}}\}S^{np_N}T^{p_N},$$
where the last multiple sum equals zero according to the relations (2.12) and (2.19). Obviously for $`n=2`$ the matrices (2.45) and (2.51),(2.52) created from them correspond, up to some phase factors, to the Pauli matrices $`\sigma _j`$ and Dirac matrices $`\gamma _\mu `$.
Obviously, from the set of matrices $`Q_{rs}`$ (with exception of $`Q_{nn}=I`$) one can easily make the $`n^21`$ generators of the fundamental representation of $`SU(n)`$ group
$$G_{rs}=a_{rs}Q_{rs}+a_{rs}^{}Q_{rs}^+,$$
(2.53)
where $`a_{rs}`$ are suitable factors. For example the choice
$$a_{kl}=\frac{1}{\sqrt{2}}\alpha ^{[kl+n(k+l1/4)]/2}$$
(2.54)
gives commutation relations
$$[G_{kl},G_{rs}]=i\mathrm{sin}\left(\pi (kslr)/n\right)$$
(2.55)
$$\{\mathrm{sg}(k+r,l+s,n)(G_{k+r,l+s}(1)^{n+k+l+r+s}G_{kr,ls})$$
$$\mathrm{sg}(kr,ls,n)(G_{kr,ls}(1)^{n+k+l+r+s}G_{rk,sl})\},$$
where
$$\mathrm{sg}(p,q,n)=(1)^{pm_q+qm_pn},m_x=\frac{x\text{mod}(x1,n)1}{n}.$$
and indices at $`G`$ (on r.h.s.) in Eq. (2.55) are understood in the sense of mod - like in the relation (2.35). One can easily check e.g. for $`n=2`$ the matrices (2.53) with the factors $`a_{rs}`$ according to Eq. (2.54) are the Pauli matrices - generators of the fundamental representation of the $`SU(2)`$ group.
## 3 Wave equations generated by the roots <br>of D’Alambertian operator $`\mathrm{}^{1/n}`$
Now, using the generalized Dirac matrices (2.51), (2.52) we shall assemble the corresponding wave equation as follows. These four matrices with normalization
$$\left(Q_0\right)^n=\left(Q_l\right)^n=I,l=1,2,3,$$
(3.1)
allow to write down the set of algebraic equations
$$\left(\mathrm{\Gamma }(p)\mu I\right)\mathrm{\Psi }(p)=0,$$
(3.2)
where
$$\mathrm{\Gamma }(p)=\underset{\lambda =0}{\overset{3}{}}\pi _\lambda Q_\lambda .$$
(3.3)
If the variables $`\mu ,\pi _\lambda `$ represent the fractional powers of the mass and momentum components
$$\mu ^n=m^2,\pi _\lambda ^n=p_\lambda ^2,$$
(3.4)
then
$$\mathrm{\Gamma }(p)^n=p_0^2p_1^2p_2^2p_3^2p^2$$
(3.5)
and after $`n1`$ times repeated application of the operator $`\mathrm{\Gamma }`$ on Eq. (3.2) one gets the set of Klein-Gordon equations in the $`p`$representation
$$\left(p^2m^2\right)\mathrm{\Psi }(p)=0.$$
(3.6)
The Eqs. (3.2) and (3.6) are the sets of $`n^2`$ equations with solution $`\mathrm{\Psi }`$ having $`n^2`$ components. Obviously, the case $`n^2=4`$ corresponds to the Dirac equation. For $`n>2`$ the Eq. (3.2) is new, more complicated and immediately invoking some questions. In the present paper we shall attempt to answer at least some of them. One can check, that the solution of the set (3.2) reads
$$\mathrm{\Psi }(p)=\left(\begin{array}{c}𝐡\\ \frac{U(p)}{\alpha \pi _0\mu }𝐡\\ \frac{U^2(p)}{(\alpha \pi _0\mu )(\alpha ^2\pi _0\mu )}𝐡\\ \\ \\ \frac{U^{n1}(p)}{(\alpha \pi _0\mu )\mathrm{}(\alpha ^{n1}\pi _0\mu )}𝐡\end{array}\right),𝐡=\left(\begin{array}{c}h_1\\ h_2\\ \\ \\ h_n\end{array}\right),$$
(3.7)
where
$$U(p)=\underset{l=1}{\overset{3}{}}\pi _lU_l,\left(U_l\right)^n=I,$$
($`U_l`$ is the triad from which the matrices $`Q_l`$ are constructed in accordance with Eqs. (2.51), (2.52)) and $`h_1,h_2,\mathrm{}h_n`$ are arbitrary functions of $`p`$. At the same time, $`\pi _\lambda `$ satisfy the constraint
$$\pi _0^n\pi _1^n\pi _2^n\pi _3^n=\mu ^n=m^2.$$
(3.8)
First of all, one can bring to notice, that in Eq. (3.2) the fractional powers of the momentum components appear, which means that the equation in the $`x`$representation will contain the fractional derivatives:
$$\pi _\lambda =(p_\lambda )^{2/n}(i_\lambda )^{2/n}.$$
(3.9)
Our primary considerations will concern $`p`$representation, but afterwards we shall show how the transition to the $`x`$representation can be realized by means of the Fourier transformation, in accordance with the approach suggested in .
Further question concerns relativistic covariance of Eq. (3.2): How to transform simultaneously the operator
$$\mathrm{\Gamma }(p)\mathrm{\Gamma }(p^{})=\mathrm{\Lambda }\mathrm{\Gamma }(p)\mathrm{\Lambda }^1$$
(3.10)
and the solution
$$\mathrm{\Psi }(p)\mathrm{\Psi }^{}(p^{})=\mathrm{\Lambda }\mathrm{\Psi }(p)$$
(3.11)
to preserve the equal form of the operator $`\mathrm{\Gamma }`$ for initial variables $`p_\lambda `$ and the boosted ones $`p_\lambda ^{}`$ ?
### 3.1 Infinitesimal transformations
First let us consider the infinitesimal transformations
$$\mathrm{\Lambda }(d\omega )=I+id\omega L_\omega ,$$
(3.12)
where $`d\omega `$ represents the infinitesimal values of the six parameters of the Lorentz group corresponding to the space rotations
$$p_i^{}=p_i+ϵ_{ijk}p_jd\phi _k,i=1,2,3$$
(3.13)
and the Lorentz transformations
$$p_i^{}=p_i+p_0d\psi _i,p_0^{}=p_0+p_id\psi _i,i=1,2,3,$$
(3.14)
where $`\mathrm{tanh}\psi _i=v_i/c\beta _i`$ is the corresponding velocity. Here, and anywhere in the next we use the convention, that in the expressions involving the antisymmetric tensor $`ϵ_{ijk}`$, the summation over indices appearing twice is done. From the infinitesimal transformations (3.13), (3.14) one can obtain the finite ones. For the three space rotations we get
$$p_1^{}=p_1\mathrm{cos}\phi _3+p_2\mathrm{sin}\phi _3,p_2^{}=p_2\mathrm{cos}\phi _3p_1\mathrm{sin}\phi _3,p_3^{}=p_3,$$
(3.15)
$$p_2^{}=p_2\mathrm{cos}\phi _1+p_3\mathrm{sin}\phi _1,p_3^{}=p_3\mathrm{cos}\phi _1p_2\mathrm{sin}\phi _1,p_1^{}=p_1,$$
(3.16)
$$p_3^{}=p_3\mathrm{cos}\phi _2+p_1\mathrm{sin}\phi _2,p_1^{}=p_1\mathrm{cos}\phi _2p_3\mathrm{sin}\phi _2,p_2^{}=p_2$$
(3.17)
and for the Lorentz transformations similarly
$$p_0^{}=p_0\mathrm{cosh}\psi _i+p_i\mathrm{sinh}\psi _i,i=1,2,3,$$
(3.18)
where
$$\mathrm{cosh}\psi _i=\frac{1}{\sqrt{1\beta _i^2}},\mathrm{sinh}\psi _i=\frac{\beta _i}{\sqrt{1\beta _i^2}}.$$
(3.19)
The definition of the six parameters implies that the corresponding infinitesimal transformations of the reference frame $`pp^{}`$ changes a function $`f(p)`$ :
$$f(p)f(p^{})=f(p+\delta p)=f(p)+\frac{df}{d\omega }d\omega ,$$
(3.20)
where $`d/d\omega `$ stands for
$$\frac{d}{d\phi _i}=ϵ_{ijk}p_j\frac{}{p_k},\frac{d}{d\psi _i}=p_0\frac{}{p_i}+p_i\frac{}{p_0},i=1,2,3.$$
(3.21)
Obviously, the equation
$$p^{}=p+\frac{dp}{d\omega }d\omega $$
(3.22)
combined with Eq. (3.21) is identical to Eqs. (3.13), (3.14). Further, with the use of formulas (3.12) and (3.21) the relations (3.10), (3.11) can be rewritten in the infinitesimal form
$$\mathrm{\Gamma }(p^{})=\mathrm{\Gamma }(p)+\frac{d\mathrm{\Gamma }(p)}{d\omega }d\omega =\left(I+id\omega L_\omega \right)\mathrm{\Gamma }(p)\left(Iid\omega L_\omega \right),$$
(3.23)
$$\mathrm{\Psi }^{}(p^{})=\mathrm{\Psi }^{}(p)+\frac{d\mathrm{\Psi }^{}(p)}{d\omega }d\omega =\left(I+id\omega L_\omega \right)\mathrm{\Psi }(p).$$
(3.24)
If we define
$$𝐋_\omega =L_\omega +i\frac{d}{d\omega },$$
(3.25)
then the relations (3.23), (3.24) imply
$$[𝐋_\omega ,\mathrm{\Gamma }]=0,$$
(3.26)
$$\mathrm{\Psi }^{}(p)=\left(I+id\omega 𝐋_\omega \right)\mathrm{\Psi }(p).$$
(3.27)
The six operators $`𝐋_\omega `$ are generators of the corresponding representation of the Lorentz group, so they have to satisfy the commutation relations
$$[𝐋_{\phi _𝐣},𝐋_{\phi _𝐤}]=iϵ_{jkl}𝐋_{\phi _𝐥},$$
(3.28)
$$[𝐋_{\psi _𝐣},𝐋_{\psi _𝐤}]=iϵ_{jkl}𝐋_{\phi _𝐥},$$
(3.29)
$$[𝐋_{\phi _𝐣},𝐋_{\psi _𝐤}]=iϵ_{jkl}𝐋_{\psi _𝐥},j,k,l=1,2,3.$$
(3.30)
How this representation looks like, in other words, what operators $`𝐋_\omega `$ satisfy Eqs.(3.28) - (3.30) and (3.26)? First, one can easily check, that for $`n>2`$ there do not exist matrices $`L_\omega `$ with constant elements representing the first term in r.h.s. of equality (3.25) and satisfying the Eq. (3.26). If one assumes, that $`L_\omega `$ consist only of constant elements, then the elements of matrix $`\frac{d}{d\omega }\mathrm{\Gamma }(p)`$ involving the terms like $`p_i^{2/n1}p_j`$ certainly cannot be expressed through the elements of the difference $`L_\omega \mathrm{\Gamma }\mathrm{\Gamma }L_\omega `$ consisting only of the elements proportional to $`p_k^{2/n}`$ \- in contradistinction to the case $`n=2,`$ i.e. the case of Dirac equation. In this way the Eq. (3.26) cannot be satisfied for $`n>2`$ and $`L_\omega `$ constant. Nevertheless, one can show, that the set of Eqs. (3.26),(3.28) - (3.30) is solvable, provided that we accept the elements of the matrices $`L_\omega `$ are not constants, but the functions of $`p_i`$. To prove this, let us first make a few preparing steps.
###### Definition 13
Let $`\mathrm{\Gamma }_1(p),\mathrm{\Gamma }_2(p)`$ and $`X`$ be the square matrices of the same dimension and
$$\mathrm{\Gamma }_1(p)^n=\mathrm{\Gamma }_2(p)^n=p^2.$$
Then for any matrix $`X`$ we define the form
$$Z(\mathrm{\Gamma }_1,X,\mathrm{\Gamma }_2)=\frac{1}{np^2}\underset{j=1}{\overset{n}{}}\mathrm{\Gamma }_1^jX\mathrm{\Gamma }_2^{nj}.$$
(3.31)
One can easily check, that the matrix $`Z`$ satisfies e.g.
$$\mathrm{\Gamma }_1Z=Z\mathrm{\Gamma }_2,$$
(3.32)
$$Z(Z(X))=Z(X)$$
(3.33)
and in particular for $`\mathrm{\Gamma }_1=\mathrm{\Gamma }_2\mathrm{\Gamma }`$
$$[\mathrm{\Gamma },Z]=0,$$
(3.34)
$$[\mathrm{\Gamma },X]=0X=Z(X).$$
(3.35)
###### Lemma 14
The Eq. (3.2) can be expressed in the diagonalized (canonical) form
$$\left(\mathrm{\Gamma }_0(p)\mu \right)\mathrm{\Psi }_0(p)=0;\mathrm{\Gamma }_0(p)\left(p^2\right)^{1/n}Q_{0,}$$
(3.36)
where $`Q_0`$ is the matrix (2.51), i.e. there exists the set of transformations $`Y`$, that
$$\mathrm{\Gamma }_0(p)=Y(p)\mathrm{\Gamma }(p)Y^1(p);Y=Z(\mathrm{\Gamma }_0,X,\mathrm{\Gamma })$$
(3.37)
and a particular form reads
$$Y=yZ(\mathrm{\Gamma }_0,I,\mathrm{\Gamma }),Y^1=yZ(\mathrm{\Gamma },I,\mathrm{\Gamma }_0),$$
(3.38)
where
$$y=\sqrt{\frac{n\left[1\left(p_0^2/p^2\right)^{1/n}\right]}{1p_0^2/p^2}}.$$
Proof:
The Eq. (3.32) implies
$$\mathrm{\Gamma }_0=Z(\mathrm{\Gamma }_0,X,\mathrm{\Gamma })\mathrm{\Gamma }Z(\mathrm{\Gamma }_0,X,\mathrm{\Gamma })^1,$$
(3.39)
therefore, if the matrix $`X`$ is chosen in such a way that $`detZ0`$, then $`Z^1`$ exists and the transformation (3.39) diagonalizes the matrix $`\mathrm{\Gamma }`$. Let us put $`X=I`$ and calculate the following product
$$C=Z(\mathrm{\Gamma }_0,I,\mathrm{\Gamma })Z(\mathrm{\Gamma },I,\mathrm{\Gamma }_0)=\frac{1}{n^2p^4}\underset{i,j=1}{\overset{n}{}}\mathrm{\Gamma }_0^i\mathrm{\Gamma }^{ni+j}\mathrm{\Gamma }_0^{nj}.$$
(3.40)
The last sum can be rearranged, instead of the summation index $`j`$ we use the new one
$$k=ij\mathrm{for}\text{ }ij,k=ij+n\mathrm{for}\text{ }i<j;k=0,\mathrm{}n1,$$
then the Eq. (3.40) reads
$$C=\frac{1}{n^2p^4}\underset{k=0}{\overset{n1}{}}\left(\underset{i=k+1}{\overset{n}{}}\mathrm{\Gamma }_0^i\mathrm{\Gamma }^{nk}\mathrm{\Gamma }_0^{n+ki}+\underset{i=1}{\overset{k}{}}\mathrm{\Gamma }_0^i\mathrm{\Gamma }^{2nk}\mathrm{\Gamma }_0^{ki}\right)$$
(3.41)
and if we take into account that $`\mathrm{\Gamma }_0^n=\mathrm{\Gamma }^n=p^2,`$ then this sum can be simplified
$$C=\underset{k=0}{\overset{n1}{}}C_k=\frac{1}{n^2p^2}\underset{k=0}{\overset{n1}{}}\underset{i=1}{\overset{n}{}}\mathrm{\Gamma }_0^i\mathrm{\Gamma }^{nk}\mathrm{\Gamma }_0^{ki}.$$
(3.42)
For the term $`k=0`$ we get
$$C_0=\frac{1}{n^2p^2}\underset{i=1}{\overset{n}{}}\mathrm{\Gamma }_0^i\mathrm{\Gamma }^n\mathrm{\Gamma }_0^i=\frac{1}{n}$$
(3.43)
and for $`k>0`$, using Eqs. (3.3), (2.51), (2.52), (3.36) and Definition 3 one obtains
$$C_k=\frac{1}{n^2p^2}\underset{i=1}{\overset{n}{}}\mathrm{\Gamma }_0^i\mathrm{\Gamma }^{nk}\mathrm{\Gamma }_0^{ki}=\frac{1}{n^2p^2}\underset{i=1}{\overset{n}{}}\mathrm{\Gamma }_0^i\left(\underset{\lambda =0}{\overset{3}{}}\pi _\lambda Q_\lambda \right)^{nk}\mathrm{\Gamma }_0^{ki}$$
(3.44)
$$=\frac{1}{n^2p^2}\underset{i=1}{\overset{n}{}}\mathrm{\Gamma }_0^i\left(\pi _0IT+\left[\underset{\lambda =1}{\overset{3}{}}\pi _\lambda U_\lambda \right]S\right)^{nk}\mathrm{\Gamma }_0^{ki}$$
$$=\frac{1}{n^2p^2}\underset{i=1}{\overset{n}{}}\mathrm{\Gamma }_0^i\left(\pi _0IT+US\right)^{nk}\mathrm{\Gamma }_0^{ki}$$
$$=\frac{1}{n^2p^2}\underset{i=1}{\overset{n}{}}\mathrm{\Gamma }_0^i\left(\underset{p=0}{\overset{nk}{}}\pi _0^pU^{nkp}\{T^p,S^{nkp}\}\right)\mathrm{\Gamma }_0^{ki}$$
$$=\frac{\left(p^2\right)^{k/n}}{n^2p^2}\underset{p=0}{\overset{nk}{}}\pi _0^pU^{nkp}\underset{i=1}{\overset{n}{}}T^i\{T^p,S^{nkp}\}T^{ki}.$$
For $`p<nkl`$ the last sum can be with the use of the relation (2.3) modified
$$\underset{i=1}{\overset{n}{}}T^i\{T^p,S^{lp}\}T^{ki}=\{T^p,S^{lp}\}T^k\underset{i=1}{\overset{n}{}}\alpha ^{i(lp)}$$
(3.45)
$$=\{T^p,S^{lp}\}T^k\alpha ^{(lp)}\frac{1\alpha ^{n(lp)}}{1\alpha ^{(lp)}}=0,$$
therefore only the term $`p=nk`$ contributes:
$$C_k=\frac{\left(p^2\right)^{k/n}}{n^2p^2}\left(p_0^2\right)^{(nk)/n}n=\frac{1}{n}\left(\frac{p_0^2}{p^2}\right)^{(nk)/n}.$$
(3.46)
So the sum (3.42) gives in total
$$C=\frac{1}{n}\left[1+\left(\frac{p_0^2}{p^2}\right)^{1/n}+\left(\frac{p_0^2}{p^2}\right)^{2/n}+\mathrm{}+\left(\frac{p_0^2}{p^2}\right)^{(n1)/n}\right]$$
(3.47)
$$=\frac{1p_0^2/p^2}{n\left[1\left(p_0^2/p^2\right)^{1/n}\right]},$$
therefore Eq. (3.37) is satisfied with $`Y,Y^1`$ given by Eq. (3.38) and the proof is completed.
Solution of the Eq. (3.36) reads
$$\mathrm{\Psi }_0(p)=\left(\begin{array}{c}\mathrm{𝟎}\\ \\ \mathrm{𝟎}\\ 𝐠\\ \mathrm{𝟎}\\ \\ \\ \mathrm{𝟎}\end{array}\right);𝐠\left(\begin{array}{c}g_1\\ g_2\\ \\ \\ g_n\end{array}\right),\mathrm{𝟎}\left(\begin{array}{c}0\\ 0\\ \\ \\ 0\end{array}\right),$$
(3.48)
i.e. the sequence of non zero components can be only in one block, whose location depends on the choice of the phase of the power $`(p^2)^{1/n}`$. The $`g_j`$ are arbitrary functions of $`p`$ and simultaneously the constraint $`p^2=m^2`$ is required. Now we shall try to find the generators satisfying the covariance condition for Eq. (3.36)
$$[𝐋_\omega ,\mathrm{\Gamma }_0(p)]=0$$
(3.49)
together with the commutation relations (3.28)-(3.30). Some hint can be obtained from the Dirac equation transformed to the diagonal form in an accordance with the relations (3.37), (3.38). We shall use the current representation of the Pauli and Dirac matrices
$$\sigma _1=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\sigma _2=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right),\sigma _3=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),$$
(3.50)
$$\gamma _0=\left(\begin{array}{cc}\mathrm{𝟏}& \mathrm{𝟎}\\ \mathrm{𝟎}& \mathrm{𝟏}\end{array}\right),\gamma _j=\left(\begin{array}{cc}\mathrm{𝟎}& \sigma _j\\ \sigma _j& \mathrm{𝟎}\end{array}\right);j=1,2,3,$$
(3.51)
where the bold 0,1 stand for zero and unit matrices $`2\times 2`$. The Dirac equation
$$\left(\mathrm{\Gamma }(p)m\right)\mathrm{\Psi }(p)=0,\mathrm{\Gamma }(p)\underset{\lambda =0}{\overset{3}{}}p_\lambda \gamma _\lambda $$
(3.52)
is covariant under the transformations generated by
$$𝐋_{\phi _𝐣}=\frac{i}{4}ϵ_{jkl}\gamma _k\gamma _l+i\frac{d}{d\phi _j}=L_{\phi _j}+i\frac{d}{d\phi _j};L_{\phi _j}=\frac{1}{2}\left(\begin{array}{cc}\sigma _j& \mathrm{𝟎}\\ \mathrm{𝟎}& \sigma _j\end{array}\right),$$
(3.53)
$$𝐋_{\psi _𝐣}=\frac{i}{2}\gamma _0\gamma _j+i\frac{d}{d\psi _j}=L_{\psi _j}+i\frac{d}{d\psi _j};L_{\psi _j}=\frac{i}{2}\left(\begin{array}{cc}\mathrm{𝟎}& \sigma _j\\ \sigma _j& \mathrm{𝟎}\end{array}\right),$$
(3.54)
where $`j,k,l=1,2,3.`$ Obviously, to preserve covariance, one has with the transformation $`\mathrm{\Gamma }\mathrm{\Gamma }_0=Y\mathrm{\Gamma }Y^1`$ perform also
$$𝐋_\omega 𝐌_\omega =Y(p)𝐋_\omega Y^1(p).$$
(3.55)
For the space rotations $`𝐋_{\phi _𝐣}`$commuting with both $`\mathrm{\Gamma }_0,\mathrm{\Gamma }`$ and with the $`Y`$ from the relation (3.38) the result is quite straightforward
$$𝐌_{\phi _𝐣}=𝐋_{\phi _𝐣}=L_{\phi _j}+i\frac{d}{d\phi _j},$$
(3.56)
i.e. the generators of the space rotations are not changed by the transformation (3.55). The similar procedure with the Lorentz transformations is slightly more complicated, nevertheless after calculation of the commutator $`[𝐋_{\psi _𝐣},\mathrm{\Gamma }_0/\sqrt{1+p_0/\sqrt{p^2}}]`$ and a few further steps one obtains
$$𝐌_{\psi _𝐣}=M_{\psi _j}(p)+i\frac{d}{d\psi _j};M_{\psi _j}(p)=ϵ_{jkl}\frac{p_kL_{\phi _l}}{p_0+\sqrt{p^2}}.$$
(3.57)
So the generators (3.56), (3.57) guarantee the covariance of the diagonalized Dirac equation obtained from Eq. (3.52) according to Lemma 14. At the same time it is obvious, that having the set of generators $`L_{\phi _j}`$ (with constant elements) of space rotations, one can according to Eq. (3.57) construct the generators of Lorentz transformations $`M_{\psi _j}(p)`$ (or $`𝐌_{\psi _𝐣}`$), which satisfy commutation relations (3.28) - (3.30). Obviously this recipe is valid for any representation of infinitesimal space rotations. Let us make a remark, that the algebra given by Eqs. (3.56), (3.57) appears in a slightly modified form already in . Now, we shall show, that if one requires a linear relation between the generators $`M_{\psi _j}`$ and $`L_{\phi _l}`$, like in Eq. (3.57), then this relation can have a more general shape, than that in Eq. (3.57).
###### Lemma 15
Let $`L_{\phi _j}`$ be matrices with constant elements satisfying commutation relations (3.28). Then the operators
$$𝐌_{\psi _𝐣}=M_{\psi _j}(p)+i\frac{d}{d\psi _j};M_{\psi _j}(p)=\frac{\kappa L_{\phi _j}+ϵ_{jkl}p_kL_{\phi _l}}{p_0+\sqrt{p^2\kappa ^2}},$$
(3.58)
where $`\kappa `$ is any complex constant, satisfy the commutation relations (3.29), (3.30).
Proof:
After insertion of the generators (3.58) into the relations (3.29), (3.30) one can check, that the commutation relations are satisfied. In fact, it is sufficient to verify e.g. the commutators $`[𝐋_{\phi _\mathrm{𝟏}},𝐋_{\psi _\mathrm{𝟐}}],[𝐋_{\phi _\mathrm{𝟏}},𝐋_{\psi _\mathrm{𝟏}}]`$ and $`[𝐋_{\psi _\mathrm{𝟏}},𝐋_{\psi _\mathrm{𝟑}}],`$ the remaining follow from the cyclic symmetry.
Let us note, the formula (3.58) covers also the limit case $`\left|\kappa \right|\mathrm{}`$, then
$$M_{\psi _j}=iL_{\phi _j}.$$
(3.59)
On the other hand, the relation (3.57) corresponds to $`\kappa =0.`$ The representations of the Lorentz group defined by the generators (3.56), (3.58) and differing only in the parameter $`\kappa `$ should be equivalent in the sense
$$𝐌_\omega (\kappa ^{})=X^1(p)𝐌_\omega (\kappa )X(p).$$
(3.60)
We shall not make a general proof of this relation, but rather we shall show, that the representations defined in the Lemma 15 and differing only in $`\kappa `$, can be classified by the same mass $`m^2=p^2`$ and spin $`𝐬^2=s(s+1).`$ First, let us note, that the six generators considered in the lemma together with the four generators $`p_\alpha `$ of the space-time translations form the set of generators of the Poincaré group. One can easily check, that the corresponding additional commutation relations are satisfied:
$$[p_\alpha ,p_\beta ]=0,[𝐌_{\phi _𝐣},p_0]=0,[p_\alpha ,\mathrm{\Gamma }_0]=0,$$
(3.61)
$$[𝐌_{\phi _𝐣},p_k]=iϵ_{jkl}p_l,[𝐌_{\psi _𝐣},p_k]=i\delta _{jk}p_0,[𝐌_{\psi _𝐣},p_0]=ip_j.$$
(3.62)
Further, the generators $`𝐌_\omega `$ can be rewritten in the covariant notation
$$𝐌_{jk}=ϵ_{jkl}𝐌_{\phi _l},𝐌_{j0}=𝐌_{\psi _j},𝐌_{\alpha \beta }=𝐌_{\beta \alpha }.$$
(3.63)
Now the Pauli - Lubanski vector can be constructed:
$$V_\alpha =ϵ_{\alpha \beta \gamma \delta }𝐌^{\beta \gamma }p^\delta /2,$$
(3.64)
which has satisfy
$$V_\alpha V^\alpha =m^2s(s+1),$$
(3.65)
where $`s`$ is the corresponding spin number. One can check, that after inserting the generators (3.63) into relations (3.64), (3.65), the result does not depend on $`\kappa `$
$$V_\alpha V^\alpha =p^2\left(M_{\phi _1}^2+M_{\phi _2}^2+M_{\phi _3}^2\right)=m^2s(s+1).$$
So the generators of the Lorentz group, which satisfy Eq. (3.49), can have the form
$$𝐑_\omega =R_\omega +i\frac{d}{d\omega };R_\omega =\left(\begin{array}{cccccc}M_\omega & & & & & \\ & M_\omega & & & & \\ & & M_\omega & & & \\ & & & & & \\ & & & & & \\ & & & & & M_\omega \end{array}\right),$$
(3.66)
where $`M_\omega `$ are the $`n\times n`$ matrices defined in accordance with the Lemma 15. There are $`n`$ such matrices on the diagonal and apparently these matrices may not be identical.
Finally, it is obvious that Eq. (3.36) is covariant also under any infinitesimal transform
$$\mathrm{\Lambda }(d\xi )=I+id\xi K_\xi ,$$
(3.67)
where the generators $`K_\xi `$ have the similar form as the generators (3.66)
$$K_\xi =\left(\begin{array}{cccccc}k_\xi & & & & & \\ & k_\xi & & & & \\ & & k_\xi & & & \\ & & & & & \\ & & & & & \\ & & & & & k_\xi \end{array}\right)$$
(3.68)
and generally their elements may depend on $`p`$. Obviously, one can put the question: If the generators $`L_\phi ,k_\xi `$ from Eqs. (3.56), (3.68) with constant elements represent the algebra of some group (containing the rotation group as a subgroup), then what linear combination $`M_{\psi _j}(p)`$ of these generators satisfy the commutation relations (3.29), (3.30) for the generators of Lorentz transformations? In the other words, what are the coefficients in summation
$$M_{\psi _j}(p)=\underset{k=1}{\overset{3}{}}c_{jk}(p)L_{\phi _k}+\underset{\xi }{}c_{j\xi }k_\xi $$
(3.69)
satisfying the commutation relations for the generators of the Lorentz transformations? In this paper we shall not discuss this more general task, for our present purpose it is sufficient, that we proved existence of the generators of infinitesimal Lorentz transformations, under which the Eq. (3.36) is covariant.
### 3.2 Finite transformations
Now, having the infinitesimal transformations, one can proceed to finite ones, corresponding to the parameters $`\omega `$ and $`\xi `$:
$$\mathrm{\Psi }_0^{}(p^{})=\mathrm{\Lambda }(\omega )\mathrm{\Psi }_0(p),\mathrm{\Psi }_0^{}(p)=\mathrm{\Lambda }(\xi )\mathrm{\Psi }_0(p),$$
(3.70)
where $`pp^{}`$ is some of the transformations (3.15) \- (3.18). The matrices $`\mathrm{\Lambda }`$ satisfy
$$\mathrm{\Lambda }(\omega +d\omega )=\mathrm{\Lambda }(\omega )\mathrm{\Lambda }(d\omega ),\mathrm{\Lambda }(\xi +d\xi )=\mathrm{\Lambda }(\xi )\mathrm{\Lambda }(d\xi ),$$
(3.71)
which for the parameters $`\phi `$ (space rotations only) and $`\xi `$ imply
$$\frac{d\mathrm{\Lambda }(\phi _j)}{d\phi _j}=i\mathrm{\Lambda }(\phi _j)R_{\phi _j},\frac{d\mathrm{\Lambda }(\xi )}{d\xi }=i\mathrm{\Lambda }(\xi )K_\xi .$$
(3.72)
Assuming the constant elements of the matrices $`R_{\phi _j}`$ and $`K_\xi `$, the solutions of the last equations can be written in the usual exponential form
$$\mathrm{\Lambda }(\phi _j)=\mathrm{exp}(i\phi _jR_{\phi _j}),\mathrm{\Lambda }(\xi )=\mathrm{exp}(i\xi K_\xi ).$$
(3.73)
The space rotation by an angle $`\phi `$ about the axis having the direction $`\stackrel{}{u},`$ $`\left|\stackrel{}{u}\right|=1`$ is represented by
$$\mathrm{\Lambda }(\phi ,\stackrel{}{u})=\mathrm{exp}\left[i\phi \left(\stackrel{}{u}\stackrel{}{R}_\phi \right)\right];\stackrel{}{R}_\phi =(R_{\phi _1},R_{\phi _2},R_{\phi _3}).$$
(3.74)
For the Lorentz transformations we get instead of Eq. (3.72)
$$\frac{d\mathrm{\Lambda }(\psi _j)}{d\psi _j}=if_j(\psi _j)\mathrm{\Lambda }(\psi _j)N_j,$$
(3.75)
where in accordance with Eqs. (3.18) and (3.58) there stand for
$$f_j(\psi _j)=\frac{1}{p_0\mathrm{cosh}\psi _j+p_j\mathrm{sinh}\psi _j+\sqrt{p^2\kappa ^2}},$$
(3.76)
$$N_j=\kappa R_{\phi _j}+ϵ_{jkl}p_kR_{\phi _l}.$$
(3.77)
The solution of Eq. (3.75) reads
$$\mathrm{\Lambda }(\psi _j)=\mathrm{exp}\left(iF(\psi _j)N_j\right);F(\psi _j)=_0^{\psi _j}f_j(\eta )𝑑\eta .$$
(3.78)
The Lorentz boost in a general direction $`\stackrel{}{u}`$ with the velocity $`\beta `$ is represented by
$$\mathrm{\Lambda }(\psi ,\stackrel{}{u})=\mathrm{exp}\left(iF(\psi )N\right),\mathrm{tanh}\psi =\beta ,$$
(3.79)
where
$$F(\psi )=_0^\psi \frac{d\eta }{p_0\mathrm{cosh}\eta +\stackrel{}{p}\stackrel{}{u}\mathrm{sinh}\eta +\sqrt{p^2\kappa ^2}},$$
(3.80)
$$N=\kappa \stackrel{}{u}\stackrel{}{R}_\phi +\left(\stackrel{}{u}\times \stackrel{}{p}\right)\stackrel{}{R}_\phi .$$
(3.81)
The corresponding integrals can be found e.g. in the handbook .
Let us note, from the technical point of view, solution of the equation
$$\frac{d\mathrm{\Lambda }(t)}{dt}=\mathrm{\Omega }(t)\mathrm{\Lambda }(t),$$
(3.82)
where $`\mathrm{\Lambda },\mathrm{\Omega }`$ are some square matrices, can be written in the exponential form
$$\mathrm{\Lambda }(t)=\mathrm{exp}\left(_0^t\mathrm{\Omega }(\eta )𝑑\eta \right)$$
(3.83)
only if the matrix $`\mathrm{\Omega }`$ satisfies
$$[\mathrm{\Omega }(t),_0^t\mathrm{\Omega }(\eta )𝑑\eta ]=0.$$
(3.84)
This condition is necessary for differentiation
$$\frac{d\mathrm{\Lambda }(t)}{dt}=\frac{d}{dt}\underset{j=0}{\overset{\mathrm{}}{}}\frac{\left(_0^t\mathrm{\Omega }(\eta )𝑑\eta \right)^j}{j!}=\mathrm{\Omega }(t)\underset{j=0}{\overset{\mathrm{}}{}}\frac{\left(_0^t\mathrm{\Omega }(\eta )𝑑\eta \right)^j}{j!}$$
(3.85)
$$=\mathrm{\Omega }(t)\mathrm{\Lambda }(t)=\mathrm{\Lambda }(t)\mathrm{\Omega }(t).$$
Obviously the condition (3.85) is satisfied for the generators of all the considered transformations, including the Lorentz ones in Eq. (3.79), since the matrix $`N`$ does not depend on $`\psi `$. ($`N`$ depends only on the momenta components perpendicular the direction of the Lorentz boost.)
### 3.3 Equivalent transformations
Now, from the symmetry of the Eq. (3.36) one can obtain the corresponding transformations for the Eq. (3.2). The generators (3.66) satisfy the relations (3.49) and (3.28) - (3.30), it follows that the generators
$$𝐑_\omega (\mathrm{\Gamma })=Y^1(p)𝐑_\omega (\mathrm{\Gamma }_0)Y(p)=R_\omega (\mathrm{\Gamma })+i\frac{d}{d\omega },$$
(3.86)
$$R_\omega (\mathrm{\Gamma })=Y^1(p)R_\omega (\mathrm{\Gamma }_0)Y(p)+iY^1(p)\frac{dY(p)}{d\omega },$$
where $`𝐑_\omega (\mathrm{\Gamma }_0),R_\omega (\mathrm{\Gamma }_0)`$ are generators (3.66) and $`Y(p)`$ is the transformation (3.38), will satisfy the same conditions, but with the relation (3.26) instead of the relation (3.49). Similarly the generators $`K_\xi (\mathrm{\Gamma }_0)`$ in relation (3.68) will be for Eq. (3.2) replaced by
$$K_\xi (\mathrm{\Gamma })=Y^1(p)K_\xi (\mathrm{\Gamma }_0)Y(p).$$
(3.87)
The finite transformations of the Eq. (3.2) and its solutions can be obtained as follows. First let us consider the transformations $`\mathrm{\Lambda }(\mathrm{\Gamma }_0,\omega ,\stackrel{}{u})`$ given by Eqs. (3.74) and (3.79). In accordance with Eq. (3.37) we have
$$\mathrm{\Gamma }(p)=Y^1(p)\mathrm{\Gamma }_0(p)Y(p),\mathrm{\Gamma }(p^{})=Y^1(p^{})\mathrm{\Gamma }_0(p^{})Y(p^{})$$
(3.88)
and correspondingly for the solutions of Eqs. (3.2), (3.36)
$$\mathrm{\Psi }(p)=Y^1(p)\mathrm{\Psi }_0(p),\mathrm{\Psi }^{}(p^{})=Y^1(p^{})\mathrm{\Psi }_0^{}(p^{}).$$
(3.89)
Since
$$\mathrm{\Psi }_0^{}(p^{})=\mathrm{\Lambda }(\mathrm{\Gamma }_0,\omega ,\stackrel{}{u})\mathrm{\Psi }_0(p),$$
(3.90)
then Eqs. (3.89) imply
$$\mathrm{\Psi }^{}(p^{})=\mathrm{\Lambda }(\mathrm{\Gamma },\omega ,\stackrel{}{u})\mathrm{\Psi }(p);\mathrm{\Lambda }(\mathrm{\Gamma },\omega ,\stackrel{}{u})=Y^1(p^{})\mathrm{\Lambda }(\mathrm{\Gamma }_0,\omega ,\stackrel{}{u})Y(p).$$
(3.91)
Similarly, the transformations $`\mathrm{\Lambda }(\mathrm{\Gamma }_0,\xi )`$ given by Eq. (3.73) are for Eq. (3.2) replaced by
$$\mathrm{\Lambda }(\mathrm{\Gamma },\xi )=Y^1(p)\mathrm{\Lambda }(\mathrm{\Gamma }_0,\xi )Y(p).$$
(3.92)
Let us note, all the symmetries of Eq. (3.2) like the transformation (3.92), which are not connected with a change of the reference frame ($`p`$), can be in accordance with the relation (3.34) expressed
$$\mathrm{\Lambda }(\mathrm{\Gamma },X)=Z(\mathrm{\Gamma },X,\mathrm{\Gamma }),$$
(3.93)
where $`Z`$ is defined by Eq. (3.31) and $`X(p)`$ is any matrix for which there exists $`Z(\mathrm{\Gamma },X,\mathrm{\Gamma })^1`$. Further, it is obvious that if we have some set of generators $`𝐑_\omega (\mathrm{\Gamma })`$ satisfying the relations (3.26) and (3.28) - (3.30), then also any set
$$\widehat{𝐑}_\omega (\mathrm{\Gamma })=Z(\mathrm{\Gamma },X,\mathrm{\Gamma })^1𝐑_\omega (\mathrm{\Gamma })Z(\mathrm{\Gamma },X,\mathrm{\Gamma })$$
(3.94)
satisfies these conditions. For the finite transformations one gets correspondingly
$$\widehat{\mathrm{\Lambda }}(\mathrm{\Gamma },\omega ,\stackrel{}{u})=Z(\mathrm{\Gamma }(p^{}),X(p^{}),\mathrm{\Gamma }(p^{}))^1\mathrm{\Lambda }(\mathrm{\Gamma },\omega ,\stackrel{}{u})Z(\mathrm{\Gamma }(p),X(p),\mathrm{\Gamma }(p)).$$
(3.95)
In the same way, the sets of equivalent generators and transformations can be obtained for the diagonalized equation (3.36).
Let us remark, according to the Lemma 14 there exists the set of transformations $`\mathrm{\Gamma }(p)\mathrm{\Gamma }_0(p)`$ given by the relation (3.37). We used its particular form (3.38), but how will the generators
$$𝐑_\omega (\mathrm{\Gamma },X_k)=Z(\mathrm{\Gamma }_0,X_k,\mathrm{\Gamma })^1𝐑_\omega (\mathrm{\Gamma }_0)Z(\mathrm{\Gamma }_0,X_k,\mathrm{\Gamma });k=1,2$$
(3.96)
differ for the two different matrices $`X_1`$ and $`X_2`$? The last relation implies
$$𝐑_\omega (\mathrm{\Gamma },X_1)=Z(\mathrm{\Gamma }_0,X_1,\mathrm{\Gamma })^1Z(\mathrm{\Gamma }_0,X_2,\mathrm{\Gamma })𝐑_\omega (\mathrm{\Gamma },X_2)Z(\mathrm{\Gamma }_0,X_2,\mathrm{\Gamma })^1Z(\mathrm{\Gamma }_0,X_1,\mathrm{\Gamma })$$
(3.97)
and according to the relation (3.32)
$$\mathrm{\Gamma }Z(\mathrm{\Gamma }_0,X_1,\mathrm{\Gamma })^1Z(\mathrm{\Gamma }_0,X_2,\mathrm{\Gamma })=Z(\mathrm{\Gamma }_0,X_1,\mathrm{\Gamma })^1\mathrm{\Gamma }_0Z(\mathrm{\Gamma }_0,X_2,\mathrm{\Gamma })$$
(3.98)
$$=Z(\mathrm{\Gamma }_0,X_1,\mathrm{\Gamma })^1Z(\mathrm{\Gamma }_0,X_2,\mathrm{\Gamma })\mathrm{\Gamma },$$
which means
$$[Z(\mathrm{\Gamma }_0,X_1,\mathrm{\Gamma })^1Z(\mathrm{\Gamma }_0,X_2,\mathrm{\Gamma }),\mathrm{\Gamma }]=0.$$
(3.99)
It follows that there must exists a matrix $`X_3`$ \[e.g. according to implication (3.35) one can put $`X_3=Z(\mathrm{\Gamma }_0,X_1,\mathrm{\Gamma })^1Z(\mathrm{\Gamma }_0,X_2,\mathrm{\Gamma })`$\] so that
$$Z(\mathrm{\Gamma },X_3,\mathrm{\Gamma })=Z(\mathrm{\Gamma }_0,X_1,\mathrm{\Gamma })^1Z(\mathrm{\Gamma }_0,X_2,\mathrm{\Gamma }),$$
(3.100)
then the relation (3.97) can be rewritten
$$𝐑_\omega (\mathrm{\Gamma },X_1)=Z(\mathrm{\Gamma },X_3,\mathrm{\Gamma })𝐑_\omega (\mathrm{\Gamma },X_2)Z(\mathrm{\Gamma },X_3,\mathrm{\Gamma })^1,$$
(3.101)
i.e. the generators $`𝐑_\omega (\mathrm{\Gamma },X_1),𝐑_\omega (\mathrm{\Gamma },X_2)`$ are equivalent in the sense of the relation (3.94).
### 3.4 Scalar product and unitary representations
###### Definition 16
The scalar product of the two functions satisfying Eq. (3.2) or (3.36) is defined:
$$(\mathrm{\Phi }(p),\mathrm{\Psi }(q))=\{\begin{array}{c}0\\ \mathrm{\Phi }^{}(p)W(p)\mathrm{\Psi }(q)\end{array}\mathrm{for}\begin{array}{c}pq\\ p=q\end{array},$$
(3.102)
where the metric $`W`$ is the matrix, which satisfies
$$W^{}(p)=W(p),$$
(3.103)
$$R_\omega ^{}(p)W(p)W(p)R_\omega (p)+i\frac{dW}{d\omega }=0,$$
(3.104)
$$K_\xi ^{}(p)W(p)W(p)K_\xi (p)=0.$$
(3.105)
The conditions (3.104), (3.105) in the above definition imply, that the scalar product is invariant under corresponding infinitesimal transformations. For example for the Lorentz group the transformed scalar product reads
$$\mathrm{\Phi }^{}(p^{})W(p^{})\mathrm{\Psi }^{}(p^{})$$
(3.106)
$$=\mathrm{\Phi }^{}(p)\left(Iid\omega R_\omega ^{}(p)\right)\left(W(p)+d\omega \frac{dW}{d\omega }\right)\left(I+id\omega R_\omega (p)\right)\mathrm{\Psi }(p)$$
and with the use of the condition (3.104) one gets
$$\mathrm{\Phi }^{}(p^{})W(p^{})\mathrm{\Psi }^{}(p^{})=\mathrm{\Phi }^{}(p)W(p)\mathrm{\Psi }(p).$$
(3.107)
According to a general definition, the transformations conserving the scalar product are unitary. In this way the Eqs. (3.104), (3.105) represent the condition of unitarity for representation of the corresponding group generated by $`𝐑_\omega `$ and $`K_\xi .`$
How to choose these generators and the matrix $`W(p)`$ to solve Eqs. (3.104), (3.105)? Similarly as in the case of solution of Eqs. (3.26) and (3.28) - (3.30) it is convenient to begin with the representation related to the canonical equation (3.36). Apparently the generators of the space rotations can be chosen Hermitian
$$R_{\phi _j}^{}(\mathrm{\Gamma }_0)=R_{\phi _j}(\mathrm{\Gamma }_0).$$
(3.108)
Then also for the Lorentz transformations one gets
$$R_{\psi _j}^{}(\mathrm{\Gamma }_0)=R_{\psi _j}(\mathrm{\Gamma }_0),$$
(3.109)
provided that the constant $`\kappa `$ in Eq. (3.58) is real and $`\left|\kappa \right|m`$. Also the generators $`K_\xi `$ can be chosen in the same way:
$$K_\xi ^{}(\mathrm{\Gamma }_0)=K_\xi (\mathrm{\Gamma }_0).$$
(3.110)
It follows, that instead of the conditions (3.104), (3.105) one can write
$$[𝐑_\omega (\mathrm{\Gamma }_0),W(\mathrm{\Gamma }_0)]=0,[K_\xi (\mathrm{\Gamma }_0),W(\mathrm{\Gamma }_0)]=0.$$
(3.111)
The structure of the generators $`𝐑_\omega (\mathrm{\Gamma }_0),K_\xi (\mathrm{\Gamma }_0)`$ given by Eqs. (3.66), (3.68) suggests, that the metric $`W`$ satisfying the condition (3.111) can have a similar structure, but in which the corresponding blocks on the diagonal are occupied by unit matrices multiplied by some constants. Nevertheless, let us note, that the condition (3.111) in general can be satisfied also for some other structures of $`W(\mathrm{\Gamma }_0)`$.
From $`W(\mathrm{\Gamma }_0)`$ we can obtain matrix $`W(\mathrm{\Gamma })`$ the metric for the scalar product of the two solutions of Eq. (3.2). One can check, that after the transformations
$$𝐑_\omega (\mathrm{\Gamma }_0)𝐑_\omega (\mathrm{\Gamma },X)=Z(\mathrm{\Gamma }_0,X,\mathrm{\Gamma })^1𝐑_\omega (\mathrm{\Gamma }_0)Z(\mathrm{\Gamma }_0,X,\mathrm{\Gamma }),$$
(3.112)
$$K_\xi (\mathrm{\Gamma }_0)K_\xi (\mathrm{\Gamma },X)=Z(\mathrm{\Gamma }_0,X,\mathrm{\Gamma })^1K_\xi (\mathrm{\Gamma }_0)Z(\mathrm{\Gamma }_0,X,\mathrm{\Gamma })$$
(3.113)
and simultaneously
$$W(\mathrm{\Gamma }_0)W(\mathrm{\Gamma },X)=Z(\mathrm{\Gamma }_0,X,\mathrm{\Gamma })^{}W(\mathrm{\Gamma }_0)Z(\mathrm{\Gamma }_0,X,\mathrm{\Gamma })$$
(3.114)
the unitarity in the sense of conditions (3.104), (3.105) is conserved - in spite of the fact that equalities (3.108) - (3.110) may not hold for $`R_\omega (\mathrm{\Gamma },X),K_\xi (\mathrm{\Gamma },X)`$.
### 3.5 Space-time representation and Green functions
If we take the solutions of the wave equation (3.2) or (3.36) in the form of the functions $`\mathrm{\Psi }(p)`$, for which there exists the Fourier picture
$$\stackrel{~}{\mathrm{\Psi }}(x)=\frac{1}{(2\pi )^4}\mathrm{\Psi }(p)\delta (p^2m^2)\mathrm{exp}(ipx)d^4p,$$
(3.115)
then the space of the functions $`\stackrel{~}{\mathrm{\Psi }}(x)`$ constitutes the $`x`$ representation of wave functions. Correspondingly, for all the operators $`D(p)`$ given in the $`p`$ representation and discussed in the previous paragraphs, one can formally define their $`x`$ representation:
$$\stackrel{~}{D}(z)=\frac{1}{(2\pi )^4}D(p)\mathrm{exp}\left(ipz\right)d^4p,$$
(3.116)
which means
$$\stackrel{~}{D}\stackrel{~}{\mathrm{\Psi }}(x)=\frac{1}{(2\pi )^4}D(p)\mathrm{\Psi }(p)\delta (p^2m^2)\mathrm{exp}(ipx)d^4p$$
(3.117)
$$=\frac{1}{(2\pi )^4}D(p)\mathrm{exp}(ipx)\stackrel{~}{\mathrm{\Psi }}(y)\mathrm{exp}(ipy)d^4yd^4p=\frac{1}{(2\pi )^4}\stackrel{~}{D}(xy)\stackrel{~}{\mathrm{\Psi }}(y)d^4y.$$
In this way we get for our operators:
$$\mathrm{\Gamma }_0(p)\stackrel{~}{\mathrm{\Gamma }}_0(z)=Q_0\frac{1}{(2\pi )^4}(p^2)^{1/n}\mathrm{exp}(ipz)d^4p;pzp_0z_0\stackrel{}{p}\stackrel{}{z},$$
(3.118)
$$\mathrm{\Gamma }(p)\stackrel{~}{\mathrm{\Gamma }}(z)=\underset{\lambda =0}{\overset{3}{}}Q_\lambda \frac{1}{(2\pi )^4}p_\lambda ^{2/n}\mathrm{exp}(ipz)d^4p,$$
(3.119)
$$𝐑_{\phi _j}(\mathrm{\Gamma }_0)\stackrel{~}{𝐑}_{\phi _j}(\mathrm{\Gamma }_0)=R_{\phi _j}(\mathrm{\Gamma }_0)+i\frac{d}{d\stackrel{~}{\phi }_j};\frac{d}{d\stackrel{~}{\phi }_j}=ϵ_{jkl}x_k\frac{}{x_l},$$
(3.120)
$$𝐑_{\psi _j}(\mathrm{\Gamma }_0)\stackrel{~}{𝐑}_{\psi _j}(z)$$
(3.121)
$$=\frac{1}{(2\pi )^4}\frac{\kappa R_{\phi _j}(\mathrm{\Gamma }_0)+ϵ_{jkl}p_kR_{\phi _l}(\mathrm{\Gamma }_0)}{p_0+\sqrt{p^2\kappa ^2}}\mathrm{exp}(ipz)d^4p+i\frac{d}{d\stackrel{~}{\psi }_j};$$
$$\frac{d}{d\stackrel{~}{\psi }_j}=x_0\frac{}{x_j}x_j\frac{}{x_0}$$
and in the same way
$$𝐑_\omega (\mathrm{\Gamma })\stackrel{~}{𝐑}_\omega (z)$$
(3.122)
$$=\frac{1}{(2\pi )^4}Z(\mathrm{\Gamma }_0,X,\mathrm{\Gamma })^1𝐑_\omega (\mathrm{\Gamma }_0)Z(\mathrm{\Gamma }_0,X,\mathrm{\Gamma })\mathrm{exp}(ipz)d^4p.$$
Apparently, the similar relations are valid also for the remaining operators $`K_\xi ,W,Z,Z^1`$ and the finite transformations $`\mathrm{\Lambda }`$ in the $`x`$representation. Concerning the translations, the usual correspondence is valid: $`p_\alpha i_\alpha `$.
Further, the solutions of the inhomogeneous version of the Eqs. (3.36), (3.2)
$$\left(\mathrm{\Gamma }_0(p)\mu \right)G_0(p)=I,\left(\mathrm{\Gamma }(p)\mu \right)G(p)=I$$
(3.123)
can be obtained with the use of formula (2.27):
$$G_0(p)=\frac{(\mathrm{\Gamma }_0\alpha \mu )(\mathrm{\Gamma }_0\alpha ^2\mu )\mathrm{}(\mathrm{\Gamma }_0\alpha ^{n1}\mu )}{p^2m^2},$$
(3.124)
$$G(p)=\frac{(\mathrm{\Gamma }\alpha \mu )(\mathrm{\Gamma }\alpha ^2\mu )\mathrm{}(\mathrm{\Gamma }\alpha ^{n1}\mu )}{p^2m^2}$$
(3.125)
and Eq. (3.37) implies
$$G(p)=Z(\mathrm{\Gamma }_0,X,\mathrm{\Gamma })^1G_0(p)Z(\mathrm{\Gamma }_0,X,\mathrm{\Gamma }).$$
(3.126)
Apparently, the functions
$$\stackrel{~}{G}_0(x)=\frac{1}{(2\pi )^4}G_0(p)\mathrm{exp}(ipx)d^4p,$$
(3.127)
$$\stackrel{~}{G}(x)=\frac{1}{(2\pi )^4}G(p)\mathrm{exp}(ipx)d^4p$$
(3.128)
formally satisfy Eqs. (3.123) in the $`x`$representation
$$\stackrel{~}{\mathrm{\Gamma }}_0(xy)\stackrel{~}{G}_0(y)d^4y\mu \stackrel{~}{G}_0(x)=\delta ^4(x),$$
(3.129)
$$\stackrel{~}{\mathrm{\Gamma }}(xy)\stackrel{~}{G}(y)d^4y\mu \stackrel{~}{G}(x)=\left[\underset{\lambda =0}{\overset{3}{}}Q_\lambda (i_\lambda )^{2/n}\mu \right]\stackrel{~}{G}(x)=\delta ^4(x).$$
(3.130)
The last equation contains the fractional derivatives defined in . Obviously, the functions $`\stackrel{~}{G}_0,\stackrel{~}{G}`$ can be identified with the Green functions related to $`x`$ representation of Eqs. (3.36), (3.2).
With the exception of the operators $`\stackrel{~}{𝐑}_{\phi _j}(\mathrm{\Gamma }_0),\stackrel{~}{W}(\mathrm{\Gamma }_0)`$ and $`i_\alpha `$ all the remaining operators considered above are pseudo-differential ones, which are in general non-local. The ways, how to deal with such operators, are suggested in ,,. A more general treatise of the pseudo-differential operators can be found e.g. in -. In our case it is significant, that the corresponding integrals will depend on the choice of passing about the singularities and the choice of the cuts of the power functions $`p^{2j/n}`$. This choice should reflect contained physics, however corresponding discussion would exceed the scope of this paper.
## 4 Summary and concluding remarks
In this paper we have first studied the algebra of the matrices $`Q_{pr}=S^pT^r`$ generated by the pair of matrices $`S,T`$ with the structure given by the Definition 1. We have proved, that for a given $`n2`$ one can in the corresponding set $`\{Q_{pr}\}`$ always find a triad, for which Eq. (2.48) is satisfied, hereat the Pauli matrices represent its particular case $`n=2`$. On this base we have got the rule, how to construct the generalized Dirac matrices \[Eqs. (2.51), (2.52)\]. Further we have shown, that there is a simple relation \[Eqs. (2.53), (2.54)\] between the set $`\{Q_{pr}\}`$ and the algebra of generators of the fundamental representation of the $`SU(n)`$ group.
In the further part, using the generalized Dirac matrices we have demonstrated, how one can from the roots of the D’Alambertian operator generate a class of relativistic equations containing the Dirac equation as a particular case. In this context we have shown, how the corresponding representations of the Lorentz group, which guarantee the covariance of these equations, can be found. At the same time we have found additional symmetry transformations on these equations. Further, we have suggested how one can define the scalar product in the space of the corresponding wave functions and make the unitary representation of the whole group of symmetry. Finally, we have suggested, how to construct the corresponding Green functions. In the $`x`$ representation the equations themselves and all the mentioned transformations are in general non-local, being represented by the fractional derivatives and pseudo-differential operators in the four space-time dimensions.
In line with the choice of the representation of the rotation group used for the construction of the unitary representation of the Lorentz group according to which the equations transform, one can ascribe to the related wave functions the corresponding spin - and further quantum numbers connected with the additional symmetries. Nevertheless it is obvious, that before more serious physical speculations, one should answer some more questions requiring a further study. Perhaps the first could be the problem how to introduce the interaction. The usual direct replacement
$$_\lambda _\lambda +igA_\lambda (x)$$
would lead to the difficulties, first of all with the rigorous definition of the terms like
$$\left(_\lambda +igA_\lambda (x)\right)^{2/n}.$$
On the end one should answer the more general question: Is it possible on the base of the discussed wave equations to build up a meaningful quantum field theory?
Acknowledgement: I would like to thank J.Tolar and M. Rausch de Traubenberg for drawing my attention to the cited articles on generalized Clifford algebras and M. Bednář for critical reading of the manuscript. |
warning/0003/cond-mat0003331.html | ar5iv | text | # Critical light scattering in liquids.
## I Introduction
Dynamical critical phenomena manifest themselves in a singular temperature dependence of hydrodynamic transport coefficients . In pure fluids these transport coefficients are the thermal conductivity and the shear viscosity, both diverging on approach of the critical point. In Ref. the field theoretic renormalization group (RG) theory has been used for a quantitative description of this non-analytic behavior and attention was given to the crossover to the analytic behavior in the background further away from $`T_c`$. The thermal conductivity can be measured in two ways, (i) by measuring the temperature difference when a heat current flows through the liquid (this is an experiment at zero wave vector $`k`$), and (ii) by light scattering experiments in the hydrodynamic region, where the wave vector and the temperature dependent correlation length $`\xi `$ fullfil the relation $`k\xi <<1`$. In the last case the thermal diffusivity $`D_T`$ is measured which is related to the thermal conductivity $`\kappa _T`$ by $`D_T=\kappa _T/(\rho C_P)`$ so that for a comparison of the two experimental data the specific heat has to be known. For the thermal diffusivity and the shear viscosity however theoretical calculations show that no other static quantity apart from the correlation length has to be known. This makes these two transport coefficients most suitable to check the dynamical renormalization calculation.
In light scattering experiments in liquids the characteristic frequency $`\omega _c`$, defined as the half width at half height of the central Rayleigh peak, provides useful additional information about the dynamical properties of the system. Far away from the critical point in the hydrodynamic region the characteristic frequency is given by $`\omega _c=D_T(T,\rho )k^2`$. Approaching the critical point a crossover from the hydrodynamic to the so called critical region ($`k\xi >>1`$) takes place and finally at ($`T_c,\rho _c`$) the characteristic frequency is a function of the wave vector alone. Asymptotically near the critical point (for very small values of $`k`$) the power law behavior $`\omega _ck^z`$ is expected with the dynamical critical exponent $`z3`$. Further away, that means for larger wave vector modulus, a crossover to the background behavior with a non singular thermal conductivity takes place. This is described by the van Hove theory where the characteristic frequency behaves as $`\omega _ck^4`$.
It is the aim of this paper to calculate the characteristic frequency in the whole ($`\xi ,k)`$-plane within the non-asymptotic RG theory in order to describe all types of crossover quantitatively. In addition the density dependence of the line width is considered. The non-universal background parameters entering the expression for the characteristic frequency are taken from other dynamical experiments, e.g measurements of the shear viscosity. Recently very precise data became available for xenon from experiments performed in microgravity . This allows also to reconsider the frequency dependence of the shear viscosity within RG theory already discussed in Ref. .
The results for pure fluids are also compared with light scattering experiments in polymer solutions and polymer blends. The non-asymptotic behavior in a mixture is not completely described by the critical model for pure fluids but the asymptotics is the same. Therefore agreement should be found as long as the non-universal dynamic parameters are near to their fixed point values.
## II The dynamic model
The dynamic order parameter correlation function for the gas-liquid transition can be described within the model H , which is a special case of the model H’ described in detail in Ref. , containing dynamic equations for the order parameter $`\varphi _0`$ (the entropy density) and the transverse momentum density $`𝒋_t`$,
$`{\displaystyle \frac{\varphi _0}{t}}`$ $`=`$ $`\stackrel{o}{\mathrm{\Gamma }}^2{\displaystyle \frac{\delta H}{\delta \varphi _0}}\stackrel{o}{g}(\mathbf{}\varphi _0){\displaystyle \frac{\delta H}{\delta 𝒋_l}}+\mathrm{\Theta }_\varphi ,`$ (1)
$`{\displaystyle \frac{𝒋_t}{t}}`$ $`=`$ $`\underset{t}{\overset{o}{\lambda }}^2{\displaystyle \frac{\delta H}{\delta 𝒋_t}}+\stackrel{o}{g}𝒯\left\{(\mathbf{}\varphi _0){\displaystyle \frac{\delta H}{\delta \varphi _0}}{\displaystyle \underset{k}{}}\left[j_k\mathbf{}{\displaystyle \frac{\delta H}{\delta j_k}}_k𝒋{\displaystyle \frac{\delta H}{\delta j_k}}\right]\right\}+𝚯_t,`$ (2)
with fast fluctuating forces $`\mathrm{\Theta }_i`$ and the projector $`𝒯`$ to the direction of the transverse momentum density. The Hamiltonian appearing in the dynamic equations is the normal Hamiltonian of a $`\varphi ^4`$-theory together with the conserved density $`𝒋_t`$ entering quadratically:
$`H`$ $`=`$ $`{\displaystyle d^dx\left\{\frac{1}{2}\stackrel{o}{\tau }\varphi _0^2+\frac{1}{2}(\mathbf{}\varphi _0)^2+\frac{\stackrel{o}{\stackrel{~}{u}}}{4!}\varphi _0^4+\frac{1}{2}a_j𝒋_t^2\right\}}.`$ (3)
As described in Ref. the dynamic equations may be transformed into a dynamic functional leading to dynamic vertex functions which can be calculated in perturbation theory. In general the dynamic scattering function is different from a Lorentzian due to fluctuation effects. This prediction of scaling theory has been observed in ferromagnets and even compared with RG calculations . The same scaling arguments as for the ferromagnet also apply for pure fluids although the deviation from a Lorentzian is expected to be smaller . Moreover it turns out that in one-loop order there are no frequency dependent contributions in the one-loop perturbation terms of the order parameter vertex functions . Therefore the shape of the dynamic correlation function is approximated by a Lorentzian and may be written as
$$\chi _{dyn}(k,\xi ,\omega )=2\chi _{st}\mathrm{}\left[\underset{\varphi \stackrel{~}{\varphi }}{\overset{1}{\stackrel{o}{\mathrm{\Gamma }}}}(k,\xi ,\omega )\right]=\frac{\chi _{st}(k,\xi )}{\omega _c(k,\xi )}\frac{2}{1+y^2}$$
(4)
in one-loop order with $`y=\omega /\omega _c`$ and the characteristic frequency $`\omega _c`$, defined as the half width at half height of the Rayleigh peak. The width is given by the vertex function $`\mathrm{\Gamma }_{\varphi \stackrel{~}{\varphi }}(k,\xi ,\omega =0)`$ so that the unrenormalized characteristic frequency reads
$$\underset{c}{\overset{o}{\omega }}(k,\xi )=\stackrel{}{\mathrm{\Gamma }}k^2(\xi ^2+k^2)(1+\frac{\underset{t}{\overset{2}{\stackrel{}{f}}}}{\xi ^{d4}}d^dp\frac{1}{1+(𝐱𝐩)^2}\frac{\mathrm{sin}^2\theta }{p^2}),$$
(5)
with $`x=k\xi `$, $`\mathrm{\Omega }=\omega /\stackrel{}{\mathrm{\Gamma }}`$ and $`\underset{t}{\overset{}{f}}=\stackrel{}{g}/\sqrt{\stackrel{}{\mathrm{\Gamma }}\underset{t}{\overset{}{\lambda }}}`$ after setting the parameter $`\underset{\varphi }{\overset{}{w}}=\stackrel{}{\mathrm{\Gamma }}/a_j\underset{t}{\overset{}{\lambda }}`$, which is irrelevant under renormalization, to zero. In full analogy to the renormalization of the transport coefficients the pole in the unrenormalized characteristic frequency may be absorbed into Z-factors using field theoretic renormalization group theory. As we get the same Z-factors (and thus the same flow equations for the Onsager coefficient and the mode coupling) as for the transport coefficients we shall skip the details here.
## III The characteristic frequency
### A General expression
After renormalization the characteristic frequency $`\omega _c`$ is finally found to be
$`\omega _c(k,\xi )=\mathrm{\Gamma }(\mathrm{})k^2(\xi ^2+k^2)\left\{1{\displaystyle \frac{f_t^2(\mathrm{})}{16}}\left[5+6x^2\mathrm{ln}(1+x^2)\right]\right\}`$ (6)
in an $`ϵ`$-expansion with $`ϵ=4d`$. The temperature dependence enters via the flow equations for the mode coupling and the Onsager coefficient,
$`f_t^2(\mathrm{})`$ $`=`$ $`f_t^2\left[1+{\displaystyle \frac{\mathrm{}}{\mathrm{}_0}}\left({\displaystyle \frac{f_t^2}{f_0^2}}1\right)\right]^1,`$ (7)
$`\mathrm{\Gamma }(\mathrm{})`$ $`=`$ $`\mathrm{\Gamma }_0\left({\displaystyle \frac{f_0^2}{f_t^2}}{\displaystyle \frac{\mathrm{}_0}{\mathrm{}}}\left[1+{\displaystyle \frac{\mathrm{}}{\mathrm{}_0}}\left({\displaystyle \frac{f_t^2}{f_0^2}}1\right)\right]\right)^{1x_\eta },`$ (8)
with the one-loop fixed point value of the mode coupling $`f_t^2=\frac{24}{19}`$ and the one-loop value of the exponent $`x_\eta =\frac{1}{19}`$. The connection between the flow parameter $`\mathrm{}`$ and the correlation length or the wave vector respectively is found from the matching condition
$$\left(\xi _0^1\mathrm{}\right)^2=\xi ^2+k^2,$$
(9)
for the Lorentzian approximation where the correlation length may be expressed in terms of the reduced temperature $`t`$ via $`\xi =\xi _0t^\nu `$ with $`\nu =0.63`$ along the critical isochore. As described in Ref. we may use the cubic model to include non-critical values of the reduced density. In (7)-(9) $`\mathrm{\Gamma }_0`$ and $`f_0`$ are the initial values of the Onsager coefficient and the mode coupling at an arbitrary reduced temperature $`t_0`$ along the critical isochore, $`\mathrm{}_0`$ is the solution of the matching condition at $`t_0`$ and $`k=0`$ and $`\xi _0`$ is the amplitude of the correlation length. Eq. (9) is the frequency independent matching condition which has been used since the vertex function $`\mathrm{\Gamma }_{\varphi \stackrel{~}{\varphi }}`$, expressing the characteristic frequency in the Lorentzian approximation, is evaluated at zero frequency.
We may rewrite (6) extracting the asymptotic expressions for the Onsager coefficient and the mode coupling,
$$\omega _c(k,x)=\mathrm{\Gamma }_{as}k^z\left(\frac{1+x^2}{x^2}\right)^{1x_\lambda /2}\left[c_{na}(k,x)\right]^{x_\lambda }f(k,x)$$
(10)
with $`x=k\xi `$, $`z=4x_\lambda `$ and the function $`f(k,x)`$ defined as
$$f(k,x)=1\frac{f_t^2}{16c_{na}(k,x)}\left[5+6x^2\mathrm{ln}(1+x^2)\right].$$
(11)
The non-asymptotic contributions are collected in
$$c_{na}(k,x)=\left[1+\frac{k}{k_0}\sqrt{\frac{1+x^2}{x^2}}\right]$$
(12)
so that the asymptotic region is characterized by $`c_{na}(k,x)=1`$. Finally the asymptotic Onsager coefficient $`\mathrm{\Gamma }_{as}`$ and the crossover wave length $`k_0`$ are given by
$$\mathrm{\Gamma }_{as}=\mathrm{\Gamma }_0\left(\frac{f_0^2\mathrm{}_0}{f_t^2\xi _0}\right)^{x_\lambda },k_0^1=\left(\frac{f_t^2}{f_0^2}1\right)\frac{\xi _0}{\mathrm{}_0}.$$
(13)
The advantage of Eq. (10) over Eq. (6) is the clear separation of the asymptotic and the non-asymptotic behavior which will make the discussion of the various limits of the characteristic frequency easier. Before we come to that point in the next section we should remark here that it is also possible to evaluate the crossover function in three dimensions instead of performing an $`ϵ`$-expansion. The characteristic frequency then reads
$`\omega _c(k,\xi )=\mathrm{\Gamma }(\mathrm{})k^2(\xi ^2+k^2)\left\{1+f_t^2(\mathrm{})\left[{\displaystyle \frac{2}{3}}\sqrt{{\displaystyle \frac{1+x^2}{x^2}}}\mathrm{arctan}x{\displaystyle \frac{3}{4}}\right]\right\}`$ (14)
As expressions (6) and (14) are almost identical after choosing the right initial values for the Onsager coefficient and the mode coupling we shall only discuss the $`ϵ`$-expansion result (also used for the evaluation of the transport coefficients in ) in the following.
### B Various limits of the characteristic frequency
First we should note that Eq. (10) yields a finite value for the characteristic frequency in the hydrodynamic limit $`x0`$,
$$\underset{x0}{lim}\omega _c(k,\xi )=\mathrm{\Gamma }_{as}k^2\xi ^{2+x_\lambda }\left[1+\frac{1}{x_0}\right]^{x_\lambda }\left\{1\frac{f_t^2}{16}\left[1+\frac{1}{x_0}\right]^1\right\},$$
(15)
with $`x_0=k_0\xi `$. Here the coefficient of $`k^2`$ is the non-asymptotic expression for the temperature dependent thermal diffusion coefficient $`D_T(\xi )`$ discussed in Ref. so that we can rewrite Eq. 15 in the well-known form $`\omega _c=D_T(\xi )k^2`$ for the hydrodynamic region. Also in the opposite critical limit $`x\mathrm{}`$ we obtain a finite value for the characteristic frequency,
$$\underset{x\mathrm{}}{lim}\omega _c(k,\xi )=\omega _c(k)=\mathrm{\Gamma }_{as}k^z\left[1+\frac{k}{k_0}\right]^{x_\lambda }\left\{1+\frac{5f_t^2}{16}\left[1+\frac{k}{k_0}\right]^1\right\},$$
(16)
which is the wave vector dependent non-asymptotic expression of the characteristic frequency. Both non-asymptotic expressions allow to discuss the crossover from the asymptotic limit $`\xi k_0\mathrm{}`$ or $`k/k_00`$ to the background limit $`\xi k_00`$ or $`k/k_0\mathrm{}`$ respectively.
In the hydrodynamic case we obtain the limits
$`\underset{\xi k_0\mathrm{}}{lim}\omega _c(k,\xi )`$ $`=`$ $`\mathrm{\Gamma }_{as}k^2\xi ^{2+x_\lambda }\left(1{\displaystyle \frac{f_t^2}{16}}\right),`$ (17)
$`\underset{\xi k_00}{lim}\omega _c(k,\xi )`$ $`=`$ $`\mathrm{\Gamma }_0k^2\xi ^2\left(1{\displaystyle \frac{f_0^2}{f_t^2}}\right)^{x_\lambda },`$ (18)
where we used the expression for $`k_0`$ given in Eq.(13) for the last limit. In the background limit our expression reaches the van Hove behavior. In the critical region we obtain
$`\underset{k/k_00}{lim}\omega _c(k)`$ $`=`$ $`\mathrm{\Gamma }_{as}k^z\left(1+{\displaystyle \frac{5f_t^2}{16}}\right),`$ (19)
$`\underset{k/k_0\mathrm{}}{lim}\omega _c(k)`$ $`=`$ $`\mathrm{\Gamma }_0k^4\left(1{\displaystyle \frac{f_0^2}{f_t^2}}\right)^{x_\lambda },`$ (20)
where again we reach the van Hove theory for large values of the ratio $`k/k_0`$. This means that our results for the characteristic frequency describe the crossover in the correlation length (from $`\xi ^{2+x_\lambda }`$ to $`\xi ^2`$) in the hydrodynamic region characterized by the limit $`x0`$ as well as the crossover in the wave vector (from $`k^z`$ to $`k^4`$) in the critical region characterized by the limit $`x\mathrm{}`$.
We have seen that with our non-asymptotic theory we always reach the van Hove behavior in the non-asymptotic limit for large values of the wave vector or small values or the correlation length respectively. This is different from the non-asymptotic mode coupling expression of Olchowy , where the characteristic frequency is given by
$$\omega _c(k,\xi )=\frac{k_BT}{6\pi \overline{\eta }^B\xi }k^2\frac{3}{4}(1+x^2)^{1/2}\left[y_D+y_\delta (1+x^2)^{1/2}\right]$$
(21)
with
$$y_D=\mathrm{arctan}x_D,y_\delta =(1+x_D^2)^{1/2}[y_D+\mathrm{arctan}(x_D(1+x_D^2)^{1/2}]$$
(22)
depending both on the non-universal parameter $`x_D=q_D\xi `$ which is similar to the parameter $`k_0`$ appearing in our non-asymptotic theory. Eq. (21) does not yield the van Hove theory in the non-asymptotic region but instead becomes negative for $`x>2x_D`$. This region of unphysical negative values of the characteristic frequency is always reached at constant correlation length when the wave vector becomes larger than the non-universal parameter $`q_D`$. On the other hand the parameter $`q_D`$ cannot be set to infinity as this limit yields an unphysical divergence in the hydrodynamic limit for $`\xi 0`$ .
### C Discussion of the crossover behavior
In the background we always reach the van Hove behavior for the characteristic frequency. This is a general feature of our non-asymptotic theory. The parameter which describes the crossover from the van Hove expression of the characteristic frequency to its asymptotic expression is in fact the value of the mode coupling $`f_0`$ which can take on values from zero to the fixed point value $`f_t^{}`$. Note that this corresponds to a crossover of $`k_0`$ from its asymptotic limit $`k_0\mathrm{}`$ to its van Hove limit $`k_00`$. The van Hove behavior for $`f_0=0`$,
$$\omega _c^{vH}(k,x)=\mathrm{\Gamma }_0k^4\left(1+x^2\right),$$
(23)
is different from the background behavior at finite $`f_0`$ so that we can define a background van Hove characteristic frequency $`\omega _c^{BvH}`$ as
$$\omega _c^{BvH}(k,x)=\mathrm{\Gamma }_0k^4\left(1\frac{f_0^2}{f_t^2}\right)^{x_\lambda }\left(1+x^2\right),$$
(24)
which we now always reach with our non-asymptotic theory in the background limit $`\xi k_00`$ or $`k/k_0\mathrm{}`$ respectively.
Now we can extract the background van Hove behavior from the full characteristic frequency given in Eq. 10,
$$\omega _c(k,x)=\omega _c^{BvH}(k,x)\left(k\xi _0\right)^{x_\lambda }\left(\frac{f_0^2}{f_t^2}\mathrm{}_0\right)^{x_\lambda }\left(\frac{1+x^2}{x^2}\right)^{x_\lambda /2}\left(\frac{c_{na}(k,x)}{1\frac{f_0^2}{f_t^2}}\right)^{x_\lambda }f(k,x),$$
(25)
and plot the ratio $`\omega _c/\omega _c^{BvH}`$ in order to demonstrate the crossover behavior of the characteristic line width. This is done in Fig. 1 from which we see that the ratio $`\omega _c/\omega _c^{BvH}`$ increases near the critical point (characterized by $`\xi \mathrm{}`$ and $`k0`$) as the characteristic frequency then approaches its asymptotic power law behavior, of course with nonuniversal amplitudes depending on value of the mode coupling $`f_0`$ in the background. This effect increases with increasing values of the mode coupling $`f_0`$. Especially we see that choosing the fixed point value $`f_0=f_t^{}`$ the surface of the characteristic frequency never reaches a flat surface (corresponding to the van-Hove behavior).
The crossover from the asymptotic power-law behavior in the critical region, where the characteristic frequency is proportional to $`k^z`$, to the van Hove behavior with $`\omega _ck^4`$ in the non-asymptotic background region can also be seen in Fig. 2 where we compare our asymptotic and non-asymptotic results with the van Hove theory and the result of Kawasaki . To do this we rewrite Eq. (10) extracting $`k^2`$ instead of $`k^z`$,
$$\omega _c^{as}(k,\xi )=\frac{\mathrm{\Gamma }_{as}}{\xi ^{1+x_\eta }}k^2(1+x^2)^{1x_\lambda /2}\left[c_{na}(k,x)\right]^{x_\lambda }f(k,x)\frac{\mathrm{\Gamma }_{as}}{\xi ^{1+x_\eta }}k^2\mathrm{\Omega }(x),$$
(26)
and compare the various results for the function $`\mathrm{\Omega }(x)/x`$ at constant correlation length instead of $`\omega _c`$ itself. Therefore we have to note that the function $`\mathrm{\Omega }(x)`$ defined in Eq. (26) is in general not only a function of $`x=k\xi `$ but also of the wave vector $`k`$ which enters via the non-asymptotic function $`c_{na}(k,x)`$. But keeping the correlation length constant as in Fig. 2 we can express $`k`$ in terms of $`x`$ so that $`\mathrm{\Omega }(x)`$ is really only a function of $`x`$.
As Kawasaki’s result is proportional to $`k^3`$ instead of $`k^z`$ the function $`\mathrm{\Omega }(x)/x`$ plotted in Fig. 2 becomes constant for large values of $`x`$ whereas our asymptotic result (characterized by $`c_{na}(k,x)=1`$) is proportional to $`x^{x_\eta }`$ and the van Hove theory to $`x`$. Our non-asymptotic results (at constant values of the correlation length $`\xi `$) behave for large values of $`x`$ like the van Hove theory and are therefore proportional to $`x`$. We also see in this figure that the setin of the crossover to the van Hove theory is determined by initial value of the mode coupling $`f_0`$ which is the only free parameter in our non-asymptotic theory. We also should note that in Kawasaki’s theory there is a different prefactor for the function $`\mathrm{\Omega }(x)`$ so that we have normalized the function $`\mathrm{\Omega }(x)/x`$ so that the curves coincide for $`x0`$.
We can also use the function $`\mathrm{\Omega }(x)`$ defined in Eq. (26) to compare our asymptotic result for the characteristic frequency with other theories: In Fig. 3 we have plotted the asymptotic result for $`\mathrm{\Omega }(x)/x`$ (which is only a function of $`x`$ as we have $`c_{na}(k,x)=1`$) as well as the theoretical results of Kawasaki and Lo , Paladin and Peliti and Burstyn et al. . As the other authors have a different prefactor for $`\mathrm{\Omega }(x)`$ we have normalized $`\mathrm{\Omega }(x)`$ so that the curves coincide for $`x0`$. Again we see that the Kawasaki result for $`\mathrm{\Omega }(x)/x`$ becomes constant whereas the other results show the correct $`x^{x_\eta }`$ behavior for large values of $`x`$. In addition to this comparison we should note that at the critical dimension $`d=4`$ our result for the characteristic frequency is identical with the result of Siggia et al. .
And finally let us mention that we can extend our theory to non-critical values of the density and calculate the crossover in the characteristic frequency when we leave the critical isochore: Using the parametric representation to connect the correlation length to the reduced temperature $`t=(TT_c)/T_c`$ and the reduced density $`\mathrm{\Delta }\rho =(\rho \rho _c)/\rho _c`$ we are able to evaluate the correlation length as a function of $`t`$, $`\mathrm{\Delta }\rho `$ and $`k`$. In Fig. 4 we have plotted the characteristic frequency in the hydrodynamic limit for $`k=0`$ as a function of $`t`$ and $`\mathrm{\Delta }\rho `$. We see that the characteristic frequency goes to zero in the critical limit $`t0`$ and $`\mathrm{\Delta }\rho 0`$ corresponding to $`\xi \mathrm{}`$.
## IV Comparison with experiments
### A Pure liquids
In Fig. 5-8 we compare our asymptotic and non-asymptotic results for the characteristic frequency $`\omega _c/k^2`$ as a function of the reduced temperature $`t`$ and the function $`\mathrm{\Omega }(x)/x`$ as a function of $`x`$ with experiments in Xe and $`\mathrm{CO}_2`$ (all non-universal parameters are given in Tab. I). As discussed in Ref. we can treat the exponent $`x_\lambda =1x_\eta `$ as an additional free parameter so that we can fit $`f_0`$ and $`x_\eta `$ from the experimental data (the initial value of the Onsager coefficient $`\mathrm{\Gamma }_0`$ is determined by the value of the shear viscosity at $`t_0`$). But this means that we need additional data for this fit. In Xe we have used the recent shear viscosity data of Berg et al. discussed in the next section. Fitting the parameter $`f_0`$ from the characteristic frequency data (the exact value of $`x_\eta `$ does hardly affect the exponent $`x_\lambda =1x_\eta `$) and the exponent $`x_\eta `$ from the shear viscosity data we find good agreement for the characteristic frequency (Fig. 5 and 6) as well as for the frequency dependent shear viscosity (Fig. 11 and 12). For $`\mathrm{CO}_2`$ we have taken $`t_0`$, $`f_0`$ and $`\mathrm{\Gamma }_0`$ (also given in Tab. I) from the comparison of the shear viscosity and the thermal diffusivity with experiments in Ref. so that the curves shown in Fig. 7 and 8 are obtained without any free parameter!
As we can see in these figures the experimental data are not described correctly by our asymptotic expressions but only by the non-asymptotic expressions which show the crossover to the van Hove theory for large values of the reduced temperature $`t`$ or small values of the variable $`x`$ respectively. Analogously any asymptotic theory fails to describe the experimental data correctly. In Ref. this problem was eliminated adding a regular background contribution of the form $`\omega _c^B=(\lambda ^B/\rho c_p)k^2(1+x^2)`$ to the critical expression for the characteristic frequency with $`\lambda ^B`$ being the regular part of the thermal conductivity and $`c_p`$ the full specific heat at constant pressure containing also critical contributions. The use of the full specific heat together with the term $`1+x^2`$ ensures the crossover to the van Hove theory for large values of the reduced temperature as well as for large values of the wave vector (the background characteristic frequency is proportional to $`k^2\xi ^2`$ for $`x0`$ and to $`k^4`$ for $`x\mathrm{}`$) so that the full characteristic frequency $`\omega _c=\omega _c^C+\omega _c^B`$ obtained by this procedure yields basically the same curves as our non-asymptotic theory (see Fig. 6 of Ref. ). In our theory however we use a different form of the background characteristic frequency: Following the discussion of the regular background added to the transport coefficients we would have to add a background of the form $`\omega ^B=D_T^B(T,\rho )k^2D_T^B(T_c,\rho _c)k^2`$ to our results with the background thermal diffusivity given by $`D_T^B=\lambda ^B/\rho c_p^B`$ and the background specific heat $`c_p^B`$ containing only the regular temperature dependence without the critical singularity. As this background term turns out to be negligibly small in the temperature range shown in Fig. 5-8 we have neglected it so that our asymptotic and non-asymptotic curves for Xe and $`\mathrm{CO}_2`$ contain only the critical contributions discussed in this paper. So the main difference between our non-asymptotic theory and the results of Ref. is that the crossover to the van Hove theory, which is clearly seen in experiments, is already contained in our expressions for the characteristic frequency and not added by an appropriate form of the background contribution!
In Fig. 6 and 8 we also see that the non-asymptotic results for $`\mathrm{\Omega }(x)/x`$ do of course not collapse on a single curve (in contrary to our asymptotic result and the theories of Ref. ) as the non-asymptotic contribution $`c_{na}`$ does not only depend on the variable $`x=k\xi `$ but also on the wave vector $`k`$ and the correlation length $`\xi `$ separately. This behavior can also be seen in the Xe and $`\mathrm{CO}_2`$ data in Fig. 6 and 8 although the experimental data are not precise enough for a true confirmation of the validity of our non-asymptotic theory.
### B Polymer solutions and blends
And finally we apply our theory for the characteristic frequency to light scattering experiments in binary polymer solutions: In Fig. 9 we compare our non-asymptotic theory for the characteristic frequency $`\omega _c/k^2`$ as a function or the reduced temperature with experimental data in a solution of polydisperse polystyrene (PDPS) in cyclohexane . For this figure the initial value of the Onsager coefficient was determined from the value of the background shear viscosity at the critical point also measured in Ref. . The amplitude of the correlation length $`\xi _0`$ as well as the exponents $`\nu =0.7`$ and $`x_\eta =0.065`$ were taken from the same paper. Therefore we have to note that the exponent $`\nu `$ found by Ref. for the polymer solution is higher than the value $`\nu =0.63`$ found for pure liquids or liquid mixtures. Fitting the initial value of the mode coupling $`f_0`$, which is the only free parameter in our theory, from the experimental data (all values given in Tab. I) we reach a satisfactory description of the experimental data although the curves for large wave vectors lie above the experimental data for small values of the reduced temperature. Nevertheless we have to note that the quality of the description cannot be compared to the one reached for Xe and $`\mathrm{CO}_2`$ as there are no detailed experimental data for the shear viscosity of this polymer solution in the vicinity of the critical point available so that an exact determination of $`\overline{\eta }_0`$ and thus of $`\mathrm{\Gamma }_0`$ was not possible and also the critical exponent had to be fixed and could not be fitted from the experiments.
However one crucial point remains: The fact the we have used the non-asymptotic theory developed for pure liquids to describe a polymer solution is of course a problem as liquids and liquid mixtures do have the same asymptotic behavior but show a slightly different crossover to the non-asymptotic behavior. But as an asymptotic theory is not able to describe the experimental data (in the same way as we were not able to describe the characteristic frequency in pure liquids with the asymptotic theory) and a non-asymptotic theory for critical light scattering in mixtures has not yet been set up we believe that the systematic errors made by applying a non-asymptotic theory for pure liquids to mixtures (which basically means setting the additional parameter $`w_3`$ found for the transport coefficients in liquid mixtures to zero) are rather small and can be tolerated. In addition we have to note that a background characteristic frequency given in Ref. was subtracted from the experimental data as well as from the non-asymptotic results for $`\omega _c/k^2`$.
In Fig. 10 we compare our asymptotic result for the function $`\mathrm{\Omega }(x)/x`$ as well as the theoretical results of Kawasaki and Burstyn et al. with experimental data in the polymer blend of polydimethylsiloxane and polyethylmethylsiloxane . As all these data are only available in a rather small range of $`x`$ we can apply the asymptotic theory and avoid the discussion of the last paragraph. The use of a non-asymptotic theory would also not be possible for a comparison with these experimental data for a second reason: All data shown in Fig. 10 were obtained for different temperatures and wave vectors but these different values of $`k`$ and $`\xi `$ were not indicated separately in the paper but only the corresponding value of $`x=k\xi `$. This was also the reason why we could not fit the initial value of the Onsager coefficient so that the only fit parameter, the prefactor of $`\mathrm{\Omega }(x)`$, was set by the choice that our result shall coincide with the result of Burstyn et al. in the limit of small values of $`x`$. In addition we have to note that the experimental values for the function $`\mathrm{\Omega }(x)`$ where obtained from the data for the characteristic frequency dividing not by the full shear viscosity depending on the correlation length but only by its constant background value so that we had to correct this multiplying our theoretical expression for the function $`\mathrm{\Omega }(x)`$ by $`x^{x_\eta }`$. In any case the experimental data shown in Fig. 10 are not precise enough to favor any of the presented theoretical expressions.
## V The frequency dependent shear viscosity
Since we have used information from the shear viscosity in the discussion of the light scattering line width we shall add an analysis of the most recent shear viscosity data for Xe to this paper. These new data allow a much more detailed analysis of the frequency dependent shear viscosity leading to slightly different parameters than the discussion of the shear viscosity of Xe in Ref. which was based on older data. In Ref. we have discussed the theoretical expression for the frequency dependent shear viscosity, which is given by
$$\overline{\eta }(t,\mathrm{\Delta }\rho ,\omega )=\frac{k_BT}{4\pi }\frac{\xi _0}{\mathrm{}f_t^2(\mathrm{})\mathrm{\Gamma }(\mathrm{})}\left[1+E_t(f_t(\mathrm{}),v(\mathrm{}),w(\mathrm{}))\right],$$
(27)
with the one-loop perturbational contribution
$`E_t(f_t(\mathrm{}),v(\mathrm{}),w(\mathrm{}))={\displaystyle \frac{f_t^2}{96}}\{1+6[i{\displaystyle \frac{v^2}{w}}\mathrm{ln}v+{\displaystyle \frac{1}{v_+v_{}}}({\displaystyle \frac{v_{}^2}{v_+}}\mathrm{ln}v_{}{\displaystyle \frac{v_+^2}{v_{}}}\mathrm{ln}v_+)]`$ (28)
$``$ $`{\displaystyle \frac{4}{(v_+v_{})^3}}\left[{\displaystyle \frac{v_+^3v_{}^3}{3}}+{\displaystyle \frac{3}{2}}(v_+v_{})(v_+^2\mathrm{ln}v_++v_{}^2\mathrm{ln}v_{})(v_+^3\mathrm{ln}v_+v_{}^3\mathrm{ln}v_{})\right]`$ (29)
$`+`$ $`{\displaystyle \frac{2}{(v_+v_{})^2}}[{\displaystyle \frac{v_+^3}{v_{}}}(1+4\mathrm{ln}v_+)+{\displaystyle \frac{v_{}^3}{v_+}}(1+4\mathrm{ln}v_{})`$ (30)
$`+`$ $`({\displaystyle \frac{1}{v_{}}}{\displaystyle \frac{2}{v_+v_{}}}){\displaystyle \frac{v_+^4\mathrm{ln}v_+v^4\mathrm{ln}v}{v_{}}}+({\displaystyle \frac{1}{v_+}}+{\displaystyle \frac{2}{v_+v_{}}}){\displaystyle \frac{v_{}^4\mathrm{ln}v_{}v^4\mathrm{ln}v}{v_+}}]\}.`$ (31)
The parameters introduced in Eq. (28) are defined as
$$v(\mathrm{})=\frac{\xi ^2(t)}{(\xi _0^1\mathrm{})^2},w(\mathrm{},\omega )=\frac{\omega }{2\mathrm{\Gamma }(\mathrm{})(\xi _0^1\mathrm{})^4},$$
(32)
$$v_\pm (\mathrm{},\omega )=\frac{v}{2}\pm \sqrt{\left(\frac{v}{2}\right)^2+iw},$$
(33)
with the mode coupling $`f_t(\mathrm{})`$ and the Onsager coefficient $`\mathrm{\Gamma }(\mathrm{})`$ given by Eqs. (7) and (8). The mode coupling parameter $`\mathrm{}`$ is now a function of the correlation length $`\xi `$ and the frequency $`\omega `$ and results from the solution of the matching condition
$$\left(\frac{\xi _0}{\xi }\right)^8+\left(\frac{2\omega }{\mathrm{\Gamma }(\mathrm{})}\right)^2=\mathrm{}^8.$$
(34)
At the moment of publication no experimental data were available to compare them to our theoretical expressions. The situation has changed meanwhile as Berg et al. performed shear viscosity experiments at small frequencies in a microgravity environment onboard a space shuttle. Comparing their experimental results with the mode coupling theory they found that they could only describe their data correctly multiplying the frequency by a factor of 2 in the theoretical expressions. They explained the introduction of this factor as a two-loop effect correcting the errors of the one-loop expression used for the frequency dependent shear viscosity. With this multiplicative factor for the frequency they were able to reproduce the experimental data for the shear viscosity very well.
In Fig. 11 and 12 we compare our theory with experimental data in microgravity and in the earth’s gravitational field fitting the exponent $`x_\eta `$ with $`f_0`$ taken from the light scattering experiments of Ref. . In doing so we found $`x_\eta =0.065`$ instead of the value $`x_\eta =0.069`$ used by Berg et al. We should note here that we can use the exponent $`x_\eta =0.069`$ (with the initial values $`f_0=0.959`$ and $`\mathrm{\Gamma }_0=8.82\times 10^{18}\mathrm{cm}^4/\mathrm{s}`$) to get exactly the same theoretical curves as shown in Fig. 11 and 12 but then we are not able to describe the characteristic frequency data correctly with this choice of $`f_0`$ and $`\mathrm{\Gamma }_0`$. This fact that the parameters $`f_0`$ and $`x_\eta `$ cannot be determined unambiguously from the shear viscosity data alone was already discussed in detail in Ref. .
In Fig. 11 and 12 it turns out that we can describe the experimental data in microgravity only if we multiply the frequency by a factor of 5, which may be justified for the same reason as the factor of 2 in the mode coupling theory . But then we are able to describe not only the microgravity data but also the earth-bound experiments very well with a single set of parameters shown in Tab I. And once again let us mention that we have used the same set of parameters to describe the characteristic frequency in Xe correctly in Fig. 5 and 6. As the experimental data shown in Fig. 12 cover a large range of reduced temperatures we had to add the regular background contribution found in Ref. , which is completely independent of the critical behavior described within our model.
Berg et al. did not only measure the real part of the shear viscosity but determined also the imaginary part of $`\overline{\eta }`$ from the phase shift. In Ref. they compared the mode coupling result for the ratio $`\mathrm{}(\overline{\eta })/\mathrm{}(\overline{\eta })`$ with their experimental data and found good agreement. Comparing our results with these experimental data we get less satisfactory results because in our theory the ratio $`\mathrm{}(\overline{\eta })/\mathrm{}(\overline{\eta })`$ approaches the finite value
$$\underset{TT_c}{lim}\frac{\mathrm{Im}(\overline{\eta })}{\mathrm{Re}(\overline{\eta })}=\frac{1}{76}\frac{\pi }{2}\left[1\frac{1}{76}\left\{3\mathrm{ln}(1/4)1/3\right\}\right]^10.0195$$
(35)
at $`T_c`$ which is different from the value $`0.0353`$ obtained from the mode coupling theory with the exponent $`x_\eta =0.069`$ which turns out to be in good agreement with the experimental data. As the limit of the ratio $`\mathrm{}(\overline{\eta })/\mathrm{}(\overline{\eta })`$ does not contain any free parameter at $`T_c`$ it cannot be improved and the deviation of our theory from the experiments may be explained by the fact that a one-loop order perturbation theory is not able to describe such small effects (the imaginary part of the shear viscosity is only about 3% of the total complex shear viscosity) and therefore a two-loop theory may be expected to yield much better agreement. In this respect we should also note that the mode coupling expression used by Berg et al. is not purely of one-loop order since it makes use of the experimental value for the exponent $`x_\eta `$ which differs significantly from its one-loop value. If we insert the one-loop value $`x_\eta =1/19`$ into the mode coupling expressions we would get a limit $`\mathrm{}(\overline{\eta })/\mathrm{}(\overline{\eta })0.0271`$ at $`T_c`$ which is also significantly lower than the measured limiting ratio. So the main difference between the mode coupling theory and our theory is, that it is not possible to introduce the true critical exponent $`x_\eta `$ in our expression for $`\mathrm{}(\overline{\eta })/\mathrm{}(\overline{\eta })`$ and therefore deviations from the one-loop order perturbation theory cannot be weakened by the use of the correct value for $`x_\eta `$.
## VI Conclusion
We were able to show that our one-loop perturbation theory result for the characteristic frequency evaluated within the field theoretical method of the renormalization group theory does not only reproduce the correct wave vector and correlation length dependence in the hydrodynamic region as well as in the critical region but is also able to describe experimental data correctly for a large range of wave vectors and reduced temperatures. In addition we showed that also the result for the shear viscosity evaluated within the same model is in good agreement with experiments if a two-loop value for the critical exponent is taken.
There are however some points which indicate the need for a two-loop analysis of the model: First we have seen that in one-loop order the dynamic correlation function is always of Lorentzian form whereas scaling theory predicts deviations for large frequencies. Second we are not able to get the experimental limiting value for the ratio of the imaginary and real part of the frequency dependent shear viscosity $`\mathrm{}(\overline{\eta })/\mathrm{}(\overline{\eta })`$ at $`T_c`$ and we have to introduce a multiplicative factor for the frequency in order to describe the experimental data correctly. This makes a profound two-loop analysis inevitable which is currently in progress.
Acknowledgment: This work was supported by the Fonds zur Förderung der wissenschaftlichen Forschung under Project No. P12422-TPH.
Figure caption
Fig. 1
Ratio of the characteristic frequency $`\omega _c`$ divided by the van Hove background expression $`\omega _c^{BvH}`$ for $`f_0=0.1`$ or $`k_0=7.98\times 10^3\AA ^1`$ respectively where the ratio becomes 1 in the background limit $`\xi 0`$ and $`k\mathrm{}`$ and for $`f_0f_t^{}`$ where the van Hove expression is never reached by the asymptotic characteristic frequency.
Fig. 2
Comparison of our asymptotic and non-asymptotic (for various values of $`f_0`$ at constant correlation length $`\xi `$) results for $`\mathrm{\Omega }(x)/x`$ with the Ornstein-Zernike theory and the theoretical result of Kawasaki .
Fig. 3
Comparison of our asymptotic result for $`\mathrm{\Omega }(x)/x`$ with the theoretical results of Kawasaki , Paladin and Peliti , and Burstyn et al. .
Fig. 4
The characteristic frequency $`\omega _c/\mathrm{\Gamma }_0k^2`$ as a function of the reduced temperature $`t`$ and the reduced density $`\mathrm{\Delta }\rho `$ in the hydrodynamic limit for $`k=0`$.
Fig. 5
Comparison of the asymptotic (full lines) and non-asymptotic (dashed lines) expressions for $`\omega _c/k^2`$ with the $`\mathrm{Xe}`$-data of Ref. .
Fig. 6
Comparison of the asymptotic (full lines) and non-asymptotic (dashed lines) expressions for $`\mathrm{\Omega }(x)/x`$ with the $`\mathrm{Xe}`$-data of Ref. .
Fig. 7
Comparison of the asymptotic (full lines) and non-asymptotic (dashed lines) expression for $`\omega _c/k^2`$ with the $`\mathrm{CO}_2`$-data of Ref. .
Fig. 8
Comparison of the asymptotic (full lines) and non-asymptotic (dashed line) expressions for $`\mathrm{\Omega }(x)/x`$ with the $`\mathrm{CO}_2`$-data of Ref. .
Fig. 9
Comparison of the non-asymptotic characteristic frequency $`\omega _c/k^2`$ with the experimental data of Ref. in a polymer solution after subtracting the regular background.
Fig. 10
Comparison of the asymptotic expression for $`\mathrm{\Omega }(x)/x`$ with the experimental data of Ref. in a polymer mixture.
Fig. 11
Comparison of the theoretical expression for the real part of the shear viscosity in microgravity at various frequencies with the experimental data of Ref. . See text for details.
Fig. 12
Comparison of the theoretical expressions for the real part of the shear viscosity in microgravity (at frequencies 0Hz and 2Hz) as well as in earthbound experiments (for two different vessel heights) with the experimental data of Ref. .
Table caption
TableI
Nonuniversal parameters of Xe. |
warning/0003/astro-ph0003408.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Coherent radio Cherenkov emission is a remarkably effective method for detecting high energy particles. The history of the effect goes back to Jelly, who first asked whether cosmic ray air showers might produce a radio signal . Askaryan subsequently predicted a net charge imbalance in air showers, and coherent radio power scaling like the energy of the shower squared. Substantial radio emission from atmospheric electromagnetic cascades was observed more than 30 years ago. . Progress in ultra-high energy air showers has sparked renewed interest, and new observations of radio pulses have been reported recently . The current pilot project RICE uses radio Cherenkov emission to detect $`100TeV`$ and higher energy neutrinos in Antarctic ice. The radio Cherenkov signal is the most efficient known mechanism for detecting neutrinos of $`100TeV`$ and above in solid media, yielding detection volumes of order $`1km^3`$ per radio detector for $`PeV`$ neutrinos on ice targets. At $`PeV`$ energies and above, the neutrino interaction cross sections offer fascinating new tests of Standard Model physics and new physics . Tomography of the Earth is also possible with $`PeV`$-scale neutrinos. Radio Cherenkov signals have also been used to search recently for neutrinos and cosmic rays with energies upwards of $`10^{20}eV`$ impinging on the Moon.
Cherenkov radiation is also an intrinsically interesting and beautiful physical phenomenon. Coherence is a basic feature of electrodynamics, and the coherent enhancement of Cherenkov radiation in the microwave region has been observed in the laboratory . The Argonne wake-field acceleration project has successfully generated extremely large microwave field-strengths by manipulating coherent radiation from an intense electron beam.
Despite a long history, the previous literature apparently does not contain a careful treatment of evolving charge distributions, such as those of electromagnetic showers in air or ice, which incorporates all important features of the problem. The problem is intricate because of a multitude of scales. When an electromagnetic shower evolves, it produces a pancake of charge with a finite thickness, a finite width, probed at a finite wavelength of radiation, and for a finite distance over which the shower is big. All this occurs at a finite distance from the detector. Results on evolving and finite-sized charge distributions are few. Tamm grappled with the problem of a charged particle on a track of limited length in the early days of the theory. Askaryan anticipated a coherence cut-off in air showers at high frequencies of order the inverse pancake size, imposed somewhat by hand. Allan gave physical arguments and order of magnitude estimates based on one of Feynman’s electrodynamic formulas. Kahn and Lerche attempted to resolve the coherence issue using superpositions of infinite tracks. McKay and Ralston and Alvarez-Muñiz and Zas ($`AZ`$) considered the influence of the $`LPM`$ effect at ultra-high energies. Zas, Halsen, and Stanev ($`ZHS`$) , and $`AZ`$ reported results from summing asymptotic far-fields track-by-track in Monte Carlo calculations of great complexity.
We present an approach which incorporates all the scales and allows a general analysis. Main results include an expression for the electric field in a factorized form. The “factorization” occurs when distance scales can be separated: the characteristic size of the moving charge distribution must be substantially smaller than the scale over which the charge develops. This condition is well satisfied for all cosmic ray applications we have examined. A form factor characterizes the moving charge distribution, which multiplies a charge-evolution integral. Not all the scales decouple: subtleties coming under the classic description of Fraunhofer and Fresnel zones need careful treatment. Finally, the generic situation can be summarized by analytic formulae. This is indispensable given the large parameter space. For example, the numerous and varied numerical plots obtained from immense Monte Carlos can be summarized by a few parameters. With the parameters fixed, predictions can be made for any number of circumstances.
## 2 Consequences of Coherence
Before beginning analysis we review a few basics. The well-known Frank-Tamm (1937) formula uses an exact solution to Maxwell’s equations for a uniformly moving charge on an infinitely long track. The solution is kinematic and can be obtained by the trick of boosting the charge to a speed faster than light in the medium. Extension to a track of finite length has pitfalls. Tamm’s 1939 finite-track formula assumes a uniform charge $`e`$ traveling at a uniform velocity $`v`$ along the $`z`$-axis for $`L/2<z<L/2`$. Tamm gives the energy loss $`dP`$ per angular frequency $`d\omega `$ per solid angle $`d\mathrm{\Omega }`$:
$$\frac{d^2P}{d\omega d\mathrm{\Omega }}=\frac{ne^2}{4\pi ^2c^3}(\omega L)^2\mathrm{sin}^2\theta \frac{\mathrm{sin}^2X}{X^2},$$
(1)
where $`X=n\omega L/(2c)(\mathrm{cos}\theta _c\mathrm{cos}\theta )`$ and $`\mathrm{cos}\theta _c=c/(nv)`$; $`n`$ is the index of refraction.<sup>1</sup><sup>1</sup>1We diverge from current practice and use the special symbol $`c`$ to represent the speed of light, otherwise known as “1”. This formula has been cited in the high energy physics literature, and used to interpret experiments observing Cherenkov radiation in the millimeter wavelength range .
Tamm’s finite-track formula includes two competing and distinct physical processes: the Cherenkov radiation of a uniformly moving charge, and bremmstrahlung or acceleration radiation from charges modeled as starting and stopping instantly at the track’s endpoints. The interference of the sudden start and stop contributions with the straight line contribution leads to strong oscillations in the angular distribution. Compared to a typical high energy process the acceleration at the endpoints is fake, that is, the Tamm model is unreliable. This is because a charge created by pair production is accompanied by an opposite charge, which coherently shields the pair from radiating until the oppositely charged partners gradually separate. This “Perkins effect” has been observed and is closely related to the coherence phenomena in QCD of color transparency . At the end of the Cherenkov processes, charges also do not stop instantly, but instead slow gradually to subluminal speeds. While the slowing has stochastic elements, it is better approximated by a uniform decelleration than by a catastrophic disappearance of charge. Of course, the evolution of a cosmic ray shower over many radiation lengths and involving billions of particles is an even smoother macroscopic process. Thus the Tamm formula and related approximations may misrepresent the physics if the artificial treatment of the endpoints play a major role. Conversely, experimental situations with conditions close to those assumed by the Tamm formula can be constructed: Takahashi et al. report on the sudden appearance of a charge in a cavity with metallic boundary conditions, leading to a strong mixing of Cherenkov and boundary-condition effects.
### 2.1 Fraunhofer versus Fresnel
Our study uncovers another, deeper problem with certain asymptotic assumptions of the Tamm-type approach. In a typical application of Cherenkov radiation in high energy physics, we might have a track of length $`L1m`$, observed at a distance $`R1m`$, and in the optical regime with $`R/\lambda L/\lambda >10^6`$. The application to radio detection in ice might have $`R1001000m`$, $`\lambda 0.11m`$, $`L10m`$, with $`L/\lambda 10100`$ with $`R/\lambda `$ even greater. Cosmic ray air showers develop and are observed over tens of kilometers. In all cases, all lengths are large in units of the wavelength. Given $`L`$ large enough for the acceleration contributions to be small, the Tamm formula, Eq. (1), might appear ideal at first sight. Indeed for $`\omega L\mathrm{}`$, the $`\mathrm{sin}^2X/X^2`$ distribution approaches $`2\pi c/(n\omega L)\delta (\mathrm{cos}\theta \mathrm{cos}\theta _c)`$. Integrating over angles we recover the well-tested Frank-Tamm result for an infinitely long track,
$$\frac{d^2P}{d\omega dL}=\omega \frac{e^2}{c^2}\mathrm{sin}^2\theta _c.$$
Yet the Tamm formula is quite inapplicable to such problems. This is evident from the formula’s prediction that the radiated energy will be concentrated in coordinate space at $`\theta =\theta _c`$, up to a small width due to diffraction. For many of the physical situations cited, the energy is actually spread rather uniformly over the length of a cylinder surrounding the charge’s trajectory. This is a broad angular distribution extending over angles $`\mathrm{\Delta }\theta L/R`$, where $`R`$ is the distance to the receiver. The Tamm formula, or any asymptotic far field approximation, does not depend on the distance $`R`$, and cannot describe this simple truth. True enough, the momenta (wave numbers) of photons have directions that may be peaked at $`\theta _k\theta _c`$, but this is not the same thing as the power density $`dP/d\mathrm{\Omega }`$ seen on a sphere surrounding the system going like $`\delta (\mathrm{cos}\theta \mathrm{cos}\theta _c)`$. There is no paradox: if one fixes $`L`$ arbitrarily large, and then moves to an asymptotically distant location $`R\mathrm{}`$, the photons traveling at the Cherenkov angle will appear to come from a point source whose angular size is diffraction-limited. The Tamm formula is derived by taking the limit $`R\mathrm{}`$; once taken, the case of finite $`L/R`$ is unavailable.
To see this breakdown from a different perspective, one can use simple dimensional analysis and geometrical reasoning. The Fourier transform of the electric field $`E_\omega `$ has dimensions of mass. Cherenkov radiation for the “long uniformly moving track” has cylindrical symmetry. At cylindrical radius $`\rho `$ from the track, the energy per length $`2\pi \rho E^2(\rho )`$ remains constant. Dimensional analysis plus the cylindrical symmetry forces the electric field to go like $`\sqrt{\omega /\rho }`$. This is nothing like the usual radiation from accelerated charges which has fields falling like $`1/R`$ in three-dimensional space. And this peculiarity applies out to arbitrarily large distance $`\rho `$, provided the track is long enough. But if $`\rho `$ is taken so large that the radiation appears to emerge from a point source, the $`E`$ field must fall like $`L\omega /\rho L\omega /R`$. (The factor of $`L`$ comes from the linear power per unit length dependence. Momentarily we will examine this in more detail.)
The breakdown of Tamm’s formula is thus due to an interchange of limits. Tamm’s formula is obtained by making the Fraunhofer approximation, which fails under a broad range of finite track-lengths. The Fresnel zone describes a complementary far-field region where the Fraunhofer approximation must be modified. The basic physics of the Fresnel zone for Cherenkov radiation is elementary but requires some care.
### 2.2 The Coherence Zone
Consider (Figure 1) a charge moving on a straight line. Let $`R(t)`$ be the instantaneous distance from the charge to the observation point. Information propagating in the medium at speed $`c_m`$ will arrive simultaneously from the track if $`R/t=c_m`$. This is the Cherenkov condition: $`R/t=v\mathrm{cos}\theta =c_m`$ for velocity $`\stackrel{}{v}`$ oriented at angle $`\theta `$ relative to the direction $`\widehat{R}`$. Note that $`R(t)`$ is the radius from the charge to the point, not the vector position.
Due to the geometry of the track and observation point, uniform motion produces acceleration of $`R(t)`$. If $`R/t`$ were constant, the fields arriving would all be in phase for the whole track length. However, the acceleration relative to the observation point produces an extra radial change of order $`\mathrm{\Delta }R=1/2(^2R/t^2)(\mathrm{\Delta }t)^2`$. Coherence of modes of wavelength $`\lambda `$ is then maintained only over a finite region of $`\mathrm{\Delta }R<\lambda .`$ Since $`^2R/t^2=v^2\mathrm{sin}^2\theta /R`$, we solve to find the condition $`\mathrm{\Delta }t_{coh}<\sqrt{R\lambda }/(v\mathrm{sin}\theta )`$. Equivalently, there is a finite spatial coherence region for given $`R`$, given wavenumber $`k=2\pi /\lambda `$, namely
$$\mathrm{\Delta }z_{coh}<\sqrt{R/(k\mathrm{sin}^2\theta )}$$
over which the “sonic boom” of radiation is built coherently. Since $`\mathrm{\Delta }z_{coh}\sqrt{R}`$, the coherence zone grows to infinite size as $`R\mathrm{}`$: but this limit cannot be taken carelessly.
### 2.3 Coherence of Evolving Charge Distributions
We now return to the emission from an electromagnetic shower or other time-evolving charge distribution. To a reasonable approximation, the number of particles in a highly relativistic shower scales like the primary energy divided by a suitable low energy threshold. The charge imbalance near the shower maximum is of order $`20\%e`$ of the total number of particles. These numbers have been confirmed over and over, with each generation of numerical simulation contributing further detail. The origin of the emitted Cherenkov power going like the shower-energy squared is basic electrodynamics: the electric field will scale like the charge, and the radiated power scales like the electric field-squared.
The evolving shower has a finite length scale $`a`$ over which it is near its maximum, and radiating copiously. This length scale, known as the “longitudinal spread” in cosmic ray physics , is akin to the length scale $`L`$ of Tamm’s approach but represents a smooth onset and decline of maximum power. The shower maximum-length scale $`a`$ is determined by the material, and is conceptually distinct from the shower’s total depth to reach the maximum (which goes like the logarithm of the energy) or the charged pancake size (which is fairly constant once the shower is developed.) A cartoon of these ideas is given in Figure 1.
There are then two characteristic limits. Suppose the longitudinal spread of the shower is “short” compared to the coherence length, $`a\mathrm{\Delta }z_{coh}`$. Then coherence is maintained over the whole range where the current is appreciable. The amplitude is proportional to the total length $`a`$ over which the current was “on”, times $`1/R`$. This is then the $`R\mathrm{}`$ limit, or Fraunhofer approximation, and looks like normal radiation. From dimensional analysis (and here we recall the discussion earlier), $`E_\omega a\omega /R`$.
However, in the limit $`a\mathrm{\Delta }z_{coh}`$, the coherence length is not as big as the longitudinal spread, and coherence only exists over the smaller of the two. Adding amplitudes only over the region $`\mathrm{\Delta }z_{coh}`$ and weighted by $`1/R`$, the $`E_\omega `$ field goes like $`\omega \mathrm{\Delta }z_{coh}/R=\sqrt{\omega /R}`$. This behavior is rather different from the previous case: indeed Cherenkov radiation is fundamentally a Fresnel-zone effect, as seen by the $`1/\sqrt{R}`$ dependence of the fields.
Both the Fresnel and Fraunhofer limits are far-field approximations, in the sense that $`kR1`$ is assumed. The subtlety lies in the dimensionless ratio
$$\eta =(a/\mathrm{\Delta }z_{coh})^2=\frac{ka^2}{R}\mathrm{sin}^2\theta $$
which controls how the limit $`R\mathrm{}`$ is taken. Confusion on this point is easy; one has $`R/\lambda 10^4`$ in the same regime, and yet $`R`$ is not large enough for a “large $`R`$” Fraunhofer approximation to apply, exactly because the term “large R” is undefined until the limit parameter $`\eta `$ is specified. In the RICE experiment one typically has $`a12m`$, $`\omega 1001000MHz`$, and $`R10^3m`$, so $`\eta <1`$ holds. Extension to closer observation points, or to energies where the $`LPM`$ effect can give a much larger $`a`$, makes $`\eta 1`$ possible. Partly due to the obscurity of the coherence criteria, the Fraunhofer approximation has received much attention in the previous literature , except for those estimates using fields with ”cylindrical” symmetry .
### 2.4 General Set-Up
Let $`\stackrel{}{E}_\omega (x)`$ be the time-Fourier transform of the $`E`$ field, with similar notation for other fields. The Maxwell equations for a dielectric medium are
$`\stackrel{}{}\stackrel{}{D}_\omega =4\pi \rho ,c\stackrel{}{}\times \stackrel{}{B}_\omega =4\pi \stackrel{}{J}_\omega i\omega \stackrel{}{D}_\omega ,\stackrel{}{}\stackrel{}{B}_\omega =0,c\stackrel{}{}\times \stackrel{}{E}_\omega =i\omega \stackrel{}{B}_\omega ,`$ where $`\stackrel{}{D}_\omega (x)=ϵ(\omega )\stackrel{}{E}_\omega (x)`$. There is a wave equation for the vector potential $`A_\omega ^\mu (x)`$, given by $`c(^2+k^2)A_\omega ^\mu (x)=4\pi J_\omega ^\mu (x)`$, with $`k=\omega \sqrt{ϵ}/c`$. Then we have
$$c\stackrel{}{A}_\omega (\stackrel{}{x})=d^3x^{}\frac{\mathrm{exp}(ik|\stackrel{}{x}\stackrel{}{x}^{}|)}{|\stackrel{}{x}\stackrel{}{x}^{}|}𝑑t^{}\mathrm{exp}(i\omega t^{})\stackrel{}{J}(t^{},\stackrel{}{x}^{}).$$
(2)
The $`4`$-potential $`A^\mu =(A^0,\stackrel{}{A})`$ has been defined in a generalized Lorentz gauge $`c\stackrel{}{}\stackrel{}{A}+ϵA^0/t=0`$ appropriate to the medium. The $`4`$-current $`J^\mu =(\rho /ϵ,\stackrel{}{J})`$, $`\rho `$. Since the components of $`\stackrel{}{J}`$ are related by $`\stackrel{}{J}=\stackrel{}{v}\rho `$, we have $`\stackrel{}{A}=A^0ϵ\stackrel{}{v}/c`$. We calculate $`\stackrel{}{A}_\omega `$ and then use $`A_\omega ^0=\stackrel{}{v}\stackrel{}{A}c/(ϵv^2)`$. For radiation problems the denominator factor $`1/|\stackrel{}{x}\stackrel{}{x}^{}|`$ is replaced by $`1/R`$. This is standard, with corrections of order $`a^2/R^2`$ or similar effects in the “near field” regime, which is not our subject.
### 2.5 Evading The Fraunhofer Approximation
The Fraunhofer approximation is the textbook expansion for the phase
$`\mathrm{exp}(ik|\stackrel{}{x}\stackrel{}{x}^{}|)\mathrm{exp}(ik|\stackrel{}{x}|ik\widehat{x}\stackrel{}{x}^{})`$, dropping terms of order $`k|\stackrel{}{x}^{}|^2/R`$. All existing simulations make this approximation for the phase, and for good reasons: the subsequent integrations become much simpler. The integrand in Eq. (2) oscillates wildly. A Monte Carlo simulation has to find the surviving phases from myriad cancellations, due to the phases generated over the length of each track, and then summed over thousands to millions of tracks moving in three dimensions. For a $`1TeV`$ shower the code of $`ZHS`$ runs in about 20 minutes on a workstation. Increasing the energy by a factor of 100, the calculational time scales up faster than linear, and computer time becomes prohibitive. For this reason various strategies to rescale the output have been used in arriving at the published values of electric fields. Even for $`TeV`$ energies, standard Monte Carlo routines such as GEANT challenge a work-station’s capacity. For cosmic rays of the highest energies the entire approach of direct numerical evaluation is unfeasible.
Unfortunately the Fraunhofer approximation also neglects terms in the phases, namely $`k|\stackrel{}{x}^{}|^2/R`$, that may be of order unity given our previous discussion of length and frequency scales. We must avoid this step. Progress is possible due to the translational features of the macroscopic current
$`J^\mu (t^{},\stackrel{}{x}^{})`$. A rather general model is
$$\stackrel{}{J}(t^{},\stackrel{}{x}^{})=\stackrel{}{v}n(z^{})f(z^{}vt^{},\stackrel{}{\rho }^{}).$$
(3)
An even more general situation will be discussed shortly. The charge packet travels with the speed $`\stackrel{}{v}`$, chosen here to be along the $`z`$-axis of the coordinate system. The function $`f(z^{}vt^{},\stackrel{}{\rho }^{})`$ represents a normalized charge density of the traveling packet, with $`\stackrel{}{\rho }^{}`$ the transverse cylindrical coordinate relative to the velocity axis. We normalize $`f`$ by $`𝑑z^{}d^2\rho ^{}f(z^{},\stackrel{}{\rho }^{})=1.`$ In ice the packet is about $`\mathrm{\Delta }z^{}=10cm`$ thick in the longitudinal direction, and $`\mathrm{\Delta }\rho ^{}=10cm`$ in radius (the Moliere radius) in the vicinity of the shower maximum. These size scales are limited because of relativistic propagation. The time evolution of these scales is negligible near the shower maximum, and indeed the Moliere radius is usually approximated by a material constant for the whole shower. Similarly, in air showers the scale of charge separation is small compared to the scale of shower longitudinal spread.
The shower’s net charge evolution appears in the factor $`n(z^{})`$. With our normalization, $`n(z^{})`$ represents the total charge crossing a plane at $`z^{}`$. The symbol $`n_{max}`$ will denote the maxiumum value of $`n(z^{})`$; later we will see that the electric field scales linearly with $`n_{max}`$. The longitudinal spread $`a`$ is a property of $`n(z^{})`$ near the shower maximum. The model neglects charge (current) left behind, and moving at less than light-speed in the medium, which does not emit Cherenkov radiation. We do not have a sharp cut-off at the beginning or end of tracks, and the function $`n(z^{})`$ will vary smoothly.
### 2.6 Factorization for the Fresnel Zone
Now while we cannot expand around $`\stackrel{}{x}^{}=0`$ (the Fraunhofer approximation), the conditions of the problem do permit an expansion around $`\stackrel{}{\rho }^{}=0`$ (the shower axis), namely for $`R(z^{})=[(zz^{})^2+\rho ^2]^{1/2}`$, that
$`|\stackrel{}{x}\stackrel{}{x}^{}|`$ $`=`$ $`[(zz^{})^2+(\stackrel{}{\rho }\stackrel{}{\rho }^{})^2]^{1/2},`$
$`=`$ $`R(z^{}){\displaystyle \frac{\stackrel{}{\rho }\stackrel{}{\rho }^{}}{R}}+𝒪({\displaystyle \frac{\rho ^2}{R}}).`$
For typical values in this problem, the second term is $`10`$ times smaller than the first, and the third is $`10^3`$ smaller than the second. For the exponent in Eq. (2), the third term does not contribute if $`k\mathrm{\Delta }\rho ^2/R1`$, that is $`\omega 250GHz`$.
Collecting terms, we have
$`cR\stackrel{}{A}_\omega =\stackrel{}{v}{\displaystyle 𝑑z^{}n(z^{})\mathrm{exp}[i(\frac{\omega }{v}z^{}+kR)]}`$
$`\times {\displaystyle }{\displaystyle }dt^{}d^2\rho ^{}\mathrm{exp}\{i[{\displaystyle \frac{\omega }{v}}(z^{}vt^{})+\stackrel{}{q}\stackrel{}{\rho }^{}]\}f(z^{}vt^{},\stackrel{}{\rho }^{}).`$ (4)
We have shifted the $`t^{}`$ integral which produces the translational phase in the $`z^{}`$ integral. This gives the factorization:
$$\stackrel{}{A}_\omega F(\stackrel{}{q})\stackrel{}{A}_\omega ^{FF}(\eta ),$$
(5)
where
$$F(\stackrel{}{q})=d^3x^{}e^{i\stackrel{}{q}\stackrel{}{x}^{}}f(\stackrel{}{x}^{}),$$
(6)
$$vcR\stackrel{}{A}_\omega ^{FF}=\stackrel{}{v}I^{FF},$$
(7)
and
$`I^{FF}(\eta ,\theta )={\displaystyle 𝑑z^{}\mathrm{exp}[\varphi (z^{})]},`$
$`\varphi (z^{})=ik(z^{}\mathrm{cos}\theta _c+R(z^{},\rho ))+\mathrm{log}n(z^{}).`$ (8)
Here $`\stackrel{}{q}=(\omega /v,\stackrel{}{q}_{})`$, $`\stackrel{}{q}_{}=k\stackrel{}{\rho }/R`$, and $`\stackrel{}{x}^{}=(z^{},\stackrel{}{\rho }^{})`$. Provided $`F(\omega )1`$ in either frequency region $`k\mathrm{\Delta }\rho ^2/R1`$ or $`k\mathrm{\Delta }z^2/R1`$, the decoupling of the integrals is excellent. Here $`\mathrm{\Delta }\rho ^{}`$ and $`\mathrm{\Delta }z^{}`$ refer to the regions over which the charge exists near the maximum. $`F(\stackrel{}{q})`$ is the form factor of the charge distribution, which happens to be defined, just as in the rest of physics, in terms of the Fourier transform of the snapshot of the distribution. From our definitions $`F(0)=1`$.
It is worth noting that the dependence on orientation of $`\stackrel{}{q}`$ is observable. For example, in a Giant Air Shower, where the mechanism of charge separation might cause an azimuthal asymmetry about the shower axis labeled by a dipole $`\stackrel{}{p}`$, then $`F(\stackrel{}{q})`$ depends on $`F(\stackrel{}{q}\stackrel{}{p})`$. The orientation of the dipole relative to the observation point thus has a strong effect on the emission. (Other numerically large effects will also be important: for example Allan incorrectly assumes the fields go like $`1/R`$ from his use of Feynman’s formula).
As a consequence of separating out the form factor, the integrations have become effectively one-dimensional.
## 3 Numerical Work
At this stage we have a formula for the vector potential which is a product of a form factor and an object $`I^{FF}(\eta ,\theta )`$ containing the information about the shower history. We will denote $`I^{FF}(\eta ,\theta )`$ the Fresnel-Fraunhofer integral because it interpolates between these regimes. In the Fraunhofer approximation it is easily shown that the factorization is an exact kinematic feature of translational symmetry as exemplified in Eq. (3). If one makes a one-dimensional approximation, the Fraunhofer integral then evaluates the Fourier transform of the current.
The factorization in the Fresnel zone is more demanding, yet should be an excellent approximation. When calculating $`I^{FF}(\eta ,\theta )`$, we cannot (as mentioned earlier) consistently expand in powers of $`z^{}/R`$ because $`\eta 1`$ will be needed.
It makes sense at this point to make a numerical comparison with previous work. Summarizing the results of an extensive Monte Carlo calculation in the Fraunhofer approximation, $`ZHS`$ gave a numerical fit to the electric field
$$\frac{R|\stackrel{}{E}_\omega ^{ZHS}(\theta )|}{F^{ZHS}(\nu )}=1.1\times 10^7\frac{\nu }{\nu _0}\frac{E_0}{TeV}\mathrm{exp}[(\frac{\theta \theta _c}{\mathrm{\Delta }\theta })^2][\frac{V}{MHz}],$$
(9)
with $`\nu _0=500MHz`$. This is the result of a global fit to many angles, energies and frequencies $`\nu =\omega /2\pi <\nu _0`$. In this convention $`\omega `$ is positive. The normalized form factor is $`F_{ZHS}(\nu )=1/[1+0.4(\nu /\nu _0)^2]`$, as discussed below. In making the calculation, results for the field were also rescaled due to computer limitations. As a result, the field reported is strictly linear in the primary energy $`E_0`$. We will comment on this shortly.
We calculated our own result proportional to $`I^{FF}(\eta =0,\theta )`$ over a range of many frequencies and angles (Figure 2). Before doing the integral we scale out the electromagnetic and dimensional factors which are obvious. In our convention $`\mathrm{}<\omega <\mathrm{}`$. The results of our numerical integration are quite well fit by
$$\frac{R|\stackrel{}{E}_\omega ^{\eta =0}(\theta )|}{F(\omega )}=\frac{e}{c^2}a\sqrt{2\pi }n_{max}\omega \mathrm{sin}\theta \mathrm{exp}[\frac{1}{2}(ka)^2(\mathrm{cos}\theta \mathrm{cos}\theta _c)^2]$$
(10)
Putting in numerical values, this gives:
$$\frac{R|\stackrel{}{E}_\omega ^{\eta =0}(\theta )|}{F(\omega )}=2.09\times 10^7\frac{a}{m}\frac{n_{max}}{1000}\frac{\nu }{GHz}\mathrm{exp}[\frac{1}{2}(\frac{\mathrm{cos}\theta \mathrm{cos}\theta _c}{\mathrm{\Delta }(\mathrm{cos}\theta )})^2][\frac{V}{MHz}],$$
(11)
where
$$\mathrm{\Delta }(\mathrm{cos}\theta )=0.048\frac{2}{\sqrt{ϵ}}\frac{m}{a}\frac{GHz}{\nu }.$$
( We have indicated that $`a`$ is in units of $`m.`$ and $`\nu `$ in $`GHz`$.) Note the linear dependence on $`a\nu `$, argued earlier to come from dimensional analysis applied to the limit $`\eta 0`$. A cursory inspection shows that this result and the Monte Carlo have the same general features.
To continue the numerical comparison we need numbers for the longitudinal spread parameter $`a`$ and the number of charges at shower maximum $`n_{max}`$. There are several ways to estimate this. Running the $`ZHS`$ code many times and fitting the output of a $`1TeV`$ shower with a cutoff of $`611KeV`$ gives $`a=1.5m`$, $`n_{max}=345`$. Using these and allowing for the factor of two in conventions gives agreement to a few percent in normalization with $`ZHS`$. However, the other way to do the calculation is to evaluate the product $`an_{max}`$ many times. This method is preferred because fluctuations in $`a`$ and $`n_{max}`$ are correlated. Doing this gives $`an_{max}=(570\pm 50)m`$ at $`1TeV`$, which would predict a normalization factor of $`(1.2\pm 0.1)\times 10^7`$ in Eq. (9). This (plus the angular dependence studied below) indicates that the factorized result is quite consistent with the Monte Carlo.
In Figure 2 we show numerically integrated values of $`I^{FF}`$ Eq. (8). These factors appear directly in the fit just cited, and serve to check the formulas. The form factor has been divided out. For the range of parameters relevant to the problem, agreement is very good, and relative error is much less than $`1\%`$.
For experimental purposes one would like independent confirmation of the parameters from another source. Net particle evolution is well described by Greissen’s classic solution, which was simplified further by Rossi to a Gaussian, $`n(z)=n_{max}\mathrm{exp}[z^2/(2a^2)]/\sqrt{2\pi }a`$. While Greissen’s $`a`$ refers to the whole shower, it should also be a reasonable description for the longitudinal spread of the charge imbalance, which tends to be a fixed fraction of the total number of particles after a few radiation lengths. (There is one caution that the Greissen formula does not explicitly include low-energy physics important for the charge imbalance.) The particular Greissen formula we consulted for the longitudinal spread in radiation lengths $`X_0`$ gives $`a/X_0\sqrt{2/3log(E_0/E)}`$ for particles in the shower with energy greater than $`E`$ and a primary with energy $`E_0`$. In that case one estimates $`a=1.8m`$ at $`E_0=1TeV`$ with $`E=611MeV,X_0=0.39m`$ in ice, which is acceptably close to the previous estimates.
At higher energies there is every reason to believe that Greissen’s
stretched $`\sqrt{\mathrm{log}E}`$ energy dependence will apply. In that case $`a=2.1m`$, $`2.3m`$ for $`100TeV`$, $`1PeV`$ showers, respectively. Note that the product $`n_{max}a`$ is relevant for the field normalization. In this case we also need $`n_{max}1/\sqrt{\mathrm{log}(E_0/E)0.33}`$ from the same Greissen approximation. Rather amazingly, the product $`n_{max}a(E_0/E)\sqrt{2/3\mathrm{log}(E_0/E)}/\sqrt{\mathrm{log}(E_0/E)0.33}E_0/E`$ at high energies. This confirms the phenomenon observed by $`ZHS`$: the normalization of the electric field (Fraunhofer approximation, $`\theta =\theta _c`$) scales precisely linearly in the primary energy. It is rather pleasing that the result can be understood from first principles . Later we will see that the parameter $`a`$ enters in a much more complicated way in the Fresnel zone, creating an extra, weak energy dependence.
Regarding the angular dependence, our work (Figures 3) indicates a general dependence on $`cos\theta cos\theta _c`$ rather than $`\theta \theta _c`$. When fitting numerical output the two functional forms are rather different, unless one has a very narrow distribution. Linearizing for small $`cos\theta cos\theta _c`$ with $`a=1.5m`$ for the comparison, we would predict the scale in the angular dependence $`\mathrm{\Delta }\theta =2.1^o(\nu _0/\nu )`$ while $`ZHS`$ have the same expression with $`2.4^o`$. We find that $`\mathrm{\Delta }\theta `$ is proportional to $`1/a`$. If $`a`$ grows slowly with energy, as Greissen’s formula indicates, then the angular width decreases, which is not seen in $`ZHS`$. Another possible explanation for the small discrepancy is the improper radiation from tracks terminating abruptly at the ends used in the Monte Carlo. We have identified these effects as responsible for the small oscillations seen in the Monte Carlo output, an effect apparently too small to measure.
When numerical output to the frequency dependence is fit, there is a slight coupling between the model for the form factor and the parameter $`a`$ one will extract for the longitudinal spread. We made our own fit to the $`ZHS`$ code’s frequency dependence including the region up to $`1GHz`$, using a Gaussian form factor because of its better analytic properties. (The $`ZHS`$ fit, which contains poles in both the upper and lower half-plane, violates causality.) Specifically, we find
$$F(\nu )=\mathrm{exp}[\nu ^2/(2\nu _{}^2)],\nu _{}=0.93GHz.$$
Using the corresponding $`a`$ value we would predict the Fraunhofer Monte Carlo $`\mathrm{\Delta }\theta 2.3^o(500MHz/\nu )`$, quite close to $`ZHS`$.
In real life, shower to shower fluctuations are highly important. We studied the statistical features<sup>2</sup><sup>2</sup>2We thank Soeb Razzaque for help with this. of the parameters $`n_{max}`$, $`a`$ by fitting individual showers many times and looking at the average and $`rms`$ fit values. The results at $`E=1TeV`$ were $`a=1.5\pm 0.2m`$, $`n_{max}=345\pm 60`$. Multiplying these and adding fluctuations in quadratures gives $`an_{max}=520(1\pm 0.22)m`$. The combination $`an_{max}`$, which is the primary variable in determining the normalization of the electric field, was found to be $`(570\pm 50)m`$. The fluctuation of $`an_{max}`$ is less than half the value that the uncorrelated fluctuations of the separate terms would give. The relative fluctuation is said to decrease with increasing energy, but there are uncertainties. For example, threshold rescaling is used in Monte Carlos, leading to loss of information about the true fluctuations. Very preliminary results of running the standard Monte Carlo GEANT show variations in average shower parameters such as $`an_{max}`$ at the 30$`\%`$ level compared to the average of $`ZHS`$. These comparisons indicate that the electrodynamics is probably determined better than the rest of the problem. Indeed, the deviations from Gaussian behavior in showers is an effect contributing to the fields at the few percent level. At the level of $`1015\%`$, many other small effects contribute. Unless one uses details about the uncertainties and errors in fits, and especially about the shower-to shower fluctuations in all relevant quantities, it is pointless to fine-tune the comparison further. We conclude from the numerical work that the factorized expression is at least as reliable as the Monte Carlos, and has the attractive feature that the parameters can be adjusted directly.
## 4 The Saddle Point Approximation
While the Rossi-Greissen Gaussian approximation to the shower is common, there are additional features which favor such an approach to the emitted radiation. The coherence is dominated by regions where the phases add constructively, greatly enhancing the peak region. In such circumstances analysis is helpful, especially when the largest contribution to $`I^{FF}`$ is dominated by saddle points. These are points where the phase is stationary, $`d\varphi (z_{})/dz_{}=0`$.
Here we describe the saddle-point method to evaluate $`I^{FF}(\eta ,\theta )`$. This is a classic, controlled approximation when the charge distribution has a single maximum and $`kR1`$. The method turns out to give the exact result in the limit of flat charge evolution, that is, the Frank-Tamm formula, where numerical evaluation is highly unstable. With the saddle-point approximation, we can extract analytic formulas which are as good as the numerical integration.
We now describe the saddle-point features. By translational symmetry the shower maximum can be located at $`z^{}=0`$. Referring to the formula of Eq. (8) cited earlier, the saddle point of the phase is given by solving
$$\mathrm{cos}\theta _c(zz_{})/R(z_{})(i/k)(d/dz_{})\mathrm{log}n(z_{})=0$$
for the point $`z_{}`$ dominating the integral. The maximum electric field is already known to occur at points $`(z,R)`$ near the Cherenkov cone. For such observation points $`z=R\mathrm{cos}\theta _c`$, and the saddle point equation has an easy solution at $`z_{}=0`$. Thus the dominant integration region is near the shower maximum, as physically expected.
As the point of observation moves off the Cherenkov cone, the saddle point moves away into the complex $`z^{}`$ plane. To find the complex saddle-point, we approximate $`\mathrm{log}n(z^{})z^2/(2a^2)`$ in the vicinity of the shower maximum; that is, we fit the top of the shower locally with a Gaussian. To re-iterate: the saddle-point approximation does not need to replace the entire shower by a Gaussian, but replaces the vicinity of the region where phases are contributing coherently by a Gaussian. The saddle-point condition gives a quartic equation which can be solved. Unfortunately the solution is impossibly complicated, thwarting a direct approach. We circumvented this by studying the saddle-point location numerically. We found the quartic solution is accurately linearized in a special variable: expand about $`\mathrm{cos}\theta `$ close to $`\mathrm{cos}\theta _c`$. We then find $`z_{}R(0)\mathrm{sin}^2\theta (\mathrm{cos}\theta \mathrm{cos}\theta _c)[1+iR(0)/ka^2\mathrm{sin}^2\theta ]^1`$. With this formula the reader inclined can repeat all the calculations. We also show the saddle-point to highlight the appearance of the ratio $`R(0)/(ka^2\mathrm{sin}^2\theta )`$ indicated by the qualitative “acceleration” argument. Finally, replace $`R(0)=R=\sqrt{z^2+\rho ^2}`$.
The rest of the calculation is standard mathematical physics , so we just quote the results. Calculating fields from the potential and keeping only the leading terms in $`1/(kR)1`$ we find:
$`\stackrel{}{E}_\omega `$ $`=`$ $`{\displaystyle \frac{i\omega }{Rc^2}}F(\stackrel{}{q})I^{FF}(\eta ,\theta )[(\mathrm{cos}\theta \mathrm{cos}\theta _c)\stackrel{}{e}_R`$ (12)
$`(1i\eta {\displaystyle \frac{\mathrm{cos}\theta _c}{sin^2\theta }}{\displaystyle \frac{\mathrm{cos}\theta \mathrm{cos}\theta _c}{1i\eta }})\mathrm{sin}\theta \stackrel{}{e}_\theta ],`$
$$\stackrel{}{B}_\omega =\frac{i\omega }{vcR\mathrm{cos}\theta _c}F(\stackrel{}{q})I^{FF}(\eta ,\theta )(1+i\eta \frac{\mathrm{cos}\theta }{sin^2\theta }\frac{\mathrm{cos}\theta \mathrm{cos}\theta _c}{1i\eta })\mathrm{sin}\theta \stackrel{}{e}_\varphi ,$$
(13)
where
$`I^{FF}(\eta ,\theta )=`$ $`e^{ikR}a\sqrt{2\pi }[1i\eta (\theta )(13i\eta {\displaystyle \frac{\mathrm{cos}\theta }{\mathrm{sin}^2\theta }}{\displaystyle \frac{\mathrm{cos}\theta \mathrm{cos}\theta _c}{1i\eta }})]^{1/2}`$ (14)
$`\mathrm{exp}[{\displaystyle \frac{1}{2}}(ka)^2{\displaystyle \frac{(\mathrm{cos}\theta \mathrm{cos}\theta _c)^2}{1i\eta }}],`$
and where
$$\eta =\frac{ka^2}{R}\mathrm{sin}^2\theta .$$
Inspection reveals that these formulas have the dependence on $`\omega `$, $`R`$, and a distance scale $`a`$ quoted earlier on physical grounds. The formulas cited earlier as summarizing the numerical work are, in fact, the saddle-point approximations evaluated at $`\eta =0`$, which fit closer than any empirical formula.
### 4.1 Remarks
The dependence of the fields on symbol $`\eta `$ summarizes a good deal of complexity. For example:
* The limit $`\eta 0`$ yields the Fraunhofer limit, with spherical wave fronts and $`E_\omega a\omega /R.`$
* The limit $`\eta \mathrm{}`$ gives the cylindrically symmetric $`E_\omega \sqrt{\omega /R}`$ fields. This field can be substantially different from the Fraunhofer approximation. In fact one must take this limit to get the Frank-Tamm formula.
A notable application is emission from ultra-high energy air showers. In such showers the $`LPM`$ effect plays a definite role in suppressing the soft emissions from the hardest charges. There is no corresponding suppression of the evolution of the low-energy regions of showers where most particles exist, however . The major effect that we find is that the showers become long kinematically: that is, the parameter $`a`$ gets big if one is working with, say, a primary of $`10^{20}eV`$. The emitted fields approach those of the Frank-Tamm formula as $`a\mathrm{}`$, via the formula cited. The effect is important numerically: a $`10^{20}eV`$ air shower does not approach the Fraunhofer limit nearer than $`300km`$. The conditions of RICE are more amenable to the limit, and at $`R1km`$, $`\omega 1GHz`$ the Fraunhofer-based estimates near the Cherenkov cone in the literature are good to $`20\%`$.
* The frequency dependence in the Fraunhofer approximation is strongly affected by “diffraction”. Independent of the form factor effect, the Fraunhofer approximation imposes an upper limit to the frequency of order $`\omega <(1/a)(\mathrm{cos}\theta \mathrm{cos}\theta _c)`$. The true behavior is substantialli different: from Eq. (14) we see that the field exists in a region
$$\omega <(1/a)[(\mathrm{cos}\theta \mathrm{cos}\theta _c)^2(a^2/R^2)\mathrm{sin}^4(\theta )]^{1/2}.$$
For large $`\omega `$ and in the angular region where the signal exists, the behaviour is much flatter. This is illustrated in Figures 4 and 5. This effect is invisible at the exact value $`\theta =\theta _c`$, where the fake Fraunhofer frequency cutoff and most of the true functional dependence in the exponent of $`I^{FF}`$ both drop out. As a result of the difference in frequency dependence, the time-structure of the electric field may be substantially different from the Fraunhofer approximation. We will return to this point in a Section Causal Features .
* Fields in the forward and backward directions, $`\mathrm{sin}^2\theta 0`$, are the fields of $`\eta 0`$, the Fraunhofer approximation, regardless of the physical values of $`k,a,R`$. We note that the experiments of Takahashi observe an extremely limited region of $`\mathrm{sin}^2\theta 0`$. Perhaps this contributes to the observed agreement with Tamm’s formula in a regime where $`a^2k/R`$ is not close to zero.
* The polarization varies considerably. From symmetry the polarization is in the plane of the charge and the observation point. Moreover, for $`\theta =\theta _c`$ the electric field is transverse to $`\stackrel{}{R}`$ for any $`\eta `$. Yet naive transversality is not true in general at any finite $`\eta `$.
* Between the various limits the dependence on every scale in the problem, namely the frequency $`\omega `$, the distance $`R`$, the length scale $`a`$, and the angle $`\theta `$, is neither that of the Fraunhofer limit nor that of the infinite track Frank-Tamm limit, but instead a smooth interpolation between the cylindrical and spherical wave regimes.
* In the finite $`\eta `$ limit, one may also include a further effect, namely that as $`\omega 0`$ one has a ‘near-zone’ Coulomb-like response at small $`R`$. (Indeed, the $`\omega 0`$ limit measures the net charge.) This effect, important below about $`10MHz`$, also has a slight effect on the time structure of pulses.
We pause to comment on the generality of the result. What if we had not made the physical, but specific ansatz (3)? The entire analysis can be repeated for an arbitrary charge distribution $`j(t^{},z^{},\stackrel{}{\rho }^{})`$. The Fraunhofer expansion of the transverse variable, and the Fourier integral of the $`t^{}`$ variable are general. Provided the $`\stackrel{}{\rho }^{}`$ extent is finite, and there exists a dominant $`z^{}`$ region, then the integrals always factor into a product of a form factor and a one-dimensional integral for $`A_\omega ^{FF}(\eta )`$. In fact nothing changes (the reader can repeat the calculation) except that when an arbitrary current is set up, the existence of a single saddle point cannot be assumed. Corrections to the local Gaussian approximation are straightforward. The slight skewness of real showers (or other arbitrary charge distributions) can also be developed as a saddle-point power series. Again: there are elements of bremmstrahlung in real showers, having a stochastic nature, which the current model has not attempted to reproduce. Detailed Monte Carlo simulations of our group have included the bend-by-bend amplitudes of tracks undergoing collisions. This goes well beyond the approximation of a single straight line track, suddenly beginning and ending, of the previous literature . The effect of all the small kinks is negligible except in the very high frequency region $`\omega 100GHz`$, while the endpoint accelerations give oscillations in the angular dependence down by orders of magnitude. As a final side remark: we explicitly studied contributions of finite tracks, just to see what would happen, in the development towards the conditions of the Tamm formula. It is straightforward to develop these pieces if one needs them for, say, the Takahashi-type experiments , in the Fresnel zone.
### 4.2 The General Case
We now turn to fields valid for any $`\eta `$. The form factor, which was extracted from the Fraunhofer calculations, is universal and need not be changed. The validity of the saddle-point approximation does not depend explicitly on the value of $`\eta `$. The procedure of linearization to locate the saddle-point happens to be good to $`cos\theta cos\theta _c1`$, so that the approximation is rather good in the entire region $`R/a1`$, $`kR1`$.
For practical applications it is useful to have a formula for the fields with quantities measured in physically motivated units. For this purpose we rewrite (12) as follows
$$R\stackrel{}{E}_\omega 2.5210^7\frac{a}{m}\frac{n_{max}}{1000}\frac{\nu }{GHz}F(\stackrel{}{q})\psi \stackrel{}{}[\frac{V}{MHz}].$$
(15)
Here $`n_{max}`$ is the excess of electrons over positrons at the shower maximum, $`R`$ is measured in meters, $`\nu `$ is measured in $`GHz`$. The rescaled field is
$`\stackrel{}{}`$ $`=`$ $`[{\displaystyle \frac{\mathrm{cos}\theta \mathrm{cos}\theta _c}{\mathrm{sin}\theta }}\stackrel{}{e}_R+(1i\eta {\displaystyle \frac{\mathrm{cos}\theta _c}{sin^2\theta }}{\displaystyle \frac{\mathrm{cos}\theta \mathrm{cos}\theta _c}{1i\eta }})\stackrel{}{e}_\theta ]`$ (16)
$`[1i\eta (13i\eta {\displaystyle \frac{\mathrm{cos}\theta }{\mathrm{sin}^2\theta }}{\displaystyle \frac{\mathrm{cos}\theta \mathrm{cos}\theta _c}{1i\eta }})]^{1/2}`$
$`\mathrm{exp}[{\displaystyle \frac{1}{2}}(ka)^2{\displaystyle \frac{(\mathrm{cos}\theta \mathrm{cos}\theta _c)^2}{1i\eta }}].`$
We have defined a kinematic factor $`\psi =i\mathrm{exp}(ikR)\mathrm{sin}\theta `$ in such a way that the rescaled field $`\stackrel{}{}`$ is normalized at $`\theta =\theta _c`$,
$$\stackrel{}{}(\theta =\theta _c)=(1i\eta )^{1/2}\stackrel{}{e}_\theta .$$
(17)
It is convenient to plot the magnitude of the rescaled field, Eq. (16). Figure 3 shows the magnitude as a function of the angle difference $`\theta \theta _c`$ in various limits. The Fraunhofer approximation is shown by a dashed curve, and our result by the solid curve. One observes that the Fraunhofer limit is approached from below. This is physically clear: The Fresnel zone fields have a wider angular spread, and conservation of energy forces them to be smaller in magnitude compared to the sharper, diffraction-limited Fraunhofer fields. As the fields evolve to infinity, they coalese into narrower and taller beams.
The frequency dependence of $`R|\stackrel{}{E}_\omega |/F(\omega )`$ is shown on Figures 4 and 5. Exactly at the Cherenkov angle the difference between the Fraunhofer approximation and our results are minor for the typical parameters of $`RICE`$. However, away from the Cherenkov angle there is a substantial difference between the two, throughout the region where the magnitude of the field is large. This effect can be masked by the form factor, so we have plotted $`RE_\omega /F(\omega )`$ to show it. This effect may have important repurcussions for the time-structure of pulses, which are also discussed in the last Section.
Figures 6-10 are contour plots of the electric field. We did not bother to remove the small region $`a/R1`$, where our result does not apply. The Fraunhofer approximation has trivial $`1/R`$ dependence on the distance to the observation point (Figure 10). The exact result is certainly different, with Figures 6, 8, and 9, in particular, illlustrating the effects of constructive interference in the region of cylindrical symmetry. A complementary view examines contour plots of constant phase. This is shown in Figures 8 and 9. In making these figures we decreased the kinematic phase $`ikR`$ to values showing several oscillations (as opposed to hundreds) across the range of the plots. The lines of constant phase can be used to illustrate the time-evolution of waves of a given frequency: that is, the Fourier transform of $`\delta (\omega \omega _{})E_\omega exp(i\omega t)`$ has wave fronts at each moment in time given by the lines of constant phase. The constant phase lines of the Fraunhofer approximation are, of course, spherical (Figure 10). The constant phase lines of the true behaviour interpolate between cylindrical and spherical symmetry. As a consequence of the Fresnel-zone behaviour, the eikonals of the expanding radiation field do not emerge radially, but actually curve due to interference effects. This is a sobering impact of very basic physics, which has a measurable effect in the signal propagation speed discussed later under the topic of causal feature.
## 5 Causal Features
With our convention that the electric field $`E(t,\stackrel{}{x})=1/(2\pi )_{\mathrm{}}^{\mathrm{}}𝑑\omega \mathrm{exp}(i\omega t)E_\omega (\stackrel{}{x})`$, causality requires $`E_\omega `$ to be analytic in the upper half-plane of $`\omega `$. Singularities in the lower half-plane determine the details of $`E(t,\stackrel{}{x})`$ and the precise causal structure.
To discuss this we consider detection of signals via an antenna-system response function $`𝒜_\omega `$. By standard arguments the detected voltage is a convolution in time, and therefore a product in $`\omega `$ space, of the antenna function and the perturbing electric field. The antenna function has the same causal analytic properties as the electric field. A proto-typical antenna or circuit function for a driven $`LRC`$ circuit is
$$𝒜_\omega =Z/(\omega ^2+\omega _𝒜^2i\mathrm{\Gamma }\omega ),$$
where $`Z`$ depends on where the amplifier is connected in the circuit and can be treated as a constant. Note that $`𝒜_\omega =𝒵/(\omega \omega _+)(\omega \omega _{})`$ where $`\omega _\pm `$ are in the lower half-plane. The dielectric function $`ϵ(\omega )`$ can then be taken as slowly varying in the region where the antenna and form factor allow a response, and also has its analytic structure in the lower half-plane if this detail needs to be included. We will also ignore the form factor for this discussion, which earlier was cited as a formula analytic in the complex plane. While nothing in our analysis depends on these idealizations, this approach to the analytic structure serves to make our point.
As a first illustration, consider the electric field fit given by $`ZHS`$, proportional to $`1/(\omega i\omega _0)(\omega +i\omega _0)`$ with $`\omega _0/(2\pi )=500MHz`$. This field has poles in both the upper and lower half-plane, violating causality. One may argue that the literal analytic behavior in the complex plane goes beyond the ambitions of the original semi-empirical fit. Nevertheless, Cauchy’s theorem applies to the subsequent numerical integrations that have been made , giving non-causal branches to numerically evaluated Fourier transforms, as well as unphysical short-time structure.
Let us compare the analytic structure of the electric field in the saddle-point approximation. This approximation does not attempt to describe the region $`\omega 0`$, which requires treatment of the near zone. However for causality we do not need $`E_\omega `$ near the origin but at large $`|\omega |`$. The saddle point approximation is good here so the results should be reliable.
Let us investigate this in more detail. In the exponent in the expression for the electric field we have $`(1/2)(ka)^2(\mathrm{cos}\theta \mathrm{cos}\theta _c)^2(1+i\eta )/(1+\eta ^2)`$. Since $`\eta `$ goes like $`\omega `$ there is a phase linear in $`\omega `$ at large $`|\omega |`$. There is also a branch-cut and pole from the prefactor which occurs at $`\eta =i`$. All singularities are in the lower half-plane and consistent with causality. The causal structure of $`E(t,\stackrel{}{x})`$ then hinges on closing the contour at infinity. For this the details of the antenna function, which generally has isolated singularities, as well as singular behavior of $`E_\omega `$ near the origin do not matter.
We close the contour at infinity avoiding the branch cuts, which may be oriented along the negative imaginary axis or as consistent. Convergence then requires
$$\underset{\omega i\mathrm{}}{lim}\mathrm{Re}[i\omega t+i\omega \sqrt{ϵ}R\frac{1}{2}\omega ^2ϵa^2(\mathrm{cos}\theta \mathrm{cos}\theta _c)^2\frac{i\eta }{1+\eta ^2}]<0.$$
Using the definition of $`\eta `$, causality implies
$$t\sqrt{ϵ}\frac{R}{c}[1\frac{(\mathrm{cos}\theta \mathrm{cos}\theta _c)^2}{2\mathrm{sin}^2\theta }]^1>0.$$
This result has a natural interpretation. While the radiation from the shower appears to come primarily from the geometric location of the maximum, $`\theta \theta _c`$, the shower actually develops and radiates earlier. Consequently the strict causal limit must correspond to an apparent propagation speed slightly faster than the naive speed of $`c/\sqrt{ϵ}`$ deduced from the location of the observation point at $`R`$. The earliest signal actually arrives at an apparent speed $`v_{app}`$ of
$$v_{app}=\frac{c}{\sqrt{ϵ}}[1\frac{(cos\theta cos\theta _c)^2}{2\mathrm{sin}^2\theta }]^1.$$
This formula is entirely geometrical, consistent with the simple picture that distance differences in the problem are causing the effect, but it also incorporates subtle features of coherence. For example, the distance scale $`a`$ cancels out. Yet $`a`$ determines the angular spread $`\mathrm{cos}\theta \mathrm{cos}\theta _c`$ over which most of the power in the wave is contained, and enters in this fashion.
The Fraunhofer limit, in which all signals originate at a single point of origin, is incapable of capturing such an effect. It is interesting to trace the origin of the discrepancy. The singularities of interest are located by the zeroes of $`1i\eta =1i\omega \sqrt{ϵ}\mathrm{sin}^2\theta /R.`$ When the limit $`R\mathrm{}`$ is taken in the first step of the Fraunhofer approximation, all the non-trivial analytic structure moves away to $`\omega i\mathrm{}`$ and is lost. This procedure does not commute with closing the contour at $`|\omega |\mathrm{}`$. The correct procedure, of course, is to first close the contour, and then take the limit of large $`R`$.
In practice, of course, not all of the signal arrives at the earliest possible moment. The time scale over which the signal is detected depends on competition between dispersion, the antenna and form factor details, and something like twice the ”advanced” time interval. This time interval is $`\mathrm{\Delta }t_{caus}=R\sqrt{ϵ}/(2c)(\mathrm{cos}\theta \mathrm{cos}\theta _c)^2/\mathrm{sin}^2\theta `$. Since the $`\mathrm{\Delta }t_{caus}`$ effect scales proportional to the distance $`R`$, it does not become negligible in any limit, exhibiting another subtle facet of breakdown of the Fraunhofer approximation.
Acknowledgments: This work was supported in part by the Department of Energy, the University of Kansas General Research Fund, the
K\*STAR programs and the Kansas Institute for Theoretical and Computational Science. We thank Jaime Alvarez-Muñiz, Enrique Zas, Doug McKay, Soeb Razzaque and Suruj Seunarine for many helpful suggestions and conversations. We especially thank Enrique for writing and sharing the ZHS code, and Soeb and Suruj for generously sharing their results from running shower codes. |
warning/0003/cond-mat0003267.html | ar5iv | text | # Stripe glasses: self generated randomness in a uniformly frustrated system
## Abstract
We show that a system with competing interactions on different length scales, as relevant for the formation of stripes in doped Mott insulators, undergoes a self-generated glass transition which is caused by the frustrated nature of the interactions and not related to the presence of quenched disorder. An exponentially large number of metastable configurations is found, leading to a slow, landscape dominated long time relaxation and a break up of the system into a disordered inhomogeneous state.
Competing interactions on different length scales are able to stabilize mesoscale phase separations and the creation of spatial inhomogeneities in a wide variety of systems. Examples are stripe formation in doped Mott insulators, as found in transition metal oxides (TMO), domains in magnetic multilayer compounds, or mesoscopic structures formed by assembling polymers in solution and amphiphiles in water-oil mixtures . In many of these cases the tendency towards a perfectly ordered array of domains, stripes etc. is undermined by frustrating long range interactions. Very often, these assemblies exhibit a long time dynamics similar to the relaxation seen in glasses. In the context of stripes it has been argued that the presence of only very few quenched impurities might already cause a strictly disordered glassy state . Furthermore, recent molecular dynamics calculations for charge ordering in TMO found an anomalous long time relaxation with a power spectrum similar to $`1/f`$-noise. Indeed, there is experimental evidence for the formation of intrinsic inhomogeneities and even a stripe-glass in high temperature superconductors and other transition metal oxides. In particular slow, activated dynamics as observed in NMR experiments exhibits a striking universality, rather independent of the details of added impurities etc. It is therefore tempting to speculate that glassiness in these systems is self generated and does not rely on the presence of quenched disorder, which may of course further stabilize a glassy state.
In this paper we show that the competition of interactions on different length scales in a uniformly frustrated systems exhibits a self generated glass transition due to the emergence of an exponentially large number of metastable states. This result is obtained using the replica approach of Ref. and by solving the corresponding many body problem using the self consistent screening approximation. Since only very few examples exists for models which exhibit self generated glassiness, all these approaches are extremely important for a better understanding of glassiness in general. Even though our findings apply to broader class of problems than stripes in TMO, we will adopt a language which is specific to that problem.
A model for a uniformly frustrated system with competition on different length scales is given by the Hamiltonian:
$``$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle d^3x\left\{r_0\phi (𝐱)^2+\left(\phi (𝐱)\right)^2+\frac{u}{2}\phi (𝐱)^4\right\}}`$ (2)
$`+{\displaystyle \frac{Q}{2}}{\displaystyle d^3xd^3x^{}\frac{\phi (𝐱)\phi (𝐱^{})}{\left|𝐱𝐱^{}\right|}}.`$
Here, $`\phi (𝐱)`$ characterizes charge degrees of freedom, with $`\phi (𝐱)>0`$ in a hole-rich region, $`\phi (𝐱)<0`$ in a hole poor region, and $`\phi (𝐱)=0`$ if the local density equals the averaged one. If $`r_0<0`$ the system tends to phase separate since we have to guarantee charge neutrality, $`\phi =0`$. The coupling constant, $`Q`$, is a measure for the strength of the Coulomb interaction and characterizes the competition between short and long range interactions. In case of strongly anisotropic, quasi two-dimensional cuprate superconductors one expects an anisotropy of the gradient term in Eq.2, which we neglect for simplicity. Despite the absence of a clean derivation of Eq.2 from the many electron Schrödinger equation, we note that it describes, on a phenomenological level, many of the major competing effects which yield in microscopic theories a rich phase diagram of inhomogeneous spin and charge structures. For $`Q=0`$ and $`r_0<0`$ we expect at low temperatures long range ordered charge modulations. As shown in Ref. , the Coulomb interaction suppresses this ordered state for all $`Q>0`$ and finite $`T`$. Instead, the system undergoes several crossovers. Most interestingly, at low temperatures, where $`\left|r(T)\right|<2\sqrt{Q}`$, a mean field analysis of Eq.2 shows that besides a correlation length, $`\xi =2\left(r+2\sqrt{Q}\right)^{1/2}`$, an additional length scale, $`l_m=4\pi \left(2\sqrt{Q}r\right)^{1/2}`$ emerges, which characterizes the spatial modulation of the field correlations, where $`r=r_0+uT\phi ^2`$. These modulations are particularly relevant for low enough $`T`$ where $`r(T)0`$, where mean field theory gives locally ordered regions with characteristic size $`l_m\xi `$.We will show that a stripe glass emerges in this temperature regime.
An essential prerequisite for the anomalous dynamical features of glassiness, like aging, memory effects and ergodicity breaking, is most certainly the occurrence of a large number of metastable states, $`𝒩_{ms}`$, separated by energy barriers which are large compared to the temperature. In viscous liquids undergoing vitrification calorimetry suggests $`𝒩_{ms}\mathrm{exp}(const.V)`$, were, $`V`$ is the system size. This observation is the heart of an ideal glass transition scenario based on random first order transitions which was originally motivated by microscopic stability analyses of structural glasses and mean field theories for random Potts-glasses. Below a cross-over temperature, $`T_A`$, a ’viscous’, energy-landscape dominated long time relaxation sets in due to the occurrence of exponentially many metastable states, i.e. the configurational entropy, $`S_c=k_B\mathrm{log}𝒩_{ms}`$, becomes extensive. Because of the large barriers between these states, the system will get stuck for extremely long times in one of the metastable states, i.e. it will freeze into a glass, at some temperature $`T_G<T_A`$ which depends for example on the cooling rate. Even though this laboratory glass transition is purely dynamical, a key ingredient of the ideal glass transition scenario is that the dynamical slowing arises from proximity to an underlying phase transition at $`T_K<T_G`$, where the configurational entropy would vanish like $`S_c(T)TT_K`$. If such an ideal transition exists, even for an infinitely slow cooling rate freezing will occur at $`T_K`$ since all the liquid degrees of freedom die out due to this ’entropy crisis’.
Detailed theoretical investigation of this scenario have concentrated on systems with quenched randomness. A major step forward for studying nonrandom systems was made in Ref. , where a new replica approach was developed. Within this approach, the configurational entropy for a model of a structural glass without quenched disorder was calculated and found to be in good agreement with computer simulations.
We will use this approach to determine $`S_c`$ for a system governed by Eq.2. The key idea is to introduce, in analogy to the theory of conventional phase transitions, an appropriate symmetry breaking field, $`\psi \left(𝐫\right)`$, and to compute the partition sum
$$\stackrel{~}{Z}\left[\psi \right]=D\phi e^{[\phi ]/T\frac{g}{2}{\scriptscriptstyle d^3x\left[\psi \left(𝐱\right)\phi \left(𝐱\right)\right]^2}},$$
(3)
where $`g0^+`$. The energy $`\stackrel{~}{f}\left[\psi \right]=T\mathrm{log}\stackrel{~}{Z}\left[\psi \right]`$ will be low if $`\psi (𝐫)`$ equals to configurations which locally minimize $``$. Sampling all configurations of the $`\psi `$-field, weighted with $`\mathrm{exp}\left(\stackrel{~}{f}\left[\psi \right]/T\right)`$, is therefore equivalent to scanning all metastable states such that
$$\stackrel{~}{F}=\underset{g0}{lim}\frac{1}{W}D\psi \text{ }\stackrel{~}{f}\left[\psi \right]\mathrm{exp}\left(\stackrel{~}{f}\left[\psi \right]/T\right)$$
(4)
is a weighted average of the free energy in the various metastable configurations, where $`W=`$ $`D\psi \mathrm{exp}\left(g/2d^3x\psi ^2\left(𝐱\right)\right)`$ is introduced for proper normalization. If there are only few local minima, the limit $`g0^+`$ behaves perturbatively and $`\stackrel{~}{F}`$ equals to the free energy, $`F`$, of the system. However, in case exponentially many local minima with large barriers between them exist, a nontrivial contributions arises from the $`\psi `$-integral even for $`g0^+`$ and the averaged free energy, $`\stackrel{~}{F}`$, differs from $`F`$. This enables us to identify the configurational entropy, $`S_c`$, via $`F=\stackrel{~}{F}TS_c`$. For an illustration of the corresponding free energy landscape, see inset of Fig.1.
An explicit expression for $`S_c`$ can be obtained within a replicated theory with
$$F\left(m\right)=\underset{g0}{lim}\frac{T}{m}\mathrm{log}\frac{1}{W}D\psi \stackrel{~}{Z}^m\left[\psi \right].$$
(5)
It follows that $`\stackrel{~}{F}=\frac{mF(m)}{m}|_{m=1}`$, which gives:
$$S_c=\frac{1}{T}\frac{F(m)}{m}|_{m=1}.$$
(6)
Inserting $`\stackrel{~}{Z}\left[\psi \right]`$ of Eq. 3 into Eq. 5 finally leads to:
$`Z(m)`$ $`=`$ $`\underset{g0}{lim}{\displaystyle }D^m\phi \mathrm{exp}({\displaystyle \underset{a=1}{\overset{m}{}}}\left[\phi ^a\right]/T`$ (8)
$`\text{ }{\displaystyle \frac{g}{2m}}{\displaystyle \underset{a,b=1}{\overset{m}{}}}{\displaystyle }d^3x\phi ^a(𝐱)\phi ^b(𝐱)),`$
with $`F\left(m\right)=\frac{T}{m}\mathrm{log}Z(m)`$. Eq.8 has a formal similarity to the action of the random field Ising model, obtained within the conventional replica approach, which allows us to use techniques, developed for this model. In the following we use the self consistent screening approximation (SCSA) of Eq.8 and determine the Green’s function, $`𝒢_{ab}\left(𝐪\right)=\phi ^a(𝐪)\phi ^b(𝐪)`$, in replica space. $`𝒢_{ab}\left(𝐪\right)`$ then determines the partition function, $`Z(m)`$, and correspondingly $`S_c`$.
The interaction between different replicas is symmetric with respect to the replica index suggesting the mean field ansatz
$$𝒢_{ab}\left(𝐪\right)=\left(𝒢\left(𝐪\right)\left(𝐪\right)\right)\delta _{ab}+\left(𝐪\right),$$
(9)
with equal diagonal elements, $`𝒢\left(𝐪\right)`$, and equal off-diagonal elements, $`\left(𝐪\right)`$.The physical interpretation of $`𝒢\left(𝐫𝐫^{}\right)=\phi (𝐫)\phi (𝐫^{})`$ as thermodynamic (instantaneous) correlation function is straightforward. On the other hand, $`\left(𝐫𝐫^{}\right)=lim_t\mathrm{}\phi (𝐫,t)\phi (𝐫^{},0)`$ can be interpreted as measuring long time correlations, arising from trapping in metastable minima which, in mean field theory, have infinite barriers between them. An analogous structure in replica space follows for the diagonal elements, $`\mathrm{\Sigma }_𝒢\left(𝐪\right)`$, and off-diagonal elements, $`\mathrm{\Sigma }_{}\left(𝐪\right)`$, of the self energy, which are given in the SCSA as:
$$\mathrm{\Sigma }_𝒜\left(𝐪\right)=2\frac{d^3p}{\left(2\pi \right)^3}𝒟_𝒜\left(𝐩\right)𝒜\left(𝐩+𝐪\right)$$
(10)
with $`𝒜\{𝒢,\}`$. The screening of the interaction is characterized by the collective propagators $`𝒟_𝒢^1\left(𝐩\right)=(uT)^1+\mathrm{\Pi }_𝒢\left(𝐩\right)`$ and $`𝒟_{}\left(𝐩\right)=\frac{uT\mathrm{\Pi }_{}\left(𝐩\right)𝒟_𝒢^2\left(𝐩\right)}{1uT𝒟_𝒢\left(𝐩\right)\mathrm{\Pi }_{}\left(𝐩\right)}`$ with polarization functions $`\mathrm{\Pi }_𝒜\left(𝐩\right)=\frac{d^3q}{\left(2\pi \right)^3}𝒜\left(𝐪+𝐩\right)𝒜\left(𝐪\right)`$. The set of equations is closed by the Dyson equation: $`𝒢^1\left(𝐤\right)=𝒢_0^1\left(𝐤\right)+\mathrm{\Sigma }_𝒢\left(𝐤\right)`$ for the diagonal elements, and $`\left(𝐤\right)=\frac{𝒢^2\left(𝐤\right)\mathrm{\Sigma }_F\left(𝐤\right)}{1𝒢\left(𝐤\right)\mathrm{\Sigma }_F\left(𝐤\right)}`$ for the off diagonal elements, respectively. $`𝒢_0^1\left(𝐪\right)=r+q^2+Qq^2`$ is the inverse Hartree propagator. Note, all momentum integrations have to be cut-off at $`\left|𝐩\right|=\mathrm{\Lambda }`$, which is of the order of an inverse lattice constant. Once the $`𝒢_{ab}`$ and $`𝒟_{ab}`$ are determined the free energy becomes
$$F(m)/(2mT)=\text{tr}\mathrm{log}𝒢^1+\text{tr}\mathrm{log}𝒟^1\text{tr}\mathrm{\Sigma }𝒢.$$
(11)
After performing the trace in replica space for arbitrary integer $`m`$ and analytical continuation to $`m1`$, $`S_c`$ follows from Eq.6. One finds immeadiately $`S_c=0`$ if $`(𝐤)`$ vanishes. In the following we discuss the numerical solution of this set of coupled integral equations.
In Fig. 1, $`S_c/V`$ is shown for two different $`Q`$-values as function of $`T`$. At $`T_A`$ the long time correlation function $`\left(k\right)`$ emerges leading to $`S_c>0`$ and a glassy dynamics sets in. The corresponding free energy landscape is schematically illustrated in the inset. Just as in mean field Potts glasses with quenched randomness, $`S_c`$ vanishes at a lower temperature, $`T_K`$. At $`T_K`$ the entropy of the amorphous stripe solid equals that of the stripe liquid. There is no entropic advantage anymore to be in a liquid state, leading to an obligatory glass transition no matter how slow the cooling rate. The laboratory glass temperature, $`T_G`$, will lie somewhere between $`T_K`$ and $`T_A`$ and cannot be determined within our theory. We also find that $`T_K`$ and $`T_A`$ are only weakly decreasing for increasing $`Q`$, see inset of Fig.2. Both temperatures remain finite for $`Q0`$. However, $`S_c\left(Q0\right)0`$, i.e. the fragility $`\frac{dS_c}{dT}`$ of the glass vanishes. In other words, the larger the modulation length, the smaller is the number of metastable states. Due to the $`1/r`$ interaction, the limit $`Q0`$ does not smoothly connect to the behavior at $`Q=0`$. If one includes a finite screening length $`l_s`$ in Eq.2, we expect that the glassy state disappears for $`Ql_s^4`$.
In Fig. 2, the instantaneous ($`𝒢\left(k\right)`$) and long-time ($`\left(k\right)`$) charge correlation functions, are shown at $`T=T_A`$. Even though no charge ordering occurs, the pronounced peaks at finite $`k`$ demonstrate that there is a modulated state with strong short range correlations, $`l_m<\xi `$. Also, since $`\left(k\right)𝒢\left(k\right)`$, the modulated state exhibits an anomalous dynamics, where long time correlations are only slightly reduced compared to instantaneous correlations. Introducing a Debye-Waller factor, $`W=\mathrm{log}\left(/𝒢\right)`$ for $`k`$ close to the peak maximum, gives $`W=0.12`$ $`(0.13)`$ for $`Q=0.01`$ $`(0.001)`$. The modulation length (inverse peak positions) is $`3.5\mathrm{\Lambda }^1`$ ($`6\mathrm{\Lambda }^1`$) and the correlation length is $`45\mathrm{\Lambda }^1`$ ($`80\mathrm{\Lambda }^1`$) for $`Q=0.01`$ $`(0.001)`$.
Due to the competing interactions in Eq.2, an entropy crisis occurs, causing a transition into a glass. This purely thermodynamic characterization of the spectrum of metastable states, is only the first important step for understanding glassiness, and the investigation of dynamical features is an even bigger challenge, because it requires going beyond mean field theory. An argument based on ”entropic droplets” explains quite well the phenomenology of viscous liquid dynamics and can even be made semi-quantitative. Here we apply these arguments to the present stripe model. The entropic droplet argument recognizes an intrinsic instability of the homogeneous metastable solutions, as characterized by Eq. 9: namely, creating a droplet of one metastable solution within another costs free energy that can at most scale as a surface energy but the exponential number of configurations gives an entropic driving force for such a droplet that scales with V. A mosaic structure hence will form. The activation free energy of turning over a single region can be computed where the entropic gain is given by $`TS_c`$. A renormalization group calculation, based on Ref., leads to a size dependent surface tension $`\sigma (R)=\sigma _0\left(R\mathrm{\Lambda }\right)^\theta `$ with $`\theta =\frac{d2}{2}`$ reflecting the fact that the interface between two states is wetted by intermediate states. This analysis leads to an characteristic energy barrier $`\mathrm{\Delta }E\left(TS_c(T)\right)^1`$ which implies a relaxation time obeying a Vogel-Fulcher law
$$\tau \mathrm{exp}\left(\frac{DT_K}{TT_K}\right).$$
(12)
An estimate for the surface tension $`\sigma _0|r_0|/(u\xi )`$ for the stripe model yields $`R_0^3\mathrm{\Lambda }^1(V\sigma _0/TS_c)^2`$ for the droplet volume and $`D=3V\sigma _0^2/(\mathrm{\Lambda }T^2T_KS_c/T|_{T_K})`$. Using our numerical results for $`S_c`$, this leads to $`D60200`$, typical for moderately fragile and strong glasses, and droplet sizes $`R_0(2550)\mathrm{\Lambda }^1(510)l_m`$. Note that this estimate is only qualitative since $`R_0\xi `$ and a real separation of scales never occurs. The droplet picture implies that the glass state breaks up into domains of different metastable states, separated by wetted surfaces, build by intermediate states. This physical picture is very similar to the conclusions made in Ref. based on NMR experiments.
In summary, we have shown that an exponentially large number of metastable configurations emerges in a system with competing interactions on different length scales, leading to a glass transition and anomalous long time dynamics. This glass state is self generated, implying that the barriers characterizing the activated dynamics are rather universal and should not depend on details like added impurities but only on the generic interactions on short and long scales, i.e. the magnetic exchange interactions and the Coulomb interaction. Furthermore, we showed that the magnitude of the frustration controls the fragility of the glass transition. Finally, following Ref., we argued that the configurational entropy causes a break up of the stripe glass into a mosaic of domains or droplets, build up by the various metastable states, allowing us to estimate time scales of motions. This causes an intrinsic inhomogeneity of all relevant correlation functions, modulation and correlation lengths etc. in the amorphous glassy state.
We gratefully acknowledge stimulating discussions at the workshop on Mesoscopic organization in soft hard and biological matter hosted by the Institute for Complex Adaptive Matter and with A. V. Chubukov, P. C. Hammel, J. Haase, D. C. Johnston, D. K. Morr, D. Pines, C. P. Slichter, R. Stern and B. P. Stojkovic. The work was supported by NSF (PGW), grant No.ChE-9530680. Ames Laboratory is operated for the U.S. Department of Energy by Iowa State University under Contract No. W-7405-Eng-82. |
warning/0003/astro-ph0003076.html | ar5iv | text | # VLA Polarimetry of Two Extended Radio Galaxies
## 1 Introduction
As part of our investigation of steep-spectrum ($`\alpha >0.5`$, S$`\nu ^\alpha `$), low-frequency-variable (LFV; $`\nu <`$ 1 GHz) sources, we have made a series of images with sub-arcsecond resolutions (Mantovani et al. 1992) of a sample of sources. These sources were selected from the papers of Cotton (1976), McAdam (1980), Spangler & Cotton (1981), Fanti et al. (1983) and Altschuler et al. (1984). The aim was to detect the high-brightness components required by the refractive scintillation model for low-frequency variability (Rickett 1986). Most of the sources in the sample showed compact features (deconvolved sizes $`<`$ 0.15 – 0.3 arcsec) both in MERLIN, 408 MHz and VLA, A-array, 5 GHz images (Mantovani et al. 1992). Further observations with VLBI show these features to contain components which are bright and compact enough to explain the variability at low frequency by propagation effects in the interstellar medium; see, for example, 3C99, Mantovani et al. (1990a). Sources such as these do not generally show any variability at high frequency (Padrielli et al. 1987).
However, there are sources like 0621+400 (3C159), which are variable at low frequencies and not at frequencies $`>2`$ GHz, where the compact components are too weak to account for the observed variability. In these cases, it is possible that the variability is caused by instrumental effects. The source 3C159 has been monitored for about 10 years at 408 MHz. This source has a steep radio spectrum and an extended double radio structure — a combination which is very unusual for a variable radio source. Browne et al. (1985) have suggested that the variations, which seem to show an annual cycle, may not be intrinsic but could arise from the combined effects of strong source linear polarization and ionospheric Faraday rotation.
Ionospheric Faraday rotation can easily reach 7–8 rad m<sup>-2</sup> (Sakurai & Spangler 1994). With the plausible estimate of 5 rad m<sup>-2</sup> as an expectable difference in the ionospheric RM, one finds that the position angle difference at 408 MHz is 2.7 radians; enough to produce the effect being discussed.
MERLIN observations at 408 MHz (Cerchiara et al. 1994) show that 3C159 is highly polarized ($``$10$`\%`$). The plane of polarization of the source emission could be rotated by changes in ionospheric Faraday rotation relative to the linearly polarized E-W arm of the Northern Cross Bologna telescope used for the monitoring program. A source with a linear polarization $`>`$6$`\%`$ could exhibit apparent variations of roughly the observed size if the ionospheric Faraday rotation changed by $``$90<sup>o</sup> between observations. The 3C159 observations were made at transit during the day in summer and the night in winter and so they were accompanied by annual changes in the ionospheric electron content.
The purported variability measured in the two extended radio sources 0235$``$197 and 1203$`+`$043 may have originated in a similar manner to that in 3C159. They were monitored at 408 MHz with a similar instrument, the Molonglo Cross. With a peak-to-peak fractional variability of $``$10%, they were classified as ‘probably variable’ by McAdam (1980).
Consequently, although it was expected that most of the sources belonging to our sample of steep-spectrum, low-frequency-variable sources would be core-dominated sources, it is likely that the sample has been contaminated by lobe-dominated, strongly-linearly-polarized, extended sources.
In order to test if ionospheric Faraday rotation is the cause of the apparent variability of 0235$``$197 and 1203$`+`$043 we have investigated the linear polarizations of these sources at 320 MHz with the VLA in the ‘A’ configuration and with the already available 5 GHz VLA C-array data. Both sources were also observed in the X (8.4 GHz) and U (15 GHz) bands. These observations allowed high resolution images of the ’hot spot’ regions to be made. The images were combined with available, high-resolution, C band observations to produce rotation measures (RMs) for the outer parts of the sources.
## 2 VLA Observations
VLA (Thompson et al. 1980) observing dates and observational parameters are summarized in Table 1. The high resolution (A-array) Total Intensity images at 5 GHz for 0235$``$197 and 1203$`+`$043 were presented in Mantovani et al. (1992) while the lower resolution C-array, total intensity map for 0235$``$197 has been published in Morganti, Killeen & Tadhunter (1993). However, because the polarization information was never presented in the earlier papers, we have summarized in the tables all the observational details and the derived parameters from those observations.
Because of narrowband interference, the P (320 MHz) band data were taken in spectral line mode. The data were edited to remove channels with interference. Bandpass corrections were determined from the calibrator source 3C286 (1331+305) and applied to the spectral line database. A new “Channel 0” database was then constructed and the data were calibrated for total intensity and polarization in the standard fashion (see, for example, Perley, Schwab & Bridle 1989). Instrumental polarization calibration was done using the calibration sources 3C48 (0134$`+`$329), 3C138 (0521$`+`$166) and 3C286 (1331+305) to get sufficient parallactic coverage across both days during which the target sources were observed. We assume that the instrument is stable between observing epochs.
The parallel hand (RR, LL) data were self-calibrated and imaged in the normal iterative manner. The complex gain corrections derived from self-calibration were also applied to the cross-hand (RL or LR) fringes. In turn, images in Stokes parameters I, Q, U and V were produced. Maps of the polarized flux density $`P=(Q^2+U^2)^{1/2}`$ and position angle $`\chi =0.5\times \mathrm{tan}^1(U/Q)`$, were then generated from the Q and U images. The Stokes V images were used to test the integrity of the calibration and self-calibration procedures and to diagnose problems due to interference.
At low frequencies, the primary beam of the antennas contains many background sources. In order to image the target sources of interest, it was necessary to image some of these background sources. The brightest ($`>20`$ mJy at L band) background sources were identified from the NRAO VLA Sky Survey (NVSS) (Condon et al. 1998). This threshold is somewhat arbitrary but gives a reasonable number of secondary fields to image at 320 MHz. We imaged a total of 26 fields for 0235–197 and 23 fields for 1203$`+`$043 in the AIPS mapping program, IMAGR — one field containing the program source and the others on the brightest background sources. We were able to account for substantially all of the flux density seen on the shortest baselines in both cases and to obtain satisfactory convergence in the iterative self-calibration and imaging loop. The final images in all Stokes parameters (I, Q, U & V) of the target sources contained no obvious artefacts due to sidelobes from nearby confusing sources. The linear polarizations of the background sources have been checked in order to test for beam squint between R and L beams. The r.m.s. noises in the final images are within a factor of 3 of the expected thermal noises; this is entirely consistent with other observing programs at 320 MHz. Together, these suggest that the images of the target sources are not greatly affected by confusion.
The low declination source 0235–197 was observed for only $`2`$ hours at transit; consequently, some of the data were corrupted by cross-talk between antennas. These data, mostly on baselines with immediately adjacent antennas, were excised from the database during the iterative self-calibartion and imaging process.
The C (5 GHz), X (8.4 GHz) and U (15 GHz) band data were calibrated in the standard way using VLA calibrators and AIPS procedures. Polarimetric images were generated in a manner similar to that for the P band data.
## 3 Sources structure and Polarimetry
### 3.1 Observational Parameters
The derived parameters for the low resolution observations at 320 MHz and at 5 GHz (VLA C-array) are listed in Table 2. Maps at higher resolution have been obtained at 8.4 and 15 GHz with the VLA in the A configuration. Values derived from the maps are listed in Table 3. Comments on the sources structure will be given in Section 4.
The contents of Tables 2 and 3 are: column 1 $``$ source name; column 2 $``$ the observing frequency in MHz; columns 3 to 5 $``$ major axis, minor axis (both in arcsec) and the PA in degrees of the restoring beam major axis; column 6 $``$ the rms noise in the total intensity map far from the source of emission; column 7 $``$ the rms noise $`\sqrt{\sigma _Q^2+\sigma _U^2}`$, where $`\sigma _Q`$ and $`\sigma _U`$ are the rms noises on the blank sky in the distributions of the Stokes parameters Q and U; column 8 $``$ component label; columns 9 and 10$``$ RA and Dec. of the component peak; column 11 $``$ peak flux density (mJy) of the component; column 12 $``$ total flux density (mJy) of the component.
In Tables 4 & 5, we give the measured position angle (PA) in degrees of the electric field vector at the peak of polarized emission ($`\pm `$1 rms error calculated from the distribution of PAs found in a small box around the peak of polarized emission); the Rotation Measure (RM$`=\mathrm{\Delta }\varphi (\lambda )\pm n\pi /\mathrm{\Delta }(\lambda ^2)`$ in $`radm^2`$ where $`\varphi `$($`\lambda `$) is the PA at wavelength $`\lambda `$ and $`n`$ an integer; when three frequencies are available, as for some of the components in Tab. 4, the ambiguity inplied by the integer $`n`$ can be resolved); the RM, corrected for the redshift; the percentage polarization; the depolarization index, defined as the ratio of the fractional polarization at longer wavelength to the fractional polarization at shorter wavelength; from the high and low resolution observations.
In order to compare the 5 GHz (C-Array) and 320 MHz images of 0235$``$197, we have produced 5 GHz maps (I,Q,U) at the resolution of the 320 MHz maps. This was done by restoring the 5 GHz images with the appropriate Gaussian beam during imaging. The polarization parameters derived from those images are reported in Table 5.
## 4 Notes on individual sources
### 4.1 0235$``$197
The 5 GHz images presented by Mantovani et al. (1992) and Morganti, Killeen & Tadhunter (1993), show the classical, double structure typical of powerful radio galaxies. The source is associated with a galaxy at $`z=0.620`$ (Tadhunter et al. 1993). There are no radio components above the detection limits ($``$0.2 mJy at 5 and 8.5 GHz; $``$0.5 mJy at 15 GHz) inside the error box of the optical position. Adopting the optical position as a reference, 0235$``$197 looks rather symmetric with a ratio of $``$0.8 between the full length of the two lobes, with the western being the longer.
0235$``$197 appears to be dominated by the outer lobes at frequencies $`>5`$ GHz. The most interesting feature is the bright hot spot at the far end of the Eastern lobe. At both 5 and 8.4 GHz (Figs. 1 to 3), the hot spot has a double structure with individual components labelled $`E_1`$ and $`E_2`$. Component $`E_1`$ also appears double when observed with higher resolution at 15 GHz (Fig. 4). The images have been convolved to 0.4 arcsecond resolution in order to calculate the hot spot spectral indices. Both components $`E_1`$ and $`E_2`$ show a large steepening in spectral index, $`\alpha `$ (S$`\nu ^\alpha `$), between 5–8.4 GHz and 8.4–15 GHz. We find values of $`\alpha =0.5`$ and $`\alpha =1.4`$ respectively for $`E_1`$ and $`\alpha =0.68`$ and $`\alpha =2.5`$ respectively for $`E_2`$. The front shock of the lobe $`W_1`$ at the opposite side is resolved in all of the images and can hardly be defined as a ‘hot spot’. The spectral index of the bright part at the far end is also steep ($`\alpha =1.2`$) in the range 5–8.4 GHz and it steepens to $`\alpha 2`$ in the range 8.4-15 GHz (since the 15 GHz flux density estimate is an upper limit). Note that the observations at 15 GHz have lower sensitivity to diffuse, extended emission.
All of the hot spots are highly polarized, with little depolarization and Faraday rotation. The magnetic field is parallel to the front shock and rather ordered in the lobe regions with weak diffuse emission. (The electric vector is shown in all of the images shown in the paper). At 320 MHz (Fig. 5), the polarized emission, if any, is below the detection limit of our observations. Polarized emission is detected over all of the source in the low resolution observations at 5 GHz (Fig. 6). The magnetic field is again ordered and, generally speaking, parallel to the source major axis, apart from the hot spot area, where the magnetic field is parallel to the front shock. There are two main regions of polarized emission in the $`E`$ lobe. The mean PA given in Table 5 should be treated with care because the position angles of the polarization vectors in the lobe actually vary greatly. The depolarization between 5 GHz and 320 MHz is very high (DP$`>`$0.02). Statistically, in these classical sources, the lobe nearest to the nucleus usually shows a steeper spectral index. Here we find similar values ($`\alpha =`$0.84) for the two lobes of 0235$``$197. There are indications of Faraday rotation in the hot spots from the high resolution maps.
### 4.2 1203$`+`$043
The images at 8.4 and 15 GHz (Figs. 7 to 8 & 9 respectively) do not add much new information about the overall source structure as derived from the 5 GHz image of Mantovani et al. (1992). This earlier image showed a long bent jet. The components found along the jet, labelled $`J_1`$ and $`J_2`$, are rather polarized. The emission from $`J_1`$ has a small Faraday rotation and depolarization between 6 cm and 4 cm of 17 $`radm^2`$ and 0.74 respectively.
However, the new observations have allowed us to identify component C with the core of the radio source. This component has an inverted spectrum which peaks at frequencies $``$ 15 GHz. 1203$`+`$043 therefore has an asymmetric structure, with a long bent jet pointing South which fades slowly and with a weak lobe of emission to the north where there is marginal evidence of an hot spot. The radio position of the core of 1203$`+`$043 does not coincide within the errors to any optical counterpart on the Palomar Sky Survey prints.
Much more interesting is the structure found at 320 MHz (Fig. 10). Together with the North-South structure which dominates at higher frequencies and which appears here as a ridge of emission, there is a region of diffuse emission with its major axis perpendicular to the main ridge. This new feature is about 50 arcseconds in extent and is comparable in width with the main North-South ridge.
## 5 Discussion
Due to the non-detection of polarized emission at 320 MHz in both 0235$``$197 and 1203$`+`$043, we cannot explain the low frequency variability observed with the Molonglo Cross (McAdam 1980) in terms of ionospheric Faraday rotation as in the case of 3C159.
Can refractive scintillation (Rickett, 1986) explain the variability of 0235$``$197? The refractive scintillation models assume a supposedly-variable radio source to have most of its flux density in a single compact component. The degree of variability is a function of the characterization of the interstellar medium (itself a function of galactic coordinates) and source size. 0235$``$197 is at galactic latitude $`|b|=65^{}`$ and was reported to vary (rms variability $`0.2`$ Jy) on a time scale of the order of 1 year (McAdam, 1980).
In the usual refractive model of interstellar turbulence (Mantovani et al. 1990b, Spangler et al. 1993, Spangler et al. 1994, Bondi et al. 1994), the relevant parameter for the observed scintillation index is $`\theta _{eff}^{7/6}\sqrt{\mathrm{sin}b/I}`$, where $`\theta _{eff}=\theta _{\mathrm{FWHM}}/2.35`$ and $`I`$ is the parameter indicator of the source structure ($`=1`$ for a gaussian structure). In such a model a source of ($`3`$ Jy) with a rms variability of 0.2 Jy and corresponding scintillation index of $`0.06`$ should have an angular diameter in the range 10–20 mas (see Spangler et al. 1993 for details).
The $`E_{1a}`$ hot spot has a spectrum which, extrapolated towards lower frequencies, gives a mean flux density of $``$3 Jy at 408 MHz. This is comparable to the peak flux density found at 320 MHz. The crucial parameter is, however, the angular size of the hot spot. The deconvolved size found for the $`E_{1a}`$ hot spot at 15 GHz is $`<`$0.2 arcsec. Even if in principle there is not contradiction, the scintillation theory requires an angular size for the hot spot in 0235$``$197 that is a factor 7–10 smaller than the size measured at 15 GHz, which is close to the sizes usually measured for the hot spots. Low frequency VLBI observations are needed to confirm the existence of such a compact component in the hot spot.
The lack of polarized emission at 320 MHz for 1203$`+`$043 and a radio structure which lacks a bright compact component rules out both of the mechanisms for low frequency variability in this source. Consequently, we conclude that this source might be a spurious case of variability. However, 1203$`+`$043 has an interesting structure at 320 MHz. It shows a pair of secondary lobes in a direction perpendicular to the main source axis, making the object one of a few known ‘X’-shaped sources. At present, only about ten sources are known to show such morphology. They are believed to have both young and old lobes. These lobes may be supplied by jets whose direction has changed with time. A change in the orientation of the central engine due to precession has been suggested by Ekers et al. (1978) for NGC326 to account for its ‘X’ shaped morphology. Such a model has been applied successfully to 0828$`+`$32 by Klein et al. (1995) but with the extra assumption that the length of the precessing beam changes with time. A merger between galaxies is thought to be the cause of the precession. However, Ulrich-Demoulin & Rönnback (1996) have reported that optical images of 0828$`+`$32 do not show the signature of a recent major merger event.
However, the structure of 1203$`+`$043 looks peculiar when compared with other ’X’ shaped sources. For example, it has an asymmetric structure with respect to the component C which is believed to be the core (Fig. 7 and Fig. 6 in Mantovani et al. 1992). The long, bent jet is clearly ’one-sided’ while, generally speaking, the members of the class show two-sided jets (at the available resolution). Moreover, the young lobes of the ’X’ shaped sources are dominated by hot spots while here the jet emission fades away from the core and the northern lobe contains only diffuse emission without any bright component. This asymmetry is reflected in the 320 MHz map where the region to the North-West is more extended and brighter than the opposite side.
## 6 Conclusions
We have conducted a program of multi-wavelength VLA observations of the suspected low frequency variable sources 0235$``$197 and 1203+043. Since 0235$``$197 is not polarized at 320 MHz, its variability cannot be accounted for by instrumental polarization effects as in the case of 3C159. 0235$``$197 may contain a low frequency component sufficiently compact and bright as required by the refractive scintillation model for low frequency variability. Our observations have insufficient resolution to test this suggestion; low frequency VLBI observations are required for this purpose. However, this component would have to have extremely unusual properties among hot spots in radio sources.
In our high frequency images of 1203+043 we have identified the core of the radio source; its location indicates that the source has a large apparent asymmetry. At 320 MHz, this source shows no polarization. However, it does have an additional, steep-spectrum component at this frequency; this previously-undetected component lies perpendicular to the main axis and predominantly to one side. However, the overall morphology of 1203+043 at low frequencies seems similar to that of the ‘X’-shaped sources like NGC326. From its morphology and component sizes, we conclude that 1203+043 is likely not variable at low frequencies and that its inclusion in such catalogs is spurious.
###### Acknowledgements.
The authors like to thank the referee, Dr. Steve Spangler for his comments to the paper and Dr. Ian Browne for a critical reading of the manuscript. FM thanks Miller Goss, Assistant Director, NRAO, Socorro, for his hospitality during period when part of the work was done. The National Radio Astronomy Observatory is operated by Associated Universities Inc., under cooperative agreement with the National Science Foundation; AIPS is NRAO’s Astronomical Image Processing System. |
warning/0003/quant-ph0003039.html | ar5iv | text | # REFERENCES
Perturbative expansion for master equation and its applications
X. X. Yi<sup>a,b</sup>,C. Li<sup>a</sup> J. C. Su<sup>c</sup>
<sup>a</sup>Institute of Theoretical Physics, Northeast Normal University,
Changchun 130024, China
<sup>b</sup>Institute of Theoretical Physics, Academia Sinica, Peking 100080, China<sup>*</sup><sup>*</sup>*Corresponding address
<sup>c</sup> Department of Physics, Jilin University, Changchun 130023, China
We construct generally applicable small-loss rate expansions for the density operator of an open system. Successive terms of those expansions yield characteristic loss rates for dissipation processes. Three applications are presented in order to give further insight into the context of those expansions. The first application, of a two-level atom coupling to a bosonic environment, shows the procedure and the advantage of the expansion, whereas the second application that consists of a single mode field in a cavity with linewidth $`\kappa `$ due to partial transmission through one mirror illustrates a practical use of those expansions in quantum measurements, and the third one, for an atom coupled to modes of a lossy cavity shows the another use of the perturbative expansion.
PACS numbers:03.65.-w,05.30.Fk,42.50.Dv
The study of open quantum systems has recently attracted the attention of physicists from various fields: cosmology, condensed matter, quantum optics\[3-7\], particle physics, quantum measurement, and quantum computation. The problem can be described generally as interest in the effective dynamics of one subsystem of several interacting subsystems. A formal framework to describe the effective dynamics of such a subsystem is set up in ref., and a short-time perturbative expansion for coherence loss has also been constructed. To some extent(for example, if we are interested in a behavior for finite time), however, time is not as good as the loss rate as a perturbative parameter. Motivated by this and recent experimental developments\[15-20\] as well as the analysis of models related to them\[21-24\], we construct generally small-loss rate expansions for dissipation. The results suggest that these are useful in many areas such as high-Q Cavity QED\[15-20\], quantum computation\[11,12,21-26\], quantum measurement, quantum optics\[3-7\] etc.
We consider an open quantum system, the total Hamiltonian describing such a system is expressed as
$$H=H_0+H_{env}+H_I,$$
(1)
where $`H_0`$ and $`H_{env}`$ indicate the free Hamiltonian of the system and of the environment, respectively. $`H_I`$ is the interaction Hamiltonian between the system and the environment. It is well known that the form of the master equation depends on the precise kind of the system-environment interaction. In order to derive a master equation for a quite general $`H_I`$, let us suppose that, in the Schrödinger picture, $`H_I`$ can be written as
$$H_I=\mathrm{}\underset{m}{}(X_m^+A_m+X_m^{}A_m^{})$$
(2)
where the $`X_m^\pm `$are eigenoperators of the system satisfying
$$[H_0,X_m^\pm ]=\pm \mathrm{}\omega _mX_m^\pm .$$
(3)
This form is quite general, since any system operator can be decomposed into eigenoperators of $`H_0`$. As shown in ref., we can write the master equation in the following form (in the Schrödinger picture)
$`\dot{\rho }(t)`$ $`=`$ $`i[H_0,\rho ]+{\displaystyle \frac{1}{2}}{\displaystyle \underset{m}{}}K_m(2X_m^{}\rho X_m^+X_m^+X_m^{}\rho \rho X_m^+X_m^{})`$ (4)
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{m}{}}G_m(2X_m^+\rho X_m^{}X_m^{}X_m^+\rho \rho X_m^{}X_m^+),`$ (5)
where
$$K_m=2Re[_0^{\mathrm{}}𝑑\tau e^{i\omega _m\tau }\mathrm{Tr}_{env}\{A_m(\tau )A_m^{}(0)\rho _{env}\}],$$
$$G_m=2Re[_0^{\mathrm{}}𝑑\tau e^{i\omega _m\tau }\mathrm{Tr}_{env}\{A_m^{}(\tau )A_m(0)\rho _{env}\}],$$
$`\rho (t)=\rho (t,K_m,G_m)`$ stands for the density operator of the system and $`\rho _{env}`$ denotes the density operator of the environment. Notice from eq.(4) that $`G_m`$ should vanish at zero temperature $`T=0`$, while $`K_m`$ should not if $`A_m`$ are indeed destruction operators of some kind. In case the constant $`G_m`$ and $`K_m`$ are smaller than any one of the internal coupling parameters of the system, the density operator may be expanded in powers of $`K_m`$ and $`G_m`$,
$`\rho (t,K_m,G_m)`$ $`=`$ $`\rho (t,0,0)+{\displaystyle \underset{m}{}}{\displaystyle \frac{\rho }{K_m}}K_m+{\displaystyle \underset{m}{}}{\displaystyle \frac{\rho }{G_m}}G_m`$ (6)
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{m,n}{}}{\displaystyle \frac{^2\rho }{K_mK_n}}K_mK_n+{\displaystyle \frac{1}{2}}{\displaystyle \underset{m,n}{}}{\displaystyle \frac{^2\rho }{G_mG_n}}G_mG_n`$ (7)
$`+`$ $`{\displaystyle \underset{m,n}{}}{\displaystyle \frac{^2\rho }{G_mK_n}}G_mK_n+\mathrm{}.`$ (8)
Substituting this expression into the master equation, we find the following set of equations
$`\dot{\rho }(t,0,0)`$ $`=`$ $`i[H_0,\rho (t,0,0)]`$ (9)
$`{\displaystyle \frac{\dot{\rho }}{G_m}}`$ $`=`$ $`i[H_0,{\displaystyle \frac{\rho }{G_m}}]+{\displaystyle \frac{1}{2}}(2X_m^+\rho (t,0,0)X_m^{}X_m^{}X_m^+\rho (t,0,0)\rho (t,0,0)X_m^{}X_m^+)`$ (10)
$`{\displaystyle \frac{\dot{\rho }}{K_m}}`$ $`=`$ $`i[H_0,{\displaystyle \frac{\rho }{K_m}}]+{\displaystyle \frac{1}{2}}(2X_m^{}\rho (t,0,0)X_m^+X_m^+X_m^{}\rho (t,0,0)\rho (t,0,0)X_m^+X_m^{})`$ (11)
$`{\displaystyle \frac{^2\dot{\rho }}{K_mK_n}}`$ $`=`$ $`i[H_0,{\displaystyle \frac{\rho }{K_mK_n}}]+{\displaystyle \frac{1}{2}}(2X_m^{}{\displaystyle \frac{\rho }{K_n}}X_m^+X_m^+X_m^{}{\displaystyle \frac{\rho }{K_n}}{\displaystyle \frac{\rho }{K_n}}X_m^+X_m^{})`$ (12)
$`+`$ $`{\displaystyle \frac{1}{2}}(2X_n^{}{\displaystyle \frac{\rho }{K_m}}X_n^+X_n^+X_n^{}{\displaystyle \frac{\rho }{K_m}}{\displaystyle \frac{\rho }{K_m}}X_n^+X_n^{})`$ (13)
$`{\displaystyle \frac{^2\dot{\rho }}{G_mG_n}}`$ $`=`$ $`i[H_0,{\displaystyle \frac{\rho }{G_mG_n}}]+{\displaystyle \frac{1}{2}}(2X_m^+{\displaystyle \frac{\rho }{G_n}}X_m^{}X_m^{}X_m^+{\displaystyle \frac{\rho }{G_n}}{\displaystyle \frac{\rho }{G_n}}X_m^{}X_m^+)`$ (14)
$`+`$ $`{\displaystyle \frac{1}{2}}(2X_n^+{\displaystyle \frac{\rho }{G_m}}X_n^{}X_n^{}X_n^+{\displaystyle \frac{\rho }{G_m}}{\displaystyle \frac{\rho }{G_m}}X_n^{}X_n^+)`$ (15)
$`{\displaystyle \frac{^2\dot{\rho }}{K_mG_n}}`$ $`=`$ $`i[H_0,{\displaystyle \frac{\rho }{K_mG_n}}]+{\displaystyle \frac{1}{2}}(2X_m^{}{\displaystyle \frac{\rho }{G_n}}X_m^+X_m^+X_m^{}{\displaystyle \frac{\rho }{G_n}}{\displaystyle \frac{\rho }{G_n}}X_m^+X_m^{})`$ (16)
$`+`$ $`{\displaystyle \frac{1}{2}}(2X_n^+{\displaystyle \frac{\rho }{K_m}}X_n^{}X_n^{}X_n^+{\displaystyle \frac{\rho }{K_m}}{\displaystyle \frac{\rho }{K_m}}X_n^{}X_n^+)`$ (17)
Generally speaking, given an initial condition, $`\rho (0,0,0)`$, we can solve eq.(6) exactly, which gives the zeroth order solution for the density operator $`\rho (t,0,0)`$. Substituting the zeroth order solution into eq.(7), $`\rho /K_m`$ or $`\rho /G_m`$ can be calculated. Following this procedure, successive terms of the expansion (5) could be worked out, though the calculation is complicated. Some words of caution are now in order. From the mathematical point of view, the expansion (5) holds if and only if the series converges. This may be satisfied easily in physics for a large number of open systems. For example in a high-Q cavity, the loss rate of the atom-cavity system is small enough to permit us to expand the density operator in powers of the loss rate. At the end of this section, we will present some discussion about this point in detail.
To illustrate the advantage of the expansion (5), we present here a simplest model, which describes a two-level atom coupling to a bose-mode environment. The master equation of such a system is given by
$$\dot{\rho }=\frac{1}{2}i\mathrm{\Omega }[\sigma _z,\rho ]+\frac{1}{2}\gamma \{2\sigma ^{}\rho \sigma ^+\rho \sigma ^+\sigma ^{}\sigma ^+\sigma ^{}\rho \}$$
(18)
with
$$\gamma =2\pi Re[_0^{\mathrm{}}e^{i\omega _m\tau }\mathrm{Tr}_{env}\{b_m(\tau )b_m^{}(0)\rho _{env}\}],$$
where $`b_m^{}(b_m)`$ stands for the creation(annihilation) operator of the m-th mode of the environment, $`\mathrm{\Omega }`$ is the Rabi frequency , and $`\sigma _z(\sigma ^+,\sigma ^{})`$ denote the Pauli matrices. To obtain the form of the master equation given in eq.(9), the environment was assumed to be in its vacuum state. According to eq.(5), $`\sigma _z(t)`$ reads:
$$\sigma _z(t)=\mathrm{Tr}(\rho (t,0,0)\sigma _z)+\gamma \mathrm{Tr}(\frac{\rho }{\gamma }\sigma _z)+\frac{1}{2}\gamma ^2\mathrm{Tr}(\frac{^2\rho }{\gamma ^2}\sigma _z)+\mathrm{}$$
(19)
The first term in eq.(10) is $`\sigma _z(0)`$. In order to calculate $`\mathrm{Tr}(\frac{\rho }{\gamma }\sigma _z)`$, we first evaluate $`\mathrm{Tr}(\frac{\dot{\rho }}{\gamma }\sigma _z)`$, it is given that
$$\mathrm{Tr}(\frac{\dot{\rho }}{\gamma }\sigma _z)=\sigma _z(0).$$
(20)
Using the same procedure as mentioned, we arrive at
$$\mathrm{Tr}(\frac{^2\dot{\rho }}{\gamma ^2}\sigma _z)=2\sigma _z(0)t.$$
(21)
The eqs.(11,12) together give
$$\sigma _z(t)=\sigma _z(0)\gamma t\sigma _z(0)+\frac{\gamma ^2}{2!}t^2\sigma _z(0)+\mathrm{}$$
(22)
Using the algebra of Pauli matrices, we obtain straightforwardly from the master equation that
$$\sigma _z(t)=\sigma _z(0)e^{\gamma t}.$$
(23)
A comparison between eq.(14) and (13) shows that for small $`\gamma `$, the expansions are a quite good approach for the two level dissipative system, and this result is quite general.
Noticing the eq.(13) is expanded in the product of loss rate and time, we present here the other example to show that these expansions are generally in powers of the loss rate, but not in the product of time and the loss rate. Consider the master equation given in eq.(9) As mentioned above, we calculate $`\sigma _z(t)`$ in order to illustrate the advantage of the expansions. The results of $`\sigma _z(t)`$ show no difference between the short-time expansions and the small loss-rate expansions. To show the difference between the two expansions, we calculate $`\sigma _x(t)`$. For simplicity, we only present the results up to first order of $`\gamma `$. It follows from eq.(5) that
$$\sigma _x(t)=\mathrm{Tr}(\rho (t,0,0)\sigma _x)+\gamma \mathrm{Tr}(\frac{\rho }{\gamma }\sigma _x)+\mathrm{}$$
(24)
It is easy to show that (setting $`\mathrm{}=1,\sigma _x(0)=1`$)
$$\mathrm{Tr}(\rho (t,0,0)\sigma _x)=\mathrm{cos}(\mathrm{\Omega }t).$$
In order to compute $`\mathrm{Tr}(\frac{\rho }{\gamma }\sigma _x)`$, we have to calculate $`\mathrm{Tr}(\frac{\dot{\rho }}{\gamma }\sigma _x)`$. Based on the expansions, we arrive at
$$\mathrm{Tr}(\frac{\dot{\rho }}{\gamma }\sigma _x)=\mathrm{Tr}(i[H_0,\frac{\rho }{\gamma }\sigma _x])+\frac{1}{2}(2\sigma ^{}\rho (t,0,0)\sigma ^+\sigma _x\rho (t,0,0)\sigma ^+\sigma ^{}\sigma _x\sigma ^+\sigma ^{}\rho (t,0,0)\sigma _x)a+b.$$
Simple calculation gives
$$a=\mathrm{\Omega }\mathrm{Tr}(\frac{\rho }{\gamma }\sigma _y).$$
As state above, in order to compute $`a`$, we have to calculate $`\mathrm{Tr}(\frac{\dot{\rho }}{\gamma }\sigma _y).`$ Using the same procedure as above, we show that
$$a=\mathrm{\Omega }^2_0^t𝑑t\mathrm{Tr}(\frac{\rho }{\gamma }\sigma _x)(\mathrm{cos}(\mathrm{\Omega }t)1),$$
and
$$b=\mathrm{cos}(\mathrm{\Omega }t).$$
The results for $`a`$ and $`b`$ together give
$$\frac{^2y(t)}{t^2}+\mathrm{\Omega }^2y(t)=2\mathrm{\Omega }\mathrm{sin}(\mathrm{\Omega }t),$$
with $`y(t)\mathrm{Tr}(\frac{\rho }{\gamma }\sigma _x)`$ and initial conditions $`y(t=0)=0,\dot{y}(t)|_{t=0}=1`$. This is a two order differential equation and can be solved easily, once $`y(t)`$ is known, $`\sigma _x(t)`$ up to first order of $`\gamma `$ is given. It is obvious that results given above are indeed different from the short time expansions, since the results given by short time expansions are in powers of time $`t`$.
The results up to first order of $`\gamma `$(15) and an exact numerical results are illustrated in Fig.1. The parameters chosen are $`\mathrm{\Omega }=2`$, and time is in units of $`\mathrm{\Omega }`$. In Fig.1 the scattering line represents the exact numerical results, whereas the dot line and the solid line show the results from the expansion. The dot line and the solid line are for different $`\gamma `$, and $`\gamma `$ for dot line is smaller than one in solid line, those curves show that the expansions (15) are indeed a good approximation to the exact solution.
We need to point out that for most open systems, average values of meaningful quantities can’t be obtained exactly in any way. Therefore the expansions of the density operator provide a practical approach to the exact solution. For example, consider a single-mode field in a lossy cavity. The density operator for that mode obeys the following master equation in the Schrödinger picture,
$$\dot{\rho }=i[\omega _fa^{}a,\rho ]+\frac{\kappa }{2}(2a\rho a^{}a^{}a\rho \rho a^{}a),$$
(25)
where $`\kappa `$ is the linewidth of the cavity mode with frequency $`\omega _f`$. In most textbooks, the solution of the master equation is given in terms of diagonal matrix elements $`n|\rho |n`$ in a stationary state. Given an initial condition for the density operator, the evolution of $`\rho `$, however, is more useful than the stationary solution. In what follows, we present a solution of the master equation in a number state (Fock state) basis.
For a high-Q cavity, the linewidth $`\kappa `$ due to partial transmission through one mirror is so small that we can expand $`\rho `$ in powers of $`\kappa `$:
$$\rho (t)=\rho ^0(t)+\frac{\rho }{\kappa }\kappa +\frac{1}{2}\frac{^2\rho }{\kappa ^2}\kappa ^2+\mathrm{},$$
(26)
In a number state basis $`\{|n,n=0,1,2,3\mathrm{}\}`$, the expansion can be written as
$$\rho (t)=\underset{m,n}{}\rho _{mn}^0(t)|mn|+\underset{m,n}{}\underset{k=1}{\overset{\mathrm{}}{}}\frac{1}{k!}\frac{^k\rho _{mn}}{\kappa ^k}|mn|\kappa ^k.$$
(27)
Here the subscripts on the density operator $`\rho _{mn}`$ indicate matrix elements of $`\rho `$ in the number basis and $`\rho _0`$ is the solution of eq.(16) with $`\kappa =0`$. With this notation, it follows from eqs.(6),(7) and (8) that
$`\rho _{mn}^0(t)`$ $`=`$ $`\rho _{mn}^0(0)e^{i(mn)\omega _ft},`$ (28)
$`{\displaystyle \frac{^k\rho _{mn}}{\kappa ^k}}`$ $`=`$ $`c_k(t)e^{i\omega _f(nm)t}.`$ (29)
This iterative equation gives the density operator expansions of the system under consideration. Here, $`\rho _{mn}(0)`$ stands for the initial condition of $`\rho `$, and
$`c_k(t)`$ $`=`$ $`{\displaystyle _0^t}F_k(t^{^{}})e^{i\omega _f(nm)t^{^{}}}𝑑t^{^{}},`$ (30)
$`F_1(t)`$ $`=`$ $`\sqrt{(n+1)(m+1)}\rho _{m+1,n+1}^0(t){\displaystyle \frac{1}{2}}m\rho _{mn}^0(t){\displaystyle \frac{1}{2}}n\rho _{mn}^0(t),`$ (31)
$`F_k(t)`$ $`=`$ $`\sqrt{(n+1)(m+1)}{\displaystyle \frac{^{k1}\rho _{m+1,n+1}}{\kappa ^{k1}}}{\displaystyle \frac{m+n}{2}}{\displaystyle \frac{^{k1}\rho _{m,n}}{\kappa ^{k1}}},k=2,3,4\mathrm{}`$ (32)
The master equation in the form (16) is widely used in field-quadrature measurement. As shown in Ref., different approximations to eq.(16) correspond to different measurement schemes, therefore the expansions(6-8) for the density operator provide a new method to develop quantum measurement theory. In contrast to the short time perturbative expansions, the expansions (6-8) hold for finite time as long as the linewidth $`\kappa `$ is small. In other words, whether the expansions hold does not depend on time $`t`$.
In addition to quantum measurement, these expansions have use in high-Q cavity QED. There are many interesting features in cavity QED. One of them is spontaneous emission. Spontaneous emission is so fundamental that it is usually regarded as an inherent property of matter. The master equation for a single atom coupling to a mode of a lossy cavity is given in the interaction picture under rotating-wave and dipole approximations by
$$\dot{\rho _I}=\frac{\gamma }{2}(2\sigma ^{}\rho _I\sigma ^+\sigma ^+\sigma ^{}\rho _I\rho _I\sigma ^+\sigma ^{})+\frac{\kappa }{2}(2a\rho _Ia^{}a^{}a\rho _I\rho _Ia^{}a),$$
(33)
where $`\rho _I`$ stands for the reduced density operator of the system that consists of an atom and a cavity mode, $`\gamma `$ denotes the linewidth of the atom, and $`\kappa `$ describes the loss rate of the cavity.This is different from eq.(9) in which the loss of the single-mode field is neglected. It is well known that the emission spectrum may be expressed in terms of average values of the atom operator $`\stackrel{}{\sigma }`$. In this sense, we may calculate the average value of the atom operator to replace computing the emission spectrum without any loss of generality. Moreover, the study of many other effects in cavity QED such as atomic dipole squeezing, population trapping, and atomic collapse-and-revival phenomenon may be reduced to calculate and analyse the average value of the atom operator. In the remainder of this paper, based on the expansion scheme, we compute the average value of an atom operator given by $`A=\lambda ^{(+)}\sigma ^++\lambda ^{()}\sigma ^{}+\lambda ^{(z)}\sigma _z`$. For this end, we first of all list the expansions of the density operator $`\rho _I(t)`$ in the interaction picture
$`\rho _I(t)`$ $`=`$ $`\rho _I^0(t)+{\displaystyle \frac{\rho _I(t)}{\kappa }}\kappa +{\displaystyle \frac{\rho _I(t)}{\gamma }}\gamma +{\displaystyle \frac{1}{2}}{\displaystyle \frac{^2\rho _I(t)}{\kappa ^2}}\kappa ^2+{\displaystyle \frac{1}{2}}{\displaystyle \frac{^2\rho _I(t)}{\gamma ^2}}\gamma ^2+{\displaystyle \frac{1}{2}}{\displaystyle \frac{^2\rho _I(t)}{\gamma \kappa }}\gamma \kappa +\mathrm{}`$ (34)
Here,
$$\dot{\rho }_I^0(t)=0,$$
$$\frac{\dot{\rho _I}(t)}{\gamma }=\frac{1}{2}(2\sigma ^{}\rho _I^0(t)\sigma ^+\sigma ^+\sigma ^{}\rho _I^0(t)\rho _I^0(t)\sigma ^+\sigma ^{}),$$
$$\frac{\dot{\rho _I}(t)}{\kappa }=\frac{1}{2}(2a\rho _I^0(t)a^{}a^{}a\rho _I^0(t)\rho _I^0(t)a^{}a),$$
$$\frac{^2\dot{\rho _I}(t)}{\gamma ^2}=\frac{1}{2}(2\sigma ^{}\frac{\rho _I}{\gamma }\sigma ^+\sigma ^+\sigma ^{}\frac{\rho _I}{\gamma }\frac{\rho _I}{\gamma }\sigma ^+\sigma ^{}),$$
$$\frac{^2\dot{\rho _I}(t)}{\kappa ^2}=\frac{1}{2}(2a\frac{\rho _I}{\kappa }a^{}a^{}a\frac{\rho _I}{\kappa }\frac{\rho _I}{\kappa }a^{}a),$$
$$\frac{^2\dot{\rho _I}(t)}{\gamma \kappa }=\frac{1}{2}(2\sigma ^{}\frac{\rho _I}{\kappa }\sigma ^+\sigma ^+\sigma ^{}\frac{\rho _I}{\kappa }\frac{\rho _I}{\kappa }\sigma ^+\sigma ^{})+\frac{1}{2}(2a\frac{\rho _I}{\gamma }a^{}a^{}a\frac{\rho _I}{\gamma }\frac{\rho _I}{\gamma }a^{}a).$$
It is easy to show that for any atom operator $`A`$,
$$\mathrm{Tr}(\frac{^n\dot{\rho }_I(t)}{\kappa ^n}A)=\mathrm{Tr}(\frac{^n\dot{\rho }_I(t)}{\kappa ^m\gamma ^{nm}}A)=0$$
(35)
for $`nm0`$, while
$`\mathrm{Tr}({\displaystyle \frac{\dot{\rho }_I}{\gamma }}A)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{Tr}(\rho ^0(t)B)`$ (36)
$`\mathrm{Tr}({\displaystyle \frac{^n\dot{\rho }_I}{\gamma ^n}}A)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{Tr}({\displaystyle \frac{^{n1}\rho _I}{\gamma ^{n1}}}B),`$ (37)
where $`B=4\lambda ^{(+)}\sigma ^+4\lambda ^{()}\sigma ^{}.`$ Eqs.(23) and (24) suggest that the average value for any atom operator can be calculated analytically as an expansion in powers of $`\kappa `$ and $`\gamma `$, provided $`\rho ^0(t)`$ (the zeroth order density operator in the Schrödinger picture) is known. Generally speaking, given a initial condition for $`\rho `$, the $`\rho ^0(t)`$ that obeys the von Neumann equation can be given readily. In the model presented above, the von Neumann equation is given by
$$\dot{\rho }^0=i[H_0,\rho ^0],$$
(38)
Where $`H_0`$ denotes the free Hamiltonian for the cavity-atom system (Jaynes-Cummings model)
$$H_0=\omega _fa^{}a+\frac{1}{2}\omega _a\sigma _z+g(a^{}\sigma ^{}+\sigma ^+a).$$
(39)
If the cavity-atom system is initially in a state $`|e,n=|e|n`$, i.e. the atom is in its excited state, while the single-mode cavity is in the number state $`|n`$, then $`\rho ^0(t)`$ reads
$$\rho ^0(t)=|\psi ^0(t)\psi ^0(t)|,$$
(40)
where
$`|\psi ^0(t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{sin}\theta _{n+1}(e^{iE_+(n+1)t}e^{iE_{}(n+1)t})|g,n+1`$ (41)
$`+`$ $`(\mathrm{sin}^2{\displaystyle \frac{\theta _{n+1}}{2}}e^{iE_+(n+1)t}+\mathrm{cos}^2{\displaystyle \frac{\theta _{n+1}}{2}}e^{iE_{}(n+1)t})|e,n.`$ (42)
$$E_\pm (n+1)=\frac{\omega _f}{2}(2n+1)\pm \sqrt{\delta ^2+g^2(n+1)},\theta _{n+1}=arctg\frac{2g\sqrt{n+1}}{\delta },\delta =\omega _f\omega _a.$$
(43)
It follows from eq.(6,7,8) that
$`\mathrm{Tr}(\rho ^0(t)A)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\mathrm{Tr}(\rho _0(t)B)+\lambda ^{(z)}{\displaystyle \underset{n}{}}\{|\mathrm{sin}^2{\displaystyle \frac{\theta _{n+1}}{2}}e^{iE_+(n+1)t}+\mathrm{cos}^2{\displaystyle \frac{\theta _{n+1}}{2}}e^{iE_{}(n+1)t}|^2`$ (44)
$``$ $`\mathrm{sin}^2\theta _{n+1}\mathrm{sin}^2(\sqrt{\delta ^2+g^2(n+1)}t)\}`$ (45)
$`\mathrm{Tr}(B\rho ^0(t))`$ $`=`$ $`{\displaystyle \underset{n}{}}\{2\lambda ^{(+)}\mathrm{sin}\theta _n(e^{iE_+(n)t}e^{iE_{}(n)t})`$ (46)
$`(\mathrm{sin}^2{\displaystyle \frac{\theta _{n+1}}{2}}e^{iE_+(n+1)t}+\mathrm{cos}^2{\displaystyle \frac{\theta _{n+1}}{2}}e^{iE_{}(n+1)t})`$ (47)
$`2\lambda ^{()}(\mathrm{sin}^2{\displaystyle \frac{\theta _{n+1}}{2}}e^{iE_+(n+1)t}+\mathrm{cos}^2{\displaystyle \frac{\theta _{n+1}}{2}}e^{iE_{}(n+1)t})`$ (48)
$`\mathrm{sin}\theta _n(e^{iE_+(n)t}e^{iE_{}(n)t})\}.`$ (49)
Then successive perturbative terms of average value up to second order of $`\gamma `$ for an atom operator are given by
$$A(t)=\mathrm{Tr}(\rho _0(t)A)+\frac{\gamma }{2}_0^t\mathrm{Tr}(\rho ^0(t^{^{}})B)𝑑t^{^{}}\frac{\gamma ^2}{2}_0^t𝑑t^{^{}}_0^t^{^{}}\mathrm{Tr}(\rho ^0(t^{^{\prime \prime }})B)𝑑t^{^{\prime \prime }}+\mathrm{}.$$
(50)
Based on the short-time expansion, $`A(t)=\mathrm{Tr}(A\rho (t))A_0+A_1t+A_2t^2+\mathrm{}`$, in powers of $`t`$, where $`A_0=\mathrm{Tr}(A\rho (0))`$, $`A_1=\mathrm{Tr}(A\frac{\rho }{t}(0))`$ and $`A_2=\frac{1}{2}\mathrm{Tr}(A\frac{^2\rho }{t^2}(0)).`$ This is quite different from the results given by eq.(31).
Although we are currently investigating the perturbative expansion for an open system, we opt here for a few qualitative comments. Mathematically, the perturbative expansion is a good approach to the exact solution of the master equation so long as loss rates $`\gamma `$ and $`\kappa `$ are smaller than all other internal coupling constants of the system. This condition holds for high-Q cavities from the physical point of view. In fact, an optical cavity of $``$ 20$`\mu m`$ diameter has $`g/2\pi 125MHz`$ and $`\kappa /2\pi 100KHz`$ for reasonable $`Q10^9`$. Thus the ratio $`g/\kappa 10^3`$. Even $`g/\kappa 10^4`$ seems feasible for microspheres. Generally, in the optical domain $`g/\gamma 10^2`$, great enough for the perturbative expansion in powers of $`\gamma `$ to hold.
In the end of this paper, we turn our attention to study the decoherence in $`N`$ two-level atoms. This problem is usually related to the register in quantum computer. A few papers have been published on this subject , but a key additional feature of the present paper is to study the decoherence from a new aspect. If the system consists of $`N`$ two-level atoms, the decoherence is due to the inevitable coupling of the $`N`$ atoms to the external environment. Generally, the environment may be treated as that consists of an infinite number of oscillators. The Hamiltonian describing such decoherence process takes the form
$`H`$ $`=`$ $`H_s+H_{env}+H_I,`$ (51)
$`H_s`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\mathrm{\Omega }_i\sigma _i^z,`$ (52)
$`H_{env}`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}\omega _kb_k^{}b_k,`$ (53)
$`H_I`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}(g_{ki}b_k^{}\sigma _i^{}+h.c.),`$ (54)
where $`\sigma _i^\alpha `$ are the spin-$`\frac{1}{2}`$ Pauli operators ($`i`$ denotes the qubit index) and $`b_k,b_k^{}`$ are the bosonic operators, $`H_s`$, $`H_{env}`$ are the free Hamiltonian of the system and the environment, respectively. And $`H_I`$ stands for the $`N`$ qubits-environment interaction. This model is closely related to the Dicke maser model. The Hamiltonian (32) is complicated so it is hard to find its exact solution though the Hilbert space associated with this model can split into invariant eigenspaces. Fortunately, with the perturbative approach created in the previous section, the complex system can be easily treated. To start with, we give the master equation of the system
$`\dot{\rho }`$ $`=`$ $`i[H,\rho ]+{\displaystyle \frac{1}{2}}\{{\displaystyle \underset{i}{}}K_i(2\sigma _i^{}\rho \sigma _i^+\rho \sigma _i^+\sigma _i^{}\sigma _i^+\sigma _i^{}\rho )\}`$ (55)
$`+`$ $`{\displaystyle \frac{1}{2}}\{{\displaystyle \underset{i}{}}G_i(2\sigma _i^+\rho \sigma _i^{}\sigma _i^{}\sigma _i^+\rho \rho \sigma _i^{}\sigma _i^+)\}`$ (56)
$`=`$ $`i[H,\rho ]+\rho .`$ (57)
Here
$$K_m=2Re𝑑\tau e^{i\mathrm{\Omega }_m\tau }\mathrm{Tr}_{env}\{A_m(\tau )A_m^{}(0)\rho _{env}\},$$
$$G_m=2Re𝑑\tau e^{i\mathrm{\Omega }_m\tau }\mathrm{Tr}_{env}\{A_m^{}(\tau )A_m(0)\rho _{env}\},$$
$$A_m(\tau )=\underset{j=1}{\overset{\mathrm{}}{}}\mathrm{}g_{mj}b_je^{i\omega _j\tau }.$$
In the case of so-called Dicke limit, $`A_m`$ does not depend on the atom index $`m`$. This holds, for example, when the typical environment wavelengths are much greater than the distances between the atoms.
In order to study the decoherence of the atoms, we assume that the initial state of the system is
$$|\psi _m=S_m^+|0+\underset{j=m+1}{\overset{N}{}}|0),$$
(58)
where $`S_m^+=_{j=1}^{mN}\sigma _j^+|0`$, and $`|0=|0_1|0_2\mathrm{}|0_N`$ stands for the lower state of the atoms. Eq.(34) indicates that there are $`m`$ atoms in the upper state $`|1`$, and the rest of the $`N`$ atoms are in their lower state. With those initial conditions, the probability of the atoms remaining in the initial state is given by $`F_m(t)`$
$$F_m(t)=1\frac{K}{\mathrm{\Gamma }_{K1}}\frac{K^2}{\mathrm{\Gamma }_{K2}}\frac{G}{\mathrm{\Gamma }_{G1}}\frac{G^2}{\mathrm{\Gamma }_{G2}}\frac{GK}{\mathrm{\Gamma }_{GK}}+O(G^3)+O(K^3),$$
(59)
with
$$\frac{1}{\mathrm{\Gamma }_{K1}}=t[2+2(2mN)],$$
$$\frac{1}{\mathrm{\Gamma }_{K2}}=t^2[2+2(2mN)],$$
$$\frac{1}{\mathrm{\Gamma }_{G1}}=t[22(2mN)],$$
$$\frac{1}{\mathrm{\Gamma }_{G2}}=t^2[22(2mN)],$$
$$\frac{1}{\mathrm{\Gamma }_{GK}}=\frac{1}{2}GKt^2[4(2mN)4].$$
Here, we suppose that all qubits are alike, so $`\mathrm{\Omega }_m=\mathrm{\Omega }`$ and $`K_m=K`$ and $`G_m=G`$. $`G`$ and $`K=G+1`$ depend on environment temperature $`T`$ through $`G=1/[exp(\mathrm{}\mathrm{\Omega }/kT)1]`$, which indicate that the probability decrease with the temperature increasing. In fact, the fidelity in the field of quantum information is nothing but an overlap between the initial and final state of the qubits(two-level system). Eq. (35) suggests that the fidelity depends on $`m`$ i.e. the number of the atoms in upper state initially. And to get the maximum of the fidelity, the variable $`m`$ should be taken as small as possible.
To sum up, in this paper, we construct the small-loss rate perturbative expansion for the density operator of an open system. The expansions provide a quite good approach to the exact solution in case the master equation of the system can not be solved exactly. As an interesting application of this expansions, we used it to calculate some average values such as $`\sigma _z`$ and $`\sigma _x`$ in dissipative two-level system, the expansions of the density operator for a single-mode field in a lossy cavity are also presented, and the dynamical property in $`N`$ two-level atom system.
In addition, the other meaningful quantities of the open system such as energy, occupation probability etc. can be expanded in the same spirit of the density operator, so long as the master equation of the system is known.
Figure captions
Fig.1: $`\sigma _x(t)`$ vs time $`t`$, the parameter chosen is $`\mathrm{\Omega }=2`$. The scattering line represent the exact numerical results, whereas the solid line and the dot line show the results from the expansion, $`\gamma `$ in the solid line and the dot line is different, dot line:$`\gamma =0.01`$,solid line:$`\gamma =0.05`$. |
warning/0003/hep-ph0003159.html | ar5iv | text | # 1 Feynman diagrams contributing to 𝑂(𝛼_{𝑒𝑤}𝑚_{𝑡(𝑏)}²/𝑚_𝑊²) Yukawa corrections to 𝑔𝑏→𝑡𝐻⁻: (𝑎) and (𝑏) are tree level diagrams; (𝑐)-(𝑣) are one-loop diagrams. The dashed lines represent 𝐻,ℎ,𝐴,𝐻^±,𝐺⁰ and 𝐺^± for diagrams (𝑐) and (𝑓); 𝐻,ℎ,𝐴 and 𝐺⁰ for diagrams (𝑚),(𝑝),(𝑡) and (𝑢); 𝑡̃,𝑏̃,𝐻,ℎ,𝐴,𝐻^±,𝐺⁰ and 𝐺^± for (𝑖) and (𝑘), where the solid lines represent charginos and neutralinos if the dashed lines represent squarks. For diagrams (𝑑) and (𝑔), the solid lines in the loop represent 𝜒̃⁰ and 𝜒̃⁺ and the dashed lines represent squarks.
1. Introduction
There has been a great deal of interest in the charged Higgs bosons appearing in the two-Higgs-doublet models(THDM), particularly the minimal supersymmetric standard model(MSSM), which predicts the existence of three neutral and two charged Higgs bosons $`h,H,A,`$ and $`H^\pm `$. When the Higgs boson of the Standard Model(SM) has a mass below 130-140 Gev and the h boson of the MSSM is in the decoupling limit (which means that $`H^\pm `$ is too heavy anyway to be possibly produced), the lightest neutral Higgs boson may be difficult to distinguish from the neutral Higgs boson of the standard model(SM). But charged Higgs bosons carry a distinctive signature of the Higgs sector in the MSSM. Therefore, the search for charged Higgs bosons is very important for probing the Higgs sector of the MSSM and, therefore, will be one of the prime objectives of the CERN Large Hadron Collider(LHC). At the LHC the integrated luminosity is expected to reach $`L=100fb^1`$ per year in the second phase. Recently, several studies of charged Higgs boson production at hadron colliders have appeared in the literature. For a relatively light charged Higgs boson, $`m_{H^\pm }<m_tm_b`$, the dominate production processes at the LHC are $`ggt\overline{t}`$ and $`q\overline{q}t\overline{t}`$ followed by the decay sequence $`tbH^+b\tau ^+\nu _\tau `$. For a heavier charged Higgs boson the dominate production process is $`gbtH^{}`$. Previous studies showed that the search for heavy charged Higgs bosons with $`m_{H^\pm }>m_t+m_b`$ at a hadron collider is seriously complicated by QCD backgrounds due to processes such as $`gbt\overline{t}b,g\overline{b}t\overline{t}\overline{b}`$, and $`ggt\overline{t}b\overline{b}`$, as well as others process. However, recent analyses indicate that the decay mode $`H^+\tau ^+\nu `$ provides an excellent signature for a heavy charged Higgs boson in searches at the LHC. The discovery region for $`H^\pm `$ is far greater than had been thought for a large range of the $`(m_{H^\pm },\mathrm{tan}\beta )`$ parameter space, extending beyond $`m_{H^\pm }1TeV`$ and down to at least $`\mathrm{tan}\beta 3`$, and potentially to $`\mathrm{tan}\beta 1.5`$, assuming the latest results for the SM parameters and parton distribution functions as well as using kinematic selection techniques and the tau polarization analysis. Of course, it is just a theoretical analysis and no experimental simulation has been performed to make the statement very reliable so far.
The one-loop radiative corrections to $`H^{}t`$ associated production have not been calculated, although this production process has been studied extensively at tree-level. In this paper we present the calculations of the $`O(\alpha _{ew}m_{t(b)}^2/m_W^2)`$ supersymmetric(SUSY) electroweak corrections to this associated $`H^{}t`$ production process at both the Fermilab Tevatron and the LHC in the MSSM. These corrections arise from the quantum effects which are induced by potentially large Yukawa couplings from the Higgs sector and the chargino-top(bottom)-sbottom(stop) couplings, neutralino- top(bottom)-stop(sbottom) couplings and charged Higgs-stop-sbottom couplings which will contribute at the $`O(\alpha _{ew}m_{t(b)}^4/m_W^4)`$ to the self-energy of the charged Higgs boson. In order to get a reliable estimate this process has to be merged with the related gluon splitting contribution $`ggH^{}t\overline{b}`$. This leads to a suppression by about $`50\%`$ at LO. However, the complete one-loop QCD corrections are probably more important, but not yet available.
2. Calculations
The tree-level amplitude for $`gbtH^{}`$ is
$$M_0=M_0^{(s)}+M_0^{(t)},$$
(1)
where $`M_0^{(s)}`$ and $`M_0^{(t)}`$ represent the amplitudes arising from diagrams in Fig.1$`(a)`$ and Fig.1$`(b)`$, respectively. Explicitly,
$`M_0^{(s)}`$ $`=`$ $`{\displaystyle \frac{igg_s}{\sqrt{2}m_W(\widehat{s}m_b^2)}}\overline{u}(p_t)[2m_t\mathrm{cot}\beta p_b^\mu P_L+2m_b\mathrm{tan}\beta p_b^\mu P_Rm_t\mathrm{cot}\beta \gamma ^\mu \mathit{}P_L`$ (2)
$`m_b\mathrm{tan}\beta \gamma ^\mu \mathit{}P_R]u(p_b)\epsilon _\mu (k)T_{ij}^a,`$
and
$`M_0^{(t)}`$ $`=`$ $`{\displaystyle \frac{igg_s}{\sqrt{2}m_W(\widehat{t}m_t^2)}}\overline{u}(p_t)[2m_t\mathrm{cot}\beta p_t^\mu P_L+2m_b\mathrm{tan}\beta p_t^\mu P_Rm_t\mathrm{cot}\beta \gamma ^\mu \mathit{}P_L`$ (3)
$`m_b\mathrm{tan}\beta \gamma ^\mu \mathit{}P_R]u(p_b)\epsilon _\mu (k)T_{ij}^a,`$
where $`T^a`$ are the $`SU(3)`$ color matrices and $`\widehat{s}`$ and $`\widehat{t}`$ are the subprocess Mandelstam variables defined by
$$\widehat{s}=(p_b+k)^2=(p_t+p_H^{})^2,$$
and
$$\widehat{t}=(p_tk)^2=(p_H^{}p_b)^2.$$
Here the Cabbibo-Kobayashi-Maskawa matrix element $`V_{CKM}[bt]`$ has been taken to be unity.
The SUSY electroweak corrections of order $`O(\alpha _{ew}m_{t(b)}^2/m_W^2)`$ and $`O(\alpha _{ew}m_{t(b)}^4/m_W^4)`$ to the process $`gbH^{}t`$ arise from the Feynman diagrams shown in Figs.1(c)-1(v) and Fig.2. We carried out the calculation in the t’Hooft-Feynman gauge and used dimensional reduction, which preserves supersymmetry, for regularization of the ultraviolet divergences in the virtual loop corrections using the on-mass-shell renormalization scheme, in which the fine-structure constant $`\alpha _{ew}`$ and physical masses are chosen to be the renormalized parameters, and finite parts of the counterterms are fixed by the renormalization conditions. The coupling constant $`g`$ is related to the input parameters $`e,m_W,`$ and $`m_Z`$ by $`g^2=e^2/s_w^2`$ and $`s_w^2=1m_w^2/m_Z^2`$. The parameter $`\beta `$ in the MSSM we are considering must also be renormalized. Following the analysis of ref., this renormalization constant was fixed by the requirement that the on-mass-shell $`H^+\overline{l}\nu _l`$ coupling remain the same form as in Eq.(2) of ref. to all orders of perturbation theory. Taking into account the $`O(\alpha _{ew}m_{t(b)}^2/m_W^2)`$ Yukawa corrections, the renormalized amplitude for the process $`gbtH^{}`$ can be written as
$`M_{ren}`$ $`=`$ $`M_0^{(s)}+M_0^{(t)}+\delta M^{V_1(s)}+\delta M^{V_1(t)}+\delta M^{s(s)}+\delta M^{s(t)}+\delta M^{V_2(s)}`$ (4)
$`+\delta M^{V_2(t)}+\delta M^{b(s)}+\delta M^{b(t)}M_0^{(s)}+M_0^{(t)}+{\displaystyle \underset{l}{}}\delta M^l,`$
where $`\delta M^{V_1(s)},\delta M^{V_1(t)},\delta M^{s(s)},\delta M^{s(t)},\delta M^{V_2(s)},\delta M^{V_2(t)},\delta M^{b(s)}`$, and $`\delta M^{b(t)}`$ represent the corrections to the tree diagrams arising, respectively, from the $`gbb`$ vertex diagram Fig.1(c)-1(d), the $`gtt`$ vertex diagram Fig.1(f)-1(g), the bottom quark self-energy diagram Fig.1(i), the top quark self-energy diagram Fig.1(k), the $`btH^{}`$ vertex diagrams Figs.1(m)-1(n) and Figs.1(p)-1(q), including their corresponding counterterms Fig.1(e), Fig.1(h), Fig.1(j), Fig.1(l), Fig.1(o), and Fig.1(r), and the box diagrams Figs.1$`(s)1(v)`$. $`_l\delta M^l`$ then represents the sum of the contributions to the Yukawa corrections from all the diagrams in Figs.1(c)-1(v). The explicit form of $`\delta M^l`$ can be expressed as
$`\delta M^l`$ $`=`$ $`{\displaystyle \frac{ig^3g_sT_{ij}^a}{4\sqrt{2}\times 16\pi ^2m_W}}C^l\overline{u}(p_t)\{f_1^l\gamma ^\mu P_L+f_2^l\gamma ^\mu P_R+f_3^lp_b^\mu P_L+f_4^lp_b^\mu P_R+f_5^lp_t^\mu P_L`$ (5)
$`+f_6^lp_t^\mu P_R+f_7^l\gamma ^\mu \mathit{}P_L+f_8^l\gamma ^\mu \mathit{}P_R+f_9^lp_b^\mu \mathit{}P_L+f_{10}^lp_b^\mu \mathit{}P_R+f_{11}^lp_t^\mu \mathit{}P_L`$
$`+f_{12}^lp_t^\mu \mathit{}P_R\}u(p_b)\epsilon _\mu (k),`$
where the $`C^l`$ are coefficients that depend on $`\widehat{s},\widehat{t}`$, and the masses, and the $`f_i^l`$ are form factors; both the coefficients $`C^l`$ and the form factors $`f_i^l`$ are given explicitly in Appendix A. The corresponding amplitude squared is
$$\overline{}|M_{ren}|^2=\overline{}|M_0^{(s)}+M_0^{(t)}|^2+2Re\overline{}[(\underset{l}{}\delta M^l)(M_0^{(s)}+M_0^{(t)})^{}],$$
(6)
where
$`\overline{{\displaystyle }}|M_0^{(s)}+M_0^{(t)}|^2`$ $`=`$ $`{\displaystyle \frac{g^2g_s^2}{2N_Cm_W^2}}\{{\displaystyle \frac{1}{(\widehat{s}m_b^2)^2}}[(m_t^2\mathrm{cot}^2\beta +m_b^2\mathrm{tan}^2\beta )(p_bkp_tk`$ (7)
$``$ $`m_b^2p_tk+2p_bkp_bp_tm_b^2p_bp_t)+2m_b^2m_t^2(p_bkm_b^2)]`$
$`+`$ $`{\displaystyle \frac{1}{(\widehat{t}m_t^2)^2}}[(m_t^2\mathrm{cot}^2\beta +m_b^2\mathrm{tan}^2\beta )(p_bkp_tk+m_t^2p_bk`$
$``$ $`m_t^2p_bp_t)+2m_b^2m_t^2(p_tkm_t^2)]+{\displaystyle \frac{1}{(\widehat{s}m_b^2)(\widehat{t}m_t^2)}}`$
$`\times `$ $`[(m_t^2\mathrm{cot}^2\beta +m_b^2\mathrm{tan}^2\beta )(2p_bkp_tk+2p_bkp_bp_t2(p_bp_t)^2`$
$``$ $`m_b^2p_tk+m_t^2p_bk)+2m_b^2m_t^2(p_tkp_bk2p_bp_t)]\},`$
$`\overline{{\displaystyle }}\delta M^l(M_0^{(s)})^{}`$ $`=`$ $`{\displaystyle \frac{g^4g_s^2}{64N_C\times 16\pi ^2m_W^2(\widehat{s}m_b^2)}}C^l{\displaystyle \underset{i=1}{\overset{12}{}}}h_i^{(s)}f_i^l,`$ (8)
and
$`\overline{{\displaystyle }}\delta M^l(M_0^{(t)})^{}`$ $`=`$ $`{\displaystyle \frac{g^4g_s^2}{64N_C\times 16\pi ^2m_W^2(\widehat{t}m_t^2)}}C^l{\displaystyle \underset{i=1}{\overset{12}{}}}h_i^{(t)}f_i^l.`$ (9)
Here the color factor $`N_C=3`$ and $`h_i^{(s)}`$ and $`h_i^{(t)}`$ are scalar functions whose explicit expressions are given in Appendix B.
The cross section for the process $`gbtH^{}`$ is
$$\widehat{\sigma }=_{\widehat{t}_{min}}^{\widehat{t}_{max}}\frac{1}{16\pi \widehat{s}^2}\overline{\mathrm{\Sigma }}|M_{ren}|^2𝑑\widehat{t}$$
(10)
with
$`\widehat{t}_{min}`$ $`=`$ $`{\displaystyle \frac{m_t^2+m_H^{}^2\widehat{s}}{2}}{\displaystyle \frac{1}{2}}\sqrt{(\widehat{s}(m_t+m_H^{})^2)(\widehat{s}(m_tm_H^{})^2)},`$
and
$`\widehat{t}_{max}`$ $`=`$ $`{\displaystyle \frac{m_t^2+m_H^{}^2\widehat{s}}{2}}+{\displaystyle \frac{1}{2}}\sqrt{(\widehat{s}(m_t+m_H^{})^2)(\widehat{s}(m_tm_H^{})^2)}.`$
The total hadronic cross section for $`ppgbtH^{}`$ can be obtained by folding the subprocess cross section $`\widehat{\sigma }`$ with the parton luminosity:
$$\sigma (s)=_{(m_t+m_H^{})/\sqrt{s}}^1𝑑z\frac{dL}{dz}\widehat{\sigma }(gbtH^{}\mathrm{at}\widehat{s}=z^2s).$$
(11)
Here $`\sqrt{s}`$ and $`\sqrt{\widehat{s}}`$ are the CM energies of the $`pp`$ and $`gb`$ states , respectively, and $`dL/dz`$ is the parton luminosity, defined as
$$\frac{dL}{dz}=2z_{z^2}^1\frac{dx}{x}f_{b/P}(x,\mu )f_{g/P}(z^2/x,\mu ),$$
(12)
where $`f_{b/P}(x,\mu )`$ and $`f_{g/P}(z^2/x,\mu )`$ are the bottom quark and gluon parton distribution functions.
3. Numerical results and conclusion
In the following we present some numerical results for charged Higgs boson production in association with a top quark at both the Tevatron and the LHC. In our numerical calculations the SM parameters were taken to be $`m_W=80.41GeV`$, $`m_Z=91.187GeV`$, $`m_t=176GeV`$, $`\alpha _s(m_Z)=0.119`$, and $`\alpha _{ew}(m_Z)=\frac{1}{128.8}`$. And we used the running b quark mass $`3GeV`$ and the one-loop relations from the MSSM between the Higgs boson masses $`m_{h,H,A,H^\pm }`$ and the parameters $`\alpha `$ and $`\beta `$, and chose $`m_{H^\pm }`$ and $`\mathrm{tan}\beta `$ as the two independent input parameters. And we used the CTEQ5M parton distributions throughout the calculations. Other MSSM parameters were determined as follows:
(i) For the parameters $`M_1,M_2`$, and $`\mu `$ in the chargino and neutralino matrix, we put $`M_2=300GeV`$ and then used the relation $`M_1=(5/3)(g^2/g^2)M_20.5M_2`$ to determine $`M_1`$. We also put $`\mu =100GeV`$ except the numerical calculations as shown in Fig.6(b), where $`\mu `$ is a variable.
(ii) For the parameters $`m_{\stackrel{~}{Q},\stackrel{~}{U},\stackrel{~}{D}}^2`$ and $`A_{t,b}`$ in squark mass matrices
$`M_{\stackrel{~}{q}}^2=\left(\begin{array}{cc}M_{LL}^2& m_qM_{LR}\\ m_qM_{RL}& M_{RR}^2\end{array}\right)`$ (15)
with
$`M_{LL}^2=m_{\stackrel{~}{Q}}^2+m_q^2+m_Z^2\mathrm{cos}2\beta (I_q^{3L}e_q\mathrm{sin}^2\theta _W),`$
$`M_{RR}^2=m_{\stackrel{~}{U},\stackrel{~}{D}}^2+m_q^2+m_Z^2\mathrm{cos}2\beta e_q\mathrm{sin}^2\theta _W,`$
$`M_{LR}=M_{RL}=\left(\begin{array}{cc}A_t\mu \mathrm{cot}\beta \hfill & (\stackrel{~}{q}=\stackrel{~}{t})\hfill \\ A_b\mu \mathrm{tan}\beta \hfill & (\stackrel{~}{q}=\stackrel{~}{b})\hfill \end{array}\right),`$ (18)
to simplify the calculation we assumed $`m_{\stackrel{~}{Q}}^2=m_{\stackrel{~}{U}}^2=m_{\stackrel{~}{D}}^2`$ and $`A_t=A_b`$, and we put $`m_{\stackrel{~}{Q}}=500GeV`$ and $`A_t=200GeV`$. But in the numerical calculations of Fig.6(a) $`A_t=A_b`$ are the variables.
Some typical numerical calculations of the tree-level total cross sections and the $`O(\alpha _{ew}m_{t(b)}^2/m_W^2)`$ SUSY electroweak corrections as the functions of the charged Higgs boson mass, $`A_t=A_b`$ and $`\mu `$, respectively, for three representative values of $`\mathrm{tan}\beta `$ are given in Figs.3-6.
Figures 3(a) and 4(a) show that the tree-level total cross sections as a function of the charged Higgs boson mass for three representative values of $`\mathrm{tan}\beta `$. For $`m_{H^\pm }=200GeV`$ the total cross sections at the Tevatron are at most only $`0.7`$ fb and $`0.1`$ fb for $`\mathrm{tan}\beta =2,30`$ and $`10`$, respectively, and for $`m_{H^\pm }=300GeV`$ the total cross sections are smaller than $`0.15`$ fb for all three values of $`\mathrm{tan}\beta `$. However, at the LHC the total cross sections are much larger: the order of thousands of fb for $`m_{H^\pm }`$ in the range $`100`$ to $`240GeV`$ and $`\mathrm{tan}\beta =2`$ and $`30`$; and they are hundreds of fb for the intermediate value $`\mathrm{tan}\beta =10`$. When the charged Higgs boson mass becomes heavy($`<500`$ GeV), the total cross sections still are larger than $`100`$ fb and $`10`$ fb for $`\mathrm{tan}\beta =2,30`$ and $`10`$, respectively. For low $`\mathrm{tan}\beta `$ the top quark contribution is enhanced while for high $`\mathrm{tan}\beta `$ the bottom quark contribution becomes large. These results are smaller than ones given in ref. because we used the running b quark mass $`3GeV`$ in the numerical calculations. We have confirmed that if we chose $`m_b=4.5GeV`$, our results will agree with ref..
In Figs. 3(b) and 4(b) we show the corrections to the total cross sections relative to the tree-level values as a function of $`m_{H^\pm }`$ for $`\mathrm{tan}\beta =2,10,`$ and $`30`$. For $`\mathrm{tan}\beta =2`$ the corrections decrease the total cross sections significantly, which exceed $`13\%`$ for $`m_{H^\pm }`$ below $`300GeV`$ at the both Tevatron and the LHC. But the corrections decrease as $`m_{H^\pm }`$ increase. For example, as shown in Fig.4(b), the corrections range between $`13\%0\%`$ when $`m_{H^\pm }`$ increase from $`300GeV`$ to $`1TeV`$ at the LHC. For high $`\mathrm{tan}\beta (=10,30)`$ these corrections become smaller, which can decrease or increase the total cross sections depending on $`\mathrm{tan}\beta `$, and the magnitude of the corrections are at most a few percent for a wide range of the charged Higgs boson mass at both the Tevatron and the LHC.
In Fig.5 we present the Yukawa correction from the Higgs sector and the genuine SUSY electroweak correction from the couplings involving the genuine SUSY particles(the chargino, neutralino and squark) for $`\mathrm{tan}\beta =30`$ at the LHC, respectively. One can see that the Yukawa correction and the genuine SUSY electroweak correction have opposite signs, and thus cancel to some extent. The former decrease the total cross sections, which can range between $`8\%4\%`$ for $`m_{H^\pm }`$ below $`300GeV`$, but the latter increase the total cross sections, which range between $`10\%7\%`$ for $`m_{H^\pm }`$ in the same range. In such a case the combined effects just are about $`2\%3\%`$.
Figs.6(a) and 6(b) give the corrections as the functions of $`A_t=A_b`$ and $`\mu `$ for $`m_{H^\pm }=300GeV`$ at the LHC, respectively, assuming $`\mathrm{tan}\beta =2,10`$ and $`30`$. From Figs.6(a) and 6(b) one sees that the corrections increase or decrease slowly with increasing $`A_t=A_b`$ and the magnitude of $`\mu `$, respectively, for $`\mathrm{tan}\beta =30,10`$, and the corrections are not very sensitive to both $`A_t=A_b`$ and $`\mu `$ for $`\mathrm{tan}\beta =2`$, where the corrections are always about $`12\%`$ and $`13\%`$, respectively. In general for large values of $`A_t`$ and small values of $`\mathrm{tan}\beta `$ or large values of $`\mu `$ and $`\mathrm{tan}\beta `$, one can get much larger corrections since the charged Higgs boson-stop-sbottom couplings become stronger. For $`\mathrm{tan}\beta =30`$, comparing Fig.4(b) with Fig.6(b), we can see that the corrections indeed become larger as the values of $`\mu `$ increase. But for $`\mathrm{tan}\beta =2`$ from Fig.4(a) and Fig.6(a) we found that the corrections almost have no change when $`A_t=A_b`$ become larger. Obviously the effects from the stronger couplings have been suppressed by the decoupling effects because with an increase of $`A_t=A_b`$ all the squark masses are still heavy, which almost is same as discussed in Ref..
In conclusion, we have calculated the $`O(\alpha _{ew}m_{t(b)}^2/m_W^2)`$ and $`O(\alpha _{ew}m_{t(b)}^4/m_W^4)`$ SUSY electroweak corrections to the cross section for the charged Higgs boson production in association with a top quark at the Tevatron and the LHC. These corrections decrease or increase the cross section depending on $`\mathrm{tan}\beta `$ but are not very sensitive to the mass of the charged Higgs boson for high $`\mathrm{tan}\beta `$. At low $`\mathrm{tan}\beta (=2)`$ the corrections decrease the total cross sections significantly, which exceed $`12\%`$ for $`m_{H^\pm }`$ below $`300GeV`$ at both the Tevatron and the LHC, but for $`m_{H^\pm }>300GeV`$ the corrections can become very small at the LHC. For high $`\mathrm{tan}\beta (=10,30)`$ these corrections can decrease or increase the total cross sections, and the magnitude of the corrections are at most a few percent at both the Tevatron and the LHC.
This work was supported in part by the National Natural Science Foundation of China, the Doctoral Program Foundation of Higher Education of China, the Post Doctoral Foundation of China, a grant from the State Commission of Science and Technology of China, and the U.S.Department of Energy, Division of High Energy Physics, under Grant No.DE-FG02-91-ER4086. S.H. Zhu also gratefully acknowledges the support of the K.C. Wong Education Foundation of Hong Kong.
Appendix A
The coefficients $`C^l`$ and form factors $`f_i^l`$ are the following:
$`C^{V_1(s)}`$ $`=`$ $`{\displaystyle \frac{m_b^2}{m_W^2(\widehat{s}m_b^2)}},C^{V_1(t)}={\displaystyle \frac{m_t^2}{m_W^2(\widehat{t}m_t^2)}},C^{s(s)}={\displaystyle \frac{m_b^2}{m_W^2(\widehat{s}m_b^2)^2}},`$
$`C^{s(t)}`$ $`=`$ $`{\displaystyle \frac{m_t^2}{m_W^2(\widehat{t}m_t^2)^2}},C^{V_2(s)}={\displaystyle \frac{m_bm_t}{m_W^2(\widehat{s}m_b^2)}},C^{V_2(t)}={\displaystyle \frac{m_bm_t}{m_W^2(\widehat{t}m_t^2)}},`$
$`C^{b(s)}`$ $`=`$ $`C^{b(t)}={\displaystyle \frac{m_tm_b}{m_W^2}},`$
$`f_1^{V_1(s)}`$ $`=`$ $`\eta ^{(1)}[m_b(g_2^{V_1(s)}g_3^{V_1(s)})2p_bkg_6^{V_1(s)}],`$
$`f_2^{V_1(s)}`$ $`=`$ $`\eta ^{(2)}[m_b(g_3^{V_1(s)}g_2^{V_1(s)})2p_bkg_7^{V_1(s)}],`$
$`f_3^{V_1(s)}`$ $`=`$ $`\eta ^{(2)}[2(g_1^{V_1(s)}+g_2^{V_1(s)})+m_b(g_4^{V_1(s)}+g_5^{V_1(s)})+2p_bkg_8^{V_1(s)}],`$
$`f_4^{V_1(s)}`$ $`=`$ $`\eta ^{(1)}[2(g_1^{V_1(s)}+g_3^{V_1(s)})+m_b(g_4^{V_1(s)}+g_5^{V_1(s)})+2p_bkg_9^{V_1(s)}],`$
$`f_7^{V_1(s)}`$ $`=`$ $`\eta ^{(2)}[(g_1^{V_1(s)}+g_2^{V_1(s)})+m_b(g_6^{V_1(s)}+g_7^{V_1(s)})],`$
$`f_8^{V_1(s)}`$ $`=`$ $`\eta ^{(1)}[(g_1^{V_1(s)}+g_3^{V_1(s)})+m_b(g_6^{V_1(s)}+g_7^{V_1(s)})],`$
$`f_9^{V_1(s)}`$ $`=`$ $`\eta ^{(1)}[g_4^{V_1(s)}+2g_6^{V_1(s)}+m_b(g_8^{V_1(s)}g_9^{V_1(s)})],`$
$`f_{10}^{V_1(s)}`$ $`=`$ $`\eta ^{(2)}[g_5^{V_1(s)}+2g_7^{V_1(s)}+m_b(g_9^{V_1(s)}g_8^{V_1(s)})],`$
$`f_1^{V_2(s)}`$ $`=`$ $`2p_bkg_3^{V_2(s)},f_2^{V_2(s)}=2p_bkg_4^{V_2(s)},`$
$`f_3^{V_2(s)}`$ $`=`$ $`2g_1^{V_2(s)}+2m_t\mathrm{cot}\beta (\delta \mathrm{\Lambda }_L^{(1)}+\delta \mathrm{\Lambda }_L^{(2)}+\delta \mathrm{\Lambda }_L^{(3)})2m_tg_3^{V_2(s)}+2m_bg_4^{V_2(s)},`$
$`f_4^{V_2(s)}`$ $`=`$ $`2g_2^{V_2(s)}+2m_b\mathrm{tan}\beta (\delta \mathrm{\Lambda }_R^{(1)}+\delta \mathrm{\Lambda }_R^{(2)}+\delta \mathrm{\Lambda }_R^{(3)})+2m_bg_3^{V_2(s)}2m_tg_4^{V_2(s)},`$
$`f_7^{V_2(s)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}f_3^{V_2(s)},f_8^{V_2(s)}={\displaystyle \frac{1}{2}}f_4^{V_2(s)},`$
$`f_1^{V_2(t)}`$ $`=`$ $`2p_tkg_3^{V_2(t)},f_2^{V_2(t)}=2p_tkg_4^{V_2(t)},`$
$`f_5^{V_2(t)}`$ $`=`$ $`2g_1^{V_2(t)}+2m_t\mathrm{cot}\beta (\delta \mathrm{\Lambda }_L^{(1)}+\delta \mathrm{\Lambda }_L^{(2)}+\delta \mathrm{\Lambda }_L^{(3)})2m_tg_3^{V_2(t)}+2m_bg_4^{V_2(t)},`$
$`f_6^{V_2(t)}`$ $`=`$ $`2g_2^{V_2(t)}+2m_b\mathrm{tan}\beta (\delta \mathrm{\Lambda }_R^{(1)}+\delta \mathrm{\Lambda }_R^{(2)}+\delta \mathrm{\Lambda }_R^{(3)})+2m_bg_3^{V_2(t)}2m_tg_4^{V_2(t)},`$
$`f_7^{V_2(t)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}f_5^{V_2(t)},f_8^{V_2(t)}={\displaystyle \frac{1}{2}}f_6^{V_2(t)},`$
$`f_1^{s(s)}`$ $`=`$ $`2\eta ^{(1)}p_bk[g_1^{s(s)}+m_b(g_2^{s(s)}+g_3^{s(s)})],`$
$`f_2^{s(s)}`$ $`=`$ $`2\eta ^{(2)}p_bk[g_5^{s(s)}+m_b(g_2^{s(s)}+g_4^{s(s)})],`$
$`f_3^{s(s)}`$ $`=`$ $`2\eta ^{(2)}[m_b(g_1^{s(s)}+g_5^{s(s)})+2(m_b^2+p_bk)g_2^{s(s)}+(m_b^2+2p_bk)g_3^{s(s)}+m_b^2g_4^{s(s)}],`$
$`f_4^{s(s)}`$ $`=`$ $`2\eta ^{(1)}[m_b(g_1^{s(s)}+g_5^{s(s)})+2(m_b^2+p_bk)g_2^{s(s)}+m_b^2g_3^{s(s)}+(m_b^2+2p_bk)g_4^{s(s)}],`$
$`f_7^{s(s)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}f_3^{s(s)},f_8^{s(s)}={\displaystyle \frac{1}{2}}f_4^{s(s)},`$
$`f_1^{b(s)}`$ $`=`$ $`{\displaystyle \underset{i=H^0,h^0,G^0,A^0}{}}\eta _i^{(3)}\{\eta ^{(2)}[2m_b(3D_{312}+(1\zeta _i)D_{27})+m_b^3(D_0+D_{12}D_{22}`$
$`D_{32})m_t^2m_b(D_{23}+2D_{39})2m_bp_bk(2D_{36}+D_{24}+\zeta _i(D_0+D_{12}))`$
$`+2m_bp_tk(D_{25}+D_{310})+2m_bp_bp_t(D_{26}+2D_{38})]+\eta ^{(1)}[2m_t(3D_{313}+(1`$
$`+\zeta _i)D_{27})m_t^3(D_{33}+(1+\zeta _i)D_{23})+m_b^2m_t(D_{13}2D_{38}+(1+\zeta _i)(D_0`$
$`D_{22}))+2m_tp_bk(D_{13}D_{310}(1+\zeta _i)(D_{12}+D_{24}))+2m_tp_tk(2D_{37}`$
$`+(1+\zeta _i)D_{25})+2m_tp_bp_t(2D_{39}+(1+\zeta _i)D_{26})]\}`$
$`(k,p_b,p_t,m_b,m_b,m_i,m_t)`$
$`{\displaystyle \frac{8\sqrt{2}m_W}{\mathrm{sin}2\beta }}{\displaystyle \underset{i,j,k}{}}N_{k4}N_{k3}^{}R_i(b)R_j(t)\sigma _{ij}D_{27}(k,p_b,p_t,m_{\stackrel{~}{b}_i},m_{\stackrel{~}{b}_i},m_{\stackrel{~}{\chi }_k^0},m_{\stackrel{~}{t}_j}),`$
$`f_2^{b(s)}`$ $`=`$ $`f_1^{b(s)}(\eta ^{(1)}\eta ^{(2)},L_lR_l,N_{kl}N_{kl}^{}),`$
$`f_3^{b(s)}`$ $`=`$ $`{\displaystyle \underset{i=H^0,h^0,G^0,A^0}{}}\eta _i^{(3)}\{\eta ^{(1)}[4D_{27}+2m_b^2(D_{22}D_0(1\zeta _i)(D_{12}+D_{22}))`$
$`+2m_t^2(D_{23}(1+\zeta _i)D_{26})+4p_tk(D_{26}D_{25})]+\eta ^{(2)}2m_tm_b(1+\zeta _i)(D_{22}`$
$`D_{12}D_{26})\}(k,p_b,p_t,m_b,m_b,m_i,m_t)`$
$`{\displaystyle \frac{8\sqrt{2}m_W}{\mathrm{sin}2\beta }}{\displaystyle \underset{i,j,k}{}}\sigma _{ij}[m_tN_{k4}N_{k3}^{}R_i(b)R_j(t)D_{26}+m_bN_{k4}^{}N_{k3}L_i(b)L_j(t)(D_{12}`$
$`+D_{22})+m_{\stackrel{~}{\chi }_k^0}N_{k4}^{}N_{k3}^{}R_i(b)L_j(t)D_{12}](k,p_b,p_t,m_{\stackrel{~}{b}_i},m_{\stackrel{~}{b}_i},m_{\stackrel{~}{\chi }_k^0},m_{\stackrel{~}{t}_j}),`$
$`f_4^{b(s)}`$ $`=`$ $`f_3^{b(s)}(\eta ^{(1)}\eta ^{(2)},L_lR_l,N_{kl}N_{kl}^{}),`$
$`f_5^{b(s)}`$ $`=`$ $`{\displaystyle \underset{i=H^0,h^0,G^0,A^0}{}}\eta _i^{(3)}\{\eta ^{(1)}[12D_{313}+2m_b^2(2D_{38}D_{13}+(1\zeta _i)(D_{13}`$
$`+D_{26}))+2m_t^2(D_{33}+(1+\zeta _i)D_{23})+4p_bk(D_{25}+D_{310})4p_tk(D_{23}`$
$`+2D_{37})4p_tp_b(D_{23}+2D_{39})]+\eta ^{(2)}2m_tm_b(1+\zeta _i)(D_{13}+D_{23}`$
$`D_{26})\}(k,p_b,p_t,m_b,m_b,m_i,m_t)`$
$`+{\displaystyle \frac{8\sqrt{2}m_W}{\mathrm{sin}2\beta }}{\displaystyle \underset{i,j,k}{}}\sigma _{ij}[m_tN_{k4}N_{k3}^{}R_i(b)R_j(t)D_{23}+m_bN_{k4}^{}N_{k3}L_i(b)L_j(t)(D_{13}`$
$`+D_{26})+m_{\stackrel{~}{\chi }_k^0}N_{k4}^{}N_{k3}^{}R_i(b)L_j(t)D_{13}](k,p_b,p_t,m_{\stackrel{~}{b}_i},m_{\stackrel{~}{b}_i},m_{\stackrel{~}{\chi }_k^0},m_{\stackrel{~}{t}_j}),`$
$`f_6^{b(s)}`$ $`=`$ $`f_5^{b(s)}(\eta ^{(1)}\eta ^{(2)},L_lR_l,N_{kl}N_{kl}^{}),`$
$`f_7^{b(s)}`$ $`=`$ $`{\displaystyle \underset{i=H^0,h^0,G^0,A^0}{}}\eta _i^{(3)}\{\eta ^{(1)}[6(D_{27}D_{311})+m_b^2(D_{11}2D_{12}2D_{22}`$
$`2D_{36}+(1+\zeta _i)(D_0+D_{12}))m_t^2(2D_{23}+2D_{37}+(1+\zeta _i)D_{13})2p_bk(D_{12}`$
$`+2D_{24}+2D_{34})+2p_tk(D_{13}+2D_{25}+2D_{35})+2p_tp_b(D_{13}+2D_{26}`$
$`+D_{310})]+\eta ^{(2)}m_tm_b(1+\zeta _i)(D_{12}D_{13}D_0)\}(k,p_b,p_t,m_b,m_b,m_i,m_t),`$
$`f_8^{b(s)}`$ $`=`$ $`f_7^{b(s)}(\eta ^{(1)}\eta ^{(2)}),`$
$`f_9^{b(s)}`$ $`=`$ $`{\displaystyle \underset{i=H^0,h^0,G^0,A^0}{}}\eta _i^{(3)}\{\eta ^{(1)}2m_t[D_{13}D_{26}+(1+\zeta _i)(D_{12}+D_{24})]`$
$`+\eta ^{(2)}2m_b[D_{22}+D_{24}+\zeta _i(D_0+2D_{12}+D_{24})]\}(k,p_b,p_t,m_b,m_b,m_i,m_t)`$
$`{\displaystyle \frac{8\sqrt{2}m_W}{\mathrm{sin}2\beta }}{\displaystyle \underset{i,j,k}{}}\sigma _{ij}N_{k4}N_{k3}^{}R_i(b)R_j(t)(D_{12}`$
$`+D_{24})(k,p_b,p_t,m_{\stackrel{~}{b}_i},m_{\stackrel{~}{b}_i},m_{\stackrel{~}{\chi }_k^0},m_{\stackrel{~}{t}_j}),`$
$`f_{10}^{b(s)}`$ $`=`$ $`f_9^{b(s)}(\eta ^{(1)}\eta ^{(2)},L_lR_l,N_{kl}N_{kl}^{}),`$
$`f_{11}^{b(s)}`$ $`=`$ $`{\displaystyle \underset{i=H^0,h^0,G^0,A^0}{}}\eta _i^{(3)}\{\eta ^{(1)}2m_t[D_{23}(1+\zeta _i)D_{25}]\eta ^{(2)}2m_b[D_{26}+D_{25}`$
$`+\zeta _i(D_{13}+D_{25})]\}(k,p_b,p_t,m_b,m_b,m_i,m_t)`$
$`+{\displaystyle \frac{8\sqrt{2}m_W}{\mathrm{sin}2\beta }}{\displaystyle \underset{i,j,k}{}}\sigma _{ij}N_{k4}N_{k3}^{}R_i(b)R_j(t)(D_{13}`$
$`+D_{25})(k,p_b,p_t,m_{\stackrel{~}{b}_i},m_{\stackrel{~}{b}_i},m_{\stackrel{~}{\chi }_k^0},m_{\stackrel{~}{t}_j}),`$
$`f_{12}^{b(s)}`$ $`=`$ $`f_{11}^{b(s)}(\eta ^{(1)}\eta ^{(2)},L_lR_l,N_{kl}N_{kl}^{}),`$
where $`D_0,D_{ij},D_{ijk}`$ are the four-point Feynman integrals . The explicit forms of $`\delta M^{V_1(t)},\delta M^{s(t)},\delta M^{b(t)}`$ can be respectively obtained from $`\delta M^{V_1(s)},\delta M^{s(s)},\delta M^{b(s)}`$ by the transformation $`U`$ defined as
$$p_bp_t,\widehat{s}\widehat{t},kk,\xi _i^{(1)}\xi _i^{(2)},\xi _i^{(3)}\xi _i^{(4)},\eta _i^{(1)}\eta _i^{(2)},$$
$$m_tm_b,\eta ^{(1)}\eta ^{(2)},\lambda _b\lambda _t,m_{\stackrel{~}{t}_i}m_{\stackrel{~}{b}_i},U_{i2}V_{i2}^{},N_{i3}N_{i4}^{},$$
$$L_i(b)L_i(t),R_i(b)R_i(t),p_b^\mu P_{L(R)}p_t^\mu P_{R(L)},\gamma ^\mu \mathit{}P_L\gamma ^\mu \mathit{}P_R.$$
All other form factors $`f_i^l`$ not listed above vanish. In the above expressions we have used the following definitions:
$`\eta ^{(1)}=m_b\mathrm{tan}\beta ,\eta ^{(2)}=m_t\mathrm{cot}\beta ,`$ $`\lambda _b={\displaystyle \frac{m_b}{\sqrt{2}m_W\mathrm{cos}\beta }},\lambda _t={\displaystyle \frac{m_t}{\sqrt{2}m_W\mathrm{sin}\beta }}`$
$`L_1(q)=\mathrm{cos}\theta _q,L_2(q)=\mathrm{sin}\theta _q,`$ $`R_1(q)=\mathrm{sin}\theta _q,R_2(q)=\mathrm{cos}\theta _q,`$
$`\eta _{H^0}^{(1)}={\displaystyle \frac{\mathrm{cos}^2\alpha }{\mathrm{cos}^2\beta }},\eta _{h^0}^{(1)}={\displaystyle \frac{\mathrm{sin}^2\alpha }{\mathrm{cos}^2\beta }},`$ $`\eta _{A^0}^{(1)}=\mathrm{tan}^2\beta ,\eta _{G^0}^{(1)}=1,`$
$`\eta _{H^0}^{(2)}={\displaystyle \frac{\mathrm{sin}^2\alpha }{\mathrm{sin}^2\beta }},\eta _{h^0}^{(2)}={\displaystyle \frac{\mathrm{cos}^2\alpha }{\mathrm{sin}^2\beta }},`$ $`\eta _{A^0}^{(2)}=\mathrm{cot}^2\beta ,\eta _{G^0}^{(2)}=1,`$
$`\eta _{H^0}^{(3)}=\eta _{h^0}^{(3)}={\displaystyle \frac{\mathrm{sin}\alpha \mathrm{cos}\alpha }{\mathrm{sin}\beta \mathrm{cos}\beta }},`$ $`\eta _{G^0}^{(3)}=\eta _{A^0}^{(3)}=1,`$
$`\xi _H^{}^{(1)}={\displaystyle \frac{m_t^2}{m_b^2}}\mathrm{cot}^2\beta ,\xi _G^{}^{(1)}={\displaystyle \frac{m_t^2}{m_b^2}},`$ $`\xi _H^{}^{(2)}={\displaystyle \frac{m_b^2}{m_t^2}}\mathrm{tan}^2\beta ,\xi _G^{}^{(2)}={\displaystyle \frac{m_b^2}{m_t^2}},`$
$`\xi _H^{}^{(3)}=\mathrm{tan}^2\beta ,\xi _G^{}^{(3)}=1,`$ $`\xi _H^{}^{(4)}=\mathrm{cot}^2\beta ,\xi _G^{}^{(4)}=1,`$
$$\zeta _{H^0}=\zeta _{h^0}=\zeta _H^{}=\zeta _{A^0}=\zeta _{G^0}=\zeta _G^{}=1,$$
$`\sigma _{ij}`$ $`=`$ $`{\displaystyle \frac{m_W}{\sqrt{2}}}(\mathrm{sin}2\beta {\displaystyle \frac{m_b^2\mathrm{tan}\beta +m_t^2\mathrm{cot}\beta }{m_W^2}})L_i(b)L_j(t)`$
$`+{\displaystyle \frac{m_tm_b}{\sqrt{2}m_W}}(\mathrm{tan}\beta +\mathrm{cot}\beta )R_i(b)R_j(t){\displaystyle \frac{m_b}{\sqrt{2}m_W}}(\mu A_b\mathrm{tan}\beta )R_i(b)L_j(t)`$
$`{\displaystyle \frac{m_t}{\sqrt{2}m_W}}(\mu A_t\mathrm{cot}\beta )L_i(b)R_j(t),`$
$`g_1^{V_1(s)}`$ $`=`$ $`{\displaystyle \underset{i=H^0,h^0,G^0,A^0}{}}\eta _i^{(1)}\{[{\displaystyle \frac{1}{2}}2\overline{C}_{24}+m_b^2(2C_{11}+C_{12}C_{21}+C_{23})\widehat{s}(C_{12}`$
$`+C_{23})](p_b,k,m_i,m_b,m_b)+[F_0+F_1+2m_b^2G_1`$
$`(1+\zeta _i)2m_b^2G_0](m_b^2,m_i,m_b)\},`$
$`g_2^{V_1(s)}`$ $`=`$ $`{\displaystyle \underset{i=H^{},G^{}}{}}2\{\xi _i^{(1)}[{\displaystyle \frac{1}{2}}2\overline{C}_{24}+m_t^2C_0+m_b^2(C_02C_{11}+C_{12}C_{21}+C_{23})`$
$`\widehat{s}(C_{12}+C_{23})](p_b,k,m_i,m_t,m_t)+[\xi _i^{(1)}(F_0+F_1)2m_t^2\zeta _iG_0`$
$`+m_b^2(\xi _i^{(1)}+\xi _i^{(3)})(G_1\zeta _iG_0)](m_b^2,m_i,m_t)\}`$
$`+{\displaystyle \frac{4m_W^2}{m_b^2}}{\displaystyle \underset{i,j}{}}\{\lambda _b^2[R_j^2(b)|N_{i3}|^2(F_0+F_1)+m_b^2|N_{i3}|^2(G_0+G_1)2m_bm_{\stackrel{~}{\chi }_i^0}`$
$`\times L_j(b)R_j(b)N_{i3}^2G_0](m_b^2,m_{\stackrel{~}{b}_j},m_{\stackrel{~}{\chi }_i^0})+[2m_bm_{\stackrel{~}{\chi }_i^+}\lambda _b\lambda _tL_j(t)R_j(t)V_{i2}^2U_{i2}^2G_0`$
$`+\lambda _t^2R_j^2(t)|V_{i2}|^2(F_0+F_1)+m_b^2(\lambda _t^2R_j^2(t)|V_{i2}|^2+\lambda _b^2L_j^2(t)|U_{i2}|^2)(G_0`$
$`+G_1)](m_b^2,m_{\stackrel{~}{t}_j},m_{\stackrel{~}{\chi }_i^+})2\lambda _b^2R_j^2(b)|N_{i3}|^2\overline{C}_{24}(p_b,k,m_{\stackrel{~}{\chi }_i^0},m_{\stackrel{~}{b}_j},m_{\stackrel{~}{b}_j})`$
$`2\lambda _t^2R_j^2(t)|V_{i2}|^2\overline{C}_{24}(p_b,k,m_{\stackrel{~}{\chi }_i^+},m_{\stackrel{~}{t}_j},m_{\stackrel{~}{t}_j})\},`$
$`g_3^{V_1(s)}`$ $`=`$ $`g_2^{V_1(s)}(\xi _i^{(1)}\xi _i^{(3)},V_{i2}U_{i2}^{},N_{i3}N_{i3}^{},L_j(b)R_j(b),\lambda _bL_j(t)\lambda _tR_j(t)),`$
$`g_4^{V_1(s)}`$ $`=`$ $`{\displaystyle \underset{i=H^0,h^0,G^0,A^0}{}}\eta _i^{(1)}2m_b[C_0+2C_{11}+C_{21}+\zeta _i(C_0+C_{11})](p_b,k,m_i,m_b,m_b)`$
$`+{\displaystyle \underset{i=H^{},G^{}}{}}4m_b[\xi _i^{(3)}(C_0+2C_{11}+C_{21})+{\displaystyle \frac{m_t^2}{m_b^2}}\zeta _i(C_0+C_{11})](p_b,k,m_i,m_t,m_t)`$
$`+{\displaystyle \frac{8m_W^2}{m_b^2}}{\displaystyle \underset{i,j}{}}\{\lambda _b^2[m_{\stackrel{~}{\chi }_i^0}L_j(b)R_j(b)N_{i3}^2(C_0+C_{11})m_bL_j^2(b)|N_{i3}|^2(C_{11}`$
$`+C_{21})](p_b,k,m_{\stackrel{~}{\chi }_i^0},m_{\stackrel{~}{b}_j},m_{\stackrel{~}{b}_j})`$
$`+[m_{\stackrel{~}{\chi }_i^+}\lambda _b\lambda _tL_j(t)R_j(t)V_{i2}^{}U_{i2}^{}(C_0+C_{11})m_b\lambda _b^2L_j^2(t)|U_{i2}|^2(C_{11}`$
$`+C_{21})](p_b,k,m_{\stackrel{~}{\chi }_i^+},m_{\stackrel{~}{t}_j},m_{\stackrel{~}{t}_j})\},`$
$`g_5^{V_1(s)}`$ $`=`$ $`g_4^{V_1(s)}(\xi _i^{(1)}\xi _i^{(3)},V_{i2}U_{i2}^{},N_{i3}N_{i3}^{},L_j(b)R_j(b),\lambda _bL_j(t)\lambda _tR_j(t)),`$
$`g_6^{V_1(s)}`$ $`=`$ $`{\displaystyle \underset{i=H^0,h^0,G^0,A^0}{}}\eta _i^{(1)}m_b(C_0+C_{11}+\zeta _iC_0)(p_b,k,m_i,m_b,m_b)`$
$`{\displaystyle \underset{i=H^{},G^{}}{}}2m_b[\xi _i^{(3)}(C_0+C_{11})+{\displaystyle \frac{m_t^2}{m_b^2}}\zeta _iC_0](p_b,k,m_i,m_t,m_t),`$
$`g_7^{V_1(s)}`$ $`=`$ $`g_6^{V_1(s)}(\xi _i^{(1)}\xi _i^{(3)}),`$
$`g_8^{V_1(s)}`$ $`=`$ $`{\displaystyle \underset{i=H^0,h^0,G^0,A^0}{}}2\eta _i^{(1)}(C_{12}+C_{23})(p_b,k,m_i,m_b,m_b)`$
$`+{\displaystyle \underset{i=H^{},G^{}}{}}4\xi _i^{(1)}(C_{12}+C_{24})(p_b,k,m_i,m_t,m_t)`$
$`{\displaystyle \frac{8m_W^2}{m_b^2}}{\displaystyle \underset{i,j}{}}\{\lambda _b^2R_j^2(b)|N_{i3}|^2(C_{12}+C_{23})(p_b,k,m_{\stackrel{~}{\chi }_i^0},m_{\stackrel{~}{b}_j},m_{\stackrel{~}{b}_j})`$
$`+\lambda _t^2R_j^2(t)|V_{i2}|^2(C_{12}+C_{23})(p_b,k,m_{\stackrel{~}{\chi }_i^+},m_{\stackrel{~}{t}_j},m_{\stackrel{~}{t}_j})\},`$
$`g_9^{V_1(s)}`$ $`=`$ $`g_8^{V_1(s)}(\xi _i^{(1)}\xi _i^{(3)},V_{i2}U_{i2}^{},N_{i3}N_{i3}^{},L_j(b)R_j(b),\lambda _bL_j(t)\lambda _tR_j(t)),`$
$`g_1^{V_2(s)}`$ $`=`$ $`{\displaystyle \underset{i=H^0,h^0,G^0,A^0}{}}\eta _i^{(3)}\{\eta ^{(1)}[{\displaystyle \frac{1}{2}}+4\overline{C}_{24}+m_t^2(C_0+2C_{11}+\zeta _i(C_0+C_{11})`$
$`+C_{21}C_{12}C_{23})+m_H^{}^2(C_{22}C_{23})+\widehat{s}(C_{12}+C_{23})]+\eta ^{(2)}m_bm_t[\zeta _iC_{11}`$
$`+(1+\zeta _i)C_0]\}(p_t,p_H^{},m_i,m_t,m_b)+{\displaystyle \frac{4\sqrt{2}m_W}{\mathrm{sin}2\beta }}{\displaystyle }_{i,j,k}[m_tR_i(b)R_j(t)N_{k3}^{}N_{k4}`$
$`\times (C_{11}+C_{12})+m_{\stackrel{~}{\chi }_k^0}L_j(t)R_i(b)N_{k3}^{}N_{k4}^{}C_0]\sigma _{ij}(p_t,p_H^{},m_{\stackrel{~}{\chi }_k^0},m_{\stackrel{~}{b}_i},m_{\stackrel{~}{t}_j}),`$
$`g_2^{V_2(s)}`$ $`=`$ $`g_1^{V_2(s)}(\eta ^{(1)}\eta ^{(2)},L_lR_l,N_{kl}N_{kl}^{}),`$
$`g_3^{V_2(s)}`$ $`=`$ $`{\displaystyle \underset{i=H^0,h^0,G^0,A^0}{}}\eta _i^{(3)}\{\eta ^{(1)}m_t[C_0+C_{11}+\zeta _i(C_0+C_{12})]+\eta ^{(2)}\zeta _im_bC_{12}\}`$
$`(p_t,p_H^{},m_i,m_t,m_b)`$
$`{\displaystyle \frac{4\sqrt{2}m_W}{\mathrm{sin}2\beta }}{\displaystyle \underset{i,j,k}{}}R_i(b)R_j(t)N_{k3}^{}N_{k4}\sigma _{ij}C_{12}(p_t,p_H^{},m_{\stackrel{~}{\chi }_k^0},m_{\stackrel{~}{b}_i},m_{\stackrel{~}{t}_j}),`$
$`g_4^{V_2(s)}`$ $`=`$ $`g_3^{V_2(s)}(\eta ^{(1)}\eta ^{(2)},L_lR_l,N_{kl}N_{kl}^{}),`$
$`g_1^{V_2(t)}`$ $`=`$ $`{\displaystyle \underset{i=H^0,h^0,G^0,A^0}{}}\eta _i^{(3)}\{\eta ^{(1)}[{\displaystyle \frac{1}{2}}+4\overline{C}_{24}+m_b^2(C_0+2C_{11}+\zeta _i(C_0+C_{11})`$
$`+C_{21}C_{12}C_{23})+m_H^{}^2(C_{22}C_{23})+\widehat{t}(C_{12}+C_{23})]`$
$`+\eta ^{(2)}m_bm_t[C_0+\zeta _i(C_0+C_{11})]\}(p_b,p_H^{},m_i,m_b,m_t)`$
$`+{\displaystyle \frac{4\sqrt{2}m_W}{\mathrm{sin}2\beta }}{\displaystyle \underset{i,j,k}{}}[m_bL_i(b)L_j(t)N_{k3}^{}N_{k4}(C_{11}+C_{12})`$
$`+m_{\stackrel{~}{\chi }_k^0}L_j(t)R_i(b)N_{k3}^{}N_{k4}^{}C_0]\sigma _{ij}(p_b,p_H^{},m_{\stackrel{~}{\chi }_k^0},m_{\stackrel{~}{b}_i},m_{\stackrel{~}{t}_j}),`$
$`g_2^{V_2(t)}`$ $`=`$ $`g_1^{V_2(t)}(\eta ^{(1)}\eta ^{(2)},L_lR_l,N_{kl}N_{kl}^{}),`$
$`g_3^{V_2(t)}`$ $`=`$ $`{\displaystyle \underset{i=H^0,h^0,G^0,A^0}{}}\eta _i^{(3)}\{\eta ^{(1)}m_b[C_0+C_{11}+\zeta _i(C_0+C_{12})]+\eta ^{(2)}\zeta _im_tC_{12}\}`$
$`(p_b,p_H^{},m_i,m_b,m_t)`$
$`+{\displaystyle \frac{4\sqrt{2}m_W}{\mathrm{sin}2\beta }}{\displaystyle \underset{i,j,k}{}}R_i(b)R_j(t)N_{k3}^{}N_{k4}\sigma _{ij}C_{12}(p_b,p_H^{},m_{\stackrel{~}{\chi }_k^0},m_{\stackrel{~}{b}_i},m_{\stackrel{~}{t}_j}),`$
$`g_4^{V_2(t)}`$ $`=`$ $`g_3^{V_2(t)}(\eta ^{(1)}\eta ^{(2)},L_lR_l,N_{kl}N_{kl}^{}),`$
$`g_1^{s(s)}`$ $`=`$ $`{\displaystyle \underset{i=H^0,h^0,G^0,A^0}{}}m_b\eta _i^{(1)}\{\zeta _iF_0(p_b+k,m_i,m_b)+[\zeta _iF_02m_b^2(1+\zeta _i)G_0`$
$`+2m_b^2G_1](m_b^2,m_i,m_b)\}+{\displaystyle }_{i=H^{},G^{}}2m_b\{{\displaystyle \frac{m_t^2}{m_b^2}}\zeta _iF_0(p_b+k,m_i,m_t)`$
$`+[2m_t^2\zeta _iG_0+m_b^2(\xi _i^{(1)}+\xi _i^{(3)})(G_1\zeta _iG_0)+\zeta _i{\displaystyle \frac{m_t^2}{m_b^2}}F_0](m_b^2,m_i,m_t)\}`$
$`+{\displaystyle \frac{4m_W^2}{m_b^2}}{\displaystyle \underset{i,j}{}}\{m_{\stackrel{~}{\chi }_i^0}\lambda _b^2L_j(b)R_j(b)N_{i3}^2F_0(p_b+k,m_{\stackrel{~}{b}_j},m_{\stackrel{~}{\chi }_i^0})`$
$`m_{\stackrel{~}{\chi }_i^+}\lambda _b\lambda _tL_j(b)R_j(b)V_{i2}^{}U_{i2}^{}F_0(p_b+k,m_{\stackrel{~}{t}_j},m_{\stackrel{~}{\chi }_i^+})+[m_b^3\lambda _b^2|N_{i3}|^2(G_0+G_1)`$
$`m_{\stackrel{~}{\chi }_i^0}\lambda _b^2L_j(b)R_j(b)N_{i3}^2(2m_b^2G_0F_0)](m_b^2,m_{\stackrel{~}{b}_j},m_{\stackrel{~}{\chi }_i^0})+[m_b^3(\lambda _b^2L_j^2(t)|U_{i2}|^2`$
$`+\lambda _t^2R_j^2(t)|V_{i2}|^2)(G_0+G_1)m_{\stackrel{~}{\chi }_i^+}\lambda _b\lambda _tL_j(t)R_j(t)V_{i2}^{}U_{i2}^{}(2m_b^2G_0`$
$`F_0)](m_b^2,m_{\stackrel{~}{t}_j},m_{\stackrel{~}{\chi }_i^+})\},`$
$`g_2^{s(s)}`$ $`=`$ $`{\displaystyle \underset{i=H^0,h^0,G^0,A^0}{}}\eta _i^{(1)}(F_0+F_1)(p_b+k,m_i,m_b),`$
$`g_3^{s(s)}`$ $`=`$ $`{\displaystyle \underset{i=H^0,h^0,G^0,A^0}{}}\eta _i^{(1)}[F_0F_12m_b^2G_1+2(1+\zeta _i)m_b^2G_0](m_b^2,m_i,m_b)`$
$`+{\displaystyle \underset{i=H^{},G^{}}{}}2\{\xi _i^{(1)}(F_0+F_1)(p_b+k,m_i,m_t)[\xi _i^{(1)}(F_0+F_1)`$
$`2\zeta _im_t^2G_0+m_b^2(\xi _i^{(1)}+\xi _i^{(3)})(G_1\zeta _iG_0)](m_b^2,m_i,m_t)\}`$
$`{\displaystyle \frac{4m_W^2}{m_b^2}}{\displaystyle \underset{i,j}{}}\{\lambda _b^2[R_j^2(b)|N_{i3}|^2(F_0+F_1)+|N_{i3}|^2m_b^2(G_0+G_1)`$
$`2m_bm_{\stackrel{~}{\chi }_i^0}L_j(b)R_j(b)N_{i3}^2G_0](m_b^2,m_{\stackrel{~}{b}_j},m_{\stackrel{~}{\chi }_i^0})`$
$`+[\lambda _t^2R_j^2(t)|V_{i2}|^2(F_0+F_1)+m_b^2(\lambda _t^2R_j^2(t)|V_{i2}|^2+\lambda _b^2L_j^2(t)|U_{i2}|^2)(G_1G_0)`$
$`2m_bm_{\stackrel{~}{\chi }_i^+}L_j(t)R_j(t)\lambda _b\lambda _tV_{i2}^{}U_{i2}^{}G_0](m_b^2,m_{\stackrel{~}{t}_j},m_{\stackrel{~}{\chi }_i^+})`$
$`\lambda _b^2R_j^2(b)|N_{i3}|^2(F_0+F_1)(p_b+k,m_{\stackrel{~}{b}_j},m_{\stackrel{~}{\chi }_i^0})`$
$`\lambda _t^2R_j^2(t)|V_{i2}|^2(F_0+F_1)(p_b+k,m_{\stackrel{~}{t}_j},m_{\stackrel{~}{\chi }_i^+})\},`$
$`g_4^{s(s)}`$ $`=`$ $`g_3^{s(s)}(\xi _i^{(1)}\xi _i^{(3)},V_{i2}U_{i2}^{},N_{i3}N_{i3}^{},L_j(b)R_j(b),\lambda _bL_j(t)\lambda _tR_j(t)),`$
$`g_5^{s(s)}`$ $`=`$ $`g_1^{s(s)}(N_{i3}^{}N_{i3},V_{i2}^{}V_{i2},U_{i2}^{}U_{i2}),`$
$`\delta \mathrm{\Lambda }_L^{(1)}`$ $`=`$ $`{\displaystyle \frac{4N_c}{3m_W^2}}(1\mathrm{cot}^2\theta _W)[2m_t^2(\mathrm{ln}{\displaystyle \frac{m_t^2}{\mu ^2}}1)+m_b^2+m_t^2{\displaystyle \frac{5}{6}}m_W^2+m_b^2F_0`$
$`+(m_b^2m_t^22m_W^2)F_1](m_W^2,m_b,m_t)+{\displaystyle \frac{4N_c}{3m_W^2}}\mathrm{cot}^2\theta _W\{{\displaystyle \frac{5}{6}}[(g_V^b)^2+(g_A^b)^2`$
$`+(g_V^t)^2+(g_A^t)^2]m_Z^2+[((g_V^t)^2+(g_A^t)^2)(2m_t^2\mathrm{ln}{\displaystyle \frac{m_t^2}{\mu ^2}}+m_t^2F_02m_Z^2F_1)`$
$`((g_V^t)^2(g_A^t)^2)3m_t^2F_0](m_Z^2,m_t,m_t)+[((g_V^b)^2+(g_A^b)^2)(2m_b^2\mathrm{ln}{\displaystyle \frac{m_b^2}{\mu ^2}}`$
$`+m_b^2F_02m_Z^2F_1)((g_V^b)^2(g_A^b)^2)3m_b^2F_0](m_Z^2,m_b,m_b)\}+{\displaystyle \frac{4N_c}{m_W^2}}[(\mathrm{cot}^2\beta `$
$`1)m_t^2F_0+(m_t^2m_b^22m_t^2\mathrm{cot}^2\beta )F_1+(m_t^2\mathrm{cot}^2\beta +m_b^2\mathrm{tan}^2\beta `$
$`+2m_b^2)m_t^2G_0(m_t^2\mathrm{cot}^2\beta +m_b^2\mathrm{tan}^2\beta )m_H^{}^2G_1](m_H^{}^2,m_t,m_b)`$
$`+{\displaystyle \underset{i=H^0,h^0,G^0,A^0}{}}{\displaystyle \frac{1}{2m_W^2}}\{m_b^2\eta _i^{(1)}[F_1F_02m_b^2(G_0+\zeta _iG_0G_1)](m_b^2,m_i,m_b)`$
$`m_t^2\eta _i^{(2)}[(1+2\zeta _i)F_0+F_1+2m_t^2(1+\zeta _i)G_02m_t^2G_1](m_t^2,m_i,m_t)\}`$
$`+{\displaystyle \underset{i=H^{},G^{}}{}}{\displaystyle \frac{1}{m_W^2}}\{m_b^2[\xi _i^{(1)}(F_0+F_1)2m_t^2\zeta _iG_0+m_b^2(\xi _i^{(1)}+\xi _i^{(3)})`$
$`\times (G_1\zeta _iG_0)](m_b^2,m_i,m_t)m_t^2[{\displaystyle \frac{2m_b^2}{m_t^2}}\zeta _iF_0+\xi _i^{(2)}(F_0+F_1)+2m_b^2\zeta _iG_0`$
$`m_t^2(\xi _i^{(2)}+\xi _i^{(4)})(G_1\zeta _iG_0)](m_t^2,m_i,m_b)\}2N_C{\displaystyle }_{i,j}\{2\sigma _{ij}\sigma _{ij}G_0`$
$`+{\displaystyle \frac{1}{m_W^2}}L_i(b)L_j(t)[L_i(b)L_j(t)({\displaystyle \frac{m_b^2}{\mathrm{cos}^2\beta }}+{\displaystyle \frac{m_t^2}{\mathrm{sin}^2\beta }})\mathrm{cos}2\beta `$
$`+R_i(b)R_j(t)m_tm_b(\mathrm{tan}^2\beta \mathrm{cot}^2\beta )]\}(m_H^{}^2,m_{\stackrel{~}{t}_j},m_{\stackrel{~}{b}_i}),`$
$`\delta \mathrm{\Lambda }_L^{(2)}`$ $`=`$ $`2{\displaystyle \underset{i,j}{}}\{\lambda _t^2[{\displaystyle \frac{2m_{\stackrel{~}{\chi }_i^0}}{m_t}}L_j(t)R_j(t)N_{i4}^2(F_0m_t^2G_0)+|N_{i4}|^2(R_j^2(t)(F_0+F_1)`$
$`m_t^2(G_0+G_1))](m_t^2,m_{\stackrel{~}{t}_j},m_{\stackrel{~}{\chi }_i^0})+[{\displaystyle \frac{2m_{\stackrel{~}{\chi }_i^+}}{m_t}}\lambda _b\lambda _tL_j(b)R_j(b)U_{i2}^{}V_{i2}^{}(F_0`$
$`m_t^2G_0)+\lambda _b^2R_j^2(b)|U_{i2}|^2(F_0+F_1)m_t^2(\lambda _t^2L_j^2(b)|V_{i2}|^2+\lambda _b^2R_j^2(b)|U_{i2}|^2)`$
$`\times (G_0+G_1)](m_t^2,m_{\stackrel{~}{b}_j},m_{\stackrel{~}{\chi }_i^+})\},`$
$`\delta \mathrm{\Lambda }_L^{(3)}`$ $`=`$ $`2{\displaystyle \underset{i,j}{}}\{\lambda _b^2[|Ni3|^2(R_j^2(b)(F_0+F_1)+m_b^2(G_0+G_1))2m_bm_{\stackrel{~}{\chi }_i^0}L_j(b)R_j(b)`$
$`\times N_{i3}^2G_0](m_b^2,m_{\stackrel{~}{b}_j},m_{\stackrel{~}{\chi }_i^0})+[2m_{\stackrel{~}{\chi }_i^+}m_b\lambda _b\lambda _tL_j(t)R_j(t)U_{i2}^{}V_{i2}^{}G_0`$
$`+\lambda _b^2L_j^2(b)|U_{i2}|^2(F_0+F_1)+m_b^2(\lambda _t^2R_j^2(t)|V_{i2}|^2+\lambda _b^2L_j^2(t)|U_{i2}|^2)`$
$`\times (G_0+G_1)](m_b^2,m_{\stackrel{~}{t}_j},m_{\stackrel{~}{\chi }_i^+})\},`$
$`\delta \mathrm{\Lambda }_R^{(1)}`$ $`=`$ $`{\displaystyle \underset{i=H^0,h^0,G^0,A^0}{}}{\displaystyle \frac{1}{2m_W^2}}\{m_t^2\eta _i^{(2)}[F_0+F_12m_t^2(G_0+\zeta _iG_0G_1)](m_t^2,m_i,m_t)`$
$`m_b^2\eta _i^{(1)}[F_0+F_12\zeta _iF_0+2m_b^2(1+\zeta _i)G_02m_b^2G_1](m_b^2,m_i,m_b)\}`$
$`+{\displaystyle \underset{i=H^{},G^{}}{}}{\displaystyle \frac{1}{m_W^2}}\{m_t^2[\xi _i^{(2)}(F_0+F_1)2m_b^2\zeta _iG_0+m_t^2(\xi _i^{(2)}+\xi _i^{(4)})(G_1`$
$`\zeta _iG_0)](m_t^2,m_i,m_b)m_b^2[{\displaystyle \frac{2m_t^2}{m_b^2}}\zeta _iF_0+\xi _i^{(1)}(F_0+F_1)+2m_t^2\zeta _iG_0`$
$`m_b^2(\xi _i^{(1)}+\xi _i^{(3)})(G_1\zeta _iG_0)](m_b^2,m_i,m_t)\},`$
$`\delta \mathrm{\Lambda }_R^{(2)}`$ $`=`$ $`\delta \mathrm{\Lambda }_L^{(2)}(U),\delta \mathrm{\Lambda }_R^{(3)}=\delta \mathrm{\Lambda }_L^{(3)}(U).`$
Here $`C_0,C_{ij}`$ are the three-point Feynman integrals and $`\overline{C}_{24}\frac{1}{4}\mathrm{\Delta }+C_{24}`$, while
$`F_n(q,m_1,m_2)`$ $`=`$ $`{\displaystyle _0^1}𝑑yy^n\mathrm{ln}[{\displaystyle \frac{q^2y(1y)+m_1^2(1y)+m_2^2y}{\mu ^2}}],`$
$`G_n(q,m_1,m_2)`$ $`=`$ $`{\displaystyle _0^1}𝑑y{\displaystyle \frac{y^{n+1}(1y)}{q^2y(1y)+m_1^2(1y)+m_2^2y}},`$
and
$`g_V^t={\displaystyle \frac{1}{2}}{\displaystyle \frac{4}{3}}\mathrm{sin}^2\theta _W,g_A^t={\displaystyle \frac{1}{2}},g_V^b={\displaystyle \frac{1}{2}}+{\displaystyle \frac{2}{3}}\mathrm{sin}^2\theta _W,g_A^b={\displaystyle \frac{1}{2}},`$
which are the SM couplings of the top and bottom quarks to the Z boson. The definitions of $`\theta _q,U_{ij},V_{ij},N_{ij},\mu ,A_q`$ can be found in ref..
Appendix B
$`h_1^{(i)}`$ $`=`$ $`4m_t\eta ^{(2)}(2p_bkp^{(i)}p_b)4m_b\eta ^{(1)}(p^{(i)}p_t+p_tk),`$
$`h_2^{(i)}`$ $`=`$ $`h_1^{(i)}(\eta ^{(1)}\eta ^{(2)}),`$
$`h_3^{(i)}`$ $`=`$ $`2\eta ^{(2)}(2p_bkp_bp_tm_b^2p_tk2p^{(i)}p_bp_bp_t)+2m_bm_t\eta ^{(1)}(p_bk`$
$`2p^{(i)}p_b),`$
$`h_4^{(i)}`$ $`=`$ $`h_3^{(i)}(\eta ^{(1)}\eta ^{(2)}),`$
$`h_5^{(i)}`$ $`=`$ $`2\eta ^{(2)}(m_t^2p_bk2p^{(i)}p_tp_bp_t)+2m_bm_tm_t\eta ^{(1)}(p_tk2p^{(i)}p_t),`$
$`h_6^{(i)}`$ $`=`$ $`h_5^{(i)}(\eta ^{(1)}\eta ^{(2)}),`$
$`h_7^{(i)}`$ $`=`$ $`4\eta ^{(2)}(p^{(i)}p_bp_tkp^{(i)}kp_bp_tp_bkp^{(i)}p_t2p_bkp_tk)`$
$`4m_bm_t\eta ^{(1)}p^{(i)}k,`$
$`h_8^{(i)}`$ $`=`$ $`h_7^{(i)}(\eta ^{(1)}\eta ^{(2)}),`$
$`h_9^{(i)}`$ $`=`$ $`4m_t\eta ^{(2)}p_bk(p_bkp^{(i)}p_b)4m_b\eta ^{(1)}p^{(i)}p_bp_tk,`$
$`h_{10}^{(i)}`$ $`=`$ $`h_9^{(i)}(\eta ^{(1)}\eta ^{(2)}),`$
$`h_{11}^{(i)}`$ $`=`$ $`4m_t\eta ^{(2)}p_bk(p_tkp^{(i)}p_t)4m_b\eta ^{(1)}p_tkp^{(i)}p_t,`$
$`h_{12}^{(i)}`$ $`=`$ $`h_{11}^{(i)}(\eta ^{(1)}\eta ^{(2)}),`$
where the index $`i`$ represents the two channels $`s`$ and $`t`$, and $`p^{(s)}=p_b`$, $`p^{(t)}=p_t`$.
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warning/0003/astro-ph0003266.html | ar5iv | text | # On Companion-Induced Off-Center Supernova-Like Explosions
## 1. Introduction
It has recently been reported that neutron stars with large luminosities could be powered by ultra-strong magnetic fields $`10^{14}G`$ (”magnetars”) rather than rapid rotation as in rotation-powered pulsars (e.g. Kouveliotou et al. 1998, 1999). These ”magnetars” are believed to be produced in supernova explosions or accretion induced collapses (Duncan & Thompson 1992).
On the other hand, the Blandford-Znajek process has been suggested as a powerful source for a variety of very energetic events ranging from active galactic nuclei (Blandford & Znajek 1977) to gamma-ray bursts (e.g. Lee et al. 2000ab and references therein). The extremely high luminosities or power outputs generally require ultra-strong magnetic fields around rapidly rotating black holes. The origin of such strong magnetic fields around the rapidly spinning black holes has not been clearly understood. It is interesting to point out that if magnetars directly collapse to black holes via rapid accretion, the resulting black holes could spin rapidly with magnetic fields left at the horizon, which are large enough to make the Blandford-Znajek power interesting for very energetic events.
Based on this possibility, we propose that the hypercritical accretion onto a neutron star with a very strong magnetic field can produce a rapidly spinning black hole with an ultra-strong field near the horizon. The hypercritical accretion has been invoked for the common envelope phase evolution of the massive stellar binaries (Bethe & Brown 1998). Unless the strong magnetic fields decay within time scales shorter than the duration of the accretion phase, $``$ 1 year, it is highly likely that the black holes with very strong fields could be produced. The Blandford-Znajek power is sufficient to power a supernova-like event with the center of explosion displaced from the companion core. The companion core, after explosion, evolves into a C/O-white dwarf or a neutron star with a second explosion. The detection of highly eccentric black-hole, C/O-white dwarf binaries or the double explosion structure in the supernova remnants, with two different centers of explosion, could be an evidence of the proposed scenario. Especially, in the later case, the neutron star can have double kicks, resulting in very high kick velocities, $`>1000`$ km s<sup>-1</sup>. If the spiral-in of the neutron star leads to the merger with the companion core, the explosion at the core would appear as an explosion from a single star. The neutron star-neutron star binary system is expected when the companion collapses before significant accretion of the neutron star occurs.
## 2. Magnetar Formation and Ultra-Strong Magnetic Fields
The NSs right after core collapse are hot and convective throughout most of the stellar interior due to a large neutrino flux. Duncan & Thompson (1992) and Thompson & Duncan (1993) describe the most likely outcome of such hot neutron stars. Right after the collapse, there exists an intense neutrino-driven convective phase with the overturn time scale of convective cell, $`\tau _{\mathrm{conv}}1`$ ms at the base of the convection zone and $`\tau _{\mathrm{conv}}0.1`$ ms at the neutrinosphere. In this case, if the neutron star spin frequency is
$`\mathrm{\Omega }\text{ }>\mathrm{\Omega }_{\mathrm{dynamo}}2\pi \tau _{\mathrm{conv}}^16\times 10^3\mathrm{s}^1,`$ (1)
the amplification of of a magnetic field by helical motions is not suppressed by turbulent diffusion, and the $`\alpha \mathrm{\Omega }`$-type dynamo could work efficiently. They found that the saturation field strength of $`B_{\mathrm{sat}}10^{16}10^{17}\mathrm{G}`$ could be reached within $`10100`$ ms. Therefore, post-collapse NSs with $`P\text{ }<1\mathrm{m}\mathrm{s}`$, within the first few seconds, can build up very strong large scale magnetic fields $`10^{15}\mathrm{G}`$ through exponential growth. In this case, the resulting ultra-strong field is independent of the fossil field strength. For strong seed fields, the small-scale fluid motion, which is required for ohmic dissipation of magnetic energy, is suppressed by the magnetic tension. The convective motion ceases when the neutrino-driven convection becomes unavailable. Such a strong field pulsar with slow spin (”magnetar”, Duncan & Thompson 1992) has recently been discovered, which strongly supports the basic idea of creation of a ultrastrong field pulsar (Kouvelioutou et al. 1998, 1999).
For NSs spinning with $`\mathrm{\Omega }\text{ }<\mathrm{\Omega }_{\mathrm{dynamo}}`$ or $`P\text{ }>1\mathrm{m}\mathrm{s}`$, the differential rotation (i.e. $`d\mathrm{\Omega }/dR0`$) could amplify the seed poloidal fields by wrapping them up into strong toroidal fields (Kluzniak & Ruderman 1998). For an initial poloidal seed field with strength $`B_p10^9\mathrm{G}`$ which would result from flux freezing of a field of $`\text{ }<3\times 10^4\mathrm{G}`$ in the pre-collapse core, is amplified to the strong toroidal field of $`B_\varphi 10^{17}\mathrm{G}`$ after $`(1/2\pi )(B_\varphi /B_p)10^7`$ rotations or time $`\tau B_\varphi /B_p\mathrm{\Omega }10^4P\mathrm{s}`$ with $`P`$ in $`\mathrm{ms}`$. The amplified field becomes buoyant after reaching the critical field strength of $`B_b10^{17}\mathrm{G}`$. Since the total energy included in the buoyant fluid cells is $`\frac{1}{8\pi }B_b^2(\frac{4\pi R^3}{3})\eta 3\times 10^{51}\eta `$ erg, where $`\eta 0.2`$ is the fractional volume of the buoyant elements, the repeated build-up and buoyant loss will occur for a number of times determined as
$`N_{\mathrm{buoyant}\mathrm{loss}}{\displaystyle \frac{(1/2)I\mathrm{\Omega }^2}{(1/8\pi )(B_b^2)(4\pi R^3/3)\eta }}{\displaystyle \frac{2}{\eta }}10,`$ (2)
where $`I\frac{2}{5}MR^2`$ with neutron star mass $`M=1.5\text{ }M_{}`$ and neutron star radius $`R=10`$ km. This process will continue until the existing poloidal seed field is exhausted as the poloidal field is not re-generated in this non-dynamo amplification process. The buoyantly rising field could dissipate on a time scale
$`\tau _{\mathrm{dissipation}}(4\pi \rho )^{1/2}{\displaystyle \frac{R}{B_b}}10^3\mathrm{s},`$ (3)
which controls the overall time scale. Therefore, we conclude that essentially all NSs with $`P`$ a few ms could achieve strong fields through either dynamo or differential rotation.
Once strong fields are created, rapid spin-down of NSs is inevitable. The electromagnetic dipole radiation luminosity is
$`L_{\mathrm{EM}}{\displaystyle \frac{2}{3c^3}}B^2R^6\mathrm{\Omega }^4=2\times 10^{50}B_{15}^2\mathrm{\Omega }_4^4\mathrm{erg}\mathrm{s}^1`$ (4)
which gives the characteristic electromagnetic spin-down time scale
$`\tau _{\mathrm{EM}}{\displaystyle \frac{I\mathrm{\Omega }^2}{2L_{\mathrm{EM}}}}3\times 10^2{\displaystyle \frac{1}{B_{15}^2\mathrm{\Omega }_4^2}}\mathrm{s}`$ (5)
where $`B_{15}=B/10^{15}\mathrm{G}`$ and $`\mathrm{\Omega }_4=\mathrm{\Omega }/10^4\mathrm{s}^1`$. Therefore, all MSPs with ms periods are likely to spin-down to near the pulsar deathline, $`P70B_{15}^{1/2}\mathrm{s}`$, on a time scale of $`10^5B_{15}^1`$ yr (e.g. Ritchings 1976). That is, these initially fast spinning MSPs would not be detected as MSPs unless they are recycled through the low mass X-ray binary stage, which is not promising (Yi & Grindlay 1998). Therefore, before the common envelope phase begins, the neutron stars are likely to be spinning with spin periods very similar to those of the magnetars and anomalous X-ray pulsars (e.g. Mereghetti 1999, Kouvelioutou 1998, 1999). That is, the initial conditions for the neutron stars which spiral-in during the common envelope phase are likely to be a slow spin and a high magnetic field.
## 3. Hypercritical Accretion
Bethe & Brown (1998) suggested that the hypercritical accretion can occur in the massive binary systems, in which both stars can go into supernova leaving the compact cores as remnants. In their scenario, the more massive star <sup>1</sup><sup>1</sup>1The progenitor masses should be different by more than 5%. If the two masses are close enough, the binary will explode nearly at the same time leaving the neutron star$``$neutron star (ns,ns) binary. In this case, there is not enough time to change neutron star into black hole by accretion. first explodes leaving a neutron in the core. If the binary is within the radius of the companion in the giant stage, the first born neutron star will spiral into the envelope of the expanding companion in the giant stage. We assume the envelope of the giant to be convective.
In the rest frame of the compact object, Bondi-Hoyle-Lyttleton accretion of the hydrogen envelope of density $`\rho _{\mathrm{}}`$ and velocity $`V`$ is
$`\dot{M}=2.23\times 10^{29}{\displaystyle \frac{M_{co,1}^2}{V_8^3}}\rho _{\mathrm{}}\mathrm{g}\mathrm{s}^1`$ (6)
where $`M_{co,1}M_{co}/1\text{ }M_{}`$ is the mass of the compact object, $`V_8`$ is the velocity in unit of 1000 km s<sup>-1</sup>, and $`\rho _{\mathrm{}}`$ is given in g cm<sup>-3</sup>. The density $`\rho _{\mathrm{}}`$ required for the hypercritical accretion, $`\dot{M}=10^8\dot{M}_{Edd}10^{26}`$ g s<sup>-1</sup>, is
$`\rho _{\mathrm{}}=0.44\times 10^3{\displaystyle \frac{V_8^3}{M_{co,1}^2}}\mathrm{g}\mathrm{cm}^3.`$ (7)
Using the Kepler law for circular orbits $`V^2=GM_{tot}/a`$, where $`M_{tot}`$ is the combined mass of the compact object and the stellar material interior to the orbit of the compact object and $`a`$ is the binary separation, one finds
$`\rho _{\mathrm{}}=2.1\times 10^5{\displaystyle \frac{1}{M_{co,1}^2}}\left({\displaystyle \frac{M_{tot,10}}{a_{12}}}\right)^{3/2}\mathrm{g}\mathrm{cm}^3,`$ (8)
where $`M_{tot,10}=M_{tot}/10^{10}M_{}`$ and $`a_{12}=a/10^{12}cm`$. Since this critical density is of the same order as the average density of the hydrogen envelope of the giant star, the hypercritical accretion $`\dot{M}\text{ }>10^8\dot{M}_{Edd}`$ can occur in this spiral-in process. This rapid accretion can add $`1\text{ }M_{}`$ to the first-born neutron star in less than about one year, which is enough to turn the first born neutron star into a black hole. In this process, the orbital radius shrinks by a factor $`50`$.
If the companion is massive enough to have a supernova explosion, the binary system will end up as a black hole$``$neutron star (bh,ns) binary instead of a (ns,ns) binary. If the companion is less massive, the binary will end up as black hole$``$CO white dwarf (bh,cowd) binaries. Brown, Lee, Portegies Zwart, & Bethe (1999) show that there are big discrepancies between the predicted and the observed (ns,ns) and (ns,cowd)<sub>c</sub> (i.e. with circular orbits) binary systems. They solved this discrepancy by arguing that in both systems the first born neutron star turned into black holes by accretion leaving behind (bh,ns) and (bh,cowd) binaries. The (bh,ns) binaries are one of the most important sources for gravitational detectors. The chirp mass of double neutron star binary with masses $`1.4\text{ }M_{}`$ is $`M_{chirp}=\mu ^{0.6}(M_1+M_2)^{0.4}=1.2\text{ }M_{}`$, while that of the (bh,ns) binary with $`M_{BH}=2.4\text{ }M_{}`$ is $`M_{chirp}=1.6\text{ }M_{}`$. The (bh,ns) binaries, if they exist, can be more easily detected by factor $`(1.6/1.2)^3=2.4`$ than the (ns,ns) binaries. Since the LIGO is based on one observation of (ns,ns) binary, their unseen (bh,ns) will enhance the detectability of LIGO by factor of $``$30 considering the effect of the chirp masses. This will be tested in near future.
One of the most strong criticisms on the hypercritical accretion was the angular momentum and possibility of a strong jet along the rotation axis. Even thouth these remain as unsolved problems, Brown, Lee, & Bethe (2000) showed that Narayan & Yi (1994) ADAF solution can be applied in this case if the accretion flows is radiation-dominated. With the hypercritical accretion, the inner boundary of accretion flow can be extended to the marginally bound orbit ($`r_{mb}=2R_{Sch}`$ for Schwarzschild black holes). Since $`r_{mb}9`$ km, the inner boundary of accretion disc can be extended to the neutron star surface $`r_{NS}\text{ }>10`$ km. By using a reasonable viscosity parameter $`\alpha =0.05`$, they found that the neutrino cooling in the disc is almost negligible $`<10^7`$ of viscosity generated energy. Brown et al. (2000) argued that in order to avoid strong polar jets or pile-up up of material above neutron star surface, the neutrinos are generated at the neutron surface to carry out the pressure and energy from the system.
## 4. Hypercritical Accretion onto a Magnetar
Despite a very strong magnetic field, the dynamic effects of the magnetic field would be negligible if the mass accretion rate satisfies
$`\dot{M}`$ $`>`$ $`{\displaystyle \frac{1}{2}}B^2\left({\displaystyle \frac{GM}{R^5}}\right)^{1/2}7\times 10^{25}B_{14}^2\left({\displaystyle \frac{M_{1.5}}{R_6^5}}\right)^{1/2}\mathrm{g}\mathrm{s}^1`$ (9)
$``$ $`4\times 10^7B_{14}^2\left({\displaystyle \frac{M_{1.5}}{R_6^7}}\right)^{1/2}\dot{M}_{\mathrm{Edd}}`$
Such a possibility is realized when the accretion is hypercritical.
During hypercritical accretion, even when the electromagnetic dipole radiation occurs for a fast spinning pulsar, the spin-up by hypercritical accretion can counter the spin-down. The spin-down torque due to the dipole radiation emission is
$`N_{em}={\displaystyle \frac{2}{3c^3}}B^2R^6\mathrm{sin}^2\theta \mathrm{\Omega }^3`$ (10)
and the spin-up torque due to accretion is
$`N_{acc}=\dot{M}\stackrel{~}{a}(GMR)^{1/2}`$ (11)
where $`\stackrel{~}{a}0.5`$ is the sub-Keplerian rotation factor in the hypercritical accretion (Brown, Lee, & Bethe 2000). The equilibrium spin is reached at the neutron star spin frequency
$`\mathrm{\Omega }_{\mathrm{NS}}1.4\times 10^3\left({\displaystyle \frac{\stackrel{~}{a}_{0.5}\dot{M}_{26}M_{1.5}^{1/2}}{\mathrm{sin}^2\theta R_6^{11/2}B_{14}^2}}\right)^{1/3}s^1.`$ (12)
By assuming $`\mathrm{sin}^2\theta 0.5`$, one gets $`\mathrm{\Omega }_{NS}1.8\times 10^3\mathrm{s}^1`$ with other quantities in unities. The collapse of the neutron star to a black hole results in a rapidly spinning black hole with a spin frequency
$`\mathrm{\Omega }_\mathrm{H}=\left({\displaystyle \frac{R_{\mathrm{NS}}}{R_\mathrm{H}}}\right)^2\mathrm{\Omega }_{\mathrm{NS}}4\mathrm{\Omega }_{\mathrm{NS}}7.2\times 10^3s^1`$ (13)
which corresponds to black hole spin of $`\stackrel{~}{a}0.1`$. The magnetic field strength
$`B_\mathrm{H}=\left({\displaystyle \frac{R_{\mathrm{NS}}}{R_\mathrm{H}}}\right)^2B_{\mathrm{NS}}4B_{\mathrm{NS}}4\times 10^{14}B_{\mathrm{NS},14}\mathrm{G}`$ (14)
where we have assumed the angular momentum conservation and flux-freezing.
If the magnetar enters the hypercritical accretion regime as a slowly spinning neutron star, the spin-up during the hypercritical accretion will take place on a time scale
$`t_{spinup}{\displaystyle \frac{I_{\mathrm{NS}}\mathrm{\Omega }_{\mathrm{NS}}}{N_{acc}}}1.6\times 10^7\mathrm{s}0.5\mathrm{yr}`$ (15)
which is comparable to the spiral-in time scale ($`1`$ yr) during the common envelope phase.
## 5. Energetics of Blandford-Znajek Process
Based on the steady pulsar periods and spin-down rates, the magnetic field of pulsars are generally estimated by assuming the dipole radiation induced by magnetic fields. The typical fresh neutron stars in the pulsar island have $`10^{12}`$ G, and the recycled ones have $`\text{ }<10^{10}`$ G. The discoveries of the magnetar and AXP suggest the possibilities of the strong magnetic field $`\text{ }>10^{14}G`$. Even though AXP and magnetars raised questions on the origin of the pulsar mechanism, we assume the dipole radiation as origin of the pulsar in these systems. The probability of pulsar with strong magnetic field is assumed to be $`10\%`$ (Kouveliotou et al. 1998, 1999).
After the discovery of GRB afterglows, the theoretical understandings of the formation of GRBs have been improved quite a lot. However, the mechanism of the central engine is still controversial. The electromagnetic extraction of the BH rotating energy by the Blandford-Znajek process (Bandford & Znajek 1977) is one of the candidate mechanisms, in which the strong magnetic field is essential to power the GRB. Typically $`B_H10^{15}`$ G can generate enough energy to power the GRB from rotating black holes. However, the existence of the strong magnetic field in the black holes is uncertain because the strong magnetic fields diffuse in very short time scales.
We suggest that the Blandford-Znajek process can operate if the collapsing neutron stars, with $`B\text{ }>10^{14}`$ G, carry the magnetic fields with them, amplifying the fields upto $`10^{15}`$ G. In this scenario, the existence of neutron stars prior to the black holes is essential.
The hypercritical accretion prior to the collapse to a black hole spins up the neutron star and the black hole is expected to be spinning at rotational angular velocity close to $`\stackrel{~}{a}0.1`$. Once black hole is formed, the spin-down torque due to the dipole emission is not working, and the spin-up torque due to accretion can increase the black hole spin further. So, our estimate given above may be a low limit.
The Blandford-Znajek power (Lee, Wijers, & Brown 2000a) from the rotating black hole with the angular momentum $`J=aM`$ and the magnetic field $`B_H`$ is given by
$`P_{\mathrm{BZ}}`$ $``$ $`10^{50}a^2B_{H,15}^2(M/M_{})^2\mathrm{erg}\mathrm{s}^1.`$ (16)
Using the values in the previous section, we have $`P_{\mathrm{BZ}}3.6\times 10^{47}`$ erg s<sup>-1</sup>. The total available rotational energy is 0.06% of the rest mass $`2\times 10^{51}`$ erg (Lee, Wijers, & Brown 2000a).
## 6. The Role of Black Hole Formation in the Spiral-In Phase
When the first born star with strong magnetic field accretes material from the companion envelope, the black hole should be formed inside the envelope. Since the black hole is formed in the process of the hypercritical accretion, the temperature near the surface of the neutron star should remain $`1`$ MeV until the collapse into a black hole. This temperature is required to generate neutrinos. In this case, the spin frequency for the r-mode instability is so high that we can neglect the effect (Lindblom & Owen 1999). If the Blandford-Znajek process works at the time of black hole formation, there is enough energy to generate GRBs. However, because of the optically thick envelope, gamma rays cannot come out through the hydrogen (or helium) envelope. Instead of GRBs, the energy will pile up in the envelope to generate supernova. The main difference from the normal SN scenarios is that the SN is not induced by the core collapse of the progenitors. The natural consequence is the strong asymmetry in the SN events. Since the companion mass is not the unique parameter, even the less massive stars (normal white-dwarf progenitors) can have SN events with high kick velocities.
The available energies powered by rotating black holes ($`>1.5\text{ }M_{}`$) with Blandford-Znajek process can be $`10^{51}`$ erg. Since the binding energies ($`0.6GM^2/R`$) of the hydrogen envelope of giant stars ($`10\text{ }M_{}`$) are $`𝒪(10^{48})`$ erg, newly formed black holes have enough energy to blow the envelope off. If only $`1\%`$ of this SN energy is available for the kinetic energy of the He core of the giant star, the available kick velocities are $`𝒪(1000`$ km s<sup>-1</sup>).
Fast moving white dwarfs from SN remnants, if exists, will be the consequences of this ”companion induced SN events”. If the black holes are formed just before the normal SN events of the companions, the double SN are also possible leaving (bh,ns) as remnants. Especially, in the later case, the neutron star can have double kicks, resulting in very high kick velocities, $`>1000`$ km s<sup>-1</sup>.
Bethe & Brown (1998) argued that the formation rate of $`(bh,ns)`$ binaries is $`10^4`$ yr<sup>-1</sup> per galaxy, where the low mass black holes ($`2.5\text{ }M_{}`$) are formed in the common envelope evolution. By assuming the flat $`q`$ distribution, Brown, Lee, Portegies Zwart, & Bethe (1999) also argued that the formation rate of circular $`(bh,cowd)_c`$ binaries is $`1.7\times 10^4`$ yr<sup>-1</sup> per galaxy. In their scenario the first born neutron star went into the common envelope evolution, where the neutron star is converted into black hole by hypercritical accretion. According to our scenario, in the formation of $`(bh,ns)`$ and $`(bh,cowd)_c`$ binaries by hypercritical accretion, about $`10\%`$ of the neutron stars have very strong magnetic field and will have the companion (newly formed black holes) induced supernova-like explosions. As a result $`10\%`$ of progenitors of $`(bh,cowd)_c`$ binaries end up as eccentric $`(bh,cowd)_e`$ binaries instead of circular ones. Also $`10\%`$ of progenitors of $`(bh,ns)`$ binaries have double explosions with two different centers of explosion. By assuming that half of the progenitors of $`(bh,cowd)_e`$ or $`(bh,ns)`$ binaries survive the explosion induced by rotating black holes, we have the formation rates of eccentric $`(bh,cowd)_e`$ binaries or double explosions
$`R10^5\mathrm{yr}^1\mathrm{per}\mathrm{galaxy}.`$ (17)
Since the cooling time of CO white dwarfs, i.e. the time required to reach luminosities $`\mathrm{log}(L/L_{})=4.5`$ with $`L_{}=`$ solar luminosity, is $`110`$ Gyr (Salaris et al. 1997), the number of eccentric $`(bh,cowd)_e`$ binaries with luminosity $`\mathrm{log}(L/L_{})>4.5`$ is
$`N10^410^5\mathrm{in}\mathrm{Galaxy}.`$ (18)
The number of non-pulsating eccentric $`(ns,cowd)_e`$ binaries with luminosity $`\mathrm{log}(L/L_{})>4.5`$ is $`10`$ times more popular than $`(bh,cowd)_e`$ systems. In order to identify $`(bh,cowd)_e`$ systems, one need well-established maximum neutron star masses. If there are mechanisms to identify the black holes from dead pulsars in addition to the mass (Balberg et al. 1999), one can detect $`(bh,cowd)_e`$ systems.
## 7. Discussion
We have looked into a possibility that a supernova-like explosion could occur during the common envelope evolution of a strongly magnetized neutron star. The ”magnetar” accretes in the form of the hypercritical accretion during the spiral-in phase and collapses to a rapidly spinning black hole with a strong magnetic field near the horizon. Then the explosion results from the sudden release of energy through the Blandford-Znajek mechanism. In this scenario, the strong magnetic field naturally occurs as it is derived from the strong field of the pre-collapse ”magnetar”.
There are a number of points to be clarified in the proposed scenario some of which could be observationally testable. First, if the magnetic field of the magnetar decays on a time scale $`10^310^4`$ yrs as is often argued for the evolutionary behavior of magnetars and anomalous X-ray pulsars (Colpi, Geppert, & Page 2000), then the common envelope phase has to follow the magnetar formation. Second, during the accretion phase, most of the emitted energy (including the dipole radiation) is ultimately radiated from the neutron star in the form of the neutrinos. Third, if the Blandford-Znajek phase occurs while the neutron star collapses into black hole in the envelope, the resulting off-center explosion will give a high eccentricity, which should be testable in terms of the highly eccentric binary systems. Fourth, there will be little metal spread from the core since there will be no core collapse. The dominant effect would be blowing-off the hydrogen envelope preceded by the expansion of the envelope.
If the merger of the companion core with the neutron star occurs before the black hole formation, then the explosion will not cause an off-center explosion. Although this possibility exists, the estimated accretion time scale on which the neutron star can accrete enough mass to collapse to a black hole is comparable to the spiral-in time scale. This implies that the off-center explosion is possible.
CHL is partly supported by the U.S. Department of Energy under grant DE-FG02-88ER40388. IY is supported in part by 1999-2000 KIAS Research Fund and KRF Research Fund KRF 1998-001-D00365. HKL is supported in part by BK21 Program of Ministry of Education and by KOSEF Grant No. 1999-2-112-003-5. |
warning/0003/quant-ph0003132.html | ar5iv | text | # Table des matières
## Chapitre 1 Introduction
A l’heure actuelle, l’informatique a pris une place considérable, si ce n’est essentielle, dans le monde de l’industrie et des sciences. Dans le but d’être toujours plus rapides et toujours plus performants, les différents utilisateurs demandent un matériel informatique adapté à leurs besoins. Ainsi, nous assistons à une incessante amélioration de la capacité des processeurs, des mémoires et des périphériques composant un ordinateur. Le nombre de transistors qu’il est possible de placer sur un chip croît exponentiellement dans le temps, un constat appelé “Loi de Moore”. Dans ce contexte, l’idée de changer le système informatique de base actuel au lieu de l’améliorer constamment est bien entendu prise en compte. L’une de ces propositions est l’ordinateur quantique.
Historiquement, c’est en 1985 que D. Deutsch démontra théoriquement qu’il était possible de réaliser une machine de Turing à l’aide des concepts de mécanique quantique . Le premier pas de la recherche pour un ordinateur quantique était fait. Cependant, l’explosion de la recherche dans ce domaine eut lieu en 1994, lorsque P. Shor proposa un algorithme pour déterminer la période d’une fonction donnée (étape essentielle pour la factorisation d’un nombre) . Ce dernier est un algorithme polynomial (le nombre d’étape varie comme $`N^\alpha `$), alors qu’il n’en existe pas d’aussi puissant lors d’une programmation classique. Par la suite d’autres algorithmes avec un nombre d’étapes réduites par rapport aux algorithmes classiques ont été démontrés. Citons l’algorithme de Grover qui permet de sélectionner un élément d’une liste et celui de Deutsch-Jozsa qui permet de déterminer si une fonction donnée $`f`$ est constante ou non . Il est donc, à première vue, intéressant de chercher à réaliser un ordinateur quantique. L’artricle de revue le plus récent est celui de C.H. Bennett et D.P. DiVincenzo
Dans cette étude, après avoir décrit le concept de l’ordinateur quantique, une analyse critique des premières réalisations expérimentales est présentée. Avant de conclure, un commentaire sur l’importance que l’ordinateur quantique a prise dans la recherche en physique depuis 1994 sera également donné.
## Chapitre 2 Le concept de l’ordinateur quantique
### 2.1 Rappel de concepts de mécanique quantique
#### 2.1.1 Description du système : l’espace de Hilbert
Lorsque nous étudions un système physique à l’aide des lois de la mécanique quantique, nous travaillons dans un espace mathématique appelé espace de Hilbert, car ce dernier possède les propriétés nécessaires (linéarité, existance d’un produit scalaire et complétude) pour satisfaire les principes de base de la théorie. Si nous prenons un système $`S`$ composé de sous-systèmes $`S_1,\mathrm{},S_n`$, l’espace de Hilbert correspondant est construit comme l’espace produit tensoriel des espaces décrivant chaque sous-système : $`=_1\mathrm{}_n`$. Ainsi, en identifiant un vecteur $`|\psi `$ à un état du système, nous nous rendons compte, de par la construction mathématique, de l’existence d’états de superposition. En effet, grâce à la linéarité de l’espace nous avons : $`a|\psi +b|\varphi ,\psi ,\varphi \text{et}a,b\text{}`$. Remarquons que même lorsque $`|\psi =|\psi _1\mathrm{}|\psi _n,|\varphi =|\varphi _1\mathrm{}|\varphi _n`$ il n’est en général pas possible d’écrire l’état de superposition comme un produit tensoriel d’états des sous-systèmes: il existe ainsi dans $``$ des états non-factorisables, appelés aussi états intriqués (entangled states). L’état intriqué décrit une situation dans laquelle, bien que l’état du système global $`S`$ soit défini, l’état de chaque sous-système $`S_i`$, ne l’est pas. Pour illustrer ce concept, prenons comme exemple un système $`S`$ composé de deux sous-systèmes $`S_1`$ et $`S_2`$ à deux niveaux (notés $`+`$ et $``$) :
$$\begin{array}{ccc}1)\hfill & |\psi =|+|+=|\left[|+|+\right]\hfill & \text{état factorisable}\hfill \\ 2)\hfill & |\psi =|+|++|\psi _1|\psi _2\hfill & \text{état intriqué.}\hfill \end{array}$$
(2.1)
Un système comportant un nombre fini $`m`$ d’états orthogonaux est décrit par un espace de Hilbert complexe de dimension $`m`$. Cet espace est donc isomorphe à $`\text{}^m`$. Comme nous le discuterons dans la prochaine section, pour des ressemblances avec le langage informatique classique il est habituel, lorsque nous traitons l’ordinateur quantique, de considérer comme système quantique typique un système formé de $`n`$ systèmes à deux niveaux, décrit donc par
$$=\underset{nfois}{\underset{}{\text{}^2\mathrm{}\text{}^2}}\text{}^{2^n}.$$
(2.2)
Il est clair que tout système décrit par $`\text{}^m`$ peut être mis en correspondance avec un sous-espace de $`\text{}^{2^n}`$ pour $`m2^n<m+1`$. Le fait que tous les systèmes quantiques possédant le même nombre d’états orthogonaux soient décrits par le même espace vectoriel a inspiré l’une des premières réflexions sur l’ordinateur quantique, due à Feynman<sup>1</sup><sup>1</sup>1Feynman envisage la possibilité de simuler le comportement d’un système quantique en utilisant non pas un ordinateur (classique), mais un autre système quantique qu’on pourrait contrôler..
#### 2.1.2 Evolution : opérateur unitaire
En mécanique quantique, l’évolution du système est régie par un opérateur unitaire, généralement noté $`U`$ (i.e. $`U^1=U^+`$). Remarquons qu’un opérateur unitaire conserve le produit scalaire :
$$\begin{array}{ccc}& U\psi |U\varphi =U^+U\psi |\varphi =U^1U\psi |\varphi =\psi |\varphi .& \end{array}$$
(2.3)
Une discussion des hypothèses qui mènent à adopter un opérateur d’évolution unitaire se trouve dans par. 8.6, p. 237. L’équation de Schrödinger est une conséquence assez immédiate de l’unitarité de l’évolution.
#### 2.1.3 Couplage avec l’environnement (“décohérence”)
Pour que la théorie rende compte de l’observation, il est nécessaire d’introduire le couplage du système avec son environnement. Par la préparation, nous plaçons le système dans un état bien défini $`|\varphi `$. L’état initial de l’environnement est $`|e`$. S’il y a couplage entre l’environnement est le système (i.e., si l’opérateur d’évolution n’est pas de la forme $`U_sU_e`$), on parle d’évolution “bruitée”. Algébriquement, cette situation se décrit de la manière suivante :
$$\begin{array}{ccc}& |\varphi |e& \\ & & \text{Evolution unitaire bruitée}\\ & _ic_i|\varphi _i|e_i.& \end{array}$$
(2.4)
Vu le nombre énorme de degrés de liberté de l’environnement, il est naturel d’admettre que — si le temps écoulé depuis la préparation est assez long — les $`|e_i`$ sont orthogonaux entre eux. Ainsi, en prenant la moyenne d’une observable du système $`𝒜=A\text{}`$ nous obtenons :
$$\begin{array}{ccc}\psi |𝒜|\psi \hfill & =\hfill & _i|c_i|^2\psi _i|A|\psi _i\underset{=1}{\underset{}{e_i|e_i}}+_{ij}c_i^{}c_j\psi _i|A|\psi _j\underset{=0}{\underset{}{e_i|e_j}}\hfill \\ & =\hfill & _i|c_i|^2\psi _i|A|\psi _i\hfill \end{array}$$
(2.5)
De par le couplage avec l’environnement, les termes croisés sont nuls. Ceci n’est évidemment plus vrai si nous considérons une évolution non bruitée (i.e. $`U=U_sU_e`$) :
$$\begin{array}{cc}|\varphi |e& \\ & \text{Evolution unitaire non bruitée}\hfill \\ \left[_ic_i|\varphi _i\right]|e^{^{}}& \\ \psi |𝒜|\psi =& _{i,j}c_i^{}c_j\psi _i|A|\psi _j\underset{=1}{\underset{}{e^{^{}}|e^{^{}}}}.\hfill \end{array}$$
(2.6)
Ainsi, dans le cas d’une évolution bruitée, nous nous retrouvons avec un état de mélange pour le système, alors que dans l’autre cas, nous gardons un état de superposition. Lorsque nous tenons compte du couplage entre le système et l’environnement, il y a une perte d’information. Comme nous l’avons montré précédemment, la perte d’information est complète lorsque les états de l’environnement $`|e_i`$ sont orthogonaux. On appelle alors temps de décohérence $`\tau _{dec}`$ le temps typique sur lequel cette perte d’information se produit. Cette décohérence peut également être décrite en ne considérant que le système et en invoquant une réduction du paquet d’onde (“collapse”). Le temps de décohérence s’interprète alors comme le temps qui s’écoule avant que le collapse n’ait lieu. La limitation principale de l’ordinateur quantique réside dans ce temps de décohérence. En effet, suite au collapse, il n’est plus possible de continuer le calcul. Il est donc nécessaire d’avoir réalisé toutes les opérations souhaitées avant que le phénomène n’ait lieu. Nous reviendrons sur cette limitation dans la section 3.2.1.
### 2.2 Réinterprétation en langage informatique
* En considérant des systèmes à deux niveaux, il est possible, au lieu de parler de spin up et down, d’hélicité $`+`$ et $``$ ou de niveau d’énergie f (fondamental) et e (excité), de parler de Q-bit (quantum binary digit) prenant la valeur 0 ou 1 :
$$\begin{array}{ccc}|& ,& |\\ |+& ,& |\\ |f& ,& |e\end{array}\}|0,|1.$$
(2.7)
Par cette simple nomenclature, nous définissons les support de l’information.
* En continuant notre nomenclature, le système physique considéré devient un ordinateur. Pour créer un ordinateur travaillant avec $`N`$ Q-bits, il suffit, par analogie au cas classique, de prendre $`N`$ systèmes à deux niveaux. Ceci est satisfaisant pour l’esprit, mais n’est pas nécessaire. En effet, un seul système possèdant $`2^N`$ niveaux convient tout aussi bien du moment qu’un adressage cohérent est effectué. Cependant, pour simplifier la suite de l’exposé, nous traiterons toujours le cas d’un ordinateur quantique composé de $`N`$ spins $`\frac{1}{2}`$, ce qui ne restreint pas la généralité des propos.
* Finalement, l’évolution du système se traduit par un calcul à l’aide de portes logiques. Les portes logiques élémentaires à l’aide desquelles tout calcul peut être réalisé, sont au nombre de deux : la rotation et l’opération XOR (cf. 2.3.1). Ainsi, le schéma de base est le suivant :
$$\begin{array}{ccccc}\text{préparation}& & \text{évolution}& & \text{détection}\\ & & & & \\ \text{entrée}& & \text{calcul}& & \text{sortie}.\end{array}$$
(2.8)
Remarque : Dans le langage des machines de Turing, le système joue le rôle de la “bande infinie” sur laquelle l’information est lue et enregistrée, et l’évolution joue naturellement le rôle du processeur.
A ce stade, nous possédons des Q-bits, un ordinateur avec une entrée et une sortie, ainsi que la possibilité d’implémenter n’importe quelle opération unitaire. Il ne manque donc plus que des algorithmes intéressants pour que l’ordinateur quantique soit théoriquement achevé.
### 2.3 L’algorithme
#### 2.3.1 Définition du concept de l’algorithme
Un algorithme est une suite convenable d’évolutions unitaires et de mesures sur le système. On veut pouvoir effectuer n’importe quelle opération unitaire sur le système de $`N`$ Q-bits. Il a été montré (Barenco et al. ) que toute opération unitaire peut être construite comme le produit de deux opérations simples :
* La rotation R. Cette opération sert à changer un seul bit indépendamment des autres. Comme par convention les spins up et down se situent sur les directions $`+z`$ et $`z`$, cette opération unitaire est de la forme :
$$\begin{array}{c}R_{\theta \varphi }|0=\mathrm{cos}\theta |0+e^{i\varphi }\mathrm{sin}\theta |1\hfill \\ R_{\theta \varphi }|1=e^{i\varphi }\mathrm{sin}\theta |0+\mathrm{cos}\theta |1\hfill \\ \text{où }\theta \text{ et }\varphi \text{ sont les angles azimutal et polaire respectivement.}\hfill \end{array}$$
(2.9)
Mise à part la difficulté de pouvoir appliquer pratiquement cette opération sur chaque Q-bit séparément, il est intéressant de remarquer que celle-ci donne lieu à un état de superposition, concept appartenant uniquement à la physique quantique. Cette opération est donc intrinsèquement non-classique.
* L’opération XOR. L’opération XOR (exclusive OR, appelée aussi CNOT : controlled-NOT) est une opération logique entre deux Q-bits : elle agit sur un Q-bit donné selon l’état d’un autre Q-bit. Soit $`|\sigma _n\sigma _m`$ l’état de deux Q-bits quelconque repérés par $`n`$ et $`m`$. L’opération XOR, notée $`C(n,m)`$, peut se définir ainsi :
$$\begin{array}{ccc}C(n,m)|00& =& |10\\ C(n,m)|10& =& |00\\ C(n,m)|01& =& |01\\ C(n,m)|11& =& |11.\end{array}$$
(2.10)
L’état du Q-bit $`n`$ est changé si et seulement si l’état du Q-bit $`m`$ est $`|0`$. L’opération XOR couple donc deux Q-bits quelconques mais, contrairement à la rotation, reste une opération classique <sup>2</sup><sup>2</sup>2Toute la logique classique peut être réalisée à partir de portes XOR.. Comme nous le voyons dans l’exemple suivant, l’effet de cette opération dans un ordinateur quantique porte sur le degré d’intrication des états :
$$\begin{array}{ccc}C(n,m)\left[|00+|11\right]\hfill & =\hfill & |10+|11\hfill \\ & =\hfill & |1\left[|0+|1\right].\hfill \end{array}$$
(2.11)
A l’aide de la rotation, nous pouvons donc changer n’importe quel Q-bit indépendamment des autres et avec l’opération XOR, nous pouvons agir sur le degré d’intrication des états. Ainsi, tout l’espace de Hilbert est atteint. Pour illustrer la manipulation de ces opérations, prenons un système à 3 spins. Nous pouvons, par exemple, à partir de l’état $`|111`$ obtenir un état d’intrication maximale:
$$\begin{array}{ccc}\left[\text{}\text{}R\left(\frac{\pi }{4}\right)\right]|111& =& \frac{1}{\sqrt{2}}\left(|111+|110\right)\\ \left[\text{}C(2,3)\right]\frac{1}{\sqrt{2}}\left(|111+|110\right)& =& \frac{1}{\sqrt{2}}\left(|111+|100\right)\\ \left[C(1,2)\text{}\right]\frac{1}{\sqrt{2}}\left(|111+|100\right)& =& \frac{1}{\sqrt{2}}\left(|111+|000\right).\end{array}$$
(2.12)
Pour d’autres exemples didactiques, il est conseillé de consulter l’article de V. Scarani .
#### 2.3.2 Avantage d’un algorithme quantique
A l’heure actuelle, quelques algorithmes quantiques plus performants que leurs homologues classiques sont connus. La particularité de tous ces algorithmes est l’utilisation d’états superposés. Par ce biais, le nombre d’opérations de base nécessaires est nettement plus faible qu’avec un algorithme classique. Ainsi, l’algorithme de Shor permet de factoriser un nombre donnée à l’aide de $`N^\alpha `$ opérations, où $`N`$ est la taille du nombre à factoriser. Les algorithmes classiques pour ce genre de problèmes nécessitent un nombre d’opérations qui croît exponentiellement avec $`N`$. L’algorithme de Grover qui permet de trouver un élément dans une liste est également moins complexe que son homologue classique. Nous nous rendons alors compte que le gain de temps ne se réalise pas sur la rapidité d’exécution d’une opération de base, mais sur le nombre d’opérations. Par conséquent, ce qui peut définir un ordinateur de quantique est la façon d’utiliser le système mais non pas le système lui-même. Ce n’est pas parce qu’un système de spins est utilisé qu’un ordinateur est quantique.
Il existe d’autres algorithmes qui, tout en étant spécifiquement quantiques, ne changent pas la complexité. Par exemple celui de Fahri et al. (détermination de la parité d’une fonction) nécessite $`N/2`$ opérations , alors que son homologue classique en nécessite $`N`$. Dans la section suivante, nous illustrons le concept d’algorithme au moyen de l’algorithme de Deutsch-Jozsa .
#### 2.3.3 Un exemple : l’algorithme de Deutsch-Jozsa
Prenons un système $`S`$ composé de deux sous-systèmes $`S_1`$ et $`S_2`$. Il est ainsi possible de construire les quatre fonctions différentes agissant sur l’ensemble $`\{0,1\}`$ :
$$\begin{array}{cccc}f_1(x)=0,& f_2(x)=1,& f_3(x)=x,& f_4(x)=NOTx\end{array}$$
(2.13)
L’algorithme de Deutsch-Jozsa permet de savoir si une fonction prise au hasard est constante $`\left(f_{1,2}\right)`$ ou non $`\left(f_{3,4}\right)`$. Pour ce faire, quatre étapes suffisent :
1. Préparation de l’état initial, en l’occurence : $`|0|0`$
2. Rotation de l’état de chaque spin, livrant ainsi un état de superposition :
$$\begin{array}{ccc}\frac{1}{2}\left(|0+|1\right)\left(|0|1\right)\hfill & =\hfill & \frac{1}{2}_{x=0}^1|x\left(|0|1\right)\hfill \end{array}$$
(2.14)
3. Appel de $`f(x)`$. Il faut appliquer $`\left(1+f(x)\right)mod2`$ sur le système $`S_2`$. Remarquons que si $`f=f_{3,4}`$, cette étape nécessite l’utilisation de l’opération XOR. Le système $`S`$ se trouve alors dans l’état:
$$\begin{array}{c}\frac{1}{2}_{x=0}^1|x\underset{=(1)^{f(x)}\left(|0|1\right)}{\underset{}{\left(|\left(0+f(x)\right)mod2|\left(1+f(x)\right)mod2\right)}}=\hfill \\ \\ =\frac{1}{2}\left(_{x=0}^1(1)^{f(x)}|x\right)(|0|1)=\{\begin{array}{cc}\pm \frac{1}{2}(|0+|1)(|0|1)\hfill & \text{si }f=f_{1,2}\hfill \\ \pm \frac{1}{2}(|0|1)(|0|1)\hfill & \text{si }f=f_{3,4}\hfill \end{array}\hfill \end{array}$$
(2.15)
4. Rotation inverse de celle décrite au point 2 pour trouver l’état final du système :
$$\begin{array}{cc}|0|0& \text{si }f=f_{1,2}\\ |1|0& \text{si }f=f_{3,4}\end{array}$$
(2.16)
La mesure du premier Q-bit nous indique ainsi si la fonction est constante ($`|0`$) ou non ($`|1`$).
Le point remarquable de cet algorithme est que la fonction $`f`$ que nous désirons tester n’est appelée qu’une seule fois, alors que dans un algorithme classique deux appels sont nécessaires. A la section 3.2.1, nous décrirons une réalisation pratique de cet algorithme à l’aide de la résonance magnétique nucléaire.
### 2.4 Conclusion
Au terme de cette section, nous remarquons donc qu’un ordinateur quantique est n’importe quel système quantique vu comme porteur d’information. Au niveau théorique, il n’y a rien de nouveau. Nous assistons seulement à une réinterprétation de concepts déjà établis en mécanique quantique. Pour pouvoir effectuer un calcul quantique, il faut se donner la possibilité d’effectuer n’importe quelle opération unitaire. Pour ce faire, deux opérations suffisent : la rotation et l’opération XOR. Cependant, le couplage du système avec l’environnement pose de graves problèmes : il faut donc réaliser toutes les opérations nécessaires pour effectuer le calcul avant l’échéance de $`\tau _{dec}`$ ($``$ limitation du nombre d’opérations) ou alors trouver un moyen de vaincre ce phénomène. Ainsi, dans toute proposition de réalisation d’un ordinateur quantique, nous trouverons trois éléments : le système physique envisagé, la manière d’implémenter les deux opérations et la manière de vaincre la décohérence.
Soulignons encore que le point crucial d’un algorithme quantique est l’utilisation de la superposition d’états, c’est-à-dire que l’information elle-même est quantique. Un système quantique qui traite de l’information classique n’est pas un ordinateur quantique. En effet, les semiconducteurs actuels possèdent déjà un caractère quantique bien que les algorithmes utilisés soient classiques. Notons également que, si le système physique est composé de sous-systèmes, alors les états intriqués seront également essentiels, car ils sont une conséquence directe de la superposition. Cependant, si nous considérons uniquement le système global, nous verrons apparaître des états superposés, mais aucun état intriqué. Ainsi, la notion d’état intriqué dépend du point de vue que l’on adopte, mais l’état physique réel auquel il fait référence demeure essentiel si nous désirons utiliser toutes les possibilités d’un ordinateur quantique.
## Chapitre 3 Réalisations pratiques
### 3.1 Les différents systèmes utilisés
Différents systèmes physiques sont considérés pour la réalisation pratique. Les principaux sont les suivants :
* Les ions piégés. Un ion piégé peut être vu comme une particule plongée dans un potentiel (problème hydrogénoïde). Il est alors possible de considérer deux configurations différentes du ion, donc deux niveaux d’énergie distincts. A l’aide d’impulsions laser nous pouvons modifier l’état de chaque ion et l’intrication des états. Par cette technique, le groupe de Boulder a obtenu une intrication de quatre Q-bits .
* Les boîtes quantiques (quantum dots). Les boîtes quantiques sont en fait une réalisation des puits de potentiel carré souvent considérés dans la théorie. Le Q-bit fait alors référence à deux niveaux d’énergie différents d’un électron placé dans un de ces puits. A ma connaissance, aucun état intriqué n’a pu être réalisé jusqu’à présent.
* Les photons. Dans le cas des photons, les deux états du Q-bit sont deux états de polarisation et les opérations sont réalisées par d’astucieuses techniques d’interférométrie. L’intrication à trois photons a été démontrée. Comme le support de l’information n’est pas solide, l’interférométrie de photons est plutôt envisagée pour la communication et la cryptographie quantiques.
La seule technique qui a réussi à implémenter des algorithmes, notamment celui de Deutsch-Jozsa et un code de correction d’erreurs, est la résonance magnétique nucléaire, que nous analysons plus en détail dans la prochaine section.
### 3.2 Réalisation par résonance magnétique nucléaire
#### 3.2.1 L’expérience de Chuang et al.
De nombreux groupes de recherche tentent de réaliser un ordinateur quantique à l’aide de la résonance magnétique nucléaire (RMN). La principale question qui se pose est la suivante : comment travailler avec des états purs <sup>1</sup><sup>1</sup>1Il ne suffit pas que le système ait une évolution bien déterminée, faut-il encore que nous puissions lire le résultat et l’interpréter. dans des expériences de RMN effectuées à température ambiante, donc avec des états distribués thermiquement ? La réponse se trouve dans la notion d’états pseudo-purs, introduite en 1997 par Cory et al. et par Gershenfeld et Chuang .
Considérons un système $`S`$ composé de deux spins. En RMN, ces spins sont plongés dans un champ magnétique extérieur. Les états propres de cette interaction sont $`|++`$, $`|+`$, $`|+`$, $`|`$. A l’équilibre, la probabilité d’occupation de ces états est donnée par la statistique de Boltzmann :
$$\begin{array}{cc}P_{\sigma \omega }=\frac{e^{\frac{E_{\sigma \omega }}{kT}}}{Z}\hfill & \sigma ,\omega =+,\hfill \\ Z=_{\sigma ,\omega }e^{\frac{E_{\sigma \omega }}{kT}}\hfill & \text{fonction de partition.}\hfill \end{array}$$
(3.1)
Nous nous rendons alors compte que si $`T\mathrm{}`$ chaque état est peuplé équitablement ($`P_{\sigma ,\omega }=\frac{1}{4},\sigma ,\omega `$). Les écarts typiques à la distribution uniforme sont de l’ordre de $`\frac{E}{kT}10^5`$ à température ambiante. Il est habituel en RMN d’écrire la matrice densité totale $`\rho _t`$ comme :
$$\rho _t=\frac{1}{4}\text{}+\rho .$$
(3.2)
L’intérêt de cette décomposition vient du fait que $`\rho `$ donne les seules contributions non-triviales à la dynamique. La matrice $`\rho `$ est appelée matrice densité réduite et nous remarquons qu’elle doit être de trace nulle puisque $`\text{Tr}(\rho _t)=1`$. Cette matrice densité réduite ne représente donc pas un état (la trace de toute matrice densité représentant un état, pur ou de mélange, vaut 1). Cependant, on a montré l’existence d’une préparation $`U`$ astucieuse (séquences d’impulsions en RMN) telle que:
$$\begin{array}{ccc}U^+\rho _tU\hfill & =\hfill & \frac{1}{4}\text{}+U^+\rho U\hfill \\ & =\hfill & \frac{1}{4}\text{}\frac{ϵ}{4}\text{}+\underset{=ϵ\rho _1}{\underset{}{\frac{ϵ}{4}\text{}+U^+\rho U}}\hfill \\ & =\hfill & \left(\frac{1ϵ}{4}\right)\text{}+ϵ\rho _1,\hfill \end{array}$$
(3.3)
avec $`\rho _1`$ la matrice densité décrivant un état pur. Le terme $`\left(\frac{1\alpha }{4}\right)\text{}`$ ne va pas influencer la dynamique du système. Ainsi, l’analyse de $`\rho _1`$ suffit. On parle alors d’états pseudo-purs. De cette manière, les chercheurs justifient leur tentative de réaliser un ordinateur quantique à l’aide de la RMN. Ce résultat se généralise facilement pour un système de $`N`$ spins :
$$\begin{array}{cc}\rho _ϵ=\frac{1ϵ}{d}\text{}_d+ϵ\rho _1,\hfill & d=2^N\hfill \end{array}$$
(3.4)
Cette justification théorique étant explicitée, nous pouvons nous intéresser à la réalisation pratique de l’algorithme de Deutsch-Jozsa (2.3.3) effectué par I.L. Chuang et al. . Les deux Q-bits sont les spins nucléaires du carbone et de l’hydrogène dans une molécule de chloroforme (CHCl<sub>3</sub>). Comme ces spins ont des fréquences de résonance différentes, il est aisé d’effectuer une rotation sur l’un des systèmes, sans influencer le second. L’opération XOR se réalise à partir d’une judicieuse combinaison d’impulsions (techniques de double résonance). Nous n’entrerons pas ici dans les détails techniques des séquences utilisées. En effectuant les quatre étapes proposées par l’algorithme, le pic de résonance apparaît selon $`+z`$ pour les fonctions $`f_{1,2}`$ et selon $`z`$ pour les fonctions $`f_{3,4}`$, comme attendu par la théorie.
Nous pourrions alors louer la réussite. Remarquons cependant que ce genre de techniques RMN sont à l’heure actuelle maîtrisées <sup>2</sup><sup>2</sup>2Les premières expériences de double irradiation ont été effectuées par Bloch en 1954 ; dès la fin des années 50, ces techniques apparaissent dans les livres et que l’algorithme était déjà établi. De plus, comme le nombre d’opérations est assez restreint pour le cas traité, les auteurs n’ont pas eu à se préoccuper du problème majeur : la décohérence. En effet, dans ce cas très simple, le couplage entre le système et l’environnement n’a pas le temps d’agir <sup>3</sup><sup>3</sup>3Dans ces expériences de RMN, il est habituel de considérer le temps de relaxation $`T_2`$ comme le temps de décohérence; pour l’heure, il s’agit d’une estimation d’ordre de grandeur plutôt que d’une idenfication conceptuelle entre $`T_2`$ et $`\tau _{dec}`$ (e-mail de R. Laflamme à V. Scarani). Cependant, dès que le problème se complique légérement, les choses deviennent catastrophiques.
Pour illustrer ces propos, nous suivons l’article de Haroche et Raimond . Rappelons que dans la section 2.1.3 nous avons défini la grandeur $`\tau _{dec}`$ comme étant le temps disponible pour effectuer nos opérations, avant que n’ait lieu la réduction du paquet d’onde. Ainsi, connaissant le temps nécessaire pour effectuer une opération ($`\tau _{op}`$), nous pouvons définir la grandeur $`M`$, le nombre de pas qu’il est possible de réaliser sans perdre d’information.
$$\begin{array}{ccc}& M=\frac{\tau _{dec}}{\tau _{op}}& \\ & M10^7& \end{array}$$
(3.5)
Pour espérer factoriser un nombre de 4 bits , $`10^6`$ opérations sont nécessaires et pour un nombre de 400 bits, $`10^{12}`$ opérations. En se rappelant que le plus grand nombre de 4 bits est 15, l’ordinateur quantique semble assez limité. Pour remédier à ce problème, des protocoles de corrections sont mis au point. L’idée théorique de ces codes est assez simple. Il suffirait de déterminer l’opérateur unitaire d’évolution bruitée, d’en calculer l’inverse, et de l’appliquer sur l’état. Ainsi, la donnée serait maintenue dans son état initial. Nous ne sommes cependant pas au bout de nos peines, car un algorithme de correction est complexe et nécessite également la superposition d’état. Un deuxième ordinateur quantique serait donc nécessaire…
Ainsi, malgré sa publication dans la prestigieuse revue Nature, l’expérience de Chuang et al. n’est pas aussi révolutionnaire qu’il n’y paraît au premier abord : tous les éléments théoriques sont connus, la technique est maîtrisée et les problèmes pratiques interéssants à résoudre sont contournés. Pour couronner le tout, un article de S.L. Braunstein et al. démontre la proposition suivante: “Tous les états utilisés jusqu’à maintenant en RMN pour l’ordinateur quantique ou pour d’autres protocoles d’information quantique sont séparables” (i.e. non intriqués). L’idée de la démonstration est la suivante. Le premier pas est une décomposition de la matrice densité définie par l’équation (3.4) dans une base surdimentionnée, en l’occurence à l’aide des matrices de Pauli. Le critère suivant est alors appliqué: si tous les coefficients de la décomposition sont non négatifs, alors la matrice considérée est séparable. Comme les coefficients négatifs proviennent uniquement de la matrice $`\rho _1`$, si le coefficient $`ϵ`$ est suffisamment petit, alors la matrice densité est séparable. En calculant une borne inférieure pour ce coefficient, ces auteurs démontrent que dans toute les tentatives de réalisation d’ordinateur quantique par RMN, les états intriqués physiques ne sont pas accessibles. L’état pseudo-pur est un simulateur d’état pur. Ainsi, bien que l’expérience réalisée par Chuang et al. utilise un algorithme quantique, ce dernier ne nécessite pas forcément l’utilisation d’états intriqués contrairement à l’algorithme de Shor. Un ordinateur réalisé par RMN à des températures “normales” ne pourra donc jamais factoriser un nombre donné avec le nombre d’opérations minimales prévues. Quelle est donc la véritable définition d’un ordinateur quantique ? A ce stade de la recherche naîssent des confusions ainsi qu’un débat de fond.
#### 3.2.2 Le débat de l’ordinateur quantique par résonance magnétique nucléaire
Ce débat est plus qu’une simple discussion entre connaissseurs du domaine. Dans ce sens, l’article de Physics Today nous dépeint le tableau actuel à l’aide d’interviews des différents opposants. Suite à sa démonstration de la séparabilité de la matrice densité des états pseudo-purs en RMN, S.L. Braunstein pose la question de base suivante : “Lorsque nous pensons à l’ordinateur quantique, nous pensons que c’est merveilleux, et qu’il utilise les fantastiques propriétés d’un système quantique. Ces propriétés sont des choses comme la superposition et l’intrication. Ainsi, quel sens ont ces machines si elles ne peuvent produire aucun état intriqué ?” Selon S. Popescu, l’intrication est l’élément de base pour le succès de l’ordinateur quantique. Ainsi, cette démonstration implique de sérieux doutes sur les réalisations d’ordinateur quantique par RMN. Les expérimentateurs qui travaillent sur ce projet se défendent d’une part par l’unitarité des transformations et d’autre part par l’évolution du système. Selon R. Laflamme, le fait d’utiliser des tranformations unitaires (à caractère quantique) permet d’obtenir la réponse cherchée de manière plus efficace. De plus, Il n’a pas été possible à ce jour de décrire l’évolution observée de manière purement classique. Cependant, Linden et Popescu ont montré que l’algorithme de factorisation de Shor nécessite l’utilisation d’états intriqués, et par conséquant, une évolution quantique ne suffit pas. Les ordinateurs réalisés par RMN ne pourront donc jamais atteindre la puissance espérée d’un ordinateur quantique. Il semble donc que ces machines ne soient qu’une “simulation” d’un véritable ordinateur quantique. C’est dans cette brêche que s’engouffre S. Lloyd en affirmant qu’il existe une multitude de façon d’implémenter des opérations, et pas seulement de façons classique ou quantique. Le fait que l’ordinateur construit par RMN se trouve à mi-chemin entre les mondes classique et quantique devient ainsi simultanément un chef d’accusation et une défense. Devant de telles discussions, il est légitime de se demander si le débat porte sur des faits scientifiques ou si nous assistons simplement à un problème de linguistique. Il ne serait pas impossible que dans un proche avenir une définition plus stricte de l’ordinateur quantique apparaisse.
Lorsque nous quittons la presse spécialisée pour des revues destinées à un plus large publique, nous nous rendons compte que le flou linguistique est encore plus prononcé. Prenons comme exemple l’article de J. Dousson paru dans le Flash Informatique de l’EPFL . Dans son introduction, elle nous explique, graphique à l’appui, que d’ici 2020, quelques atomes suffiront pour stocker un bit. Peut-être. Alors continue-t-elle, sur de tels systèmes, seules les lois de la mécanique quantique sont valables. Sûrement. Par consèquent, nous travaillerons avec de l’informatique quantique. Sûrement pas. Nous retrouvons ici de façon flagrante la confusion entre un système quantique et de l’information quantique. En dehors de ces confusions, il est quand même intéressant de souligner que ces articles de vulgarisation se multiplient dans les différentes revues (par ex. l’Ordinateur Individuel) et même dans les quotidiens (par ex. La Liberté).
Au terme de ce chapitre, nous nous rendons compte que la réalisation pratique est non seulement difficile, mais qu’une fois un prototype réalisé, il n’est pas évident de savoir (ou d’admettre) s’il remplit toutes les exigences requises pour que nous puissions parler d’ordinateur quantique.
## Chapitre 4 Aspect social de l’ordinateur quantique
### 4.1 Les publications
#### 4.1.1 Explosion de la recherche
Comme nous l’avons déjà vu en introduction, c’est en 1994 que P. Shor trouva le premier algorithme intéressant pour l’ordinateur quantique. Par cette découverte, l’idée qu’un ordinateur quantique permet de résoudre des problèmes plus rapidement qu’un ordinateur classique était née et par la même occasion, une explosion de la recherche dans ce domaine.
A l’aide de la base de données INSPEC (Intranet de l’EPFL : biberl.epfl.ch/cgi-bin/webspirs.cmd), le nombre d’articles publiés par année et qui contiennent l’un des syntagmes “quantum computer”, “quantum computing”, “quantum computation” a été évalué. Comme indicateur de la qualité des articles, une seconde courbe ne comprenant que les articles parus dans Physical Review a été ajoutée. Les résultats sont présentés sur la figure 4.1. Il est à signaler que, lors d’un survol, environ 5% de ces articles ne traitent pas de l’ordinateur quantique, le syntagme “quantum computation” pouvant être utilisé pour décrire un calcul numérique d’un système quantique. Sur la figure réalisée, il est aisé de constater une explosion de la recherche dans ce domaine dès 1994, donc après l’article de P. Shor. Nous nous rendons également compte que le quart des articles ont été publiés dans Physical Review.
Par la lecture des abstracts, nous constatons que les années 1997-98 correspondent à l’achèvement de la théorie et à quelques idées de réalisations pratiques. En effet, les scientifques reconsidèrent les algorithmes et s’intéressent aux corrections nécessaires pour minimiser la perte d’information par la décohérence (cf. 2.1.3). En 1999, les notions théoriques étant quasiment complètes, ce sont des tentatives pratiques qui dominent la recherche.
Quelles peuvent être alors les raisons de cet engoûment pour l’ordinateur quantique ? Avec 210 articles parus en 1999, donc plus de 4 par semaine, alors que le sujet n’était quasiment pas traité 5 ans auparavant, nous pouvons nous demander si l’ordinateur quantique est aussi prometteur que la figure 4.1 nous le suggère, ou si nous sommes en présence d’un effet de mode.
#### 4.1.2 Raisons
Nous proposons ici une analyse des différentes raisons qui peuvent conduire un groupe de scientifiques à se lancer dans cette recherche .
La première est tout simplement la passion du sujet. Dès son plus jeune âge, la mécanique quantique a suscité de grands débats. La discrétisation des niveaux d’énergie, la dualité onde-corpuscule ont provoqué un doute même chez les plus grands de l’époque comme Einstein. Dans cette continuité, un calcul sur un ordinateur quantique débute avec un état connu pour terminer avec un état bien défini, en passant par un état de superposition. Le point intriguant de ce schéma réside dans le fait qu’à un certain moment du calcul, nous ne pouvons pas connaître l’état dans lequel se trouve le système. En effet, considérons uniquement deux spins et supposons que le système global se trouve dans l’état $`\frac{1}{\sqrt{2}}\left[|11+|10\right]`$. Nous avons une probabilité $`\frac{1}{2}`$ que, lors d’une mesure, le système soit dans l’état $`|11`$ et également $`\frac{1}{2}`$ qu’il soit dans l’état $`|10`$. De plus, si nous effectuons la mesure, nous avons une modification de l’état, donc une perte d’information. L’ordinateur quantique entre donc dans la lignée de ces éléments surprenants, mais toujours passionnants, de la mécanique quantique.
La seconde est l’utilité de l’application finale. L’ordinateur quantique, mais aussi la cryptographie ou la théorie de l’information quantiques, semblent être les premières démonstrations de l’utilité de la superposition d’états, notion qui relevait jusqu’à très récemment des débats conceptuels. En effet, comme cette dernière permet d’établir des algorithmes plus rapides, nous pouvons espérer découvrir la solution de problèmes jusqu’alors qualifiés d’irrésolvables ou d’approximatifs de par la limitation de nos outils informatiques. Par exemple, la cryptographie classique, qui se base sur la factorisation d’un nombre donné en deux nombres premiers, deviendrait désuète.
Une troisième raison pourrait s’appeler le recyclage scientifique. Depuis 1960, de nombreuses personnes se sont intéressées aux fondements de la mécanique quantique, notamment à certains concepts comme la non-localité ou l’existence de variables cachées . Cependant, à l’heure actuelle, les expériences réalisables sont de plus en plus nombreuses et en bon accord avec les prédictions de la théorie. Une réinterprétation de la mécanique quantique permet alors d’éviter de mettre un point final prématurément à son élaboration.
Finalement, comme presque toujours, l’argent est une excellente raison. Une entreprise qui commercialiserait un ordinateur surpassant tous les modèles existants réaliseraient certainement de gigantesques bénéfices. Ainsi, en lui faisant miroiter une telle application, les scientifiques du domaine s’assurent les fonds nécessaires pour leur recherche. Soulignons également le fait qu’une partie importante du financement provient d’organismes nationaux ou supranationaux. En effet, de nos jours, la possibilité d’application est un critère de financement non seulement pour les privés, mais aussi pour ces organismes. Par exemple, le projet TOP NANO 21 du Fonds National Suisse exige un partenariat industriel. L’ordinateur quantique peut ainsi satisfaire un critère de rentabilité demandé par les privés mais malheureusement aussi par les organismes nationaux.
Toutes les raisons mentionnées — la passion, l’utilité, le recyclage scientifique et l’argent — me paraissent être de bonnes motivations pour se lancer dans cette recherche. Cependant, comme nous l’avons vu dans la section précédente, la grande limitation pratique et le peu d’algorithmes existants me poussent à croire que ce sont les deux dernières les principales.
### 4.2 L’ordinateur quantique en Suisse et dans le monde
Peu de groupes de recherche en Suisse s’occupent de l’ordinateur quantique. Le seul qui travaille activement dans le domaine (ordinateur quantique à l’aide de “quantum dots”) est D. Loss, à Bâle (theorie5.physik.unibas.ch). Le groupe de N. Gisin, à Genève, est plus spécialisé en cryptographie et théorie de l’information quantiques qu’en ordinateur quantique (www.gap-optique.unige.ch). Une brève recherche sur Internet montre que les autres sites où figure l’ordinateur quantique ne traitent pas du sujet, mais le placent comme application directe de leur recherche. Dans ce contexte, nous retrouvons notamment les personnes s’occupant de la physique mésoscopiques, par exemple J. Faist, à Neuchâtel (www.unine.ch/uer/uer$`\mathrm{\_}`$physique.htm). L’EPFL n’y échappe pas non plus: l’IMO qui travaille sur les “quantum dots” a récemment invoqué l’ordinateur quantique comme application possible (Polyrama 112, décembre 1999, p. 36-38).
Pour les gens qui s’intéressent à l’ordinateur quantique dans le monde, citons l’adresse internet du Center for Quantum Computation, situé à Oxford : www.qubit.org. Ce site propose la liste des principaux groupes de recherche dans le monde et en Europe. Notons qu’il est sponsorisé par différentes entreprises et organismes, dont, comme par hasard, Hewlett Packard.
## Chapitre 5 Conclusion
D’après le nombre d’articles publiés par année, l’ordinateur quantique semble être une application très utile et la recherche avancée. Cependant, la réalisation pratique est difficile de par la limitation du nombre d’opérations par la décohérence. Comme nous l’avons vu, $`10^7`$ opérations sont réalisables avant que la décohérence n’agissent, et $`10^6`$ (resp. $`10^{12}`$) sont nécessaires pour factoriser un nombre de 4 (resp. 400) bits. La réalisation d’états intriqués avec un grand nombre de Q-bits est également un grand défi. La factorisation d’un nombre donné $`n`$ par l’algorithme de Shor exige l’entrée du nombre $`\frac{n}{2}`$. Ainsi, la factorisation de 30 nécessite 4 Q-bits, celle de 100, 9 Q-bits, celle de $`10^6`$, 20 Q-bits et pour les applicatons visées, c’est la factorisation de nombres proches de $`10^{100}`$, donc environ 330 Q-bits, qui est intéressante. Cependant, à l’heure actuelle, le meilleur résultat est la réalisation d’un état intriqué de 4 Q-bits. De plus, nous nous rendons intuitivement compte que plus le nombre de Q-bits est élevé, plus le temps de décohérence est court. Ces ordres de grandeur me semblent alors accablantes pour un réalisation pratique de l’ordinateur quantique. Les optimistes me répondront que, si quelqu’un avait affirmé à Jules César qu’un jour l’homme marcherait sur la lune, il se serait fait dévorer par les lions. D’accord, mais entre le char de Ben-Hur et la Formule 1 d’aujourd’hui, l’humanité a assisté à quelques révolutions technologiques. A mon avis, une révolution supplémentaire est à l’heure actuelle nécessaire pour une réalisation de l’ordinateur quantique.
De plus, en écartant l’exploit que représente la réalisation d’un ordinateur quantique, les possibilités prévues d’une telle machine ne sont pas si révolutionnaires que l’on veut nous faire croire. En effet, en dépit de la recherche intensive des cinq dernières années, on a trouvé seulement trois algorithmes (détermination de la période d’une fonction, recherche d’un élément d’une liste, caractérisation d’une fonction) qui sont théoriquement plus performants que les algorithmes classiques. Les comparaisons entre un ordinateur classique et un ordinateur quantique ne sont donc pas équitables. Lorsqu’on nous parle de ce dernier, on nous vante à chaque fois le gain d’opérations que nous pouvons obtenir avec l’algorithme de Shor, mais on oublie volontairement de nous rappeler toutes les autres possibilités qu’offre un ordinateur classique. Quelques scientifiques, notamment E. Farhi et al. , ont cependant démontré que dans bien des cas, l’ordinateur quantique ne surpasserait pas son homologue classique. Ainsi, avant de s’investir pleinement dans une tentative de réalisation, une recherche plus approfondie d’algorithmes me paraît nécessaire.
Il me semble donc que les tentatives de réalisations d’un ordinateur quantique sont prématurées et sa puissance théorique exagérée. Les personnes concernées, par leurs arguments et leurs premiers résultats, adoptent le langage de la demi-vérité. Cependant, comme dans notre société le progrès et le profit ne sont plus de simples mots mais des buts en soi, les propos de ces chercheurs ravissent le public et les industries. Ces dernières fournissent alors les fonds nécessaires aux premiers, ravis à leur tour. Dans cet engrenage vicieux, il est naturel que monsieur-tout-le-monde attende désespérément la promotion des ordinateurs quantiques chez Interdiscount.
Je pense donc que l’ordinateur quantique est considéré, à l’heure actuelle, comme une “simple” application de la mécanique quantique, alors qu’il est un sujet de recherche fondamentale. Les quelques éléments prometteurs ont suggéré certaines capacités au concept de l’ordinateur quantique, mais ont surtout fait miroiter un marché juteux. Ainsi, de nombreux chercheurs et industriels se sont engouffrés tête baissée dans cette brêche. Un véritable effet de mode s’en est suivi. A l’heure actuelle, l’ordinateur quantique devient en grande partie un prétexte de publication, un moyen d’obtention de fonds, voire une justification de recherche, ceci en oubliant le premier but du concept. En consultant la littérature, nous observons également un éclatement du sujet. Je ne serais donc pas étonné si, d’ici quelques années, le nombre d’articles publiés annuellement devrait diminuer au profit d’études apparentées comme par exemple la cryptographie quantique.
De par le peu d’algorithmes existants, l’ordinateur quantique ne me semble pas être aussi puissant que le monde se l’imagine et un manque de technologie adéquate me fait penser qu’il est actuellement irréalisable. Ainsi, les optimistes qui prévoient un raz-de-marée d’ordinateurs quantiques en 2020 me paraissent plutôt être des rêveurs ou des “politiciens” que des scientifiques. En attendant le progrès technique suffisant pour une réalisation, espérons que l’ordinateur quantique, replacé dans son contexte de recherche fondamentale, nous élargira l’esprit vers de nouveaux horizons intéressants.
Au terme de ce travail, je voudrais remercier le Dr. Valerio Scarani pour ses explications et son aide, notamment pour le chapitre 2. |
warning/0003/cond-mat0003338.html | ar5iv | text | # Coulomb repulsion versus Hubbard repulsion in a disordered chain
## 1 Introduction
In low dimensions ($`d2`$) disorder always yields aalr a finite localization length $`L_1`$ when the particles do not interact and there is no spin-orbit scattering. When one wants to study the role of electron electron interaction, a first issue is to know what kind of interaction is appropriate. When there are many carriers inside a large length $`L_1`$ (large density and weak disorder), it looks reasonable to assume weakly interacting Landau quasi-particles and to take the usual short range screened Coulomb repulsion. But, for low carrier densities ($`10^{10}`$-$`10^{11}`$ carriers per $`cm^2`$ is nowadays achieved kravchenko in two dimensional heterostructures) the screening of the charges is somewhat problematic and one may find safer to consider bare long range Coulomb repulsion. One has in this case a system having charge crystallization as a natural limit when kinetic energy becomes negligible compared to Coulomb energy. The range of the interaction can also be varied by metallic gates located in the vicinity of the electron gas, as it is often done for having a tunable carrier density. This gives us the motivation to study the difference between on site Hubbard like repulsion and long range Coulomb repulsion in a simple limit: two electrons in a disordered chain. Since on-site interaction plays a role only if the orbital part of the wave function is symmetric, we restrict our study to the case with opposite spins.
The problem of two interacting particles (TIP) in one dimension has been mainly studied with on site Hubbard repulsion of strength $`U`$. As proposed by Shepelyansky shepelyansky it has been numerically proven moriond ; wmgpf ; Voppen ; Frahm that interaction delocalizes a certain number of TIP states over a length $`L_2>>L_1`$. For a “contact” interaction as Hubbard repulsion, the TIP system exhibits remarkable properties tip1 ; tip2 ; tip3 ; tip4 . The mixing of the one body states inside a scale $`L=L_1`$ and the associated delocalization effect for sizes $`LL_1`$ is maximum for $`UU_c`$, where $`U_c`$ is the fixed point of a duality transformation tip2 mapping the weak $`U/t`$ limit onto the weak $`t/U`$ limit, $`t`$ being the kinetic energy scale. For the one body problem, the spectral statistics contains important information: an Anderson insulator has uncorrelated levels (Poisson statistics) whereas a disordered metal displays Wigner-Dyson rigidity characterizing quantum chaos. At the mobility edge, lies a scale invariant critical statistics shapiro , which exhibits a weaker spectral rigidity associated to weak critical chaos. For the TIP system with Hubbard interaction, the spectrum is Poissonian when $`U0`$ and $`U\mathrm{}`$ and becomes tip2 more rigid when $`UU_c`$. However, the maximum possible rigidity does not correspond to Wigner-Dyson rigidity, but to an intermediate rigidity analogous to those characterizing the one body spectrum at a mobility edge (critical statistics). It was noticed in Ref. tip3 that those intermediate statistics are also related to very slow interaction induced TIP diffusion at scales $`L_1LL_2`$. We show in this study that, in contrast to Hubbard repulsion, Coulomb repulsion can drive the TIP one dimensional system to full quantum ergodicity with Wigner-Dyson statistics. From a statistical study of the interaction matrix elements coupling two free particle (2FP) states (i.e. the TIP eigenstates at $`U=0`$), one finds that Coulomb repulsion mainly favors hopping terms between 2FP states nearby in energy, with energy separation of the order of the TIP level spacing $`\mathrm{\Delta }_2L^2`$. Hubbard repulsion mainly induces tip2 hopping terms between 2FP states separated by a larger energy $`\mathrm{\Delta }_2^{eff}>\mathrm{\Delta }_2`$, and the measure of the coupled 2FP states is multifractal tip1 . However, the generic behavior of a TIP system with either Hubbard or Coulomb repulsions can be summarized by three regimes: a free particle limit dominated by Anderson localization; a large interaction limit dominated again by Anderson localization for Hubbard and by charge crystallization for Coulomb;and between those two Poissonian limits lies an intermediate regime characterized by a maximum mixing of the one particle states and a maximum delocalization effect. The intermediate regime in both cases is located around the interaction strengths $`U_c`$ for which the TIP system has participation ratios of same order in both preferential eigenbases characterizing the weak and strong interaction limits.
## 2 TIP Hamiltonian
The TIP Hamiltonian $``$ is given by the sum of two terms: the first $`_\mathcal{0}`$ gives the kinetic energy (parameter $`t`$) and the random potentials (parameter $`W`$) in which the two particles can move,
$$_\mathcal{0}=t\underset{\{i,j\}}{}c_i^+c_j^{}+W\underset{i}{}v_in_i.$$
(1)
$`v_i`$ is randomly taken in the interval $`[\frac{1}{2},\frac{1}{2}]`$, $`c_i^+`$ creates a particle on the site i and $`n_i=c_i^+c_i`$. The second term $`𝒰`$ is the two body repulsion, which can be either on site Hubbard repulsion:
$$𝒰=U\underset{i}{}n_i(n_i1)$$
(2)
or long range Coulomb repulsion:
$$𝒰=U\underset{i}{}n_i(n_i1)+\frac{U}{2}\underset{\stackrel{i,j=1}{|\mathrm{i}\mathrm{j}|\mathrm{L}/2}}{\overset{L}{}}\frac{n_in_j}{|ij|}$$
(3)
The convention in this work is that two particles at the same site cost an energy $`2U`$ (and not $`U`$ as assumed in previous refs. tip1 ; tip2 ; tip3 ; tip4 ). An additional cost of energy $`U/p`$ has to be paid by two particles separated by a distance $`p`$ with $`L/2p1`$ when there is Coulomb repulsion. The boundary conditions (BCs) are taken periodic.
Let us denote ($`ϵ_\alpha ,\psi _\alpha `$) and $`(E_{\alpha \beta },\psi _{\alpha \beta })`$ the eigenenergies and eigenfunctions of the one particle state $`|\alpha `$ and of the 2FP state $`|\alpha |\beta =|\alpha \beta `$ respectively. One has $`E_{\alpha \beta }=ϵ_\alpha +ϵ_\beta `$. The 2FP level spacing is $`\mathrm{\Delta }_22B/L(L+1)`$ where the band width $`B8t+2W`$. The one particle localization length $`L_1`$ is defined from the weak disorder formula $`L_1=100/W^2`$. Hereafter, the energies will be given in units of the kinetic energy hopping term $`t`$ restricted to nearest neighbors.
## 3 TIP density of states
When one compares the two repulsions, a first difference appears in the density of states $`\rho _2(E)`$. In the limit $`U\mathrm{}`$, Hubbard repulsion splits tip2 the TIP band in two parts: a small band of $`L`$ “molecular states” of high energy $`U+2Wv_i`$ corresponding to two electrons localized on the same site $`i`$ and a main band of $`L(L1)/2`$ “hard core boson” states which remain at the same small energies for $`U\mathrm{}`$ and $`U0`$. The “hard core boson” states are given by the resymmetrization silvestrov of Slater determinants corresponding to electrons in the one body state $`|\alpha >`$ and $`\beta >`$ respectively. Those states do not feel on site interaction, are not coupled to one another and become decoupled from the molecular states of much larger energies when $`U\mathrm{}`$. When one takes rigid BCs, the resymmetrization is simple and one has exactly $`E_{\alpha \beta }(U0)=E_{\alpha \beta }(U\mathrm{})`$ for $`\alpha \beta `$. For two electrons in a ring enclosing a flux $`\varphi `$ the resymmetrization is more subtle and $`E_{\alpha \beta }(\varphi +\varphi _0/2,U0)=E_{\alpha \beta }(\varphi ,U\mathrm{})`$. For periodic BCs, $`E_{\alpha \beta }`$ ($`\alpha \beta `$) goes to the corresponding 2FP eigenenergy with anti periodic BCs. In contrast to Hubbard repulsion where the majority of the TIP energy levels does not feel the interaction when $`U\mathrm{}`$, excepted $`L`$ “molecular” states, Coulomb repulsion eventually crystallizes all the TIP states as two particle “molecules”. When $`L`$ is even, the sizes of the molecules are $`d=0,\mathrm{},L/2`$ and the spectrum is split in $`L/2`$ subbands of $`L`$ states ($`dL/2`$) and one subband of $`L/2`$ states ($`d=L/2`$), each of them centered around an energy $`U/d`$. When $`N`$ is odd, one has $`(L+1)/2`$ subbbands of $`L`$ states. Without disorder, each subband shrinks onto a single $`L`$-fold degenerate state obtained from successive translations of the “molecules” by one lattice spacing, a degeneracy which is broken by the random potentials. The different densities $`\rho _2(E)`$ induced by the two repulsions are illustrated in Fig. 1, for a chain of size $`L=L_1=100`$ and various interaction strengths $`U`$.
## 4 Crossover between two preferential eigenbases
When one turns on Hubbard repulsion, it has been detailed in Ref. tip2 how the TIP system goes from the 2FP basis towards the “hard core boson basis” when $`U\mathrm{}`$. The interaction threshold $`U_c=(24)^{1/4}t/2`$ was defined (for energies near the band center) as the fixed point of the duality transformation mapping the distribution of the interaction matrix elements $`U/t`$ which couple the 2FP states onto the distribution of the kinetic energy matrix elements $`t/U`$ which couple the hard core boson states. At $`UU_c`$, the 2FP basis ceases to be preferential compared to the hard core boson basis. Being unable to extend this duality argument for Coulomb interaction, we study the participation ratios $`PR_0`$ of the TIP wavefunctions $`|\mathrm{\Psi }`$ onto the 2FP eigenbasis and $`PR_{\mathrm{}}`$ onto the basis built out from symmetrized products of site orbitals (site basis). This later basis describes the correlated “molecules” created when $`U\mathrm{}`$.
$`PR_0`$ $`=`$ $`({\displaystyle \underset{\alpha \beta }{}}|\alpha \beta |\mathrm{\Psi }|^4)^1`$
$`PR_{\mathrm{}}`$ $`=`$ $`({\displaystyle \underset{ij}{}}|ij|\mathrm{\Psi }|^4)^1`$ (4)
This allows us to extend for Coulomb repulsion the concept of a crossover threshold $`U_c`$ where the $`U=0`$ eigenbasis ceases to be preferential compared to the $`U=\mathrm{}`$ eigenbasis. As shown in Fig. 2, one has $`U_c120`$ when $`L=L_1=50`$ for Coulomb repulsion.
## 5 Interaction matrix elements in the 2FP basis
Before discussing the TIP spectral statistics, it is useful to study the structure of the interaction matrix elements coupling the 2FP states. The 2FP wave function have components $`\mathrm{\Psi }_{\alpha ,\beta }(n,m)`$ on the sites $`|nm>`$ given by
$$<\alpha \beta |nm>=\frac{\psi _\alpha (n)\psi _\beta (m)+\psi _\alpha (m)\psi _\beta (n)}{\sqrt{2}}.$$
We denote
$$Q_{\alpha \beta }^{\gamma \delta }(0)=\underset{n=1}{\overset{L}{}}\psi _\alpha ^{}(n)\psi _\beta ^{}(n)\psi _\gamma (n)\psi _\delta (n)$$
(5)
and
$$Q_{\alpha \beta }^{\gamma \delta }(p)=\underset{n=1}{\overset{L}{}}\frac{\psi _\alpha ^{}(n)\psi _\beta ^{}(n+p)\psi _\gamma (n)\psi _\delta (n+p)}{|p|}+perm$$
where $`perm`$ means the terms obtained after permuting $`(\alpha \beta ),(\gamma \delta ,)`$ and $`(\alpha \beta ,\gamma \delta )`$ for $`p0`$.
In the 2FP eigenbasis, the interaction matrix elements are
$$\alpha \beta |𝒰|\gamma \delta =4UQ_{\alpha \beta }^{\gamma \delta }(0)=UH_{\alpha \beta }^{\gamma \delta }$$
(6)
for Hubbard repulsion, and
$$\alpha \beta |𝒰|\gamma \delta =U(4Q_{\alpha \beta }^{\gamma \delta }(0)+2\underset{p=1}{\overset{L/2}{}}Q_{\alpha \beta }^{\gamma \delta }(p))=UC_{\alpha \beta }^{\gamma \delta }$$
(7)
for Coulomb repulsion.
In the absence of disorder and with periodic BCs, the one body states are plane waves $`\psi _\alpha (n)=(\mathrm{exp}ik_\alpha n)/\sqrt{L}`$ and the interaction matrix elements $`𝒰_{\alpha \beta }^{\gamma \delta }`$ only couple 2FP states of same momentum $`K=K_{\alpha \beta }=k_\alpha +k_\beta =K_{\gamma \delta }`$. For Hubbard, one has when $`L_1\mathrm{}`$
$$H_{\alpha \beta }^{\gamma \delta }\frac{4}{L}\delta _{K_{\alpha \beta },K_{\gamma \delta }}.$$
(8)
The interaction matrix has a block diagonal form. If $`L`$ is odd, one has $`L`$ blocks of size $`N_s=(L+1)/2`$. If $`N`$ is even, one has $`L/2`$ blocks of size $`N_s=L/2`$ and $`L/2`$ other blocks of size $`N_s=L/2+1`$. The $`N_s`$ TIP eigenenergies $`E_n(K)`$ of same momentum $`K`$ are given by the $`N_s`$ solutions of:
$$\underset{\gamma \delta }{}\frac{1}{E_n(K)E_{\gamma \delta }}=\frac{L}{4U}$$
(9)
where $`E_{\gamma \delta }=2\mathrm{cos}k_\gamma +2\mathrm{cos}k_\delta `$. The $`E_n(K)`$ alternate with the 2FP energies $`E_{\gamma \delta }`$ of same momentum. The TIP spectrum is the uncorrelated sum of $`L`$ such series of different momenta $`K`$.
For Coulomb, the previous block diagonal structure is preserved when $`L_1\mathrm{}`$, but each block becomes more complex:
$`C_{\alpha \beta }^{\gamma \delta }{\displaystyle \frac{4}{L}}\delta _{K_{\alpha \beta },K_{\gamma \delta }}\times `$ (10)
$`(1+{\displaystyle \underset{p=1}{\overset{L/2}{}}}{\displaystyle \frac{\mathrm{exp}i[(k_\beta k_\gamma )p]}{2p}}+{\displaystyle \underset{p=1}{\overset{L/2}{}}}{\displaystyle \frac{\mathrm{exp}i[(k_\beta k_\delta )p]}{2p}}`$
$`+{\displaystyle \underset{p=1}{\overset{L/2}{}}}{\displaystyle \frac{\mathrm{exp}i[(k_\alpha k_\gamma )p]}{2p}}+{\displaystyle \underset{p=1}{\overset{L/2}{}}}{\displaystyle \frac{\mathrm{exp}i[(k_\alpha k_\delta )p]}{2p}})`$
In the presence of a random potential, the absolute values of the matrix elements can be given using a linearly graduated grey scale in the plane $`(\gamma ,\delta )`$ for a given 2FP state $`|\alpha \beta >`$, the one body states $`|\gamma >`$ and $`|\delta >`$ being ordered by increasing energies. When $`W`$ is sufficiently small, TIP momentum remains almost conserved and one can see in the plane $`(\gamma ,\delta )`$ a white cross made by the two diagonals if $`\alpha =\beta `$. The first diagonal corresponds to coupling to other states $`|\gamma \delta >`$ with $`|\gamma >|\delta >`$, the second to coupling to 2FP states $`|\gamma \delta >`$ close in energy ($`E_{\alpha \beta }E_{\gamma \delta }`$). When $`W`$ is larger, Coulomb and Hubbard give rise to a different pattern in the $`(\gamma \delta )`$ plane: energy-momentum conservation remains partially preserved by Coulomb repulsion (the second diagonal persists) and is lost by Hubbard repulsion, as shown in Fig. 3 for $`L2L_1=100`$. 2FP states $`|\alpha \alpha >`$ with $`ϵ_\alpha 0`$ are considered in Fig. 3, but similar conclusions can be drawn from arbitrary 2FP states $`|\alpha \beta >`$.
To explain why $`C_{\alpha \beta }^{\gamma \delta }`$ continues to mainly couple $`|\alpha \beta >`$ to states $`|\gamma \delta >`$ nearby in energy, we note that when $`L_1`$ is finite, disorder smears the sharp delta function into a broader gaussian peak of width $`\sigma L_1^1`$. $`\delta _{K_{\alpha \beta },K_{\gamma \delta }}\mathrm{exp}(K_{\alpha \beta }K_{\gamma \delta })^2/(2\sigma ^2)`$ and $`\mathrm{exp}i(k_\alpha k_\beta )p`$
$`\mathrm{exp}(k_\alpha k_\beta )^2/(2\sigma ^2)`$, one gets respectively
$$H_{\alpha \beta }^{\gamma \delta }\mathrm{exp}\frac{(K_{\alpha \beta }K_{\gamma \delta })^2}{2\sigma ^2}$$
and
$`C_{\alpha \beta }^{\gamma \delta }`$ $`\mathrm{exp}({\displaystyle \frac{(K_{\alpha \beta }K_{\gamma \delta })^2}{2\sigma ^2}})\times `$
$`[2\mathrm{ln}({\displaystyle \frac{|k_\beta k_\gamma |}{\sigma }}\mathrm{ln}({\displaystyle \frac{|k_\beta k_\delta |}{\sigma }})`$
$`\mathrm{ln}({\displaystyle \frac{|k_\alpha k_\gamma |}{\sigma }})\mathrm{ln}({\displaystyle \frac{|k_\alpha k_\delta |}{\sigma }})]`$
This explains the behavior shown in Fig. 4 where the disorder averaged amplitude of the hopping terms are given as a function of the energy difference. The $`L`$ hopping terms between the 2FP states $`|\gamma \delta >`$ nearby in energy $`(\mathrm{\Delta }E\mathrm{\Delta }_2L/2)`$ are much smaller for Hubbard than for Coulomb. For larger energy separation $`\mathrm{\Delta }E`$, there is a $`\mathrm{ln}(|\mathrm{\Delta }E/\mathrm{\Delta }_2|)`$ decay which is more pronounced for Coulomb than for Hubbard.
## 6 TIP diffusion and 2FP lifetime
Coulomb repulsion couples a density $`\rho _2`$ of 2FP states nearby in energy. Hubbard repulsion effectively couples a smaller density $`\rho _2^{eff}<\rho _2`$, as explained in Ref. tip1 . Let us review three consequences of this difference.
(i) TIP diffusion: In the first studies shepelyansky ; imry ; fmgp of the TIP problem, a density $`\rho _2(L_1)`$ of 2FP states coupled by the interaction was assumed for $`LL_1`$. Under this assumption, it was predicted that the TIP dynamics should exhibit interaction assisted diffusion on scales $`L_1<L<L_2`$, the time evolution of the TIP center of mass $`R_2`$ being given fmgp by:
$$R_2(t)\sqrt{D_2(t)t}$$
(12)
with $`D_2(t)`$ is roughly constant, up to $`log(t)`$ corrections fmgp . In Ref. tip3 , a much slower propagation $`R_2\mathrm{log}t`$ was observed for Hubbard repulsion, attributed to the weak density $`\rho _2^{eff}(L_1)`$ of effectively coupled states. One expects that the original prediction will be at least partially restored for Coulomb repulsion.
(ii) TIP localization: The interaction assisted propagation stops at a scale $`L_2`$ characteristic of TIP localization. Assuming $`|<\alpha \beta |𝒰|\gamma \delta >|^2U^2/L_1^3`$ for $`LL_1`$, the enhancement factor $`L_2/L_1`$ was originally given by the estimate:
$$\frac{L_2}{L_1}|<\alpha \beta |𝒰|\gamma \delta >|^2\rho _2(L_1)L_1$$
(13)
It was pointed out in Ref. tip1 that, since $`\rho _2(L_1)L_1^2\rho _2^{eff}(L_1)L_1^{1.75}`$ in the presence of Hubbard repulsion, the enhancement factor should be weaker ($`L_2/L_1\sqrt{L_1}`$), a prediction confirmed by numerical calculations. It is likely that $`L_2/L_1L_1`$ will be a better estimate for Coulomb repulsion.
(iii) 2FP lifetime: The inverse lifetime $`(\mathrm{\Gamma }_{\alpha \beta })`$ of a 2FP state $`|\alpha \beta `$ is given by Fermi Golden rule:
$$\mathrm{\Gamma }_{\alpha \beta }\underset{\gamma \delta }{}|<\alpha \beta |𝒰|\gamma \delta >|^2\delta (E_{\alpha \beta }E_{\gamma \delta })$$
(14)
If the density of coupled states is $`\rho _2(L_1)L_1^2`$, a reasonable estimate for Coulomb repulsion, one gets
$$\mathrm{\Gamma }_{\alpha \beta }(L=L_1)|<\alpha \beta |𝒰|\gamma \delta >|^2\rho _2(L_1)L_1^1.$$
In contrast, the lifetime should be longer for Hubbard repulsion:
$$\mathrm{\Gamma }_{\alpha \beta }(L=L_1)|<\alpha \beta |𝒰|\gamma \delta >|^2\rho _2^{eff}(L_1)L_1^{1.65}$$
since the density tip1 of coupled states (by the second moment of $`|<\alpha \beta |𝒰|\gamma \delta >|`$) is $`\rho _2^{eff}L_1^{1.35}`$ for a 2FP state $`|\alpha \alpha >`$.
We have checked this prediction. In Fig. 5, the inverse lifetime calculated in a chain of length $`L=420`$ is shown as a function of $`L_1`$. When $`L_1<L`$, one has $`\mathrm{\Gamma }_{\alpha \alpha }L_1^{1.65}`$ for Hubbard repulsion while $`\mathrm{\Gamma }_{\alpha \alpha }L_1^{1.05}`$ for Coulomb repulsion, close to a simple decay $`L_1^1`$. When $`L`$ becomes larger than $`L_1`$, the lifetime becomes $`L`$-independent for Hubbard, and may continue to weakly decay as a function of $`L`$ for Coulomb. Fig. 5 shows us that this possible decay remains negligible.
## 7 Two particle spectral statistics
We now study how the spectral statistics depend on the interaction strength $`U`$ for the two repulsions, using the distribution $`P(s)`$ of energy spacings between consecutive levels and the variance $`\mathrm{\Sigma }_2(E)`$ of the number of levels inside an energy window of width E. We consider energy levels in the bulk of the low energy sub-band: $`E0`$ for Hubbard repulsion and $`E2U/L`$ for Coulomb repulsion (see Fig. 1). We take $`L=L_1`$ for having the largest interaction matrix elements, and hence the maximum mixing of the 2FP states.
After unfolding the spectra, one expects
$$P(s)P_W(s)\frac{\pi s}{2}\mathrm{exp}(\frac{\pi }{4}s^2)$$
(15)
and
$$\mathrm{\Sigma }_2(E)\mathrm{\Sigma }_2^W(E)\frac{2}{\pi ^2}(ln(2\pi E)+\gamma +1\frac{\pi ^2}{8})$$
(16)
($`\gamma `$ being the Euler constant) for correlated levels having Wigner-Dyson statistics, whereas one should have Poisson statistics with
$$P_P(s)=\mathrm{exp}(s)$$
(17)
and
$$\mathrm{\Sigma }_2^P(E)=E$$
(18)
for uncorrelated levels.
The TIP spectra are well described by Wigner-Dyson statistics for an intermediate Coulomb repulsion. In Fig. 6 and in Fig. 7, one can see indeed that the level spacing distribution and the number variance $`\mathrm{\Sigma }_2`$ are well described by the Wigner surmise $`P_W(s)`$ and $`\mathrm{\Sigma }_2^W(E)`$ respectively near the intermediate interaction strength $`U_c120`$ for which the TIP system does not have a preferential eigenbasis (see Fig. 2). As shown by $`\mathrm{\Sigma }_2(E)`$ (Fig. 7), Wigner-Dyson spectral rigidity is established over an energy interval containing a few levels. In contrast, Hubbard repulsion can only yield tip2 at the corresponding $`U_c`$ ($`1`$)
$$P_{SP}4s\mathrm{exp}(2s)$$
(19)
and
$$\mathrm{\Sigma }_20.16+0.41E$$
(20)
respectively.
To study how the spacing distribution depends on $`U`$, we use the spectral parameter $`\eta `$ defined by
$$\eta (P,U)=\frac{_0^b𝑑s[P(s)P_W(s)]}{_0^b𝑑s[P_P(s)P_W(s)]}$$
(21)
with $`b=0.4729`$. For $`U=0`$, the consecutive levels are essentially uncorrelated and $`\eta 1`$. The curves $`\eta (U)`$ given in Fig. 9 show us a striking difference between the spectral statistics yielded by the two repulsions.
For $`L=L_1=100`$, let us estimate the interaction threshold $`U^{}`$ for which the interaction matrix elements coupling consecutive 2FP states becomes of the order of their energy separation $`\mathrm{\Delta }_2`$. $`\rho _20.25`$ (density normalized to $`1`$) gives $`\mathrm{\Delta }_2(\rho _2L(L+1)/2)^14/5050`$ (see insert of Fig. 1) for the energy spacing between consecutive 2FP levels around $`E=0`$. The off-diagonal interaction matrix elements $`Q=H_{\alpha \beta }^{\gamma \delta }`$ (Hubbard) or $`Q=C_{\alpha \beta }^{\gamma \delta }`$ (Coulomb) coupling consecutive 2FP levels are normally distributed, as shown in Fig. 8. The root mean square of $`Q`$ is small for Hubbard ($`0.0038`$) and larger for Coulomb ($`0.01`$). This gives $`U_C^{}0.08`$ for Coulomb and $`U_H^{}0.2`$ for Hubbard.
When $`U`$ increases but remains lower than $`U^{}`$, there is first a perturbative regime discussed in Refs. wp ; wpi where the interaction yields Rabi oscillations between consecutive 2FP states at a frequency given by the absolute value of the coupling interaction matrix element. Moreover, the energy range $`E_U`$ under which one has level repulsion is given wp by this Rabi frequency. When $`U`$ becomes equal to $`U^{}`$, one has a transition from this perturbative regime towards an effective Fermi golden rule decay of the 2FP states and the characteristic range $`E_U`$ over which Wigner-Dyson rigidity occurs becomes wp ; wpi proportional to the square of the amplitude of the coupling matrix elements: $`E_UU`$ when $`U<U^{}`$ and $`E_UU^2`$ when $`U>U^{}`$. $`U=U^{}`$ is the interaction threshold where the spectral statistics are intermediate between Poisson and Wigner ($`\eta 0.39`$ when $`P(s)P_{SP}(s)`$). Looking at Fig. 9, one can see that $`\eta `$ decreases as a function of $`U`$ down to the characteristic value $`\eta ^{}0.39`$ reached when $`UU_C^{}0.08`$ for Coulomb repulsion, and when $`UU_H^{}0.2`$ for Hubbard repulsion.
For Coulomb repulsion, $`\eta `$ continues to decrease when $`UU_C^{}`$ down to $`\eta =0`$ with the slower $`U`$ dependence characteristic of the effective golden rule decay of the 2FP states. The Wigner-Dyson distribution $`P_W(s)`$ is fully established at $`UU_c`$, i.e. when the TIP system is exactly as far from the $`U=0`$ eigenbasis than from the $`U=\mathrm{}`$ eigenbasis. When $`UU_c`$, integrability is slowly restored and $`\eta 1`$ as $`U\mathrm{}`$.
For Hubbard repulsion, the spectral rigidity does not continue to increase above $`U_H^{}`$, but saturates to the intermediate critical rigidity characterized by Eq. 19 and Eq. 20 for $`P(s)`$ and $`\mathrm{\Sigma }_2(E)`$ respectively. Above the fixed point $`U_c1`$ of the duality transformation, the TIP system becomes closer to the $`U=\mathrm{}`$ eigenbasis and the levels become statistically uncorrelated. This critical Hubbard regime is a complicated issue where the multifractal character of the interaction matrix described in ref. tip1 is very likely relevant. However, one can do the following remark. When $`L_1\mathrm{}`$, the interaction matrix is block diagonal, a block corresponding to a pair momentum $`K`$. For Hubbard, the $`N_s`$ TIP level $`E_n(K)`$ of momentum $`K`$ are located near the 2FP levels of same momentum when $`4U/L0`$ or $`4U/L\mathrm{}`$. One can assume that they should be near the middle of consecutive 2FP levels of same $`K`$ for $`UU_c`$. The distribution $`P(s)`$ of levels located in the middle of levels with spacing distribution $`P_P(s)`$ is the semi-Poisson distribution $`P_{SP}(s)`$. If the sequence of 2FP levels of momentum $`K`$ were randomly distributed, the spacing distribution of the $`N_s`$ TIP levels should be given by $`P_{SP}(s)`$ without disorder. On the contrary, the interaction matrix being more random for Coulomb, one can expect that the $`L`$ series of $`N_s`$ levels of momentum $`K`$ will be driven towards Wigner-Dyson statistics as $`U`$ increases. However, though the argument may give hints for the existence of the Hubbard $`P_{SP}(S)`$, it does not explains the behavior of $`\mathrm{\Sigma }_2(E)`$. The breakdown of momentum conservation by the disorder, and the associated mixing of the $`L`$ independent series of $`N_s`$ levels characterizing the clean limit plays a complex role.
The insert of Fig. 6 shows how the spectrum becomes more rigid at $`U=70`$ when the range $`p`$ of the interaction is increased. $`\mathrm{\Sigma }_2(E)`$ displays a similar information in Fig. 7.
## 8 Quantum melting for intermediate Coulomb repulsions
The intermediate Wigner-Dyson regime yielded by Coul- omb repulsion corresponds, inside a scale $`L_1`$, to a complete melting of the localized 2FP states previous to crystallization. To show this, we introduce two parameters $`\gamma `$ and $`\xi `$. For a TIP wavefunction $`|\mathrm{\Psi }`$, we calculate the density $`\rho _i=\mathrm{\Psi }|c_i^{}c_i|\mathrm{\Psi }`$ at site $`i`$ and the density density correlation function $`C(r)=(1/2)_i\rho _i\rho _{ir}`$. The participation ratio $`\xi `$ (i. e. the number of sites occupied by a TIP state) is given by $`2C(0)^1`$. The crystallization parameter $`\gamma `$ is given by the difference $`MaxC(r)MinC(r)`$, where all translations $`r`$ (including $`r=0`$) are considered. If the electron density is homogeneous (as for an extended liquid state) $`\gamma 0`$ whereas $`\gamma 1`$ if the two charges are mainly located on two different sites, a situation occuring when $`U\mathrm{}`$ (Coulomb “molecule”) and when $`U0`$ and $`L_10`$ (two electrons located in two minima of the potential). The variations of $`\xi `$ and $`\gamma `$ are given in Fig. 10 when Coulomb repulsion $`U`$ increases. The curves $`\gamma (U)`$ and $`\xi (U)`$ are correlated to the curve $`\eta (U)`$ shown in Fig. 9. The curve $`\xi (U)`$ shows us that the TIP states occupy a fraction of the chain without repulsion before being uniformly spread over a scale $`L_1`$ at interaction strength $`UU_c`$ for which there is a crossing of the two curves $`PR(U)`$ given in Fig. 2. For $`L=L_1`$, one has an interaction induced quantum melting of the non interacting glass, yielding quantum ergodicity with “ more liquid” and extended wavefunctions and Wigner-Dyson spectral statistics.
## 9 Conclusions
We have studied one of the simplest problems where quantum localization and two body interaction are in competition. We have seen that the range of the interaction makes important differences. One of them is revealed by the study of the spectral statistics, providing an intriguing puzzle for quantum ergodicity: the Wigner-Dyson statistics shrinking to intermediate statistics when the two body repulsion becomes local. However, the generic behavior of the TIP system can be summarized by three regimes, independently of the interaction range. There is the free particle limit dominated by quantum localization, the Coulomb limit dominated by the pinning of a correlated system of charges (“Coulomb molecule, Wigner crystal”) or by quantum localization again (Hubbard repulsion). Between these two limits, there is an intermediate regime where one has a maximum mixing of the one body states, making the states more extended and the spectrum more rigid. In one dimension, this yields a partial delocalization effect ($`L_2`$ could be large but remains finite). Similar conclusions have been reached from a study of the many body ground state sjwp of one dimensional spinless fermions at half filling. In two dimensions, there are experimental evidences kravchenko that localization may disappear for intermediate Coulomb energy to Fermi energy ratios $`r_s`$. A new quantum regime has been observed bwp for the ground state of two dimensional spinless fermions at intermediate factors $`r_s`$ where the metallic phase is observed. The nature of the states, when the system is far from the free particle (Fermi glass) or the Coulomb (Wigner crystal) bases, and the associated transport mechanism, remain to be understood. This simple one dimensional study draws our attention to the important role played by the range of electron-electron repulsions. To know how is or is not screened Coulomb repulsion in low density electron systems is then an important issue. |
warning/0003/hep-ph0003276.html | ar5iv | text | # Could a Slepton be the Zee boson?
## 1 Zee Neutrino Mass Model
An economical way to generate small neutrino masses with a phenomenologically favorable texture is given by the Zee model, which generates masses via one-loop diagrams. The model consists of a charged singlet scalar $`h_{\text{Zee}}^\text{-}`$, the Zee scalar, which couples to lepton doublets $`\psi _{Lj}`$ via the interaction
$$f^{ij}\left(\psi _{Li}^\alpha 𝒞\psi _{Lj}^\beta \right)ϵ_{\alpha \beta }h_{\text{Zee}}^\text{-},$$
(1)
where $`\alpha ,\beta `$ are the $`SU(2)`$ indices, $`i,j`$ are the generation indices, $`𝒞`$ is the charge-conjugation matrix, and $`f^{ij}`$ are Yukawa couplings antisymmetric in $`i`$ and $`j`$. The latter fact is a result of the $`SU(2)`$ product rule and is central to the favorable texture obtained. Another ingredient of the Zee model is an extra Higgs doublet (in addition to the one that gives masses to charged leptons) that develops a vacuum expectation value (VEV) and thus provides a mass mixing between the charged Higgs boson and the Zee scalar boson. The corresponding coupling, together with the $`f^{ij}`$’s, enforces lepton number violation.
A recent analysis by Frampton and Glashow (see also Ref.) showed that the Zee mass matrix of the following texture
$$\left(\begin{array}{ccc}0& m_{e\mu }& m_{e\tau }\\ m_{e\mu }& 0& ϵ\\ m_{e\tau }& ϵ& 0\end{array}\right),$$
(2)
where $`ϵ`$ is small compared with $`m_{e\mu }`$ and $`m_{e\tau }`$, is able to provide a compatible mass pattern that explains the atmospheric and solar neutrino data. The generic Zee model guarantees the vanishing of the diagonal elements, while the suppression of the $`m_{\mu \tau }`$ entry, here denoted by the small parameter $`ϵ`$, has to be otherwise enforced. Moreover, $`m_{e\mu }m_{e\tau }`$ is required to give the maximal mixing solution for the atmospheric neutrino data.
Setting $`ϵ`$ to zero in the above mass matrix gives the following (zeroth order) result: one linear combination of $`\nu _\mu `$ and $`\nu _\tau `$ remains massless while the orthogonal state forms a Dirac pair with $`\nu _e`$. Oscillation between the Dirac pair and the massless Majorana state could explain the atmospheric neutrino data. If we restore a nonzero $`ϵ`$, or, for that matter, put in some other perturbation to the above mass matrix instead, we have the first order result: namely, we have an additonal pseudo-Dirac splitting between the massive states that could explain the solar neutrino data. To be exact, the original massless state would then also gain a tiny mass. This generalized Zee mass texture is what we will aim at in our discussion of supersymmetric version(s) of the Zee model below.
## 2 Zee Model in a Supersymmetric Framework
To take the Zee model seriously, one have to put it together with other aspects of beyond-standard-model physics. Hence it worths considering putting the model in a supersymmetric framework. The topic is studied in our recent paper, upon which the present report is based.
A naive idea along the line would be to add what is needed in the Zee mechanism to the minimal supersymmetric standard model (MSSM). However a right-handed slepton ($`\stackrel{~}{\mathrm{}}_R`$) actually has the same gauge quantum number as the Zee scalar boson, hence the question of our title — if the slepton could take the role. It is clear that, lepton number and, generically R-parity, has to be violated here. The MSSM spectrum even provides the second Higgs doublet needed. A careful study shows that all we need to complete the SUSY-Zee diagram is the following minimal set of only three R-parity violating (RPV) couplings:
$$\{\lambda _{12k},\lambda _{13k},\mu _k\}$$
where family index $`k`$ can be chosen arbitrarily with $`\stackrel{~}{\mathrm{}}_{R_k}`$ assuming the role of the Zee boson. The corresponding diagram is given in Fig.1. Note that our RPV parameters are defined in a flavor basis where the so-called sneutrino VEV’s are rotated away (see Refs. for more details).
It is interesting to note that the SUSY-Zee contribution to neutrino masses involves RPV parameters of both the bilinear and trilinear type. The existence of such a contribution under SUSY without R-parity has not been realized before our work. In fact, there is another contribution of the kind that exists even within the present minimal framework and, in that case, contributes to the same neutrino mass entries ($`m_{e\mu }`$ and $`m_{e\tau }`$) also first identified in Ref.. This is shown in Fig.2.
However, the above minimal set of RPV couplings also gives rise to other contributions that could potentially spoil the Zee mass texture. These more widely studied contributions tend to give diagonal mass matrix entries at least the same strength as the off-diagonal ones, as illustrated in Ref.. To retain the successful flavor of the Zee model, one has to go to a region of the parameter space where the SUSY-Zee contribution dominates hence giving the zeroth order texture with the sub-dominating contributions fitting into a successful final result.
The best scenario here is for $`k=3`$, i.e. making the right-handed stau the Zee boson. We give here the neutrino mass results and summarize the required conditions as follows:
$`m_{ee}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi ^2}}\left(A_\tau ^E\mu \mathrm{tan}\beta \right)f(M_{\stackrel{~}{\tau }_L}^2,M_{\stackrel{~}{\tau }_R}^2)m_\tau ^2\lambda _{133}^2;`$ (3)
$`m_{\tau \tau }`$ $`=`$ $`{\displaystyle \frac{v^2\mathrm{cos}^2\beta \left(g^2M_1+g^{}_{}{}^{}2M_2\right)}{2\mu \left[2\mu M_1M_2v^2\mathrm{sin}\beta \mathrm{cos}\beta \left(g^2M_1+g^{}_{}{}^{}2M_2\right)\right]}}\mu _3^2;`$ (4)
$`m_{e\mu }`$ $`=`$ $`{\displaystyle \frac{1}{16\pi ^2}}{\displaystyle \frac{\sqrt{2}\mathrm{tan}\beta }{v\mathrm{cos}\beta }}f(M_{h_1^\text{-}}^2,M_{\stackrel{~}{\tau }_R}^2)m_\tau \left(m_\mu ^2m_e^2\right)\mu _3\lambda _{123}`$ (5)
$`{\displaystyle \frac{1}{16\pi ^2}}{\displaystyle \frac{v\mathrm{sin}\beta }{\sqrt{2}}}f(M_{\stackrel{~}{e}_L}^2,M_{\stackrel{~}{\tau }_R}^2)m_\tau \mu _3\lambda _{123}\lambda _{133}^2;`$
$`m_{e\tau }`$ $`=`$ $`{\displaystyle \frac{1}{16\pi ^2}}{\displaystyle \frac{\sqrt{2}\mathrm{tan}\beta }{v\mathrm{cos}\beta }}f(M_{h_1^\text{-}}^2,M_{\stackrel{~}{\tau }_R}^2)m_\tau \left(m_\tau ^2m_e^2\right)\mu _3\lambda _{133}`$ (6)
$`{\displaystyle \frac{1}{16\pi ^2}}{\displaystyle \frac{v\mathrm{sin}\beta }{\sqrt{2}}}f(M_{\stackrel{~}{e}_L}^2,M_{\stackrel{~}{\tau }_R}^2)m_\tau \mu _3\lambda _{133}^3;`$
where $`f(X,Y)=\frac{1}{XY}\mathrm{log}\frac{X}{Y}`$; and with $`m_{\mu \mu }`$ and $`m_{\mu \tau }`$ being zero.
The conditions for the success of the scenario are given by :-
$`\lambda _{133}`$ $``$ $`{\displaystyle \frac{m_\mu ^2}{m_\tau ^2}}\lambda _{123};`$ (7)
$`(\mu _3\mathrm{cos}\beta )\lambda _{123}`$ $``$ $`\mathrm{cos}^3\beta \text{Max}(M_{h_1^\text{-}}^2,M_{\stackrel{~}{\tau }_R}^2)(710^5\text{GeV}^1);`$ (8)
$`\mu _3^2\mathrm{cos}^2\beta `$ $``$ $`\mu ^2M_1(110^{14}\mathrm{GeV}^1);`$ (9)
$`\lambda _{133}^2`$ $``$ $`{\displaystyle \frac{\text{Max}(M_{\stackrel{~}{\tau }_L}^2,M_{\stackrel{~}{\tau }_R}^2)}{(A_\tau ^E\mu \mathrm{tan}\beta )}}(2.510^9\text{GeV}^1);`$ (10)
$`{\displaystyle \frac{\lambda _{123}^2}{M_{\stackrel{~}{\tau }_R}^2}}`$ $``$ $`10^8\text{GeV}^2.`$ (11)
The feasibility of the scenario is marginal. Alternative versions with some extra chiral superfields, introduced and briefly discussed in Ref., are much less contrained though. |
warning/0003/gr-qc0003055.html | ar5iv | text | # Cosmological Time in (2+1)-Gravity
## 1 Introduction.
We shall be mainly concerned with maximal globally hyperbolic, matter-free spacetimes $`M`$ of topological type $`S\times `$, where $`S`$ is a compact closed oriented surface of genus $`g>1`$. The (2+1)-dimensional Einstein equation with vanishing cosmological constant actually implies that $`M`$ is (Riemann) flat.
After \[D-J-’t H\] and \[W\], a large amount of literature has grown up about this $`(2+1)`$-gravity topic, regarded as a useful toy-model for the higher dimensional case. Two main kinds of description have been experimented. A “cosmological” approach points to characterize the spacetimes in terms of some distinguished global time; for instance the constant mean curvature CMC time has been widely studied \[A-M-T\], \[Mo\]. A “geometric” time-free approach eventually identifies each flat spacetime by means of its $`ISO(2,1)`$-valued holonomy \[W\], \[Me\]. With the exception of the case with toric space ($`g=1`$), there is not a clear correspondence between the results obtained in these two approaches.
The aim of this paper is to show that this gap can be filled by using the canonical Cosmological Time CT, that is “the length of time that the events of $`M`$ have been in existence” (see \[A-G-H\]). It turns out that this is a global time which reveals the fundamental properties of spacetime. It is canonically defined by means of the very basic spacetime’s structures: its casual structure and the Lorentz distance. The cosmological time $`\tau `$ is invariant under diffeomorphisms, therefore the $`\tau =a`$ level surfaces $`S_a`$ provide a gauge-invariant description of space evolution in $`M`$. Both the intrinsic and extrinsic geometry of the surfaces $`S_a`$, as well as their past/future asymptotic states, are intrinsic features of spacetime. The asymptotic states are defined by the evolution of the observables associated to the length of closed geodesic curves on the surfaces $`S_a`$. Remarkably, they recover and decouple the linear and the translational parts of the holonomy. The study of the asymptotic states also leads to understand the initial singularity (we will always assume that the space is future expanding) and the way how the classical geometry degenerates, but does not completely disappear. The initial singularity can be interpreted as the isometric action of the fundamental group of $`S`$ on a suitable “real tree”. Differently to the case of the CMC time (for instance), the level surfaces $`S_a`$ of the CT are in general only $`C^1`$-embedded into the spacetime $`M`$. This lack of smoothness takes place on a “geodesic lamination” on $`S_a`$ and is a observable large scale manifestation of the intrinsic geometry of the initial singularity. Thus the initial singularity admits two complementary descriptions: one, in terms of real trees and, the second, in terms of geodesic laminations. The existence of a duality relation between real trees and laminations was already known in the context of Thurston theory of the boundary of the Teichmüller space. It is remarkable that Einstein theory of (2+1)-gravity sheds new light on this subject and puts duality in concrete form.
In \[B-G2\] we have also used the cosmological time in order to study certain interesting families of $`(2+1)`$-spacetimes coupled to particles.
Our main purpose consists of elucidating the central role of the cosmological time and its asymptotic states in the description of spacetimes. The cosmological time perspective provides a new interpretation of several facts spread in the literature which are related to Thurston work. More precisely, the present article is based on, and could be considered a complement of, Mess’s fundamental paper \[Me\].
## 2 The Cosmological Time Function.
For the basic notions of Lorentzian geometry and causality we refer for instance to \[B-E\], \[H-E\]. Let $`N`$ be any time oriented Lorentzian manifold of dimension $`n+1`$. The cosmological time function, $`\tau :N(0,\mathrm{}]`$, is defined as follows. Let $`C^{}(q)`$ be the set of past-directed causal curves in $`N`$ that start at $`qN`$, then
$$\tau (q)=\mathrm{sup}\{L(c):cC^{}(q)\}$$
where $`L(c)`$ denotes the Lorentzian length of the curve $`c`$:
$$L(c)=_c(\mathrm{proper}\mathrm{time}).$$
$`\tau (q)`$ can be interpreted as the length of time the event $`q`$ has been in existence in $`N`$. For example, if $`N`$ is the standard flat Minkowski space $`𝕄^{n+1}`$, $`\tau `$ is the constant $`\mathrm{}`$-valued function, so in this case it is not very interesting. In \[A-G-H\] (see also \[W-Y\]) one studies the properties of a manifold $`N`$ with regular cosmological time function. Recall that $`\tau `$ is regular if:
1) $`\tau (q)`$ is finite valued for every $`qN`$;
2) $`\tau 0`$ along every past directed inextensible causal curve.
The existence of a regular cosmological time function has strong consequences on the structure of $`N`$ and of the constant-$`\tau `$ surfaces \[A-G-H\]. In particular when $`\tau `$ is regular, $`\tau :N(0,\mathrm{}[`$ is a continuous function, which is twice differentiable almost everywhere, giving a global time on $`N`$ denoted by CT. Each $`\tau `$ level surface is a future Cauchy surface, so that $`N`$ is globally hyperbolic. For each $`qN`$ there exists a future-directed time-like unit speed geodesic ray $`\gamma _q:(0,\tau (q)]N`$ such that:
$$\gamma _q(\tau (q))=q,\tau (\gamma _q(t))=t.$$
The union of the past asymptotic end-points of these rays can be regarded as the initial singularity of $`N`$.
The cosmological time function is not related to any specific choice of coordinates in $`N`$; it is “gauge-invariant” and so it represents an intrinsic feature of spacetime. Thus, when the cosmological time is regular, the $`\tau `$-constant level surfaces and their properties have a direct physical meaning as they are observables.
We present now two basic examples of spacetime with regular cosmological time, which shall be important throughout all the paper. To fix the notations, the standard Minkowski space $`𝕄^{2+1}`$ is endowed with coordinates $`x=(x^1,x^2,x^3)`$, so that the metric is given by $`ds^2=(dx^1)^2+(dx^2)^2(dx^3)^2`$. $`𝕄^{2+1}`$ is oriented and time-oriented in the usual way.
Example 1. Consider the chronological future of the origin $`0𝕄^{2+1}`$
$$I^+(0)=\{x𝕄^{2+1}:(x^1)^2+(x^2)^2(x^3)^2<0,x^3>0\}.$$
Its cosmological time, $`\tau :I^+(0)(0,\mathrm{})`$, is a smooth submersion; the constant-time $`\{\tau =a\}`$ surfaces are the (upper) hyperboloids
$$𝕀(a)=\{x𝕄^{2+1}:(x^1)^2+(x^2)^2(x^3)^2=a^2,x^3>0\}.$$
Hence $`𝕀(a)`$ is a complete space of constant Gaussian curvature equal to $`1/a^2`$, and of constant extrinsic mean curvature $`1/a`$. The Lorentzian length of the time-like geodesic arc connecting any $`pI^+(0)`$ with $`0`$ equals $`\tau (p)`$; $`0`$ is the initial singularity. Note that $`𝕀(a)`$ can be obtained from $`𝕀(1)`$ by means of a dilatation in $`𝕄^{2+1}`$ with constant factor $`a`$; shortly we write $`𝕀(a)=a𝕀(1)`$. We shall denote by $`SO(2,1)`$ the group of oriented Lorentz transformations acting on $`𝕄^{2+1}`$ and by $`ISO(2,1)`$ the Poincaré group. $`SO^+(2,1)`$ denotes the subgroup of $`SO(2,1)`$ transformations which keep $`I^+(0)`$ and each $`𝕀(a)`$ invariant. $`ISO^+(2,1)`$ is the corresponding subgroup of $`ISO(2,1)`$.
Example 2. Let us denote by $`I^+(1,3)`$ the chronological future in $`𝕄^{2+1}`$ of the line $`\{x^1=x^3=0\}`$
$$I^+(1,3)=\{x𝕄^{2+1}:(x^1)^2(x^3)^2<0,x^3>0\}.$$
The Lorentzian length of the time-like geodesic arc connecting any $`p=(x^1,x^2,x^3)I^+(1,3)`$ with $`q=(0,x^2,0)`$ equals the cosmological time $`\tau (p)`$. The level surfaces are
$$𝕀(1,3,a)=\{x𝕄^{2+1}:(x^1)^2(x^3)^2=a^2,x^3>0\}$$
and have constant extrinsic mean curvature equal to $`(1/2a)`$. Each surface $`𝕀(1,3,a)`$ is isometric to the flat plane $`^2`$. To make this manifest, it is useful to consider the following change of coordinates. Let $`\mathrm{\Pi }^{2+1}=\{(u,y,\tau )^{2+1}:\tau >0\}`$ be endowed with the metric $`ds^2=\tau ^2du^2+dy^2d\tau ^2`$. Then, $`x^1=\tau sh(u),x^2=y,x^3=\tau ch(u)`$, is an isometry between $`\mathrm{\Pi }^{2+1}`$ and $`I^+(1,3)`$. The level set $`\{\tau =a\}`$ of $`\mathrm{\Pi }^{2+1}`$ goes isometrically onto $`𝕀(1,3,a)`$, so this is intrinsically flat. Note that the group of oriented isometries of $`\mathrm{\Pi }^{2+1}`$ is generated by the translations parallel to the planes $`\{\tau =a\}`$, and the rotation of angle $`\pi `$ of the $`(u,y)`$ coordinates.
We are going to show that any maximal globally hyperbolic, matter-free $`(2+1)`$-spacetime $`M`$, with compact space $`S`$, actually has regular cosmological time, and its initial singularity can be accurately described.
## 3 Flat (2+1)-spacetimes
A flat spacetime is, by definition, locally isometric to the Minkowski space $`𝕄^{2+1}`$. We assume that our maximal hyperbolic flat spacetimes are time-oriented and future expanding, and that these orientations locally agree with the usual ones on $`𝕄^{2+1}`$. The spacetime structures on $`S\times `$ are regarded up to oriented isometry homotopic to the identity.
### 3.1 Minkowskian suspensions
We introduce here the simplest $`(2+1)`$ spacetimes with compact space $`S`$ of genus $`g>1`$.
Recall that the upper hyperboloid $`𝕀(1)𝕄^{2+1}`$, mentioned in the previous section, is a classical model for the hyperbolic plane $`^2`$ (see \[B-P\] for this and other models); the Poincaré disk is another model which can be obtained from $`𝕀(1)`$ by means of the stereographic projection shown in figure 1. We shall use the Poincaré model in section 4.
Take any hyperbolic surface $`F=^2/\mathrm{\Gamma }`$ homeomorphic to $`S`$. $`\mathrm{\Gamma }`$ is a subgroup of $`SO^+(2,1)`$ which acts freely and properly discontinuously on $`^2𝕀(1)`$. $`𝕀(1)`$ can be identified with the universal covering of $`S`$ and $`\mathrm{\Gamma }`$ with the fundamental group $`\pi _1(S)`$. $`\mathrm{\Gamma }`$ can be thought also as a group of isometries of the spacetime $`I^+(0)`$ and $`M(F)=I^+(0)/\mathrm{\Gamma }`$ is the required spacetime with compact space homeomorphic to $`S`$. We call it the Minkowskian suspension of $`F`$. This construction is well-known; sometimes $`M(F)`$ is also called the Lorentzian cone over $`F`$ or the Löbell spacetime based on $`F`$. $`I^+(0)`$ can be regarded as the universal covering of $`M(F)`$. Let us now consider the cosmological time of $`M(F)`$. The CT of $`I^+(0)`$ naturally induces the CT of $`M(F)`$. Indeed, each level surface $`S_a`$ of $`M(F)`$ has $`𝕀(a)`$ as universal covering; moreover, $`S_1=F`$ and $`S_a=aF`$. In this case, the CT coincides with the CMC time and each level surface $`S_a`$ smoothly embeds into $`M(F)`$. The initial singularity “trivially” consists of one point.
Notation. Let $`Y`$ be any subset of $`𝕀(1)`$, we shall denote by $`\widehat{Y}`$ its “suspension” in $`I^+(0)`$ which is defined by $`\widehat{Y}=_{a(0,\mathrm{}[}aY`$.
### 3.2 (2+1) spacetimes as deformed Minkowskian suspensions
It has been shown in \[Me\] that any maximal globally hyperbolic, future expanding flat spacetime $`M`$ with compact space homeomorphic to $`S`$, as above, can be regarded as a “deformation” of some Minkowskian suspension (see also \[W\]). In fact $`M`$ is of the form $`M=U(M)/\mathrm{\Gamma }^{}`$, where:
1) The domain $`U(M)`$ of $`𝕄^{2+1}`$ is a convex set
$$U(M)=\{x𝕄^{2+1};x^3>f(x)\}$$
where $`f:\{x^3=0\}[0,\mathrm{}[`$ is a convex function.
2) $`\mathrm{\Gamma }^{}`$ is a subgroup of $`ISO^+(2,1)`$ (also called the holonomy group of $`M`$) acting freely and properly discontinuously on $`U(M)`$. Hence $`U(M)`$ is the universal covering of $`M`$ and $`\mathrm{\Gamma }^{}`$ is isomorphic to $`\pi _1(M)\pi _1(S)`$.
3) The “linear part” $`\mathrm{\Gamma }`$ of $`\mathrm{\Gamma }^{}`$ is a subgroup of $`SO^+(2,1)`$ which is isomorphic to $`\pi _1(S)`$ and acts freely and properly discontinuously on $`𝕀(1)^2`$. This is a non trivial fact which follows from a result of Goldman \[Go\]. Each element $`\gamma ^{}\mathrm{\Gamma }^{}`$ is of the form $`\gamma ^{}=\gamma +t(\gamma )`$, where $`\gamma \mathrm{\Gamma }`$ and $`t(\gamma )^3`$ is a translation. $`t:\mathrm{\Gamma }^3`$ is a cocycle representing an element of $`H^1(\mathrm{\Gamma },^3)`$. If $`t^{}=\lambda t`$, $`\lambda ^{}`$, then $`U(M^{})`$ differs from $`U(M)`$ by : $`U(M)=\lambda ^1U(M^{})`$. When $`\lambda `$ is “small”, $`U(M^{})`$ is “close” to $`I^+(0)`$ ($`M^{}`$ is “close” to $`M(F)`$, $`F=^2/\mathrm{\Gamma }`$).
Note that $`\mathrm{\Gamma }^{}`$, whence $`U(M)`$ and $`t`$, are defined up to inner automorphism of $`ISO^+(2,1)`$.
### 3.3 Spacetimes of simplicial type
In this section, we shall consider the flat spacetimes that can be obtained from Minkowskian suspensions by means of particular deformations. These spacetimes will be called of simplicial type, the origin of this name is related to the material presented in section 4. Spacetimes of simplicial type are important because they are “dense” in the set of all spacetimes; the shape of any spacetime and of its CT can be arbitrarily well approximated by some spacetime of simplicial type (see proposition 4.23). So, it is enough to understand these examples in order to have a rather complete qualitative picture of our general presentation. Moreover, all the statements of this paper can easily be checked in a spacetime of simplicial type.
Start with a Minkowskian suspension $`M(F)`$. Assume that a weighted multi-curve $``$ on $`F`$ is given. $``$ is the union of a finite number of disjoint simple closed geodesics on $`F`$, each one endowed with a strictly positive real weight. $``$ governs a specific deformation of $`M(F)`$ producing a required flat spacetime denoted by $`M(F,)`$. A particular spacetime deformation is associated to each component of $``$ and can be obtained by means of an appropriate surgery operation in Minkowski space. As the deformations associated to the components of $``$ act locally and independently from each other, we may assume for simplicity that $``$ just consists of one component $`c`$, with weight $`r`$ and length $`s`$.
Elementary deformation. In order to illustrate the deformation associated with one geodesic $`c`$ with weight $`r`$, we shall now introduce a simple hyperbolic surface $`F_0`$ which can be understood as a local model for the general surface $`S`$. Let $`\mathrm{\Gamma }_0`$ be an infinite cyclic subgroup of $`SO^+(2,1)`$ generated by an element $`g_0`$ acting on $`𝕀(1)`$ as an isometry of hyperbolic type (see for instance \[B-P\] for the classification of the isometries of $`^2`$). We can assume that $`g_0`$ is a Lorentz transformation corresponding to a boost along the $`x^1`$-direction, so that the $`g_0`$-invariant geodesic line on $`𝕀(1)`$ is the line $`\sigma _0=𝕀(1)\{x^2=0\}`$. The hyperbolic surface $`F_0=𝕀(1)/\mathrm{\Gamma }_0`$ is homeomorphic to the non-compact annulus $`S^1\times `$ and its area is not finite. The image in $`F_0`$ of the axis of $`g_0`$ is a simple closed geodesic $`c`$ of a certain length $`s`$; give it the positive weight $`r`$. So we dispose of a one-component weighted multi-curve $`_0`$ on $`F_0`$, as illustrated in figure 2.
The suspension $`M(F_0)=I^+(0)/\mathrm{\Gamma }_0`$ is a flat spacetime. Let us now construct $`M_0=M(F_0,_0)`$ which represents the deformation of $`M(F_0)`$ associated to the weighted multi-curve $`_0`$. We shall use the spacetimes $`I^+(0)`$, $`I^+(1,3)`$ and $`\mathrm{\Pi }^{2+1}`$ that we have introduced in section 2. The universal covering $`U(M_0)`$ of $`M_0`$ will be the union of three domains of $`𝕄^{2+1}`$: $`U(M_0)=ABC`$, where $`A=I^+(0)\{x^20\}`$, $`B=I^+(1,3)\{0x^2r\}`$, $`C=C^{}+r(0,1,0)`$ and $`C^{}=I^+(0)\{x^20\}`$. In our notations, $`C^{}+r(0,1,0)`$ denotes the set of points in $`𝕄^{2+1}`$ which can be obtained from $`C^{}`$ by means of a translation of length $`r`$ along the unit vector $`(0,1,0)`$. It is important to note that the cosmological times of the different pieces $`A`$, $`B`$ and $`C`$ fit well together; in fact, the CT level surfaces $`\stackrel{~}{S}_a`$ of $`U(M_0)`$ are
$$\stackrel{~}{S}_a=(a𝕀(1)\{x^2<0\})(𝕀(1,3,a)\{0x^2r\})(a𝕀(1)\{x^2>0\}+r(0,1,0)).$$
As shown in figure 3, each surface $`\stackrel{~}{S}_a𝕄^{2+1}`$ can be obtained by cutting the hyperboloid $`𝕀(a)`$ along $`a\sigma _0`$ (which is the intersection of $`𝕀(a)`$ with the $`\{x^2=0\}`$-plane) and then by inserting a band of $`𝕀(1,3,a)`$ of depth $`r`$.
The surfaces $`\stackrel{~}{S}_a`$ are only $`C^1`$-embedded into $`U(M_0)`$. The initial singularity of $`U(M_0)`$ is the segment $`J_0=\{x^1=x^3=\mathrm{0\hspace{0.17em}\; 0}x^2r\}`$.
###### Remark 3.1
We have the following characterization of $`J_0`$. The interior points of this segment make the subset of $`U(M_0)`$ (boundary of the convex set $`U(M_0)`$) of the points with exactly two null supporting planes; the end-points make the subset of $`U(M)`$ with more than two null supporting planes. Recall that a supporting plane at $`xU(M_0)`$ is a plane $`P`$ such that $`xP`$ and $`U(M_0)P=\mathrm{}`$. $`P`$ is null if it contains some null-lines.
The covering $`U(M_0)𝕄^{2+1}`$ is flat. To get $`M_0`$, we only need to specify the action of $`\pi _1=\pi _1(F_0)`$ on $`U(M_0)`$.
Action of the fundamental group. $`\pi _1`$ acts on $`A`$ by the restriction of the action of $`\mathrm{\Gamma }_0`$ on $`𝕀(1)`$. The domain $`B`$ corresponds (via the isometry established in section 2) to $`B^{}=\{(u,y,\tau )\mathrm{\Pi }^{2+1};0yr\}`$, so that the action of $`\pi _1`$ on $`B`$ transported on $`B^{}`$ is just given by the translation $`(u,y,\tau )(u+s,y,\tau )`$. Finally, if $`\alpha `$ is the translation $`(x^1,x^2,x^3)(x^1,x^2+r,x^3)`$ on $`𝕄^{2+1}`$, then the action of $`\pi _1`$ on $`C`$ is just the conjugation of $`\mathrm{\Gamma }_0`$ by $`\alpha `$.
The CT of the covering $`U(M_0)`$ passes to the quotient $`M_0=U(M_0)/\pi _1`$; each level surface $`S_a`$ is only $`C^1`$-embedded into $`M_0`$, so that it is endowed with an induced $`C^1`$-Riemannian metric. This allows anyway to define the length of curves traced on the surface $`S_a`$ and the derived length-space distance. Let $`𝒜=A/\pi _1`$, $`=B/\pi _1`$ and $`𝒞=C/\pi _1`$. Then, $`S_a`$ is a flat annulus of depth $`r`$ and parallel geodesic boundary components of length $`as`$; $`S_a(𝒜𝒞)`$ can be isometrically embedded into $`aF_0`$, and has geodesic boundary curves of length $`as`$. As shown in figure 4, $`S_a`$ can be obtained by cutting $`F_0`$ along $`c`$ and by inserting a annulus of depth $`r`$.
###### Remark 3.2
If $`g`$ is an element of $`ISO^+(2,1)`$ acting on $`X𝕄^{2+1}`$ as $`g(X)=QX+w`$, the transformed domain $`g(U(M_0))=Q(ABC)+w`$ is, of course, an isometric copy of the universal covering in $`𝕄^{2+1}`$. The curve $`\sigma =Q(\sigma _0)`$ is a geodesic line of $`𝕀(1)`$; $`\sigma `$ is the intersection of $`𝕀(1)`$ with a suitable hyperplane passing at the origin of $`𝕄^{2+1}`$. Let us denote by $`\widehat{\sigma }`$ the suspension of $`\sigma `$; then
$$Q(B)=_{\lambda [0,r]}\{\widehat{\sigma }+\lambda v\}$$
where $`v`$ is the unitary (in the Minkowski norm) vector tangent to $`𝕀(1)`$, normal to $`\sigma `$, and pointing towards $`Q(C^{})`$. We also denote $`Q(B)=B(\sigma ,v,r)`$. The shape of the CT level surfaces in $`g(U(M_0))`$ is shown in figure 5. The initial singularity of $`g(U(M_0))`$ is given by the space-like segment $`J=Q(J_0)+w`$.
Simplicial type deformation. $`M_0`$ represents a local model of the deformation $`M=M(F,)`$ we are interested in. In fact, there exists a neighborhood $`𝒲`$ of $``$ in $`M_0`$ which embeds isometrically into $`M`$, respecting the cosmological time. Let us denote by $`𝒲^{}`$ the image of $`𝒲`$ in $`M`$. Then $`M𝒲^{}`$ embeds isometrically into the Minkowskian suspension $`M(F)`$, respecting again the cosmological time.
We describe now the universal covering $`U(M)𝕄^{2+1}`$ and a cocycle $`t:\mathrm{\Gamma }^3`$ which leads to $`\mathrm{\Gamma }^{}ISO^+(2,1)`$ such that $`M=U(M)/\mathrm{\Gamma }^{}`$. The inverse image of $`cF=𝕀(1)/\mathrm{\Gamma }`$ into the covering $`𝕀(1)`$ is an infinite and locally finite set $`\stackrel{~}{}`$ of disjoint complete geodesic lines. Given any geodesic $`\sigma _0\stackrel{~}{}`$, then $`\stackrel{~}{}=\{\sigma =\gamma (\sigma _0);\gamma \mathrm{\Gamma }\}`$. Let $`\widehat{}I^+(0)`$ be the suspension of the geodesic lines of $`\stackrel{~}{}`$. The set $`𝕀(1)\stackrel{~}{}`$ is the union of an infinite number of connected components. Denote by $`R`$ any such a component, and by $`\widehat{R}`$ its suspension, which is a component of $`I^+(0)\widehat{}`$. Every $`R`$ covers a component $`F_R`$ of $`F`$; more precisely, if $`\mathrm{\Gamma }_R`$ is the subgroup of $`\mathrm{\Gamma }`$ which keeps $`R`$ invariant, then $`F_R=R/\mathrm{\Gamma }_R`$.
Now, fix one base component $`R_0`$ and take in it one base point $`x_0`$. For each $`\gamma \mathrm{\Gamma }`$, let $`\gamma (x_0)`$ be the point in $`𝕀(1)`$ which is defined by the action of $`\gamma `$ on $`x_0`$. The geodesic arc in $`𝕀(1)`$ connecting $`x_0`$ with $`\gamma (x_0)`$ crosses a finite number of lines $`\{\sigma _i\}`$ belonging to $`\stackrel{~}{}`$. At each crossing consider the unitary (in the norm of $`𝕄^{2+1}`$) vector $`v_i`$ tangent to $`𝕀(1)`$ and normal to $`\sigma _i`$, pointing far from $`x_0`$. Then, the required cocycle $`t(\gamma )^3`$ is given by
$$t(\gamma )=\underset{i}{}rv_i.$$
Note that if $`\gamma _1(x_0)`$ and $`\gamma _2(x_0)`$ belong to the same component $`R`$, then $`t(\gamma _1)=t(\gamma _2)`$, whence also $`t(R)=t(\gamma )`$ for any $`\gamma `$ such that $`\gamma (x_0)R`$, is well defined. $`U(M)`$ is tiled by tiles of two types: (i) “$`\widehat{R}+w`$”, (ii) “$`B(\sigma ,v,r)+w`$”, for some translation vector $`w^3`$. More precisely, the tiles of the first type make the open subset of $`U(M)`$
$$=_R\{\widehat{R}+t(R)\}.$$
Each line $`\sigma \stackrel{~}{}`$ is in the boundary of two regions $`R_\sigma `$, $`R_\sigma ^{}`$ and we assume that $`R_\sigma `$ is closer to $`x_0`$ than $`R_\sigma ^{}`$. Set $`v_\sigma `$ the unitary (in the norm of $`𝕄^{2+1}`$) vector tangent to $`𝕀(1)`$ and normal to $`\sigma `$, pointing towards $`R_\sigma ^{}`$. The two regions $`\widehat{R}_\sigma +t(R_\sigma )`$ and $`\widehat{R}_\sigma ^{}+t(R_\sigma ^{})`$ are connected by the tile of the second type $`B(\sigma ,v_\sigma ,r)+t(R_\sigma )`$, so that
$$U(M)=_{\sigma \stackrel{~}{}}(B(\sigma ,v_\sigma ,r)+t(R_\sigma )).$$
Note that each tile has its own CT; all the cosmological times fit well together and define the CT of $`U(M)`$ which passes to the quotient $`M`$.
###### Remark 3.3
The construction of $`M(F)`$ and of $`M(F,)`$ can be performed for any hyperbolic surface $`F`$, not necessarily compact nor of finite area. Similarly, by starting from any locally finite family of weighted geodesic lines in $`𝕀(1)`$, the simplicial deformation that we have just described produces a globally hyperbolic spacetime structure on $`^2\times `$ with a regular cosmological time.
###### Remark 3.4
When $`F`$ is compact, every region $`R`$ (defined above) is bounded by infinitely many lines of $`\stackrel{~}{}`$. In fact, as $`F`$ is compact, every $`\gamma \mathrm{\Gamma }`$ with $`\gamma 1`$ is of hyperbolic type \[B-P\]. Consequently, for any $`\sigma \stackrel{~}{}`$ and for every $`\gamma \mathrm{\Gamma }`$ with $`\gamma (\sigma )\sigma `$, $`\sigma `$ and $`\gamma (\sigma )`$ are “ultra-parallel”. This means that the hyperbolic distance satisfies $`d(\sigma ,\gamma (\sigma ))>0`$; moreover, $`\sigma `$ and $`\gamma (\sigma )`$ have a common orthogonal geodesic line. Suppose now that a region $`R`$ is bounded by finitely many lines in $`\stackrel{~}{}`$. In this case, $`R`$ contains a band $`E`$ of infinite diameter, bounded by two half-lines contained in $`\stackrel{~}{}`$. As $`F`$ is compact, $`^2`$ is tiled by tiles of the form $`\gamma (𝒟)`$, where $`𝒟`$ is a fundamental domain for $`\mathrm{\Gamma }`$ of finite diameter. So, one (at least) tile $`\gamma (𝒟)`$ must be contained in $`E`$. But clearly $`\gamma (𝒟)\stackrel{~}{}\mathrm{}`$ and this contradicts the fact that $`R`$ is a region of $`^2\stackrel{~}{}`$.
The same conclusion holds if $`F`$ is of finite area. If $`F`$ is of infinite area, we can eventually have different behaviours. For instance, in the example $`F_0`$ above, $`\stackrel{~}{}_0`$ just consists of one component which divides $`^2`$ into two regions. Other examples will be presented in subsection 4.1.
## 4 The cosmological time of (2+1)-spacetimes
In this section we describe the main properties of the CT for an arbitrary spacetime $`M`$. We adopt the notations of the previous sections; in particular, $`M`$ is assumed to be an expanding matter-free spacetime of topological type $`S\times `$ with compact surface $`S`$ of genus $`g>1`$. The validity of the following statements can be quite easily checked for spacetimes of simplicial type. We shall try to point out the main ideas; we postpone a commentary on the proofs, with references to the existing literature.
###### Proposition 4.1
The cosmological time function, $`\tau :M(0,\mathrm{}[`$, is surjective and regular, so that it defines a global time on $`M`$. It lifts to a regular cosmological time on the covering, $`\stackrel{~}{\tau }:U(M)(0,\mathrm{}[`$. Each level surface $`\stackrel{~}{S}_a`$ of $`U(M)`$ maps onto $`S_a`$ of $`M`$ and is its universal covering. In other words, the action of $`\pi _1(S)`$ on $`U(M)`$ restricts to a free, properly discontinuous isometric action on $`\stackrel{~}{S}_a`$ such that $`S_a=\stackrel{~}{S}_a/\pi _1(S)`$. Each $`\stackrel{~}{S}_a`$ ($`S_a`$) is a future Cauchy surface.
### 4.1 Initial singularity - external view
Let us give a description of the initial singularity of $`M`$ as it appears “from the exterior” point of view, that is, from the Minkowski space in which the universal covering $`U(M)`$ is placed. In subsection 4.5 we shall show how the initial singularity can also be characterized in terms of the observables associated with the CT asymptotic states.
Let us first give a definition.
###### Definition 4.2
An $``$-tree (also called a real tree) is a metric space $`(𝒯,d)`$ such that for each couple of points $`pq𝒯`$ there exists a unique arc in $`𝒯`$ with $`p`$ and $`q`$ as end-points and this arc is isometric to the interval $`[0,d(p,q)]`$. This arc is called a segment of $`𝒯`$ and is denoted $`[p,q]`$.
###### Remark 4.3
The so-called simplicial trees are the simplest examples of real trees. A simplicial tree is a real tree covered by a countable family of elementary segments, called the “edges” of the tree, in such a way that: a) whenever two edges intersect, then they just have one common endpoint; b) the edge-lengths take values in a finite set of strictly positive numbers. Any endpoint of any edge is called a “vertex” of the tree. The distance is the natural length-space distance. Note that a simplicial tree is not necessarily locally finite; in other words, vertices of infinite “valence” may occur. In general, a real tree is more complicated than a simplicial tree because one might find, for instance, a segment containing a Cantor set made by the endpoints of other segments.
###### Proposition 4.4
For any $`pU(M)𝕄^{2+1}`$ there is a unique time-like geodesic arc $`a(p)`$ contained in $`U(M)`$, which starts at $`p`$ and is directed in the past of $`p`$, such that the Lorentzian length of $`a(p)`$ equals $`\stackrel{~}{\tau }(p)`$. The other end-point of $`a(p)`$, denoted by $`i(p)`$, belongs to the boundary $`U(M)`$ of $`U(M)`$ in $`𝕄^{2+1}`$. If $`p`$ and $`q`$ are identified in $`M`$ by the action of $`\pi _1(S)`$, so are $`a(p)`$ and $`a(q)`$. The union of the initial points $`𝒯=\{i(p);pU(M)\}`$ is an $``$-tree. More precisely, each segment of $`𝒯`$ is a rectifiable space-like curve in $`U(M)`$ with its own length. There is a natural isometric action of the fundamental group $`\pi _1(S)`$ on $`𝒯`$. The quotient space $`i(M)=𝒯/\pi _1(S)`$ can be thought as the initial singularity of $`M`$.
###### Remark 4.5
We have already encountered several examples of spaces of the form $`X=\stackrel{~}{X}/\pi _1(S)`$ for some action of $`\pi _1`$ on $`\stackrel{~}{X}`$: for instance, $`F=^2/\mathrm{\Gamma }`$, $`M=U(M)/\mathrm{\Gamma }^{}`$, $`S_a=\stackrel{~}{S}_a/\mathrm{\Gamma }^{}`$. Now, the initial singularity of spacetime also is a quotient $`i(M)=𝒯/\pi _1(S)`$. Instead of the bare topological quotient space, it is more interesting to consider $`\stackrel{~}{X}`$ endowed with the action of $`\pi _1`$.
###### Remark 4.6
When $`M`$ is of simplicial type, the corresponding real tree $`𝒯`$ is actually a simplicial real tree. This justifies the name we have given to these special spacetimes. In this case, the set of edges of $`𝒯`$ consists of the union of the space-like segments which form the initial singularity of the different tiles of the form $`B(\sigma ,v_\sigma ,r)+t(R_\sigma )`$ (see section 3). The points of $`𝒯`$ can also be characterized by the properties discussed in remark 3.1. A homeomorphic (not isometric) copy of $`𝒯`$ can easily be embedded into $`𝕀(1)`$. Select one point in each region of $`𝕀(1)\stackrel{~}{}`$ and consider the set made by the union of all these points. Connect two points of this set by a geodesic segment of $`𝕀(1)`$ if and only if they belong to adjacent regions. In this way we get the required tree. This tree is manifestly “dual” of $`\stackrel{~}{}`$; in fact, the regions of $`𝕀(1)\stackrel{~}{}`$ correspond to the vertices of $`𝒯`$ and the lines of $`\stackrel{~}{}`$ correspond to the edges of $`𝒯`$. We shall return on this duality in section 4.3. Note that, as demonstrated in remark 3.4, all the vertices of $`𝒯`$ are of infinite valence.
Examples of real trees. A hyperbolic surface $`F=^2/\mathrm{\Gamma }`$ is represented in figure 6. $`F`$ is of infinite area and is homeomorphic to a torus with one puncture; two simple closed geodesics $`c`$ and $`a`$ on $`F`$ are depicted. The geodesic $`c`$ cuts open $`F`$ into a compact surface and an infinite area annulus. By using the Poincaré disk model, a fundamental domain of $`\mathrm{\Gamma }`$ in $`^2`$ is also shown in figure 6. This domain is delimited by four pair-wise ultra-parallel geodesic lines. The inverse images of $`c`$ and $`a`$ are represented on this domain. The first two terms of a sequence of partial tilings of $`^2`$, made by a finite number of copies of the fundamental domain, are shown in figure 7. The first partial tiling just contains one fundamental domain. The second is made by the union of $`1+4=5`$ copies of the fundamental domain. The next partial tiling of this sequence, which is not shown in the figure, contains $`1+4+12=17`$ tiles, and so on. For each partial tiling of $`^2`$ one can determine a corresponding partial lifting of the curves $`c`$ and $`a`$. Figure 8 shows the first two partial liftings $`\stackrel{~}{c}`$ of $`c`$ and the structure of the associated partial dual trees. In the limit of the complete infinite tiling of $`^2`$, the complete lifting of $`c`$ contains an infinite number of geodesics and the associated real tree is infinite. In this case, $`^2\stackrel{~}{c}`$ has exactly one component with infinitely many boundary lines (the associated vertex of the dual tree has infinite valence), whereas all the remaining components have one boundary line. The first three partial liftings $`\stackrel{~}{a}`$ of $`a`$ are shown in figure 9; the corresponding partial dual trees are also represented. Note that these figures are just evocative, as they are not geometrically exact.
###### Remark 4.7
The $``$-trees and the associated $`\pi _1(S)`$-actions which occur in proposition 4.4 are not arbitrary (see \[O\] pag. 32). In fact, one can prove that the $`\pi _1(S)`$-action is minimal with small edge-stabilizers. This means that there is no non-empty strictly sub-tree which is invariant for this action, and that, for each segment in the tree, the subgroup of $`\pi _1(S)`$ which keeps the segment invariant is virtually Abelian. We shortly say that a real tree which admits such a kind of $`\pi _1(S)`$-action, is geometric.
### 4.2 Intrinsic and extrinsic geometry of the constant CT surfaces
In order to describe the geometric properties of the surfaces of constant cosmological time, it is natural to introduce the notion of geodesic lamination.
###### Definition 4.8
Let $`G`$ be a surface endowed with a $`C^1`$-Riemannian metric. As usual, this induces a length-space distance on $`G`$ and the notion of geodesic arc (line) makes sense. A geodesic lamination of $`G`$ is a closed subset $`K`$ of $`G`$, also called the support of the lamination, which is the disjoint union of complete and simple geodesics, also called the leaves of the lamination. “Complete” means that we dispose of arc-length parametrization defined on the whole real line $``$; “simple” means that the geodesic has no self-normal crossing in $`G`$. In other words, each leaf is either a simple closed geodesic or a simple geodesic which is an isometric copy of $``$ embedded in $`G`$. When $`G`$ is compact, such a non compact leaf is not a closed subset of $`G`$.
###### Remark 4.9
A finite union of disjoint simple closed geodesics is called a multi-curve and is the simplest example of geodesic lamination. We have already introduced multi-curves in section 3.3. A generic geodesic lamination $`K`$ can be more complicated than a multi-curve; in fact, if $`\alpha `$ is an arc embedded in $`G`$ which is transverse to the leaves of $`K`$, typically $`\alpha K`$ is a Cantor set.
###### Proposition 4.10
For every $`a(0,\mathrm{}[`$:
1) $`\stackrel{~}{S}_a`$ is the graph of a positive convex function defined on the plane $`\{x^3=0\}`$ in $`𝕄^{2+1}`$.
2) $`\stackrel{~}{S}_a`$ is only $`C^1`$-embedded into $`U(M)`$, so that it carries an induced $`C^1`$-Riemannian metric. $`\stackrel{~}{S}_a`$ is geodesically complete and for each $`pq\stackrel{~}{S}_a`$, there is a unique geodesic arc connecting $`p`$ and $`q`$.
3) The locus $`\stackrel{~}{L}_a`$ at which the embedding of $`\stackrel{~}{S}_a`$ into $`U(M)`$ is no longer $`C^2`$ is a geodesic lamination of $`\stackrel{~}{S}_a`$. $`\stackrel{~}{L}_a`$ is in fact the pull-back of a geodesic lamination $`L_a`$ of $`S_a`$.
###### Remark 4.11
If $`M`$ is the spacetime of simplicial type which corresponds to the weighted multi-curve $``$ on the surface $`F`$, then $`L_a`$ is just made by the boundary components of the flat annular components embedded into $`S_a`$, which are associated to the components of $``$.
The content of the last remark generalizes as follows. Recall that $`\pi _1(S)`$ acts as $`\mathrm{\Gamma }`$ on each $`𝕀(a)`$. For every $`a(0,\mathrm{}[`$, let us consider the map
$$p_a:\stackrel{~}{S}_a𝕀(a)$$
defined as follows: $`p_a(x)`$ is the unique point of $`𝕀(a)`$ such that the tangent plane to $`\stackrel{~}{S}_a`$ at $`x`$ is parallel to the tangent plane to $`𝕀(a)`$ at $`p_a(x)`$. This map is well-defined, surjective and $`\pi _1(S)`$-equivariant. By taking the union of the $`p_a`$’s we get a $`\pi _1(S)`$-equivariant map $`p:U(M)I^+(0)`$, respecting the CT. This induces a map $`p^{}:MM(F)`$ respecting the CT.
###### Proposition 4.12
1) There exists a geodesic lamination $``$ on $`F=𝕀(1)/\mathrm{\Gamma }`$, which lifts to a geodesic lamination $`\stackrel{~}{}`$ on $`𝕀(1)`$, such that, for every $`a`$, one has $`p_a(\stackrel{~}{L}_a)=a\stackrel{~}{}`$ and any leaf of $`\stackrel{~}{L}_a`$ is isometrically mapped onto a leaf of $`a\stackrel{~}{}`$. That is, the union of $`p_a(\stackrel{~}{L}_a)`$’s covers the suspension $`\widehat{}`$ of $`\stackrel{~}{}`$.
2) $``$ is the disjoint union of two sublaminations
$$=^{}$$
where $``$ is the maximal multi-curve sublamination of $``$. Note that either $``$ or $`^{}`$ may be empty. Then
a) $`p`$ embeds $`U(M)p^1(\widehat{})`$ isometrically into $`I^+(0)`$ respecting the CT;
b) $`p`$ embeds $`U(M)p^1(\widehat{})`$ continuously into $`I^+(0)`$ respecting the CT.
3) The set $`p^1(\widehat{})`$ is the union of components of the type $`B(\sigma ,v_\sigma ,r)+w`$, so that $`(p^{})^1()S_a`$ is the disjoint union of flat annular components of $`S_a`$, like in the case of a spacetime of simplicial type.
We have an immediate corollary concerning the intrinsic and extrinsic geometry of the constant CT surfaces.
###### Corollary 4.13
$`\stackrel{~}{W}_a=\stackrel{~}{S}_a\stackrel{~}{L}_a`$ is an open dense set of $`\stackrel{~}{S}_a`$. Each component of $`\stackrel{~}{W}_a`$ is either isometric to an open set of $`𝕀(a)`$ or is a flat band which embeds into $`𝕀(1,3,a)`$, and projects onto an annulus of $`S_a`$. Flat annuli do occur only if $``$ is non empty.
### 4.3 CT duality
To sum up, two geometric structures are naturally associated to the spacetime $`M`$: the real tree $`𝒯`$ (the initial singularity) and the geodesic lamination $``$ on $`F=𝕀(1)/\mathrm{\Gamma }`$ which reflects the lack of smoothness of the embedding of $`S_a`$ into $`M`$. We have already noted that for a spacetime of simplicial type these two objects are “dual” to each other. Here we want to strengthen and generalize this point.
If $``$ is non empty, we extend the lamination $`L_a`$ on $`S_a`$ to a lamination $`L_a^{}`$, by foliating the flat annular regions of $`S_a`$ by closed geodesics parallel to the boundary components. As usually $`\stackrel{~}{L}_a^{}`$ denotes the lifted lamination to $`\stackrel{~}{S}_a`$. The above map $`p_a`$ sends $`\stackrel{~}{L}_a^{}`$ onto $`a\stackrel{~}{}`$.
We have a natural continuous surjective map $`i_a:\stackrel{~}{S}_a𝒯`$ which associates to $`x`$ the initial point on the arc $`a(x)`$. So $`𝒯^{}=\{i_a^1(x)\}_{x𝒯}`$ is a partition of $`\stackrel{~}{S}_a`$ by closed subsets. $`\pi _1(S)`$ acts also on $`𝒯^{}`$ and, clearly, $`i_a`$ induces an $`\pi _1`$-equivariant identification between $`𝒯^{}`$ and $`𝒯`$.
###### Proposition 4.14
For every $`a`$, each closed set $`E`$ of the partition $`𝒯^{}`$ of $`\stackrel{~}{S}_a`$ is:
1) either the closure of a component of $`\stackrel{~}{S}_a\stackrel{~}{L}_a^{}`$;
2) or a leaf of the foliation of some band component of $`\stackrel{~}{L}_a^{}`$ which projects onto a flat annular region of $`S_a`$ .
We describe how the distance $`d`$ on the real tree $`𝒯`$ can be encoded, in dual terms, by equipping the geodesic laminations $`\stackrel{~}{}`$, $``$, with suitable transverse invariant measures.
###### Definition 4.15
A measured geodesic lamination on $`F`$ is a couple $`(,\mu )`$, where $``$ is a geodesic lamination and $`\mu `$ is a transverse invariant measure which consists of a Borel measure $`\mu _J`$ on each embedded interval $`J[0,1]`$ in $`F`$, transverse to the leaves of $``$ such that
(1) the support of $`\mu _J`$ coincides with $`J`$;
(2) if $`J,J^{}`$ are arcs, homotopic through arcs which are transverse to the leaves of $``$, keeping the endpoints either on the same leaf or in the same connected components of $`F`$, then $`\mu _J(J)=\mu _J^{}(J^{})`$ . $`(,\mu )`$ lifts to $`(\stackrel{~}{},\stackrel{~}{\mu })`$ which is $`\pi _1`$-equivariant.
###### Remark 4.16
The simplest measured geodesic laminations of $`F`$ are the weighted multi-curves.
Let $`J`$ be an arc in $`𝕀(1)`$ transverse to the leaves of $`\stackrel{~}{}`$. The map $`p_a`$ lifts $`J`$ to an arc $`J^{}`$ in $`\stackrel{~}{S}_a`$ transverse to the leaves of $`\stackrel{~}{L}_a^{}`$. On the other hand, by means of the map $`i_a`$, the distance $`d`$ on $`𝒯`$ lifts to a measure $`\stackrel{~}{\mu }_J^{}`$ on $`J^{}`$ which finally gives us the required ($`\pi _1`$-equivariant) transverse measure on $`\stackrel{~}{}`$.
One can invert the above construction and associate to each measured geodesic lamination $`(,\mu )`$ of the hyperbolic surface $`F`$ a suitable geometric $``$-tree.
From geodesic laminations to real trees. Take the measured lamination $`(,\mu )`$ of the surface $`F`$. $``$ is in general the disjoint union of two sublaminations
$$=^{\prime \prime }$$
where $``$ is the maximal weighted multi-curve sublamination of $``$. Note that either $``$ or $`^{\prime \prime }`$ may be empty. $`F`$ consists of a finite number of connected components, the metric completion of any such a component is isometric to a compact hyperbolic surface with geodesic polygonal boundary. If $``$ is non-empty, let us consider the spacetime of simplicial type associated to $``$, and let $`F^{}`$ be the $`\tau =1`$ level surface of this spacetime. Let us denote by $`^{}`$ the lamination on $`F^{}`$ which coincides with $``$ outside the flat annuli of $`F^{}`$ and is defined as $`L_1^{}`$ above on these annuli. If $``$ is empty, set $`F^{}=F`$. The measured lamination $`(,\mu )`$ “extends” to a measured lamination $`(^{},\mu ^{})`$ on $`F^{}`$. The flat annular components are foliated by closed geodesics parallel to the boundary components. These annuli are endowed with a plain transverse measure of total mass equal to the corresponding annulus depth. Take the universal covering $`\stackrel{~}{F}^{}`$ of $`F^{}`$ with the lifted ($`\pi _1`$-equivariant) measured geodesic lamination $`(\stackrel{~}{^{}},\stackrel{~}{\mu }^{})`$. Now define a partition of $`\stackrel{~}{F}^{}`$ by closed subsets in the very same way we have defined above the the partition $`𝒯^{}`$ of $`\stackrel{~}{S}_1`$, with respect to the lamination $`\stackrel{~}{L}_1^{}`$. Call again this partition $`𝒯^{}`$. We can give it a distance $`d`$ which makes it an $``$-tree. If $`E`$ and $`E^{}`$ are the closure of two components of the complement of the lamination, take two points $`x`$ and $`x^{}`$ in these closed sets such that the geodesic segment $`[x,x^{}]`$ of $`\stackrel{~}{F}^{}`$ is transverse to the leaves of the lamination. By integration, the transverse measure induces a distance on the subset of $`𝒯^{}`$ made by the closed sets intersecting $`[x,x^{}]`$. In fact, by the “invariance” of the measure, this distance doesn’t depend on the segment $`[x,x^{}]`$. Finally one verifies that in this way one can actually define a distance between any two points of $`𝒯^{}`$ and that the resulting $`(𝒯^{},d)`$ is a geometric real tree.
###### Remark 4.17
Clearly, weighted multi-curves on the surface $`F`$ dually correspond to geometric simplicial real trees; the spacetimes of simplicial type do materialize this duality.
### 4.4 Reconstruction of $`M=U(M)/\mathrm{\Gamma }^{}`$
Starting from $`(F=𝕀(1)/\mathrm{\Gamma },𝒯)`$ or, equivalently, from $`(F=𝕀(1)/\mathrm{\Gamma },(,\mu ))`$, one can reconstruct a cocycle $`t`$, whence $`M=U(M)/\mathrm{\Gamma }^{}`$. This generalizes what we have done for a spacetime of simplicial type in subsection 3.3.
With the notations introduced at the end of the previous subsection, consider $`(\stackrel{~}{}^{},\stackrel{~}{\mu }^{})`$ on $`\stackrel{~}{F}^{}`$. To recover a cocycle $`t`$ do as follows: fix one base point $`p_0^{}`$ on $`\stackrel{~}{F}^{}`$ out from the support of the lamination. Let $`p_0`$ be its image on $`F^{}`$. If $`\sigma `$ is a loop in $`F^{}`$ based on $`p_0`$, which represents an element $`[\sigma ]`$ of $`\pi (F^{},p_0)`$, lift it to the oriented arc $`\sigma ^{}`$ in $`\stackrel{~}{F}^{}`$ which starts at $`p_0^{}`$; up to homotopy we can assume that $`\sigma ^{}`$ is transverse to the leaves of the lamination. Let $`f`$ be any continuous $`^3`$-valued function on $`\sigma ^{}`$ which coincides with the unit normal to the leaves of the lamination, tangent to $`\stackrel{~}{F}^{}`$, and oriented in agreement with $`\sigma ^{}`$. Now we can integrate $`f`$ along $`\sigma ^{}`$ by using the transverse measure getting a vector $`t([\sigma ])`$. By varying $`[\sigma ]`$, one gets such a cocycle $`t`$.
### 4.5 CT asymptotic states
The above discussion tells us that any spacetime $`M=U(M)/\mathrm{\Gamma }^{}`$ is completely determined by the linear part $`\mathrm{\Gamma }`$ of its holonomy $`\mathrm{\Gamma }^{}`$ (or equivalently by the surface $`F=𝕀(1)/\mathrm{\Gamma }`$) and by its initial singularity $`i(M)=𝒯/\pi _1(S)`$. The aim of this subsection is to recover these geometric objects from the “internal point of view” by “working inside the spacetime”. More precisely, we would like to show that $`F`$ and $`𝒯`$ can be interpreted as the future and past asymptotic states for the geometry of the CT level surfaces. To this aim we shall consider the observables defined by the lengths of the curves on the CT level surfaces. It is convenient to introduce the concept of Marked Spectrum associated with a metric space $`(\stackrel{~}{X},d)`$ which is endowed with an action $`\alpha `$ of the surface’s fundamental group $`\pi _1(S)`$, so that $`X=\stackrel{~}{X}/\pi _1(S)`$. Whenever we shall refer to $`X`$, we shall actually refer to the triple $`(\stackrel{~}{X},d,\alpha )`$ (see remark 4.5).
Let us denote by $`𝒞`$ the set of conjugation classes of $`\pi _1(S)\{1\}`$ which coincide with the homotopy classes of non-contractible continuous maps $`f:S^1S`$. Each marked spectrum is a point of the functional space $`(_0)^𝒞`$, endowed with the natural product topology. The Marked Spectrum $`s_X`$ of $`X`$ (denoted also by $`s_{\stackrel{~}{X}}`$), is defined as follows: for any $`c=[\gamma ]𝒞`$, $`\gamma \pi _1(S)`$, $`\gamma 1`$,
$$s_X(c)=\underset{p\stackrel{~}{X}}{inf}d(p,\alpha _\gamma (p)).$$
The spectrum is “marked” because one takes track of the map in addition to its image.
When $`X=F`$ or $`S_a`$, $`s_X(c)`$ is just the length of a closed geodesic curve (not necessarily simple; that is, self-crossings could possibly occur in $`c`$) which minimizes the length among the curves in that homotopy class. For this reason, in such a case, $`s_X`$ is called the Marked Length Spectrum and is denoted by $`l_X`$. When $`\stackrel{~}{X}=𝒯`$, $`s_𝒯`$ can be expressed, in dual terms, as the Marked Measure Spectrum of the corresponding measured geodesic lamination $``$ on $`F`$; usually this is denoted $`I_{}`$. $`I_{}(c)`$ is just the minimal transverse measure realized by the curves in that homotopy class. When $`𝒯`$ is simplicial, that is when $``$ is a weighted multi-curve $``$ of $`F`$, $`I_{}(c)`$ is easily expressed in terms of the geometric intersection number (this also justifies the notation): assume that all the weights are equal to $`1`$ (that is the length of all edges of $`𝒯`$ is equal to $`1`$); then it is easy to see that $`I_{}(c)`$ is just the minimum number of intersection points between $``$ and any curve belonging to $`c`$ and transverse to the components of the lamination. For arbitrary weights one just takes multiples of the contribution of each component of $``$.
###### Remark 4.18
Instead of the whole $`𝒞`$, one could prefer to use the subset $`𝒮𝒞`$ of isotopy classes of simple closed curve in $`S`$, and take the corresponding (restricted) marked spectra. The discussion should proceed without any substantial modification.
On the boundary of the Teichmüller space. It is convenient, at this stage, to recall the fundamental facts about the role that the marked spectra play in the study of the Teichmüller space and in Thurston’s theory of its natural boundary. Let us denote by $`T_g`$ the Teichmüller space of the hyperbolic structures on $`S`$ up to isometry isotopic to the identity. It is well-known (see \[T3\], \[B-P\], \[F-L-P\]) that the map
$$l:T_g(_0)^𝒞$$
defined by $`l(F=^2/\mathrm{\Gamma })=l_F`$, realizes a meaningful embedding of $`T_g`$ onto a subset of $`(_0)^𝒞`$ homeomorphic to the finite dimensional open ball $`B^{6g6}`$. We shall identify $`T_g`$ with $`l(T_g)`$. In fact $`T_g`$ is in a natural way a real analytic submanifold of $`(_0)^𝒞`$.
Fix any such a hyperbolic structure $`FT_g`$ on $`S`$. Let us denote by $`𝒢(F)`$ the set of measured geodesic laminations on $`F`$. Let us denote by $`𝒢𝒯(S)`$ the set of all $`\pi _1(S)`$-geometric $``$-trees (remark 4.7). At the end of subsection 4.3, we have outlined a construction which associates to each $`𝒢(F)`$ a dual $``$-tree say $`\mathrm{\Delta }()𝒢𝒯(S)`$. Note that this construction did not use the fact that $`F`$ was associated to a spacetime $`M`$.
###### Proposition 4.19
$`\mathrm{\Delta }:𝒢(F)𝒢𝒯(S)`$ is a bijection, that is it can be naturally inverted. For each $`r>0`$, $`\mathrm{\Delta }(r)=r\mathrm{\Delta }(`$; here we take either the $`r`$-multiple of the measure or the $`r`$-multiple of the distance. We can shortly say that “$`\mathrm{\Delta }`$ respects the positive rays”.
###### Proposition 4.20
Consider the maps, $`I:𝒢(F)(_0)^𝒞`$ and $`s:𝒢𝒯(S)(_0)^𝒞`$, obtained by taking the corresponding marked spectra. Then $`I=s\mathrm{\Delta }`$ and is injective. The image in $`(_0)^𝒞`$ is a positive cone based on the origin and positive rays go onto positive rays, in a obvious sense. Moreover $`T_g`$ and the image Im(I) are disjoint subsets of $`(_0)^𝒞`$
###### Remark 4.21
These spectra represent the actual “physical” observables in our discussion. The last two propositions specify the meaning of the duality between laminations and real trees. As the spectra coincide, they reveal the same physical content.
Similarly to $`T_g`$, we identify $`𝒢(F)`$ and $`𝒢𝒯(S)`$ with the image $`Im(I)(_0)^𝒞`$, endowed with the subspace topology.
Set $`𝒫^+(𝒢(F))=𝒫^+(𝒢𝒯(S))=𝒫^+(Im(I))`$ the projective quotient space, obtained by identifying to one point each positive ray in $`Im(I)\{0\}`$. Similarly $`T_g𝒫^+(Im(I))`$ has a natural quotient topology.
###### Proposition 4.22
The pair $`(\overline{T}_g,\overline{T}_g)=(T_g𝒫^+(Im(I)),𝒫^+(Im(I)))`$ is homeomorphic to the pair $`(\overline{B}^{6g6},S^{6g7})`$, where $`\overline{B}^{6g6}`$ is the closed ball and $`S^{6g7}`$ is its boundary sphere. The natural action on $`T_g`$ of the mapping class group $`Mod_g`$ of $`S`$ extends to an action on the compactification $`\overline{T}_g`$. This is called the Thurston’s natural compactification and $`\overline{T}_g`$ is the natural boundary of the Teichmüller space.
We can state precisely how the simplicial trees are dense, as we claimed in section 3. Let us denote $`𝒮𝒯(S)`$ the subset of $`𝒢𝒯(S)`$ made by the simplicial real trees.
###### Proposition 4.23
$`𝒮𝒯(S)`$ is dense in $`𝒢𝒯(S)`$ in the induced topology by $`(_0)^𝒞`$.
###### Remark 4.24
In this remark we collect a few technical complements concerning the marked spectra and the geometric meaning of spectra convergence.
1) The natural compactification of $`T_g`$ is formally similar to the the natural compactification of $`^2`$ in the hyperboloid model $`𝕀(1)`$ where $`S_{\mathrm{}}^1`$ is obtained by adding to $`𝕀(1)`$ the endpoints of the rays of the future light cone.
2) Let $`F_n`$ be a sequence in $`T_g`$ considered as a sequence of actions of $`\pi _1(S)`$ on $`^2`$. The meaning of the compactification is the following; up to passing to a subsequence (still denoted by $`F_n`$), one of the following situations occur: for every $`c𝒞`$,
(a) $`l_{F_n}(c)l_{F_0}(c)`$, for some $`F_0T_g`$.
(b) There exist a geometric real tree $`𝒯𝒢𝒯(S)`$ and a positive sequence $`ϵ_n0`$, such that $`ϵ_nl_{F_n}(c)s_𝒯(c)`$. This is also called the Morgan-Shalen convergence of the sequence of actions. This can be reformulated in a similar, equivalent, dual way as the convergence (up to positive multiples) of a sequence of marked length spectra of hyperbolic structures on $`S`$ to the marked measure spectrum of a measured geodesic lamination on a fixed base $`F_0`$.
3) The convergence of marked spectra has a deep geometric content. This can be expressed in terms of the Gromov convergence. Given two metric spaces $`(Y,d)`$ and $`(Y^{},d^{})`$ and $`ϵ>0`$, an $`ϵ`$-relation is a set $`RY\times Y^{}`$ (i.e. a relation between the two spaces) such that:
(i) the two projections of $`R`$ to $`Y`$ and $`Y^{}`$ are both surjective;
(ii) if $`(y,y^{}),(z,z^{})R`$ then $`|d(y,z)d(y^{},z^{})|<ϵ`$.
Let $`G`$ be a group, and $`\{G\times Y_nY_n\}_{n1}`$ be a sequence of isometric actions of $`G`$ on the metric spaces $`Y_n`$. We say that $`(G\times Y_nY_n)(G\times Y_0Y_0)`$ in the sense of Gromov, if for every compact subset $`K_0Y_0`$, for every $`ϵ>0`$ and for every finite subset $`P`$ of $`G`$, if $`n`$ is big enough, there are compact subsets $`K_nY_n`$ and $`ϵ`$-relations $`R_n`$ between $`K_n`$ and $`K_0`$ which are $`P`$-equivariant; this means that: if $`xK_0`$, $`gP`$, $`g(x)K_0`$, $`x_n,y_nK_n`$ and $`(x,x_n),(g(x),y_n)R_n`$, then $`d_n(g(x_n),y_n)ϵ`$.
It turns out that in case $`(a)`$ above we actually have the convergence in the Gromov sense of the sequence of actions on $`^2`$ to an interior point of $`T_g`$. In case $`(b)`$, the Morgan-Shalen convergence is equivalent to the Gromov convergence for the sequence of actions on $`ϵ_n^2`$.
4) Note that $`𝒢𝒯(S)`$ is defined by using only the topology of $`S`$ (its fundamental group indeed) while in order to adopt the dual view point we have to fix (in an arbitrary way) a base hyperbolic surface $`F_0T_g`$. In fact, the dual view point can be developed by using the marked measure spectra of the measured (singular) foliations on $`S`$ (instead of the measured geodesic laminations on $`F_0`$), which only depend on the differential structure of $`S`$ (see \[F-L-P\]). On the other hand, let us consider $`T_g`$ as a space of complex holomorphic structures on $`S`$ (thanks to the classical Uniformization Theorem). By fixing any such structure $`F_0`$, one can realize such a spectrum as the measure spectrum of the horizontal measured foliation of a unique quadratic differential $`\omega `$ on $`F_0`$. These “rigidifications” (via geodesic laminations or quadratic differentials) of softer objects (the measured foliations) is reminiscent of the role of Hodge theory with respect to De Rham Cohomology.
CT asymptotic states as limit spectra. After this somewhat long but necessary digression, let us come back to the CT asymptotic states.
###### Proposition 4.25
(a) $`\underset{a0}{lim}l_{S_a}=s_𝒯`$; (b) $`\underset{a\mathrm{}}{lim}l_{S_a}/a=l_F`$.
###### Remark 4.26
This means, in particular, that, in a far CT future the spacetime looks like the Minkowskian suspension $`M(F)`$. In order to detect the dual effect of the initial singularity on the embedding of $`S_a`$ into $`M`$, for large value of the cosmological time one needs to increase the accuracy in the measurement of geometric quantities. Nevertheless this effect is, in principle, observable for any finite value $`a`$ of the CT.
###### Proposition 4.27
For every $`a(0,\mathrm{}[`$, $`l_{S_a}/a`$ belongs to $`T_g(_0)^𝒞`$. Hence, the cosmological time determines a curve $`\gamma _M:(0,\mathrm{}[T_g`$. This is a real analytic curve which connects $`FT_g`$ with the point on the natural boundary $`[𝒯]T_g`$ (here $`[.]`$ denotes the projective class).
###### Remark 4.28
Consider a spacetime of simplicial type. To prove proposition 4.25 in this case, one has to note that the depth of the annular regions is constant on each $`S_a`$. When $`a0`$, the contribution (to the length of any curve on $`S_a`$) of the part contained in the non-annular components becomes negligible, the length of the annuli boundaries goes linearly to zero, so that only the transverse crossing of the annuli becomes dominant. When $`a\mathrm{}`$, the annuli depth goes to zero because of the rescaling by $`1/a`$, and the length spectrum converge to the spectrum of $`F`$. The general case follows by using the density stated in proposition 4.23. Concerning proposition 4.27, in the special case of a spacetime of simplicial type, the curve in $`T_g`$ is just given by the Fenchel-Nielsen flow obtained by “twisting” the hyperbolic surface $`F`$ along the closed geodesic of the multi-curve (see \[T3\] and also \[B-P\]).
### 4.6 A commentary on the proofs
The identification between cocycles of a spacetime $`M`$ with measured geodesic laminations on $`F=𝕀(1)/\mathrm{\Gamma }`$ is due to Mess \[Me\]. In fact one can find other examples of such a construction of “cocycles” from measured laminations in the contest of Thurston’s theory of “bending” or “earthquakes” (see for instance \[E-M\]).
Measured geodesic laminations emerged in the original Thurston’s approach to the natural compactification of $`T_g`$ \[T\] \[T2\] \[T3\]. See also \[F-L-P\] for the alternative approach by using the measured foliations (see remark 4.24 4). For the claim about the quadratic differentials in remark 4.24 4) see \[Ke\]. The dual approach via real trees is due to Morgan-Shalen \[M-S\] \[M-S2\]. This approach does apply to more general, higher dimensional situations. The monography \[O\] contains a rather exhaustive introduction to this matter and we mostly refer to it (and to its bibliography) for all the details. In particular one can find in \[O\] a complete proof of the duality (see proposition 4.19 and proposition 4.20). The delicate point is just the inversion of the map $`\mathrm{\Delta }`$ we have described above. The geometric interpretation of the Morgan-Shalen convergence (see remark 4.24 3) is due to Paulin and Bestvina (c.f. the bibliography of \[O\]).
It is an amazing fact that the spacetimes “materialize” this subtle duality in the way we have seen. Note also that, in the spacetime setting, the choice of the base hyperbolic surface $`F`$ (see remark 4.24 4) is fixed by the linear part of the holonomy of $`M`$, that is by its future asymptotic state.
Concerning proposition 4.27, the Fenchel-Nielsen flow generalizes to the earthquake flow (one uses again the density 4.23) with initial data $`(F,)`$ which has real analytic orbits \[Ke2\].
## 5 Complements
In this section we add a few comments about the flat spacetimes with compact space of genus $`g=1`$, and about the spacetimes with negative cosmological constant. Finally we discuss a conjecture relating the CT and the CMC time.
### 5.1 Toric space ($`g=1`$)
The case in which the surface $`S`$ is a torus is a particular example of the so called Teichmüller spacetimes which we have analysed in \[B-G\]. So we simply remind the main points. Each non static spacetime determines a curve $`\gamma :(0,\mathrm{}[T_1^{}`$, $`\gamma (a)=(w(a),\omega (a))`$, where $`T_1^{}`$ is the cotangent bundle of the Teichmüller space $`T_1`$ of conformal structures on the torus up to isomorphism isotopic to the identity. Let us recall that $`T_1`$ is isometric with the Poincaré disk. The cotangent vectors $`\omega (a)`$ at the point $`w(a)T_1`$ is a quadratic differential on a Riemann surface representing $`w(a)`$. It is not hard to verify that $`\gamma `$ is just the complete orbit of the Teichmüller flow with initial data $`(w(1),\omega (1))`$ (see \[Ab\],\[B-G\]). In particular the projection of $`\gamma `$ onto $`T_1`$ is a complete geodesic connecting two boundary points. These points can also be interpreted in terms of marked spectra. Let us denote by $``$ and by $`𝒱`$ the horizontal and vertical measured foliations of $`w(1)`$. Then: $`\underset{a\mathrm{}}{lim}l_{S_a}/a=I_{}`$ and $`\underset{a0}{lim}l_{S_a}=I_𝒱`$.
### 5.2 Spacetime with negative cosmological constant
The above discussion on CT for flat spacetimes (i.e. with cosmological constant $`\mathrm{\Lambda }=0`$) can be adapted to the case of negative $`\mathrm{\Lambda }`$ which we normalize to be $`\mathrm{\Lambda }=1`$. We denote by $`𝕏^{2+1}`$ the Universal anti de Sitter Spacetime of dimension $`2+1`$. Each spacetime is now locally isometric to $`𝕏^{2+1}`$. The role played by $`I^+(0)`$ in the flat case, is played now by the diamond-shaped domain $`D(2)`$ (see \[H-E\] pag. 132) isometric to $`B^2\times (\pi /2,\pi /2)`$ with metric, in coordinates $`(y^1,y^2,t)`$, $`ds^2=(\mathrm{cos}^2t)h_2dt^2`$, where $`h_2`$ is the usual Poincaré hyperbolic metric on the open disk $`B^2`$.
Anti-de Sitter Suspensions. If $`F=^2/\mathrm{\Gamma }`$ is a hyperbolic surface of genus $`g>1`$, then $`\mathrm{\Gamma }`$ isometrically acts also on $`D(2)`$ and $`D(F)=D(2)/\mathrm{\Gamma }`$ is the anti-de Sitter suspension of $`F`$. Up to a translation, the function $`t`$ gives the CT and it has many qualitative properties similar to the CT of the Minkowskian suspensions, but we have now both an initial and a final singularity, both reduced to one point. In a sense, $`D(F)`$ can be obtained by the Minkowskian suspension $`M(F)`$ by a procedure of warping and doubling; $`D(F)`$ and $`M(F)`$ have the same initial singularity; the future asymptotic state of $`M(F)`$ “becomes” the level surface of the CT on $`D(F)`$ where the expansion ends and the collapsing begins. Also the anti-de Sitter analogous of $`I^+(1,3)`$ is easy to figure out.
Deforming anti-de Sitter suspensions. We want to generalize the above “warping and doubling” construction. Let $`M=U(M)/\mathrm{\Gamma }^{}`$ as in the former flat-spacetime discussion, $`\mathrm{\Gamma }^{}=\mathrm{\Gamma }+t(\mathrm{\Gamma })`$. $`F=𝕀(1)/\mathrm{\Gamma }`$ as usually. For $`t(\pi /2,0)`$, $`\tau (0,\mathrm{})`$, set $`t=(\pi /2)e^\tau `$. Denote $`h(a)`$ the spatial metric on $`S_a`$. On the manifold $`F\times (\pi /2,0)`$ consider the metric $`ds^2=\mathrm{cos}^2(t)h(\tau )/\tau ^2dt^2`$, getting a spacetime $`𝒟^{}(M)`$. Similarly, take $`M^{}`$ and $`𝒟^{}(M^{})`$, where $`M^{}=U(M^{})/\mathrm{\Gamma }^{}`$, $`\mathrm{\Gamma }^{}=\mathrm{\Gamma }t(\mathrm{\Gamma })`$, $`𝒟^{}(M^{})`$ is obtained from $`𝒟^{}(M^{})`$ by reversing the time and the orientation. Finally set $`𝒟(M)=𝒟^{}(M)𝒟^{}(M^{})`$, by gluing the two pieces at $`t=0`$. $`𝒟(M)`$ is locally anti de Sitter; up to a translation, $`t`$ gives the CT. The asymptotic state for $`t\pi /2`$ (i.e. the initial singularity) is equal to the initial singularity of $`M`$. The final singularity ($`t\pi /2`$) coincides with the initial singularity of $`M^{}`$. The future asymptotic states of $`M`$ and $`M^{}`$ “glue” at the level surface $`\{t=0\}`$ of the CT where the expansion ends and the collapse begins. The orbit of $`𝒟(M)`$ in $`T_g`$ is given by the union of two earthquake rays associated to $`M`$ (pointing to the future) and to $`M^{}`$ (towards the past); note that the qualitative behaviour is similar to what we have remarked for $`g=1`$. $`𝒟(M)`$ is the quotient of a domain $`D(2)_M𝕏^{2+1}`$, which is a “deformation” of the diamond-shaped domain $`D(2)`$. Also in this case the spacetimes with simplicial asymptotic singularities are significant and particularly simple to be understood.
### 5.3 CT versus CMC
Assume again that the space $`S`$ is of genus $`g>1`$, and that the spacetimes are flat. Given any global time on a spacetime $`M=U(M)/\mathrm{\Gamma }^{}`$, the asymptotic behaviour of the geometry of the corresponding level surfaces reflects in general a property of the specific time and not of the spacetime. On the other hand, we have seen that the asymptotic states of the cosmological time are intrinsic features of the spacetime. In this sense, we can say that a global time is “good” when it has the same asymptotic states of the CT. The CMC time, $`\rho `$ say, is a widely studied global time. A natural question is whether $`\rho `$ is a good global time. We conjecture that this is the case. Let us denote by $`W_a`$ the $`\{\rho =a\}`$ level surfaces of the CMC time.
###### Conjecture 5.1
(a) $`\underset{a\mathrm{}}{lim}l_{W_a}=s_𝒯`$; (b) $`\underset{a0}{lim}l_{W_a}/a=l_F`$.
There are some strong evidences supporting the conjecture; in particular by \[A-M-T\] we know that:
(1) $`\rho `$ is a global time function with image the interval $`(0,+\mathrm{})`$.
(2) If $`\gamma :(0,\mathrm{})T_g`$ is any $`\rho `$-orbit in $`T_g`$ (here $`T_g`$ is intended as a space of conformal structures) then:
(i) The $`\underset{\rho 0}{lim}\gamma `$ exists in $`T_g`$.
(ii) $`\gamma `$ is proper, that is it goes out from any compact set of $`T_g`$, roughly it “goes to $`\mathrm{}`$”.
An idea to prove the conjecture, should be to confine each $`W_a`$ between two barriers made by CT-level surfaces $`S_a^{}`$, $`S_{a^{\prime \prime }}`$, in such a way that $`a^{}`$ and $`a^{\prime \prime }`$ depend nicely on $`a`$ and, when $`a\mathrm{}`$ or $`a0`$, $`S_a^{}`$ and $`S_{a^{\prime \prime }}`$ become more and more “geometrically” close to each other. In a recent conversation, L. Andersson confirmed that this should actually work at least for a spacetime with simplicial initial singularity. A similar conjecture can be formulated in the anti-de Sitter context. |
warning/0003/astro-ph0003113.html | ar5iv | text | # The recent star formation history of the Hipparcos solar neighbourhood
## 1 Introduction
The problem of deducing the star formation rate history, $`SFR(t)`$, of the Milky Way has been generally attempted in terms of indirect inferences mostly through chemical evolution models. The validity of these methods relies on the soundness of the assignation of a “chemical age” to each of the studied stars. Generally a metallicity indicator is chosen, and used to measure the metal content of a number of stars which are then binned into age groups through the use of an age-metallicity relation derived from a chemical evolution model. For example, Rocha-Pinto & Maciel (1997) take a variety of age-metallicity relations (AMRs) from the literature and use a closed box chemical evolution model to translate AMRs into $`SFR(t)s`$, allowing for an intrinsic Gaussian spread in the AMR assumed to be constant in time.
The advantages of these methods are that large samples of stars both in the solar neighbourhood and further away within the disk of the galaxy can be studied. An inferred $`SFR(t)`$ can be constructed over an ample time range and spatial extent within the Galaxy which is consistent with the metallicities of the sample studied, and the chemical evolution model proposed. However, the validity of the AMR assumption can not be checked independently of the proposed chemical evolution model, and is necessarily dependent on what is chosen for the mixing physics of the ISM, the possible infall of primordial non enriched gas, and the still largely unknown galactic formation scenario in general. This last problem also affects attempts at inferring the $`SFR(t)`$ from stellar kinematics (e.g. Gomez et al. 1990, Marsakov et al., 1990).
With the recent availability of the Hipparcos satellite catalogue (ESA 1997) we are now in a position to attempt recovery of the local $`SFR(t)`$ directly, without the need of any model dependent assumptions. Previous direct approaches have been undertaken through the binning of observed stars into age groups according to the degree of chromospheric activity as measured through selected emission line features, with conflicting results depending on the assumed age-activity relation (e.g. Barry, 1988, and Soderblom et al., 1991). Using this technique, Rocha-Pinto et al. (1999) have derived a star formation history from the chromospheric activity-age distribution of a larger sample comprising 552 stars, founding evidence for intermittency in the $`SFR(t)`$ over 14 Gyr. The Hipparcos catalogue offers high quality photometric data for a large number of stars in the solar neighbourhood, which can be used to construct a colour-magnitude diagram (CMD) for this region. Once a CMD is available, it is in principle possible to recover the $`SFR(t)`$ which gave rise to the observed distribution of stars, assuming only a stellar evolutionary model in terms of a set of stellar tracks. In practice the most common approach to inverting CMDs has been to propose a certain parameterization for the $`SFR(t)`$, which is used to construct synthetic CMDs, which are statistically compared to the observed ones to select the values of the parameters which result in a best match CMD. Examples of the above are Chiosi et al. (1989), Aparicio et al. (1990) and Mould et al. (1997) using Magellanic and local star clusters, and Mighell & Butcher (1992), Smecker-Hane et al. (1994), Tolstoy & Saha (1996), Aparicio & Gallart (1995) and Mighell (1997) using local dSph’s.
We have extended these methods in Hernandez et. al. (1999) (henceforth paper I) by combining a rigorous maximum likelihood statistical approach, analogous to what was introduced by Tolstoy & Saha (1996), with a variational calculus treatment. This allows a totally non-parametric solution of the problem, where no a priori assumptions are introduced. This method was applied by Hernandez et al. (2000) (paper II) to a set of HST CMDs of local dSph galaxies to infer the $`SFR(t)`$ of these interesting systems. The result differs from what can be obtained from a chemical evolution model in that a direct answer is available, with a time resolution which depends only on the accuracy of the observations.
Limitations on the applicability of our method to the Hipparcos data appear in connection to the selection function of the catalogue. The need to work only with complete volume-limited samples limits the age range over which we can recover the $`SFR(t)`$ to 0–3 Gyr, with a resolution of $`0.05`$ Gyr. This makes it impossible to compare our results with those of chemical evolution models which typically sample ages of 0–14 Gyr, with resolutions of 0.5–1.5 Gyr.
In Section 2 we give a summarized review of the method introduced in paper I, the sample selection and results are discussed in section 3. Section 4 presents a careful statistical testing of our results, and Section 5 our conclusions.
## 2 The method
In this section we give a summary description of our HR diagram inversion method, which was described extensively in our paper I. In contrast with other statistical methods, we do not need to construct synthetic colour magnitude diagrams for each of the possible star formation histories being considered. Rather we use a direct approach which solves for the best $`SFR(t)`$ compatible with the stellar evolutionary models assumed and the observations used.
The evolutionary model consists of an isochrone library, and an IMF. Our results are largely insensitive to the details of the latter, for which we use:
$$\rho (m)\{\begin{array}{cc}\hfill m^{1.3}& 0.08M_{}<m0.5M_{}\hfill \\ \hfill m^{2.2}& 0.5M_{}<m1.0M_{}\hfill \\ \hfill m^{2.7}& 1.0M_{}<m\hfill \end{array}$$
(1)
The above fit was derived by Kroupa et al. (1993) for a large sample towards both galactic poles and all the solar neighbourhood, and therefore applies to the Hipparcos data.
As we shall be treating here only data from the solar neighbourhood derived by the Hipparcos satellite, we shall assume for the observed stars a fixed metallicity of $`[Fe/H]=0`$. This assumption is valid as we will only be treating stars within a short distance from the Sun, having a small spread in ages. Once the metallicity is fixed we use the latest Padova isochrones (Fagotto et al., 1994, Girardi et al., 1996) together with a detailed constant phase interpolation scheme using only stars at constant evolutionary phase, to construct an isochrone library having a chosen temporal resolution.
To transform the isochrones from the theoretical HR diagram to the observed colour-magnitude diagrams, we used the transformations provided by the calibrations of Lejeune et al. (1997) which are appropriate for the solar metallicities considered here. Using the updated calibrations given by Bessell et al. (1998) does not change the transformations significantly in the regime used here, unlike the case for giant and AGB stars (see Weiss & Salaris, 1999).
In this case we implement the method with a formal resolution of 15 Myr, compatible with the high resolution of the Hipparcos observations. It is one of the advantages of the method that this resolution can be increased arbitrarily (up to the stellar model resolution) with computation times scaling only linearly with it.
Our only other inputs are the positions of, say, $`n`$ observed stars in the HR diagram, each having a colour and luminosity, $`c_i`$ and $`l_i`$. Starting from a full likelihood model, we first construct the probability that the $`n`$ observed stars resulted from a certain $`SFR(t)`$. This will be given by:
$$=\underset{i=1}{\overset{n}{}}\left(_{t_0}^{t_1}SFR(t)G_i(t)𝑑t\right),$$
(2)
where
$`G_i(t)`$ $`=`$ $`{\displaystyle _{m_0}^{m_1}}{\displaystyle \frac{\rho (m;t)}{2\pi \sigma (l_i)\sigma (c_i)}}\times `$
$`\times \mathrm{exp}\left({\displaystyle \frac{D(l_i;t,m)^2}{2\sigma ^2(l_i)}}\right)\mathrm{exp}\left({\displaystyle \frac{D(c_i;t,m)^2}{2\sigma ^2(c_i)}}\right)dm`$
In the above expression $`\rho (m;t)`$ is the density of points along the isochrone of age $`t`$, around the star of mass $`m`$, and is determined by the assumed IMF together with the duration of the differential phase around the star of mass $`m`$. The ages $`t_0`$ and $`t_1`$ are a minimum and a maximum age needed to be considered, as $`m_0`$ and $`m_1`$ are a minimum and a maximum mass considered along each isochrone, e.g. 0.6 and 20 $`M_{}`$. $`\sigma (l_i)`$ and $`\sigma (c_i)`$ are the amplitudes of the observational errors in the luminosity and colour of the $`i`$th star. These values are supplied by the particular observational sample one is analysing. Note that for the sample we have selected (Section 3), there is no correlation between these errors. Finally, $`D(l_i;t,m)`$ and $`D(c_i;t,m)`$ are the differences in luminosity and colour, respectively, between the $`i`$th observed star and a general star of age and mass $`(m,t)`$. We shall refer to $`G_i(t)`$ as the likelihood matrix, since each element represents the probability that a given star, $`i`$, was actually formed at time $`t`$ with any mass.
A similar version was introduced in paper I, but restricted to the case of observational errors only in one variable, which was adequate to the problem of studying the $`SFR(t)`$ of local dSph galaxies treated in paper II. Equation (2) is essentially the extension from the case of a discretised $`SFR(t)`$ used by Tolstoy & Saha (1996), to the case of a continuous function in the construction of the likelihood. The challenge now is to find the optimum $`SFR(t)`$ without evaluating equation (2) i.e. without introducing a fixed set of test $`SFR(t)`$ cases from which one is selected.
The condition that $`(SFR)`$ has an extremal can be written as
$$\delta (SFR)=0,$$
and a variational calculus treatment of the problem applied. Firstly, we develop the product over $`i`$ using the chain rule for the variational derivative, and divide the resulting sum by $``$ to obtain:
$$\underset{i=1}{\overset{n}{}}\left(\frac{\delta _{t_0}^{t_1}SFR(t)G_i(t)𝑑t}{_{t_0}^{t_1}SFR(t)G_i(t)𝑑t}\right)=0$$
(3)
Introducing the new variable $`Y(t)`$ defined as:
$$Y(t)=\sqrt{SFR(t)}𝑑tSFR(t)=\left(\frac{dY(t)}{dt}\right)^2$$
and introducing the above expression into equation (3) we can develop the Euler equation to yield
$$\frac{d^2Y(t)}{dt^2}\underset{i=1}{\overset{n}{}}\left(\frac{G_i(t)}{I(i)}\right)=\frac{dY(t)}{dt}\underset{i=1}{\overset{n}{}}\left(\frac{dG_i/dt}{I(i)}\right)$$
(4)
where
$$I(i)=_{t_0}^{t_1}SFR(t)G_i(t)𝑑t$$
This in effect has transformed the problem from one of searching for a function which maximizes a product of integrals (equation 2) to one of solving an integro-differential equation (equation 4). We solve this equation iteratively, with the boundary condition $`SFR(t_1)=0`$.
Details of the numerical procedure required to ensure convergence to the maximum likelihood $`SFR(t)`$ can be found in our paper I, where a more complete development of the method is also found. In paper I the method was tested extensively using synthetic HR diagrams, obtaining very satisfactory results. Equation (4) will be satisfied by any stationary point in the likelihood, not just the global maximum, it is therefore important to check that the numerical algorithm implemented converges to the true $`SFR(t)`$. This was shown in our paper I, using synthetic HR diagrams extensively, and testing the robustness of the method to changes in the initialization condition of the algorithm. An independent test of the validity of the results was also implemented, and is described in section 4.
The main advantages of our method over other maximum likelihood schemes are (1) the totally non parametric approach the variational calculus treatment allows, and (2) the efficient computational procedure, where no time consuming repeated comparisons between synthetic and observational CMD are necessary, as the optimal $`SFR(t)`$ is solved for directly.
### 2.1 Two tests
We now present two examples of the method’s performance, in cases similar to the Hipparcos samples CMDs. The left panel of Figure (1) shows a synthetic CMD produced from the first input $`SFR(t)`$, resulting in a number of stars similar to what the Hipparcos samples yield for small errors in $`VI`$ ($`<0.12`$ mag) and $`M_V`$ $`(<0.02`$ mag). The positions of the simulated stars are then used to construct the likelihood matrix,which is used to recover the inferred $`SFR(t)`$, through an iterative numerical procedure (see paper I). The right panel of Figure (1) shows the last 3 iterations of the method (solid curves) and the input $`SFR(t)`$, a three burst $`SFR(t)`$ (dotted curve). It can be seen that the main features of the input $`SFR(t)`$ are accurately recovered. The age, duration and shape of the input $`SFR(t)`$ are clearly well represented by the final inferred $`SFR(t)`$. As the difference between successive isochrones diminishes with age, since the errors remain constant, the accuracy of the recovery procedure diminishes with the age of the stellar populations being treated. This is seen in that the first burst is very accurately recovered, whilst the last one appears somewhat spread out.
The last example is shown in Figure (2), which is analogous to Figure (1). Here a $`SFR(t)`$ which is constant over a large period is treated. The HR diagram of this case appears by sight almost identical to that of the previous case, however, given the extremely small errors assumed (typical of the Hipparcos data) the method is capable of distinguish and accurately recover the input $`SFR(t)`$ of these two cases. The small number of stars $`(450)`$ result in a degree of shot noise, which has to be artificially suppressed using a smoothing procedure, the result of which is seen in the residual short period oscillations of the inferred $`SFR(t)`$. This smoothing procedure reduces the effective resolution of the method to $`50`$ Myr. Note that as in the previous example, the inversion method successfully recovers the main features of the input $`SFR(t)`$. In these two tests only stars bluewards of $`VI=0.7`$ where considered in the inversion procedure (see below).
## 3 Sample selection and results
In order to apply the method described in the preceding section to the Hipparcos data, we would like to construct a volume-limited sample, where no biases appear between stars of different ages. Further, such a sample should contain a sufficient number of stars coming from all age groups being considered i.e. it must go down in magnitude to the turnoff point corresponding to the oldest age being considered. Although the Hipparcos satellite produced a catalogue having very well understood errors and highly accurate magnitude and colour determinations for a large number of stars, the sample has to be reduced through several cuts before it complies with the restrictions required by our method.
The Hipparcos catalogue provides an almost complete sample of stars in the solar neighbourhood. The limiting magnitude depends both on spectral type and galactic latitude (ESA SP-1200, Volume 1, page 131). For the types earlier than G5 which we consider here, the limiting V magnitude is given by $`V_{\mathrm{lim}}=7.9+1.1\mathrm{sin}|b|`$. To avoid unnecessary complications, we consider cuts of the type $`V`$=const. at all latitudes with $`V<7.9`$
Figure (3) shows a graph of $`M_V`$ vs. distance for all stars in the Hipparcos catalogue having distance errors smaller than $`20\%`$, and an apparent magnitude $`m_V<7.9`$. The solid lines show the inclusion criteria for one possible volume limited sample, complete to $`M_V<3.15,m<7.25`$ . As it can be seen, the maximum age which can be considered will not be very large, as the number of stars in a volume-limited sample complete to absolute magnitudes greater than 4 rapidly dwindles. After experimenting with synthetic CMDs of known $`SFR(t)`$ produced using our isochrone grid and constructed to have the same numbers of stars as a function of lower magnitude limit as in Figure (3), and recovering the $`SFR(t)`$ using our method, we identified 3 Gyr as the maximum age we can accurately treat with the data at hand. This fixes $`t_0=0,t_1=3`$ Gyr as the temporal limits in equation (2), were the use of 200 isochrones establishes the formal resolution of the method to be $`15`$ Myr.
Although the absolute magnitude errors correlate tightly with the distance, the colour errors correlate more strongly with the apparent magnitude, and can actually represent the dominant error in inverting the CMDs. The solid curve shows the $`m_V<7.25`$ completion limit, which implies errors similar to those used in Figures (1) and (2). It will be with complete volume-limited samples having this apparent magnitude limit that we will be dealing.
As the limit in $`M_V`$ is moved to dimmer stars, the structure of the $`SFR(t)`$at older ages is better recovered by the inversion method, but the number of younger stars diminishes (see Figure 1) and the $`SFR(t)`$ of the younger period is under represented in the recovered $`SFR(t)`$. We constructed a variety of somewhat independent Hipparcos CMDs for different absolute magnitude limits in the range $`3.0<M_V<3.5`$, and obtained highly compatible answers.
Also, all our samples include a certain number of contaminating stars having ages greater than the 3 Gyr limit considered by the inversion method, mostly in the RGB region, as their turn off points appear at magnitudes dimmer than the minimum ones considered. To avoid these stars, the inversion method considers only stars bluewards of $`VI=0.7`$, as was done in the synthetic examples of Figures (1) and (2).
As the fraction of stars produced by the $`SFR(t)`$ which live into the observed CMD diminishes with age, in inverting a well populated CMD the older regions of the $`SFR(t)`$ are under estimated by the recovery method. This is compensated by a correction factor given by the assumed IMF, and the mass at the tip of the RGB as a function of time, as discussed in paper I.
Once the IMF, metallicity, positions of the observed stars in the CMD and observational errors in both coordinates (also supplied by the Hipparcos catalogue) are given, they are used to construct the likelihood matrix $`G_i(t)`$, which is the only input given to the inversion method. The small number of stars present in any volume-limited sample ($`450`$) leaves us insensitive to the existence of small features in the $`SFR(t)`$ producing only a few stars. The limited numbers of stars also reduces the resolution of the method, as a smoothing function has to be applied to suppress instabilities in solving the integro-differential equation of the problem. This final smoothing reduces the effective temporal resolution of the method to 50 Myr, still much higher than that of any indirect chemical evolution inference of $`SFR(t)`$.
### 3.1 Kinematic and geometric corrections
Once the apparent and absolute magnitudes of the sample have been chosen, the set of stars to be studied is fully specified. The positions of which in the CMD are compared to the assumed isochrones to construct the likelihood matrix, which is the only input required by the numerical implementation of the method, as described above. The resulting $`SFR(t)`$ will be representative of the stars which where used in the construction of the likelihood matrix.
To normalize the various inferred $`SFRs(t)`$ from samples having different $`M_V`$ limits, and hence complete out to different distances, we apply the following kinematic and geometric corrections.
Let $`F(v,h)`$ be the fraction of the time a star having vertical velocity at the disk plane $`v`$ spends between heights $`h`$ and $`+h`$. Then, for a cylindrical sample complete to height $`h`$ above and below the disk plane,
$$N(t)=\frac{N_o(t)}{\sqrt{2\pi }\sigma (t)}_{\mathrm{}}^{\mathrm{}}\frac{e^{v^2/2\sigma (t)^2}}{F(v,h)}𝑑v$$
(5)
where $`N(t)`$ is the number of stars a stellar population of age $`t`$ contains, $`N_o(t)`$ the number of stars of age $`t`$ observed and $`\sigma (t)`$ the time dependent velocity dispersion of the several populations. As volume-limited samples are generally spherical around the Sun, a further geometric factor is required, giving:
$$N(t)=\frac{3N_o(t)}{2\sqrt{2\pi }\sigma (t)R^3}_0^Rrh_{\mathrm{}}^{\mathrm{}}\frac{e^{v^2/2\sigma (t)^2}}{F(v,h)}𝑑v𝑑r$$
(6)
where $`R`$ is the radius of the observed spherical volume limited sample, $`r`$ is a radial coordinate and $`h^2=R^2r^2`$. To estimate $`F(v,h)`$ one requires the detailed vertical force law of the galactic disk at the solar neighbourhood. The best direct estimate of this function remains that of Kuijken and Gilmore (1989), who show this function to deviate from that of a harmonic potential to a large degree. This detailed force law we integrate numerically to obtain $`F(v,h)`$. We use $`\sigma (t)=\mathrm{\hspace{0.33em}20}\mathrm{km}/\mathrm{s}`$ which is appropriate for the metallicity and age ranges we are studying (Edvardsson et al. 1993, Wyse & Gilmore 1995). Note also that the scatter in metallicity within 80 pc is rather small (Garnett and Kobulnicky, 1999) and will not change significantly our results.
In this way, assuming a Gaussian distribution for the vertical velocities of the stars, and a given $`\sigma (t)`$, an observed $`N_o(t)`$ can be transformed into a total $`N(t)`$, which is equal to the total projected disk quantity.
In our case the $`SFR(t)`$ given by the method takes the place of $`N_o(t)`$, and equation (6) is used to obtain a final star formation history, which accounts for the kinematic and geometric factors described. This function is then normalized through the total number of stars in the relevant sample, to give the deduced $`SFR(t)`$ in units of $`M_{}`$Myr<sup>-1</sup>kpc<sup>-2</sup>.
### 3.2 Results
Figure (4) shows the CMD corresponding to a volume-limited sample complete to $`M_V<3.15`$ for stars in the Hipparcos catalogue having errors in parallax of less than $`20\%`$ and $`m_V<7.25`$ (left panel). The right panel of this figure shows the result of applying our inversion method to this CMD, solid curve. The dotted envelope encloses several alternative reconstructions arising from different $`M_V`$ cuts, and gives an estimate of the errors likely to be present in our result, which can be seen to increase with time. The reconstruction based on the $`(M_V,BV)`$ diagram gives essentially the same results.
A certain level of constant star formation activity can be seen, superimposed onto a strong, quasi-periodic component having a period close to $`0.5`$ Gyr, as encoded in the positions of the observed stars in the CMD. The sharp feature seen towards $`t=3`$ Gyr could be the beginning of a fifth cycle, truncated by the boundary condition $`SFR(3)=0`$. We have performed tests with synthetic CMDs having the same numbers of stars and magnitude limits as in Figure (4), and having a variety of $`SFR(t)`$. The method efficiently discriminates between constant and periodic input $`SFR(t)`$, and correctly recovers features such as those found in the inferred $`SFR(t)`$ of Figure (4). We conclude that as far as the Padova isochrones at solar metallicity are representative of the observational properties of the stars in the CMDs, the $`SFR(t)`$ of the solar neighbourhood over the last 3 Gyr has been that shown in Figure (4). The unprecedented time resolution of our $`SFR`$ reconstruction makes it difficult to compare with the results derived from chromospheric activity studies (Rocha-Pinto et al. 1999), although qualitatively we do find the same activity at both 0.5 and above 2 Gyr, but not the decrease between 1 and 2 Gyr.
One possible interpretation of a cyclic component in the $`SFR(t)`$ of the solar neighbourhood can be found in the density wave hypothesis (Lin and Shu, 1964) for the presence of spiral arms in late type galaxies. As the pattern speed and the circular velocity are in general different, a given region of the disk (e.g. the solar neighbourhood), periodically crosses an arm region where the increased local gravitational potential might possibly trigger an episode of star formation. In the simplest version of this scenario, we can take the pattern angular frequency $`\mathrm{\Omega }_p`$ equal to twice the circular frequency $`\mathrm{\Omega }`$ at the Sun’s position (Binney & Tremaine 1987), valid within a flat rotation curve region. The time interval $`\mathrm{\Delta }t`$ between encounters with an arm at the solar neighbourhood will be given in general by
$$\mathrm{\Delta }t=\frac{2\pi }{m|\mathrm{\Omega }\mathrm{\Omega }_p|}$$
that is,
$$\mathrm{\Delta }t=\frac{0.22\mathrm{Gyr}}{m}\left(\frac{\mathrm{\Omega }}{29\mathrm{km}\mathrm{s}^1\mathrm{kpc}^1}\right)^1\left|\frac{\mathrm{\Omega }_p}{\mathrm{\Omega }}1\right|^1$$
where $`m`$ is the number of arms in the spiral pattern. The classical value of the pattern speed, $`\mathrm{\Omega }_p=0.5\mathrm{\Omega }14.5`$ km s<sup>-1</sup> kpc<sup>-1</sup> would imply that the interaction with a single arm ($`m=1`$) would be enough to account for the observed regularity in the recent SFR history. However, more recent determinations tend to point to much larger values (e.g. Mishkurov et al. 1979, Avedisova 1989, Amaral and Lépine 1997) close to $`\mathrm{\Omega }_p2324`$ km s<sup>-1</sup> kpc<sup>-1</sup>, which would then imply that the regularity present in the reconstructed $`SFR(t)`$ would be consistent with a scenario where the interaction of the solar neighbourhood with a two-armed spiral pattern would have induced the star formation episodes we detect. These arms are clearly detected in, for instance, the distribution of free electrons in the galactic plane (Taylor and Cordes 1993). This is reminiscent of the explanations put forward to account for the inhomogeneities observed in the velocity distribution function, where well-defined branches associated with moving groups of different ages (Chereul et al. 1999, Skuljan et al. 1999, Asiain et al. 1999) could perhaps be also associated with an interaction with spiral arm(s), although in this case the time scales are much smaller.
Alternatively, if the solar neighbourhood is closer to the corotation radius, the galactic bar could have triggered star formation in the solar neighbourhood with episodes separated by about 0.5 Gyr if the pattern speed of the bar is larger than about 40 km s<sup>-1</sup> kpc<sup>-1</sup> (Dehnen 1999).
Of course, other explanations are possible, for example the cloud formation, collision and stellar feedback models of Vazquez & Scalo (1989) predict a phase of oscillatory $`SFR(t)`$ behaviour as a result of a self-regulated star formation régime. Close encounters with the Magellanic Clouds have also been suggested to explain the intermittent nature of the SFR on longer time scales (Rocha-Pinto et al. 2000).
We have tested the ability of our method to accurately distinguish oscillatory components in the $`SFR(t)`$ with tests such as those shown in Figures (1) and (2). The oscillatory component in the case shown in Figure (1) is successfully recovered by the method. In the case shown in Figure(2) however, although the main features are also accurately recovered, a level of small amplitude fluctuations spuriously appears. This last feature is of such a small level, that if a CMD where produced from the method answer of Figure (1) having the same total number of stars, each small fluctuation would produce a number of stars of order 1.
Our answer shown in Figure (4) shows not only a large scale oscillatory component, but superimposed onto this, a certain level of small amplitude fluctuations. Given the total number of stars present in our sample, we can not rule out the possibility (quit possibly in fact the case) that these small fluctuations are numerical, as they are of amplitude similar to the ones discussed appearing in Figure(2), and are actually buried within the error envelope. The main oscillatory component having a period of 0.5 Gyr however, involves a sufficiently large number of stars to be objectively identified.
A larger and independent data set from which to derive $`SFR(t)`$ would be necessary in order to extend our results to a broader age range, and a more extensive region of the Galactic disk. Increasing the number of stars available for the HR inversion procedure would also allow to recover finer features, and reduce the small numerical fluctuations discussed above.
## 4 Testing the results
In our Paper I we tested this method using synthetic CMDs produced from known star formation histories, with which we could assess the accuracy of the result of the inversion procedure, as shown in Figures (1) and (2). In working with real data, we require the introduction of an independent method of comparing our final result to the starting CMD, in order to check that the answer our inversion procedure gives is a good answer. From our paper I we know that when the stars being used in the inversion procedure were indeed produced from the isochrones and metallicity used to construct the likelihood matrix, the inversion method gives accurate results. The introduction of an independent comparison between our answer and the data is hence a way of checking the accuracy of the input physics used in the inversion procedure, i.e. the IMF, metallicity and observational parameters.
The most common procedure of comparing a certain $`SFR(t)`$ with an observed CMD is to use the $`SFR(t)`$ to generate a synthetic CMD, and compare this to the observations using a statistical test to determine the degree of similarity between the two.
The disadvantage however is that one is not comparing the $`SFR(t)`$ with the data, but rather a particular realisation of the $`SFR(t)`$ with the data. The distinction becomes arbitrary when large numbers of stars are found in all regions of the CMD, which is generally not the case. Following a Bayesian approach, we prefer to adopt the $`W`$ statistic presented by Saha (1998), essentially
$$W=\underset{i=1}{\overset{B}{}}\frac{(m_i+s_i)!}{m_i!s_i!}$$
where B is the number of cells into which the CMD is split, and $`m_i`$ and $`s_i`$ are the numbers of points two distributions being compared have in each cell. This asks for the probability that two distinct data sets are random realisations of the same underling distribution. In implementing this test we first produce a large number (500) of random realisations of our inferred $`SFR(t)`$, and compute the $`W`$ statistic between pairs in this sample of CMD’s. This gives a distribution which is used to determine a range of values of $`W`$ which are expected to arise in random realisations of the $`SFR(t)`$ being tested. Next the $`W`$ statistic is computed between the observed data set, and a new large number of random realisations of $`SFR(t)`$ (also 500), this gives a new distribution of $`W`$ which can be objectively compared to the one arising from the model-model comparison to assess whether both data and modeled CMD’s are compatible with a unique underling distribution.
Figure (5) shows a synthetic CMD produced from our inferred $`SFR(t)`$ for the solar neighbourhood, down to $`M_V=3.15`$. This can be compared to the Hipparcos CMD complete to the same $`M_V`$ limit of Figure (4). A visual inspection reveals approximately equal numbers of stars in each of the distinct regions of the diagram, a more rigorous statistical comparison is also included. The right panel of Figure (5) shows a histogram of the values of the $`W`$ statistic for 500 random realisations of our inferred $`SFR(t)`$ in a model-model comparison. This gives the range of values of the $`W`$ statistic likely to appear in comparisons of two CMD diagrams arising from the same underlying $`SFR(t)`$, our inferred $`SFR(t)`$. The dashed histogram shows the results of 500 synthetic CMDs vs. the observed Hipparcos data set. The compatibility of both sets of $`W`$ values is clear. In this way the observed stars in the volume limited sample complete to $`M_V=3.15`$ and the isochrones, IMF and metallicity used in estimating the inferred $`SFR(t)`$ shown in Figure (4) are shown to be compatible with each other. Similar tests were also performed changing the limiting $`M_V`$ in the range $`3.03.5`$, and comparing against the corresponding Hipparcos sample. This produces alternative CMDs which contain different stars (see Figure 3), which were compared to synthetic CMDs coming always from the same central inferred $`SFR(t)`$. The results were always equal or better than what is shown in Figure 5, model-model and data-model distributions of $`W`$ having mean values well within $`1\sigma `$ of each other, where $`\sigma `$ refers to either the model-model or the data-model $`W`$ distributions.
## 5 Conclusion
We have applied the method developed in our paper I to the data of the Hipparcos catalogue. An objective answer for the $`SFR(t)`$ of the solar neighbourhood over the last 3 Gyr was found, which can be shown to be consistent with the complete volume-limited Hipparcos samples relevant to this age range. A structured $`SFR(t)`$ is obtained showing a cyclic pattern having a period of about 0.5 Gyr, superimposed on some degree of underlying star formation activity which increases slightly with age. No random bursting behaviour was found at the time resolution of 0.05 Gyr of our method. A first order density wave model for the repeated encounter of galactic arms could explain the observed regularity. |
warning/0003/hep-ph0003248.html | ar5iv | text | # MADPH-00-1164 hep-ph/0003248 March, 2000 Effective potential calculation of the MSSM lightest CP-even Higgs boson massContribution to PASCOS99: 7th International Symposium on Particles, Strings and Cosmology, Granlibakken, Tahoe City, California, 10-16 Dec 1999.
## Abstract
I summarize results of two-loop effective potential calculations of the lightest CP-even Higgs boson mass in the minimal supersymmetric standard model.
Computing the lightest CP-even Higgs boson mass is the most important loop calculation in the minimal supersymmetric standard model because of the paramount importance of a precise $`m_{h^0}`$ value to the Higgs boson experimental discovery. Tree-level supersymmetry relations require that the Higgs field quartic coupling be related to the electroweak gauge couplings; therefore they impose a strict upper bound $`m_{h^0}m_Z`$, which is already in conflict with the current lower limit from LEP 2.
It is well-known that this tree-level limit can be drastically changed by radiative corrections. One-loop calculations show that incomplete cancellations of the top and stop loops give positive corrections of the form
$$\mathrm{\Delta }m_{h^0}^2=\frac{3h_t^2m_t^2}{4\pi ^2}\mathrm{ln}\frac{m_{\stackrel{~}{t}}^2}{m_t^2},$$
(1)
where $`m_t`$ and $`m_{\stackrel{~}{t}}`$ are top and stop masses respectively. This formula, however, suffers from an ambiguity in the definition of $`m_t`$. Numerically, using running or on-shell top-quark mass can amount to about $`20\%`$ difference in $`\mathrm{\Delta }m_{h^0}^2`$. The problem can only be resolved by an explicit two-loop calculation.
Two-loop calculations in the existing literature have used two different approaches: (a) a renormalization group resummation approach, and (b) a two-loop diagrammatic approach. In the first approach, leading and next-to-leading logarithmic corrections are calculated by integrating one- and two-loop renormalization group equations. However, two-loop non-logarithmic finite corrections are not calculable in principle. The second approach was initiated by Hempfling and Hoang using an effective potential method; they restricted their calculation to specific choice of supersymmetry parameters: i.e. large $`\mathrm{tan}\beta \mathrm{}`$ and zero left-right stop mixing. Two-loop QCD corrections were later computed at more general cases in the effective potential approach. $`m_{h^0}`$ to the same two-loop QCD order was also computed using an explicit diagrammatic method. These calculations incorporate both two-loop logarithmic and non-logarithmic finite corrections. In the following, I shall concentrate on the effective potential approach.
The general way of calculating corrections to CP-even Higgs boson mass is to compute Higgs self-energy and tadpole diagrams to the required loop order. In an effective potential approach, these diagrams can be derived from a generating functional, i.e. the effective potential, by taking explicit derivatives with respect to the Higgs fields. These quantities then enter the MSSM CP-even Higgs boson mass-squared matrix as follows
$$_h^2=\left[\begin{array}{cc}m_Z^2c_\beta ^2+m_{A^0}^2s_\beta ^2+\mathrm{\Delta }_{11}^2& (m_Z^2+m_{A^0}^2)s_\beta c_\beta +\mathrm{\Delta }_{12}^2\\ (m_Z^2+m_{A^0}^2)s_\beta c_\beta +\mathrm{\Delta }_{21}^2& m_Z^2s_\beta +m_{A^0}^2c_\beta +\mathrm{\Delta }_{22}^2\end{array}\right],$$
(2)
where $`\mathrm{\Delta }_{ij}^2`$ represents radiative corrections to the $`ij`$-entry. We note that all these corrections are computed at the zero external momentum limit; sometimes it is necessary to calculate self-energy diagrams directly to account for corrections from non-zero external momenta.
The CP-even Higgs boson masses can be calculated by diagonalizing the above matrix in eq. (2). This computation is tedious but can be greatly simplified when one considers the case $`m_{A^0}m_Z`$, where $`m_{A^0}`$ is the mass of the pseudoscalar $`A^0`$. In this case, we find the corrections to $`m_{h^0}^2`$ is
$$\mathrm{\Delta }m_{h^0}^2=\frac{4m_t^4}{v^2}\left(\frac{d}{dm_t^2}\right)^2V\mathrm{Re}\mathrm{\Pi }_{hh}(m_{h^0}^2)+\mathrm{Re}\mathrm{\Pi }_{hh}(0).$$
(3)
where $`V`$ is the effective potential, $`v`$ the Higgs field VEV, and the last two terms account for non-zero external momentum corrections.
We have carried out this calculation procedure to the two-loop order including leading QCD and top Yukawa corrections. To illustrate our analysis, we present an approximation formula which is derived under the following assumptions: the soft masses for left and right stops, gluino, heavy Higgs bosons and Higgsinos have a common mass $`M_S`$, where $`M_S`$ can be identified as the supersymmetry scale. The two eigenvalues and mixing angle of stops are then accordingly $`m_{\stackrel{~}{t}_1}^2=m_{\stackrel{~}{t}}^2+m_tX_t`$, $`m_{\stackrel{~}{t}_2}^2=m_{\stackrel{~}{t}}^2m_tX_t`$ and $`s_t=c_t=\frac{1}{\sqrt{2}}`$, where the average top-squark mass $`m_{\stackrel{~}{t}}^2=M_S^2+m_t^2`$, and $`X_t=A_t+\mu /\mathrm{tan}\beta `$ is the left-right stop mixing parameter.
We find the approximation formula for two-loop QCD+top Yukawa corrections is (in terms of on-shell mass parameters)
$`\mathrm{\Delta }m_{h^0}^2={\displaystyle \frac{3m_t^4}{2\pi ^2v^2}}\left(\mathrm{ln}{\displaystyle \frac{m_{\stackrel{~}{t}}^2}{m_t^2}}+\widehat{X}_t^2{\displaystyle \frac{\widehat{X}_t^4}{12}}\right)`$ (4)
$`+`$ $`{\displaystyle \frac{\alpha _sm_t^4}{\pi ^3v^2}}\left(3\mathrm{ln}^2{\displaystyle \frac{m_{\stackrel{~}{t}}^2}{m_t^2}}6\mathrm{ln}{\displaystyle \frac{m_{\stackrel{~}{t}}^2}{m_t^2}}+6\widehat{X}_t3\widehat{X}_t^2\mathrm{ln}{\displaystyle \frac{m_{\stackrel{~}{t}}^2}{m_t^2}}{\displaystyle \frac{3\widehat{X}_t^4}{4}}\right)`$
$`+`$ $`{\displaystyle \frac{3\alpha _tm_t^4}{16\pi ^3v^2}}\{s_\beta ^2(3\mathrm{ln}^2{\displaystyle \frac{M_S^2}{m_t^2}}+13\mathrm{ln}{\displaystyle \frac{M_S^2}{m_t^2}})1{\displaystyle \frac{\pi ^2}{3}}+c_\beta ^2(60K+{\displaystyle \frac{13}{2}}+{\displaystyle \frac{4\pi ^2}{3}})`$
$`+`$ $`\left[3s_\beta ^2\mathrm{ln}{\displaystyle \frac{M_S^2}{m_t^2}}c_\beta ^2\left({\displaystyle \frac{69}{2}}+24K\right)+41\right]\widehat{X}_t^2\left(1+{\displaystyle \frac{61}{12}}s_\beta ^2\right)\widehat{X}_t^4+{\displaystyle \frac{s_\beta ^2}{2}}\widehat{X}_t^6`$
$`+`$ $`c_\beta ^2[(316K\sqrt{3}\pi )(4\widehat{X}_t\widehat{Y}_t+\widehat{Y}_t^2)+(16K+{\displaystyle \frac{2\pi }{\sqrt{3}}})\widehat{X}_t^3\widehat{Y}_t`$
$`+`$ $`({\displaystyle \frac{4}{3}}+24K+\sqrt{3}\pi )\widehat{X}_t^2\widehat{Y}_t^2({\displaystyle \frac{7}{12}}+8K+{\displaystyle \frac{\pi }{2\sqrt{3}}})\widehat{X}_t^4\widehat{Y}_t^2]\},`$
where the constant $`K0.195`$. We note that two-loop QCD corrections depend only on $`\widehat{X}_t=X_t/m_{\stackrel{~}{t}}`$ while the top Yukawa corrections depend on $`\widehat{Y}_t=(A_t\mu \mathrm{tan}\beta )/m_{\stackrel{~}{t}}`$ as well. This approximation formula is good to a level of 0.5 GeV for most of the parameter space.
Fig. 1 shows the Higgs boson mass $`m_{h^0}`$ vs. the stop mixing parameter $`\widehat{X}_t`$, for different choices of $`M_S`$, $`\mu `$ and $`\mathrm{tan}\beta `$. The two-loop QCD corrections agree well with other approaches. They generally decrease $`m_{h^0}`$ from their one-loop values by $`1020`$ GeV depending on the parameter choice. Two-loop Yukawa corrections are sizeable for large stop mixings, in particular, for $`\widehat{X}_t\pm 2`$ two-loop Yukawa corrections can increase $`m_{h^0}`$ by about $`5`$ GeV.
Another interesting feature observed in the literature is that two-loop corrections shift the maximal mixing peaks. At the one-loop level, these peaks are at $`\widehat{X}_t=\pm \sqrt{6}`$. It is easy to see from eq. (4) that the size of shifts is about $`10\%`$, i.e. the peaks move to $`\widehat{X}_t\pm 2`$. This is confirmed by Fig. 1.
Finally, renormalization group resummation technique can be used to derive a particularly nice mass correction formula which has clearer physical interpretations. We find eq. (4) can be transformed into the following form by using solutions to the renormalization group equations
$$\mathrm{\Delta }m_{h^0}^2=\frac{3\overline{m}_t^4(Q_t)}{2\pi ^2\overline{v}^2(Q_1^{})}\mathrm{ln}\frac{m_{\stackrel{~}{t}}^2(Q_{\mathrm{th}})}{\overline{m}_t^2(Q_t^{})}+\frac{3\overline{m}_t^4(Q_{\mathrm{th}})}{2\pi ^2\overline{v}^2(Q_2^{})}\left[\widehat{X}_t^2(Q_{\mathrm{th}})\frac{\widehat{X}_t^4(Q_{\mathrm{th}})}{12}\right]+\mathrm{\Delta }_{\mathrm{th}}^{(2)},$$
(5)
where $`Q_1^{}=e^{1/3}m_t`$, $`Q_2^{}=e^{1/3}m_t`$, $`Q_t=\sqrt{m_tm_{\stackrel{~}{t}}}`$, $`Q_t^{}=(m_tm_{\stackrel{~}{t}}^2)^{1/3}`$ and $`Q_{\mathrm{th}}=m_{\stackrel{~}{t}}`$, $`\overline{v}`$ and $`\overline{m}`$ are the Standard Model $`\overline{\mathrm{MS}}`$ parameters. These choices of scales for evaluating one-loop corrections automatically take into account two-loop leading and next-to-leading logarithmic effects. The leftover finite correction term $`\mathrm{\Delta }_{\mathrm{th}}^{(2)}`$ is understood as two-loop threshold corrections and numerically small; its detail form can be found in a forthcoming paper.
I thank J. R. Espinosa for collaborations. This work was supported in part by a DOE grant No. DE-FG02-95ER40896 and in part by the Wisconsin Alumni Research Foundation. |
warning/0003/cond-mat0003140.html | ar5iv | text | # Electronic structure of self-assembled quantum dots: comparison between density functional theory and diffusion quantum Monte Carlo
## I Introduction
The electronic structure of quantum dots resembles that of atoms. Both systems display three dimensional electronic confinement leading to level degeneracies, shell structure, and spin correlation to only mention the most studied atomic properties of quantum dots. Direct as well as exchange Coulomb interactions are also present in quantum dot systems, the importance of each depending on the dot dimensions. That is why some of the theoretical methods used to calculate the electronic structure of quantum dots are the same as the ones used to study atoms and molecules. One of these methods, local spin density approximation (LSDA) density functional theory (DFT) within the effective mass approximation (EMA), has been widely used to calculate the electron-electron interaction in dots because of its simple implementation and low computer demand. LSDA is an approximate theory, well known to predict incorrectly the physical properties of some molecules and solids, while performing well on many other systems. However, experience with LSDA on atoms and molecules may not carry over to quantum dots since the confining potentials and electron densities can be very different. As a general rule, density functional theory tends to work better for high effective density, and fails for low density systems, where correlation effects become important. There have been several investigations of LSDA and exact treatment of electron interactions in parabolic quantum dots, with particular emphasis on low density systems and effects of external magnetic fields. The electron interactions in parabolic dots have been treated exactly using diffusion Monte Carlo (DMC) as well as exact diagonalization methods. The parabolic potentials are popular models because they are computationally convenient and, with two adjustable parameters, have had relative success in explaining the influence of many-body effects on the electronic properties of dots.
In this paper we compare EMA results obtained with LSDA in self-assembled InAs/GaAs quantum dots against an exact treatment of the EMA many-body interactions obtained with diffusion quantum Monte Carlo. These small dots, which can hold about six electrons, have complicated confining potentials that differ considerably from the parabolic dots considered in previous DMC and exact diagonalization calculations. The single particle states of realistically modeled pyramidal dots are significantly modified from the states in a parabolic dot, so many-body interactions should also differ from those in parabolic models. The purpose of this comparison is to quantify errors due to LSDA for calculations on realistically modeled InAs/GaAs dots. Although error of LSDA is expected to be small in these systems, it is large enough to be a limitation to the comparison of model results to experiments. For the system reported here, we consider LSDA errors of up to 10 meV to be acceptable given the current precision of the models and experiments. Determination of whether LSDA meets this acceptability criterion requires careful comparison with many-body calculations on realistically modeled potentials.
In these calculations we consider the ground state energy as a function of the number of electrons in the dot, which is the quantity for which DFT should apply. Electron addition and removal energies are experimentally measurable. These may be rigorously defined as differences in the total energies of the ground states of the dot which differ by one electron. A common approach to calculating electron addition and removal energies is to use the eigenvalues of the DFT equations, which are well known not to be interpretable as electron addition and removal energies. We have addressed this issue in a previous paper where we have shown that addition and removal energies are to a very good approximation given by the “Slater transition rule”, which uses the eigenvalues at one-half occupation. Thus the present work, which compares total energies, is a rigorous test of the LSDA, and combined with our previous work is a test of the accuracy of the addition and removal energies obtained using LSDA eigenvalues and the Slater transition rules.
The ground state energy is the lowest eigenvalue of the many-body Schrödinger equation, which we take to be
$$\left\{\underset{i=1}{\overset{N}{}}\frac{\mathrm{}^2}{2}_i\left[M^1_i\right]+V_{\mathrm{cb}}(R)+V_{\mathrm{ee}}(R)\right\}\mathrm{\Psi }(R)=E\mathrm{\Psi }(R),$$
(1)
where $`\mathrm{\Psi }(R)=\mathrm{\Psi }(𝐫_1,𝐫_2,\mathrm{}𝐫_N)`$ is the $`N`$-electron wave function, $`M`$ is the electron effective mass tensor and the potential energy terms $`V_{\mathrm{cb}}(R)`$ and $`V_{\mathrm{ee}}(R)`$ are given by
$`V_{\mathrm{cb}}(R)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}v_{\mathrm{cb}}(𝐫_i)`$ (2)
$`V_{\mathrm{ee}}(R)`$ $`=`$ $`{\displaystyle \underset{i<j}{}}{\displaystyle \frac{e^2}{ϵ|𝐫_i𝐫_j|}}`$ (3)
where $`v_{\mathrm{cb}}(𝐫)`$ is the potential for a single electron due to the offset and strain potential of the conduction band, and $`V_{\mathrm{ee}}(R)`$ is the Coulomb interaction of the conduction electrons with charge $`e`$ in a dielectric characterized by the dielectric constant $`ϵ`$. The offset of the conduction band edge causes a step potential at the edge of the dot, which is reduced by strain, resulting in a potential $`v_{\mathrm{cb}}`$ that is zero outside the dot and a roughly constant depth of several tenths of an eV in the interior of the dot, which typically has dimensions $`10\mathrm{nm}`$. Strain causes the electron effective mass to become anisotropic leading to a mass tensor given by $`\mathrm{diag}(M)=(m_{\mathrm{xx}}m_{\mathrm{yy}}m_{\mathrm{zz}})`$ and zero off-diagonal terms. In the usual case of sample growth along the crystal direction (0 0 1), the electron masses along the plane perpendicular to the growth direction are equal, i.e. $`m_{\mathrm{xx}}=m_{\mathrm{yy}}`$.
## II Density Function Theory and the Local Spin Density Approximation
The DFT approach to obtaining the ground state energy is to replace the rather complex N-electron ground state wave function and the associated Schrödinger equation by the much simpler ground state electron density $`\rho (𝐫)`$ and the corresponding functional forms $`T[\rho ]`$ and $`V[\rho ]`$ of the kinetic and potential energy operators, respectively. However, those functional forms are unknown, and approximations are necessary. In Kohn-Sham theory the density is given by a sum of densities of single particle orbitals,
$$\rho (𝐫)=\underset{\sigma }{}\underset{i=1}{\overset{N}{}}n_i|\psi _i(𝐫,\sigma )|^2.$$
(4)
The total energy can then written as
$$E=E_{\mathrm{kin}}+E_{\mathrm{cb}}[\rho ]+E_\mathrm{H}[\rho ]+E_\mathrm{x}[\rho ]+E_\mathrm{c}[\rho ],$$
(5)
where
$$E_{\mathrm{kin}}=\frac{\mathrm{}^2}{2}\underset{\sigma ,i=1}{\overset{N}{}}\psi _i^{}(𝐫)(M^1)\psi _i(𝐫,\sigma )𝑑𝐫$$
(6)
is the kinetic energy of the Kohn-Sham orbitals,
$$E_{\mathrm{cb}}[\rho ]=\rho (𝐫)v_{\mathrm{cb}}(𝐫)𝑑𝐫,$$
(7)
is the potential energy from the conduction band offset and strain,
$$E_\mathrm{H}[\rho ]=e\rho (𝐫)\varphi (𝐫)𝑑𝐫,$$
(8)
is the Hartree energy, and $`E_\mathrm{x}[\rho ]`$ and $`E_\mathrm{c}[\rho ]`$ are functionals of the electron density. The self-consistent electrostatic potential $`\varphi (𝐫)`$ is obtained as a solution to the Poisson equation
$$^2\varphi (𝐫)=\frac{e}{ϵ}\rho (𝐫)$$
(9)
with the boundary condition $`\varphi (𝐫)0`$ as $`|𝐫|\mathrm{}`$. Minimization of the total energy defined by Eqs. (4 \- 9) yields a one particle Schrödinger equation which may be solved self-consistently with the expression for the density and the Poisson equation to determine the exact many-body ground state energy. However, the nonlocal functionals $`E_\mathrm{x}[\rho ]`$ and $`E_\mathrm{c}[\rho ]`$ are in general not known, so an approximation must be made in order to proceed with the Kohn-Sham approach.
In the local spin density approximation (LSDA) the exchange and correlation functionals are taken to be functions of the local charge and spin density, which are then uniquely defined by the requirement that the approximation be exact for the homogeneous electron gas. These functions have been fit to the exchange and correlation of the uniform electron gas. The validity of LSDA has been a subject of much research in atomic, crystalline, molecular, and parabolic quantum dot systems, but has not been checked for self-assembled dot potentials, which are qualitatively like finite square wells with rather high electron density. The density gradients and shell structures of self-assembled dots cause the exchange and correlation energy to differ from that of a uniform gas. The purpose of this paper is to quantify this difference as a function of the number of electrons in a realistically modeled dot.
## III Quantum Monte Carlo Methods
The simplest QMC method is variational Monte Carlo (VMC), in which expectation values for a trial wave function are evaluated using a Metropolis algorithm. Using the variational principle of the ground state energy, parameters in the wave function can be optimized to minimize the total energy, thus providing an approximation to the ground state energy and wave function. In our calculation we use VMC to optimize the guiding wave function for our diffusion Monte Carlo algorithm, described below.
We also use Monte Carlo integration to evaluate the exact exchange energy. In the Kohn-Sham formalism, the exchange energy is defined as the difference
$$E_\mathrm{x}[n]=\mathrm{\Psi }_{\mathrm{KS}}|V_{\mathrm{ee}}|\mathrm{\Psi }_{\mathrm{KS}}E_\mathrm{H}[\rho ],$$
(10)
where $`|\mathrm{\Psi }_{\mathrm{KS}}`$ is a Slater determinant of the true Kohn-Sham orbitals; those that give the exact density. We assume that LSDA is accurate enough to provide approximations to the Kohn-Sham orbitals, i. e. that the LSDA density is close to the exact density. We then evaluate the integral $`\mathrm{\Psi }_{\mathrm{KS}}|V_{\mathrm{ee}}|\mathrm{\Psi }_{\mathrm{KS}}`$ using a Slater determinant of our calculated LSDA wave functions.
Diffusion Monte Carlo is a stochastic method that is able to project the N-electron ground state wave function $`\mathrm{\Phi }_0`$ from a trial wave function $`\mathrm{\Psi }_\mathrm{T}`$. A position basis is used to described the state of the system. A good trial wave function is important because the variance of the Monte Carlo sampling decreases as the trial wave function approaches the true ground state wave function. We chose the trial wave function to be a product of Slater determinants for spin up and spin down electrons with a Jastrow factor that introduces correlation,
$$\mathrm{\Psi }_\mathrm{T}(R)=det[\theta _\mathrm{k}(𝐫_\mathrm{i})]\times \mathrm{exp}[\underset{\mathrm{i}<\mathrm{j}}{}u(r_{\mathrm{ij}})].$$
(11)
For convenience, we take the single particle wave functions $`\theta _\mathrm{k}(𝐫)`$ to be the non-interacting eigenstates for the external potential, and choose the Jastrow factor $`u(r)=ar/(1+br)`$ with $`a=m/ϵ`$ and $`b=1/2`$.
Following the method described in Ref. \[\], we used importance sampling from a function $`\mathrm{\Phi }(R)=\mathrm{\Psi }_\mathrm{T}(R)\mathrm{\Phi }_0(R)`$. The anisotropic mass is taken into account using the mass tensor, $`M`$. The projection of a state $`R`$ to the ground state distribution $`\mathrm{\Psi }_\mathrm{T}\mathrm{\Phi }_0`$ is accomplished by repeatedly applying the projection operator $`e^{H\tau }`$, each time sampling a new configuration. Here $`H`$ is the Hamiltonian and $`\tau `$ is a parameter chosen to be small so that the projection operator is approximated by
$`R|`$ $`\mathrm{\Psi }_\mathrm{T}(R)e^{H\tau }\mathrm{\Psi }_\mathrm{T}^1(R^{})|R^{}`$ (14)
$`={\displaystyle \frac{\mathrm{exp}\{[E_\mathrm{L}(R)E_0]\tau \}}{[(2\pi \tau )^3detM^1]^{\frac{n}{2}}}}\times `$
$`{\displaystyle \underset{i=1}{\overset{n}{}}}\mathrm{exp}\left[{\displaystyle \frac{(𝐫_\mathrm{i}𝐫_\mathrm{i}\tau 𝐯_{\mathrm{d},\mathrm{i}})M(𝐫_\mathrm{i}𝐫_\mathrm{i}\tau 𝐯_{\mathrm{d},\mathrm{i}})}{2\tau }}\right],`$
where $`E_0`$ is a constant parameter to control the population, $`E_\mathrm{L}(R)=\mathrm{\Psi }_\mathrm{T}(R)^1H\mathrm{\Psi }_\mathrm{T}(R)`$ is the local energy, and $`𝐯_\mathrm{d}=M^1\mathrm{log}\mathrm{\Psi }_\mathrm{T}`$. The algorithm is thus: (1) Start with an ensemble of configurations (walkers) distributed by $`|\mathrm{\Psi }_\mathrm{T}|^2`$, (2) propagate each configuration with a drift $`𝐯_\mathrm{d}`$ and Gaussian displacement with covariance matrix $`\sigma _{\mathrm{ij}}^2=M_{\mathrm{ij}}^1\tau `$, (3) reweight each configuration by a factor $`\mathrm{exp}\{(E_\mathrm{L}E_0)\tau \}`$, (4) repeat, collecting statistics once the steady state distribution is reached. The exact ground state energy is obtained by sampling the local energy. We use branching at each step to improve the efficiency of the process.
A complication known as the fermion sign problem arises when applying diffusion Monte Carlo to fermions. Electron exchange introduces negative signs into the projection operator, which decreases the efficiency of the Monte Carlo sampling. In the present discussion the short time approximation, Eq. (14), breaks down when walkers cross nodes of $`\mathrm{\Psi }_\mathrm{T}`$. We handle this problem with the fixed node approximation, in which walkers are given a weight of zero when they cross a node of $`\mathrm{\Psi }_\mathrm{T}`$. This approximation has a lower bound property, so the true ground state energy can only be lower than the fixed node energy, and the exact choice of $`\mathrm{\Psi }_\mathrm{T}=\mathrm{\Phi }_0`$ would give the exact Fermionic ground state energy.
The correlation energy can be deduced from a knowledge of all other energies,
$$E_c[\rho ]=E_{\mathrm{tot}}E_{\mathrm{kin}}E_{\mathrm{cb}}E_\mathrm{H}E_\mathrm{x},$$
(15)
or equivalently
$$E_c[\rho ]=\mathrm{\Phi }_0|V_{\mathrm{ee}}|\mathrm{\Psi }_0\mathrm{\Psi }_{\mathrm{KS}}|V_{\mathrm{ee}}|\mathrm{\Psi }_{\mathrm{KS}}.$$
(16)
## IV Numerical Comparison
Using DFT and DMC we have calculated the total ground state energy of the system of electrons, defined for DFT by Eq. 5, and defined within DMC as the fixed node ground state energy, as defined in the previous section. Figure 1 shows the system we have considered, a pyramidal shaped InAs quantum dot in a GaAs matrix. However, our method of analysis is not restricted to this particular shape, and we expect our results to be valid for truncated pyramidal or lens shaped dots. We have introduced several simplifications in the model. Since our DMC scheme does not allow for a position dependent dielectric constant, we have used only one value, $`\overline{ϵ}=14ϵ_0`$, for both GaAs and for InAs. In the absence of strain the GaAs - InAs conduction band offset between the two materials is 770 meV. To partially take into account strain effects, we have assumed a constant shift of 400 meV in the InAs conduction band, and no strain effect in the GaAs conduction band, resulting in a conduction band offset of $`770\mathrm{m}\mathrm{e}\mathrm{V}400\mathrm{m}\mathrm{e}\mathrm{V}=370\mathrm{m}\mathrm{e}\mathrm{V}`$. We also take the electron effective masses to be constant in each material.
In order to study a realistic potential, we have used a nonuniform grid basis for the LSDA calculation and the solution of the Poisson equation. For the Schrödinger equation we set the wave functions to zero at the edge of the grid, which is reasonable due to the exponential decay of bound states, but the Poisson equation requires more care. We used a multipole expansion up to the quadrupole term to simulate the boundary conditions at infinity. This difficulty with the boundary conditions of the Poisson equation is one source of error in our comparison. We used cubic interpolation to map the gridded wave functions and conduction band potential to continuous functions for Monte Carlo evaluation of the exchange energy, for determination of the trial nodes in the DMC calculation, and for guiding wave functions in the DMC calculation. Although the cubic interpolation of the potential allows us to compare LSDA and DMC for similar external potentials $`v_{\mathrm{cb}}`$, errors in the kinetic energy operator are more difficult. The use of a finite difference approximation to the Laplacian in the computation of the LSDA solution creates an operator that has no simple counterpart in a continuous formulation of the problem. The error between the physical Laplacian operator and the artificial finite difference operator can be made arbitrarily small by the use of finer grids, but computer memory and CPU time constraints caused this to be a significant source of error in our comparison. The size of the errors are several meV, which are comparable to the rather small LSDA errors, especially in the correlation energy. Below we discuss a way we estimated this error so that we can compare correlation energy.
For the set of calculations we describe now, we have taken all masses to have the isotropic value of $`0.05m_e`$, for both the InAs and GaAs regions. The case of anisotropic masses will be discussed later. We have also accounted for numerical discrepancies between the two calculations due to grid interpolation error. To estimate this error we have compared eigenvalues of single electron states with both methods, and found that the DMC single particle eigenvalues typically lie about 0.25 meV above the eigenvalues computed by the grid method used in the DFT calculation. We attribute that difference to a systematic error in the finite difference kinetic energy operator and shift up our LSDA calculations to compensate for this.
Figure 2 shows the values of the different components of the total energy as a function of dot occupation. Differences between LSDA and QMC are not apparent on the scale of the figure. This figure clearly shows that the external potential energy and kinetic energy are much larger than the interaction energies. In other words the effects of interactions enter as a perturbative correction to the non-interacting system. The reason for this can be seen from the scaled electron density. The energy and length scales for the electron interaction are scaled by the dielectric constant and mass, so that the effective Bohr radius is $`a_0^{}=ϵ/m^{}a_0150\mathrm{\AA }`$ and the effective Hartree is $`\mathrm{Ha}^{}=ϵ^2m^{}\mathrm{Ha}7\mathrm{m}\mathrm{e}\mathrm{V}`$. If we estimate the density by assuming the $`N`$ electrons uniformly occupy the interior of the dot, we find an effective conduction electron density of $`r_s0.46N^{1/3}`$. To see the consequences of such a high effective density, consider the uniform electron gas at $`r_s=0.6`$. The electron gas has a ground state energy expansion for small $`r_s`$,
$$E=2.2099r_s^20.9163r_s^10.094+0.0622\mathrm{ln}(r_s)+\mathrm{},$$
(17)
where the first term is the kinetic energy, the next term is the exchange, and remaining terms are correlation energy. For the case of six electrons in the dot $`r_s0.25`$ and the expansion gives $`E_{\mathrm{kin}}1440\mathrm{m}\mathrm{e}\mathrm{V}`$, $`E_\mathrm{x}150\mathrm{m}\mathrm{e}\mathrm{V}`$, and $`E_\mathrm{c}0.6\mathrm{meV}`$. Although the comparison between these very different electronic systems cannot be pushed too far, this does show that our exchange and correlation energies are reasonable for this electron density. The leading effect of the interactions is the Hartree energy, with small corrections for exchange and very small correlation corrections.
The most important contribution of LSDA in this system is the exchange term, and we plot the comparison of the LSDA exchange energy in Figure 3. For finite systems, the definition for exchange includes the correction for self-interaction in the Hartree energy, which we also show in the figure. Merely correcting for self-interaction will recover at least 75% of the exchange energy, and LSDA is able to do better, recovering about 90% for an error in the exchange energy of less than 5 meV. As is well known in many other systems , the error in the local approximation to the exchange is the largest error in LSDA.
By comparison the correlation energy is much smaller and its accuracy is difficult to assess due to grid errors. Correlation is defined as the difference in total energies, Eq. (16), and in this case the total energy difference is less than 1% . Thus our estimates of the exact correlation energy are fairly uncertain. In Figure 4 we show our best attempt at checking the LSDA correlation in the quantum dot. Since the LSDA wave functions are not the true Kohn-Sham orbitals (due mostly to grid errors) the correlation is overestimated due to relaxation of our interpolated LSDA states. As stated earlier, we correct for this by performing single particle calculations to estimate the relaxation of our LSDA states, and have removed this contribution from our calculated correlation energy. These results are shown as “QMC” and “QMC with corrections” in the figure. We find that LSDA gives a reasonable estimate of the magnitude of the correlation energy in this system.
Our comparison of the energy errors from the LSDA are summarized in Figure 5. The largest error is due to LSDA underestimating the magnitude of the exchange energy, $`E_\mathrm{x}`$, which is a 6 meV error for six electrons in the dot. There is a small error in the correlation energy, $`E_\mathrm{c}`$, which slightly compensates for the error in the exchange. The error in the Hartree energy $`E_{\mathrm{Ha}}`$ is due primarily to truncation of the multipole expansion at quadrupole terms for determination of the boundary conditions of the Poison equation. There are also errors in $`E_{\mathrm{kin}}`$ and $`E_{\mathrm{cb}}`$ (not shown in the figure) due to the grid, in particular the discretization of the kinetic energy operator discussed earlier. We therefore find that LSDA leads to an error in the total energy of up to 6 meV, which in our calculation is partially compensated for by grid errors in other energy terms, leading to an overall error in the total energy of no more than 2 meV.
Because of strain, the electron effective mass in the plane parallel to the base of the dot is usually different from the mass perpendicular to it. The LSDA does not explicitly account for mass anisotropy in the exchange and correlation functionals, thus for the calculation of the exchange-correlation potential $`V_{\mathrm{xc}}`$ we have assumed a constant average mass given by
$$m_{\mathrm{xc}}=(m_{}^2m_{})^{1/3},$$
(18)
where $`m_{}`$ is the in-plane mass, $`m_{}=m_{\mathrm{xx}}=m_{\mathrm{yy}}`$, and $`m_{}`$ is the perpendicular mass, $`m_{}=m_{\mathrm{zz}}`$. The choice of the average mass $`m_{\mathrm{xc}}`$ is somewhat arbitrary. We have performed a test of this approximation by calculating the exchange and correlation energy for an anisotropic system using (i) LSDA with $`m_{\mathrm{xc}}`$ as given above, (ii) LSDA with $`m_{\mathrm{xc}}`$ optically averaged, $`m_{\mathrm{xc}}^1=(2/m_{}+1/m_{})/3`$, and (iii) DMC, which can explicitly treat anisotropic mass. Figure 6 shows the calculated exchange and correlation energies for several electron mass ratios, which we have again fixed to constant across the InAs and GaAs. We find that either choice is acceptable, as both give errors comparable to the isotropic mass case.
## V Conclusions
In conclusion, we have shown that DFT offers an accurate approximation for the ground state energy including many-body interactions in small self-assembled quantum dots, while providing a simple and fast means of modeling systems containing several electrons. This is in large part due to the small size of the self-assembled quantum dots in relation to the effective Bohr radius. In this regime the largest error is in the local approximation to the exchange, and we have verified that the errors due to this approximation are small. This is quantified in Figures 2 and 3 where we see that the exchange energy for six electrons in the dot is 72 meV with an LSDA error of 6 meV, compared to a total energy of 1490 meV. We expect this applicability of LSDA to hold in general for small quantum dots of various shapes, but we emphasize that LSDA is expected to become progressively worse if the dot becomes much larger than the effective Bohr radius. Also, the use of the EMA Schrödinger equation, Eq. (1), although reasonable for the chosen problem, is predicted to fail for very small dots. Although our model assumes pure InAs dots surrounded by GaAs, in some systems indium and gallium may in fact intermix. Again, we expect the relative magnetude of exchange and correlation energies and the acceptability of LSDA to be unchanged by intermixing or alloying of the dots.
Based on this comparison we conclude that the addition and removal energies found by EMA-LSDA calculations in realistic dots are accurate to 1-2 meV per electron as far as the many-body interactions within the EMA are concerned. Thus, comparison of such calculated charging energies with experiment can be considered a direct test of the models and the uncertainties in the analysis of experiments. Furthermore, we have carefully tested the grid errors and concluded that energies are accurate to $`\pm 10`$ meV including all errors.
We would like to point out that application of LSDA to a system of several coupled self-assembled dots may have difficulties. Systems that have weak intersite transitions along with intrasite interactions are a well-known many-body problem where LSDA is known to fail. The problem is closely related to the $`\mathrm{H}_2`$ molecule, where LSDA describes well the electronic energy at equilibrium distance, but gives an incorrect broken symmetry solution if the atoms are pulled apart beyond a critical distance. The problem of weakly coupled dots would be of particular importance for future studies.
## Acknowledgments
This work was supported by CRI from the University of Illinois and NSF Grants No. ECS 95-09751 and No. DMR 94-22496 and computer resources at NCSA. |
warning/0003/hep-th0003079.html | ar5iv | text | # Dynamics of Multiparticle Systems with non – Abelian Symmetry
## 1 Introduction and Overview
The consideration of the influence of internal symmetries on the final state of a many body system begun with the pioneering work of Bethe . Much of the subsequent interest in the subject arises from the realization that in the study of hadronic interactions and in particular in studies involving quark confinement, these constraints may be of decisive importance. An important progress in treating equilibrium systems was made employing group projection techniques. This allowed for a consistent treatment of abelian and nonabelian symmetries of compact groups and a consistent formulation of thermodynamics of many particle systems with internal symmetries taken into account . Application of these methods to specific processes demonstrated in which circumstances the presence of symmetry is of physical relevance .
However, it is not fully understood how the symmetry-modified properties of the equilibrium system arise from kinetic formulation of the dynamical evolution. When an internal symmetry is not at work, Boltzmann’s H theorem in principle assures that the statistical Bose/Fermi/Boltzmann distributions are the asymptotic (equilibrium) distributions, irrespective of the nature of microscopic interaction. However, in presence of exact symmetries the equilibrium distributions are modified, see e.g. . This implies that symmetry constraints introduce effective interactions of potentially far more complex nature than is the usual two body Boltzmann collision term. In fact it can be argued that quantum symmetry constraints are the heart of the nonlocality of quantum physics. However, in the limit of classical Boltzmann equation evolution these are implemented by a strictly local (though non-linear) consideration of Fermi blocking and Bose enhancement in phase space evolution. Our aim in this work is to make a step towards understanding how the microscopic nonabelian symmetry constrains operate within the kinetic master equation description of the time evolution, leading on to the symmetry modified (constrained) macroscopic many particle equilibrium state.
It is first important to convince oneself that an underlying symmetry of microscopic interactions does not lead in general to the desired symmetry properties of a (macroscopic) many body interacting system. To do this, we consider the high energy nuclear (heavy ion) collisions and specifically here two symmetry examples:
a) SU(2) Isospin symmetry: The initial state transforms under a given representation of the isospin $`SU(2)`$ group. All elementary high energy interactions are governed by the strong interaction, which preserves the isotopic symmetry. A final state results as a multiparticle state formed by many individual hadron – hadron collisions. In any of such microprocesses the isospin is conserved. However, proceeding ‘as usual’ without symmetry constrained treatment of local interactions does not assure that the final multiparticle state (macrostate) transforms under the same representation of the isospin group as the initial state, which is required for symmetry reasons.
b) SU(3) Colour symmetry: A similar situation appears in the context of the quark-gluon interactions, especially in case that local deconfinement occurs. The initial state is a colour singlet state, and quark-gluon and gluon-gluon interaction, although invariant under the colour $`SU(3)_c`$ symmetry group do not assure that during its evolution a (macroscopic) many particle state, once a singlet, always remains a singlet colour state, which, however, it must do because of exact colour symmetry of strong interactions.
As these examples show, in a dynamical (quantum) transport theory description of the approach to equilibrium there must exist a subsidiary condition which should be taken into account by corresponding kinetic equations governing the evolution. This condition is independent from other constraints related to dynamical gauged internal symmetries. For classical fluid dynamics a dynamical evolution equation addressing gauge symmetry has been proposed by Wong . In many current studies of the dynamics of classical non-abelian fields this proposal continues till today to attract considerable interest . However, these are constraints which have no relation to the intrinsic non-locality of the quantum system which we address here.
We first note that in the case of an abelian symmetry there are no additional constraints to consider. Quantum number conservation on a microscopic level is fully equivalent to preservation of all symmetry properties on the macroscopic level. This is easily seen considering the $`U(1)`$ symmetry related to microscopic particle-antiparticle formation: since the microscopic mechanisms produce equal number of particles and antiparticles (pair production), initial particle-antiparticle number difference is exactly preserved in the macroscopic many body state.
Thus only presence of a nonabelian symmetry poses a true challenge. A suitable mathematical method how to approach this problem is identified considering the previously treated statistical equilibrium case. Here one decomposes a general ‘macrostate’ consisting of many particles into possible irreducible representations of the symmetry group. Then a projection technique exploiting character function properties of the group is used to constrain the final state. We will here use this approach in order to describe multiparticle evolution applying microscopic kinetic theory scheme. We will show that the behaviour of particle phase - space distribution functions will depend not only on properties of basic interaction, but that it also depends on global properties of the macroscopic system. Those global properties provide subsidiary constraints needed, so that the asymptotic equilibrium state has properties consistent with the non-abelian constraints.
## 2 The Projection Method
Let $`G`$ be a compact internal symmetry group of our system consisting of particles (objects) transforming under irreducible representations of the symmetry group. These representation are denoted as $`\alpha _i`$ with corresponding dimensions $`d(\alpha _i)`$. One denotes $`f_{(\zeta )}^{(\alpha _i,\nu _i)}(\mathrm{\Gamma },\stackrel{}{r},t)`$ a distribution function of the particle which belongs to the multiplet $`\alpha _i`$ of the symmetry group. Members of this multiplet are numbered by indexes $`\nu _i(\nu _i=1,\mathrm{},d(\alpha _i))`$ which correspond to given values of charges related to the symmetry group. A subscript $`\zeta `$ denotes other quantum numbers characterizing different multiplets of the same representation $`\alpha `$. The variables $`(\mathrm{\Gamma },\stackrel{}{r})`$ denotes a set of the phase - space variables such as $`(\stackrel{}{p},\stackrel{}{r})`$ and $`t`$ is time.
The number of particles of the specie $`\{\alpha ,\nu _\alpha ,\zeta \}`$ is:
$$N_{\nu _\alpha ;(\zeta )}^{(\alpha )}(t)=𝑑V𝑑\mathrm{\Gamma }f_{(\zeta )}^{(\alpha ,\nu _\alpha )}(\mathrm{\Gamma },\stackrel{}{r},t);$$
(1)
We consider a system of $`\{N_{\alpha _1,\nu _{\alpha _1}}^{(\zeta _1)}(t),\mathrm{},N_{\alpha _n,\nu _{\alpha _n}}^{(\zeta _n)}(t)\}`$ particles at time $`t`$ . The distribution functions fulfill the generalized Vlasov \- Boltzmann kinetic equations, which can be written in the general form:
$`{\displaystyle \frac{f_{(\zeta _i)}^{(\alpha _i,\nu _i)}(\mathrm{\Gamma }_i,\stackrel{}{r},t)}{t}}`$ $`+`$ $`\stackrel{}{v}f_{(\zeta _i)}^{(\alpha _i,\nu _i)}(\mathrm{\Gamma }_i,\stackrel{}{r},t)`$ (2)
$`=`$ $`{\displaystyle \underset{\alpha _j,\alpha _k,\alpha _l}{}}{\displaystyle \underset{\nu _j,\nu _k,\nu _l}{}}{\displaystyle \underset{\zeta _j,\zeta _k,\zeta _l}{}}{\displaystyle 𝑑\mathrm{\Gamma }_j𝑑\mathrm{\Gamma }_k𝑑\mathrm{\Gamma }_l𝒲_{\nu _i\nu _j;\nu _k\nu _l}^{(\zeta _i,\zeta _j;\zeta _k,\zeta _l)}(\mathrm{\Gamma }_k,\mathrm{\Gamma }_l;\mathrm{\Gamma }_j,\mathrm{\Gamma }_i)}`$
$`[_{(\zeta _i)}^{(\alpha _i,\nu _i)}(\mathrm{\Gamma }_i,\stackrel{}{r},t)_{(\zeta _j)}^{(\alpha _j,\nu _j)}(\mathrm{\Gamma }_j,\stackrel{}{r},t)f_{(\zeta _k)}^{(\alpha _k,\nu _k)}(\mathrm{\Gamma }_k,\stackrel{}{r},t)f_{(\zeta _l)}^{(\alpha _l,\nu _l)}(\mathrm{\Gamma }_l,\stackrel{}{r},t)`$
$`_{(\zeta _k)}^{(\alpha _k,\nu _k)}(\mathrm{\Gamma }_k,\stackrel{}{r},t)_{(\zeta _l)}^{(\alpha _l,\nu _l)}(\mathrm{\Gamma }_l,\stackrel{}{r},t)f_{(\zeta _i)}^{(\alpha _i,\nu _i)}(\mathrm{\Gamma }_i,\stackrel{}{r},t)f_{(\zeta _j)}^{(\alpha _j,\nu _j)}(\mathrm{\Gamma }_j,\stackrel{}{r},t)];`$
Factors $`_{(\zeta )}^{(\alpha ,\nu )}(\mathrm{\Gamma },\stackrel{}{r},t)`$ are related to quantum statistics and they are equal to $`1`$ for classical particles, and equal to $`[1\pm f_{(\zeta )}^{(\alpha ,\nu )}(\mathrm{\Gamma },\stackrel{}{r},t)]`$ for bosons/fermions correspondingly. Since by assumption the whole system transforms under given representation $`\mathrm{\Lambda }`$ of an exact symmetry group, the system under consideration must preserve its transformations properties during its time evolution, provided that it is governed by a symmetry invariant interaction.
We now focus on the case of a quantum system and consider state vectors in particle number representation: $`|N_{\nu _{\alpha _1}}^{(\alpha _1)},\mathrm{},N_{\nu _{\alpha _n}}^{(\alpha _n)}.`$ These vectors describe symmetry properties of our systems and all other variables, related to phase-space properties of the system are suppressed here. They transform as a direct product representation of the symmetry group $`G`$. This representation is of the form:
$$\alpha _1^{N^{(\alpha _1)}}\alpha _2^{N^{(\alpha _2)}}\mathrm{}\alpha _n^{N^{(\alpha _n)}};$$
(3)
A multiplicity $`N^{(\alpha _j)}`$ of the representation $`\alpha _j`$ in this product is equal to a number of particles which transform under this representation:
$$N^{(\alpha _j)}=\underset{j}{}\left(\underset{\zeta _j}{}N_{\nu _{\alpha _j};(\zeta _j)}^{(\alpha _j)}\right)=\underset{j}{}N_{\nu _{\alpha _j}}^{(\alpha _j)};$$
(4)
The representation given by Eq. (3) can be decomposed into direct sum of irreducible representations $`\mathrm{\Lambda }_k`$. Corresponding states are denoted as $`|\mathrm{\Lambda }_k,\lambda _{\mathrm{\Lambda }_k};𝒩`$ where $`\lambda _{\mathrm{\Lambda }_k}`$ is an index numbering members of the representation $`\mathrm{\Lambda }`$ and $`𝒩`$ is a total number of particles
$$𝒩=\underset{k}{}N_{\nu _{\alpha _k}}^{(\alpha _k)};$$
(5)
Each physical state can be decomposed into irreducible representation base states with amplitudes depending on phase space variables $`\mathrm{\Gamma }`$:
$`|N_{\nu _{\alpha _1}}^{(\alpha _1)},\mathrm{},N_{\nu _{\alpha _n}}^{(\alpha _n)};\mathrm{\Gamma }={\displaystyle \underset{k}{\overset{}{}}}{\displaystyle \underset{\xi _{\mathrm{\Lambda }_k}}{\overset{}{}}}|\mathrm{\Lambda }_k,\lambda _{\mathrm{\Lambda }_k};𝒩;\xi _{\mathrm{\Lambda }_k}a_{\{N_{\nu _{\alpha _1}}^{(\alpha _1)},\mathrm{},N_{\nu _{\alpha _n}}^{(\alpha _n)}\}}^{\mathrm{\Lambda },\lambda _\mathrm{\Lambda }}(\xi _{\mathrm{\Lambda }_k};\mathrm{\Gamma });`$ (6)
Here appear new variables $`\xi _\mathrm{\Lambda }`$ which are degeneracy parameters required for the full description of a state in the ”symmetry space”.
Let us define an average weight
$$\overline{P_{\{N_{\nu _{\alpha _1}}^{(\alpha _1)},\mathrm{},N_{\nu _{\alpha _n}}^{(\alpha _n)}\}}^{\mathrm{\Lambda },\lambda _\mathrm{\Lambda }}}=\frac{\underset{\xi _\mathrm{\Lambda }}{}|a_{\{N_{\nu _{\alpha _1}}^{(\alpha _1)},\mathrm{},N_{\nu _{\alpha _n}}^{(\alpha _n)}\}}^{\mathrm{\Lambda },\lambda _\mathrm{\Lambda }}(\xi _\mathrm{\Lambda };\mathrm{\Gamma })|^2}{\underset{N_{\nu _{\alpha _1}}^{(\alpha _1)}+\mathrm{}+N_{\nu _{\alpha _n}}^{(\alpha _n)}=𝒩}{}\underset{\xi _\mathrm{\Lambda }}{}|a_{\{N_{\nu _{\alpha _1}}^{(\alpha _1)},\mathrm{},N_{\nu _{\alpha _n}}^{(\alpha _n)}\}}^{\mathrm{\Lambda },\lambda _\mathrm{\Lambda }}(\xi _\mathrm{\Lambda };\mathrm{\Gamma })|^2};$$
(7)
This expression gives the probability that $`N_{\nu _{\alpha _1}}^{(\alpha _1)},\mathrm{},N_{\nu _{\alpha _n}}^{(\alpha _n)}`$ particles transforming under the symmetry group representations $`\alpha _1,\mathrm{},\alpha _n`$ combine into $`𝒩`$ particle state transforming under representation $`\mathrm{\Lambda }`$ of the symmetry group.
We make the statistical hypothesis that average weights (7) do not depend on phase - space variables and can be calculated alone on basis of symmetry group consideration. This also can be proved under the stronger assumption that in Eq. (6) any state with fixed $`\mathrm{\Lambda },\lambda _\mathrm{\Lambda }`$ has the same weight (see e.g. )
Let us consider a projection operator $`𝒫^\mathrm{\Lambda }`$ on the subspace spanned by all states transforming under representation $`\mathrm{\Lambda }`$.
$$𝒫^\mathrm{\Lambda }|N_{\nu _{\alpha _1}}^{(\alpha _1)},\mathrm{},N_{\nu _{\alpha _n}}^{(\alpha _n)}=\underset{\xi _\mathrm{\Lambda }}{\overset{}{}}|\mathrm{\Lambda },\lambda _\mathrm{\Lambda };\xi _\mathrm{\Lambda }𝒞_{\{N_{\nu _{\alpha _1}}^{(\alpha _1)},\mathrm{},N_{\nu _{\alpha _n}}^{(\alpha _n)}\}}^{\mathrm{\Lambda },\lambda _\mathrm{\Lambda }}(\xi _\mathrm{\Lambda });$$
(8)
This operator has the generic form (see e.g. ):
$$𝒫^\mathrm{\Lambda }=d(\mathrm{\Lambda })\underset{G}{}𝑑\mu (g)\overline{\chi }^{(\mathrm{\Lambda })}(g)U(g);$$
(9)
Here $`\chi ^{(\mathrm{\Lambda })}`$ is the character of the representation $`\mathrm{\Lambda }`$, $`d\mu (g)`$ is the invariant Haar measure on the group, and $`U(g)`$ is an operator transforming a state under consideration. We will use the matrix representation:
$`U(g)|N_{\nu _{\alpha _1}}^{(\alpha _1)},\mathrm{},N_{\nu _{\alpha _n}}^{(\alpha _n)}`$ (10)
$`=`$ $`{\displaystyle \underset{\nu _1^{(1)},\mathrm{},\nu _n^{(N_{\nu _n})}}{}}D_{\nu _1^{(1)}\nu _1}^{(\alpha _1)}\mathrm{}D_{\nu _1^{(N_{\nu _1})}\nu _1}^{(\alpha _1)}\mathrm{}D_{\nu _n^{(1)}\nu _n}^{(\alpha _n)}\mathrm{}D_{\nu _n^{(N_{\nu _n})}\nu _n}^{(\alpha _n)}|N_{\nu _{\alpha _1}}^{(\alpha _1)},\mathrm{},N_{\nu _{\alpha _n}}^{(\alpha _n)};`$
$`D_{\nu ,\nu }^{(\alpha _n)}`$ is a matrix elements of the group element $`g`$ corresponding to the representation $`\alpha `$. Notation convention in Eq. (10) arises since there are $`N_{\nu _{\alpha _j}}^{(\alpha _j)}`$ states transforming under representation $`\alpha _j`$ and having quantum numbers of the $`\nu _{\alpha _j}`$-th member of a given multiplet.
The statistical hypothesis identifies the average weight $`\overline{P_{\{N_{\nu _{\alpha _1}}^{(\alpha _1)},\mathrm{},N_{\nu _{\alpha _n}}^{(\alpha _n)}\}}^{\mathrm{\Lambda },\lambda _\mathrm{\Lambda }}}`$ with a norm of the vector $`𝒫^\mathrm{\Lambda }|N_{\nu _{\alpha _1}}^{(\alpha _1)},\mathrm{},N_{\nu _{\alpha _n}}^{(\alpha _n)}`$. This norm can be written as
$`N_{\nu _{\alpha _1}}^{(\alpha _1)},\mathrm{},N_{\nu _{\alpha _n}}^{(\alpha _n)}\left|𝒫^\mathrm{\Lambda }\right|N_{\nu _{\alpha _1}}^{(\alpha _1)},\mathrm{},N_{\nu _{\alpha _n}}^{(\alpha _n)}={\displaystyle \underset{\xi _\mathrm{\Lambda }}{}}|𝒞_{\{N_{\nu _{\alpha _1}}^{(\alpha _1)},\mathrm{},N_{\nu _{\alpha _n}}^{(\alpha _n)}\}}^{\mathrm{\Lambda },\lambda _\mathrm{\Lambda }}(\xi _\mathrm{\Lambda })|^2;`$ (11)
where the relation $`(𝒫^\mathrm{\Lambda })^2=𝒫^\mathrm{\Lambda }`$ was used.
Left hand side of this equation can be calculated directly from Eqs.(9) and (10). One gets finally
$$\overline{P_{\{N_{\nu _{\alpha _1}}^{(\alpha _1)},\mathrm{},N_{\nu _{\alpha _n}}^{(\alpha _n)}\}}^{\mathrm{\Lambda },\lambda _\mathrm{\Lambda }}}=𝒜^{\{𝒩\}}d(\mathrm{\Lambda })\underset{G}{}𝑑\mu (g)\overline{\chi }^{(\mathrm{\Lambda })}(g)[D_{\nu _1\nu _1}^{(\alpha _1)}]^{N_{\nu _{\alpha _1}}^{(\alpha _1)}}\mathrm{}[D_{\nu _n\nu _n}^{(\alpha _n)}]^{N_{\nu _{\alpha _n}}^{(\alpha _n)}};$$
(12)
where $`𝒜^{\{𝒩\}}`$ is a permutation normalization factor. For particles of the kind $`\{\alpha ,\zeta \}`$ we included in Eq. (12) the permutation factor:
$$𝒜_{(\zeta )}^\alpha =\frac{𝒩_{(\zeta )}^{(\alpha )}!}{\underset{\nu _\alpha }{}𝒩_{\nu _\alpha ;(\zeta )}^{(\alpha )}!};$$
(13)
The permutation factor $`𝒜^{\{𝒩\}}`$ is a product of all ”partial” factors
$$𝒜^{\{𝒩\}}=\underset{j}{}\underset{\zeta _j}{}𝒜_{(\zeta _j)}^{\alpha _j};$$
(14)
The permutation factor assures the normalization of state vectors:
$$N_{\nu _{\alpha _1}}^{(\alpha _1)},\mathrm{},N_{\nu _{\alpha _n}}^{(\alpha _n)}|N_{\nu _{\alpha _1}}^{(\alpha _1)},\mathrm{},N_{\nu _{\alpha _n}}^{(\alpha _n)}=𝒜^{\{𝒩\}};$$
(15)
This normalization reflects an invariance of the state vector with respect to permutations which shuffle indistinguishable particles.
## 3 Incorporation of Symmetry
The expression Eq. (12) is a starting point for further considerations. It provides together with Eq. (1) and Eq. (2) subsidiary constraints on distribution functions $`f^{(\alpha _i,\nu _i)}`$. These conditions assure that in a dynamical evolution the symmetry of the system is preserved. When symmetry is conserved, then all weights in Eq. (12) are constant in time. In a case of strong interaction and colour symmetry, all weights, except for the weight corresponding to the singlet state, must remain zero.
We now convert the global constraint into a time evolution condition and consider:
$$\frac{d}{dt}\overline{P_{\{N_{\nu _{\alpha _1}}^{(\alpha _1)},\mathrm{},N_{\nu _{\alpha _n}}^{(\alpha _n)}\}}^{\mathrm{\Lambda },\lambda _\mathrm{\Lambda }}}=0;$$
(16)
Introducing here the result of Eq. (12) one obtains:
$`0`$ $`=`$ $`{\displaystyle \frac{d𝒜^{\{𝒩\}}}{dt}}d(\mathrm{\Lambda }){\displaystyle \underset{G}{}}𝑑\mu (g)\overline{\chi }^{(\mathrm{\Lambda })}(g)[D_{\nu _1\nu _1}^{(\alpha _1)}]^{N_{\nu _{\alpha _1}}^{(\alpha _1)}}\mathrm{}[D_{\nu _n\nu _n}^{(\alpha _n)}]^{N_{\nu _{\alpha _n}}^{(\alpha _n)}}`$ (17)
$`+`$ $`{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \underset{\nu _{\alpha _j}}{}}{\displaystyle \frac{dN_{\nu _{\alpha _j}}^{(\alpha _j)}}{dt}}𝒜^{\{𝒩\}}d(\mathrm{\Lambda }){\displaystyle \underset{G}{}}𝑑\mu (g)\overline{\chi }^{(\mathrm{\Lambda })}(g)[D_{\nu _1\nu _1}^{(\alpha _1)}]^{N_{\nu _{\alpha _1}}^{(\alpha _1)}}\mathrm{}[D_{\nu _n\nu _n}^{(\alpha _n)}]^{N_{\nu _{\alpha _n}}^{(\alpha _n)}}\mathrm{log}[D_{\nu _j\nu _j}^{(\alpha _j)}];`$
All integrals which appear in Eq. (12) and Eq. (17) can be expressed explicitly in an analytic form for any compact symmetry group.
To write an expression for the time derivative of the normalization factor $`𝒜^{\{𝒩\}}`$ we perform analytic continuation from integer to continuous values of variables $`N_{\nu _{\alpha _n}}^{(\alpha _n)}.`$ Thus we replace all factorials by the $`\mathrm{\Gamma }`$–function of corresponding arguments. We encounter here also the digamma function $`\psi `$ :
$$\psi (x)=\frac{d\mathrm{log}\mathrm{\Gamma }(x)}{dx};$$
(18)
This allows to write for Eq. (17):
$$\frac{d𝒜^{\{𝒩\}}}{dt}=𝒜^{\{𝒩\}}\underset{j}{}\underset{\zeta _j}{}\left[\frac{d𝒩_{(\zeta _j)}^{(\alpha _j)}}{dt}\psi (𝒩_{(\zeta _j)}^{(\alpha _j)}+1)\underset{\nu _{\alpha _j}}{}\frac{d𝒩_{\nu _{\alpha _j};(\zeta _j)}^{(\alpha _j)}}{dt}\psi (𝒩_{\nu _\alpha ;(\zeta _j)}^{(\alpha _j)}+1)\right];$$
(19)
To get a consistent analytical continuation in the number of particles one should define the time derivatives $`d𝒩_{\nu _\alpha ;(\zeta )}^{(\alpha )}/dt`$. We define these rates of particle number change from the integrated Boltzmann kinetic equation, Eq. (2), explicitly
$`{\displaystyle \frac{dN_{\nu _{\alpha _i}}^{(\alpha _i)}}{dt}}`$ $`=`$ $`{\displaystyle \underset{\alpha _j,\alpha _k,\alpha _l}{}}{\displaystyle \underset{\nu _j,\nu _k,\nu _l}{}}{\displaystyle \underset{\zeta _j,\zeta _k,\zeta _l}{}}{\displaystyle 𝑑V𝑑\mathrm{\Gamma }_j𝑑\mathrm{\Gamma }_k𝑑\mathrm{\Gamma }_l𝑑\mathrm{\Gamma }_i𝒲_{\nu _i\nu _j;\nu _k\nu _l}^{(\zeta _i,\zeta _j;\zeta _k,\zeta _l)}(\mathrm{\Gamma }_k,\mathrm{\Gamma }_l;\mathrm{\Gamma }_j,\mathrm{\Gamma }_i)}`$ (20)
$`[_{(\zeta _i)}^{(\alpha _i,\nu _i)}(\mathrm{\Gamma }_i,\stackrel{}{r},t)_{(\zeta _j)}^{(\alpha _j,\nu _j)}(\mathrm{\Gamma }_j,\stackrel{}{r},t)f_{(\zeta _k)}^{(\alpha _k,\nu _k)}(\mathrm{\Gamma }_k,\stackrel{}{r},t)f_{(\zeta _l)}^{(\alpha _l,\nu _l)}(\mathrm{\Gamma }_l,\stackrel{}{r},t)`$
$`_{(\zeta _k)}^{(\alpha _k,\nu _k)}(\mathrm{\Gamma }_k,\stackrel{}{r},t)_{(\zeta _l)}^{(\alpha _l,\nu _l)}(\mathrm{\Gamma }_l,\stackrel{}{r},t)f_{(\zeta _i)}^{(\alpha _i,\nu _i)}(\mathrm{\Gamma }_i,\stackrel{}{r},t)f_{(\zeta _j)}^{(\alpha _j,\nu _j)}(\mathrm{\Gamma }_j,\stackrel{}{r},t)];`$
Contributions from gradient terms of Eq. (2) vanish due to Gauss law. These terms are transformed in surface integrals and beyond the volume occupied by the system all distribution functions are equal to zero.
Eqs (17,19,20) in fact constitute the global subsidiary condition which should be fulfilled by the microscopic kinetic equations Eq. (2). These are the necessary conditions for preserving the internal symmetry on the macroscopic level. Rates of change $`d𝒩_{\nu _\alpha ;(\zeta )}^{(\alpha )}/dt`$ are related to “macrocurrents”, which are counterparts of “microcurrents” related directly to a symmetry on a microscopic level via the Noether theorem. Eq. (20) can be considered as a set of conditions on macrocurrents to provide consistency with the overall symmetry of the system. Therefore we believe that this equation can also be used as a starting point for multicomponent hydrodynamic equations with internal symmetry properties taken into account. One should notice that this subsidiary condition takes into account also surface effects for the finite volume systems. This is due to the space variables integration which is performed in Eqs (1) and (20).
One easily sees that for the case of abelian symmetry the two constraints Eq. (12) and Eq. (17) do not lead to new results: first we recall that all irreducible representations of abelian group are one-dimensional. Next, let basic particles have “charges” $`q_1,\mathrm{},q_n`$ , and let the global charge be $`Q`$. Then the only consequence of Eq. (12) follows for nonvanishing weight $`Q=N_1q_1+\mathrm{}+N_nq_n`$, which is a rather obvious result. New results appear only for nonabelian symmetries.
## 4 Example: Isospin
We now consider as an example the case of the $`SU(2)`$ symmetry with basic particles transforming under spinor $`(\frac{\mathrm{𝟏}}{\mathrm{𝟐}})`$ (fundamental) and vector $`(\mathrm{𝟏})`$ (adjoint) representations. This example can be realized by a gas mixture of nucleons and pions. To describe all group elements the three group’s parameters $`\alpha ,\beta ,\gamma `$ are chosen in such a way that diagonal matrix elements have the well known form :
i.) for the fundamental representation $`(\frac{\mathrm{𝟏}}{\mathrm{𝟐}})`$:
$`D_{mm}^{(1/2)}(\alpha ,\beta ,\gamma )=e^{im(\alpha +\gamma )}\mathrm{cos}{\displaystyle \frac{\beta }{2}};m=\pm {\displaystyle \frac{1}{2}};`$ (21)
ii.) and for the adjoint representation $`(\mathrm{𝟏})`$:
$`D_{\pm 1,\pm 1}^{(1)}(\alpha ,\beta ,\gamma )`$ $`=`$ $`{\displaystyle \frac{1}{2}}e^{\pm im(\alpha +\gamma )}(1+\mathrm{cos}\beta );`$ (22)
$`D_{0,0}^{(1)}(\alpha ,\beta ,\gamma )`$ $`=`$ $`\mathrm{cos}\beta ;`$ (23)
The Haar measure for the $`SU(2)`$ group in this parametrization has the form
$`{\displaystyle 𝑑\mu (g)f[g]}={\displaystyle \frac{1}{8\pi ^2}}{\displaystyle \underset{0}{\overset{2\pi }{}}}𝑑\alpha {\displaystyle \underset{0}{\overset{2\pi }{}}}𝑑\gamma {\displaystyle \underset{0}{\overset{\pi }{}}}𝑑\beta \mathrm{sin}\beta f[g(\alpha ,\beta ,\gamma )];`$ (24)
Any ‘macrostate’ is made of an arbitrary number: $`n_n,n_p,n_{},n_0,n_+`$; where subscripts refer to members of the fundamental representation neutrons, protons; and members of adjoint representations, $`\pi ^{},\pi ^0,\pi ^+`$, correspondingly. Let us consider the special case when the macrostate is a $`SU(2)`$ singlet. The weight of the singlet state is according to Eq. (12):
$`\overline{P_{\{n_n,n_p,n_{},n_0,n_+\}}^{0,0}}=`$ (25)
$`𝒜^{\{𝒩\}}{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle \underset{0}{\overset{2\pi }{}}}𝑑\alpha {\displaystyle \underset{0}{\overset{2\pi }{}}}𝑑\gamma {\displaystyle \underset{0}{\overset{\pi }{}}}𝑑\beta \mathrm{sin}\beta e^{\frac{i}{2}(n_nn_p+2n_{}2n_+)(\alpha +\gamma )}\mathrm{cos}^{}{\displaystyle \frac{\beta }{2}}\mathrm{cos}^{n_0}\beta `$
$`𝒜^{\{𝒩\}}\stackrel{~}{P}_{\{n_n,n_p,n_{},n_0,n_+\}}^{0,0};`$
where
$$=n_n+n_p+2n_{}+2n_+;$$
(26)
The permutation normalization factor is here:
$$𝒜^{\{𝒩\}}=\frac{(n_{}+n_0+n_+)!(n_n+n_p)!}{n_{}!n_0!n_+!n_n!n_p!};$$
(27)
The real nonzero values of the weight is obtained only when the argument of the exponent in Eq.(̇25) vanishes:
$$n_nn_p+2n_{}2n_+=0;$$
(28)
This is equivalent to the conservation of the third component of the isospin.
Novel behaviour is obtained only when one considers time evolution of the system. Presence of an exact symmetry means that the corresponding weight Eq. (12) is constant, here we consider the expression:
$$\frac{d}{dt}\overline{P_{\{n_n,n_p,n_{},n_0,n_+\}}^{0,0}}=0;$$
(29)
We note that an appropriate analytical continuation should be made. First we evaluate the integral appearing in Eq. (25):
$$\frac{1}{𝒜^{\{𝒩\}}}\overline{P_{\{n_n,n_p,n_{},n_0,n_+\}}^{0,0}}=(1)^{n_0}\underset{i=0}{\overset{n_0}{}}(2)^i\left(\genfrac{}{}{0pt}{}{n_0}{i}\right)\frac{1}{+1+i};$$
(30)
This discrete form is not allowing an analytic continuation which would allow for all necessary differentiations. However, we can write the integral also as the hypergeometric $`{}_{2}{}^{}F_{1}^{}`$ function :
$$\frac{1}{𝒜^{\{𝒩\}}}\overline{P_{\{n_n,n_p,n_{},n_0,n_+\}}^{0,0}}=(1)^{n_0}\frac{1}{+1}{}_{2}{}^{}F_{1}^{}(n_0,+1,+2;2);$$
(31)
where $``$ is as defined in Eq. (26). We so obtain:
$`\overline{P_{\{n_n,n_p,n_{},n_0,n_+\}}^{0,0}}`$
$`={\displaystyle \frac{\mathrm{\Gamma }(n_{}+n_0+n_++1)\mathrm{\Gamma }(n_n+n_p+1)}{\mathrm{\Gamma }(n_{}+1)\mathrm{\Gamma }(n_0+1)\mathrm{\Gamma }(n_++1)\mathrm{\Gamma }(n_n+1)\mathrm{\Gamma }(n_p+1)}}{\displaystyle \frac{\mathrm{cos}n_0\pi }{\mathrm{\Gamma }(n_0)}}{\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{\Gamma }(n_0+i)}{+1+i}}{\displaystyle \frac{2^i}{i!}};`$ (32)
Eqs. (25), (29) and (4), together with the condition (28) result in:
$`0`$ $`=`$ $`2{\displaystyle \frac{\stackrel{~}{P}_{\{n_n,n_p,n_{},n_0,n_+\}}^{0,0}}{}}\left({\displaystyle \frac{dn_n}{dt}}+2{\displaystyle \frac{dn_{}}{dt}}\right)+{\displaystyle \frac{\stackrel{~}{P}_{\{n_n,n_p,n_{},n_0,n_+\}}^{0,0}}{n_0}}{\displaystyle \frac{dn_0}{dt}}`$ (33)
$`+`$ $`{\displaystyle \frac{d\mathrm{log}𝒜^{(𝒩)}}{dt}}\stackrel{~}{P}_{\{n_n,n_p,n_{},n_0,n_+\}}^{0,0};`$
The projection integrals determine the coefficients which are, explicitly:
$`{\displaystyle \frac{\stackrel{~}{P}_{\{n_n,n_p,n_{},n_0,n_+\}}^{0,0}}{n_0}}`$ $`=`$ $`\pi {\displaystyle \frac{\mathrm{sin}n_0\pi }{\mathrm{\Gamma }(n_0)}}{\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{\Gamma }(n_0+i)}{+1+i}}{\displaystyle \frac{2^i}{i!}}`$
$`+{\displaystyle \frac{\mathrm{cos}n_0\pi }{\mathrm{\Gamma }(n_0)}}{\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{\Gamma }(n_0+i)[\psi (n_0)\psi (n_0+i)]}{+1+i}}{\displaystyle \frac{2^i}{i!}}`$
$`=`$ $`(1)^{n_0}{\displaystyle \underset{i=1}{\overset{n_0}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n_0}{i}}\right){\displaystyle \frac{[\psi (1+n_0)\psi (1+n_0i)]}{+1+i}}(2)^i;`$
$`{\displaystyle \frac{\stackrel{~}{P}_{\{n_n,n_p,n_{},n_0,n_+\}}^{0,0}}{}}`$ $`=`$ $`(1)^{n_0+1}{\displaystyle \underset{i=0}{\overset{n_0}{}}}(2)^i\left({\displaystyle \genfrac{}{}{0pt}{}{n_0}{i}}\right){\displaystyle \frac{1}{(+1+i)^2}};`$ (35)
and
$`{\displaystyle \frac{d\mathrm{log}𝒜^{(𝒩)}}{dt}}`$ $`=`$ $`{\displaystyle \frac{dn_n}{dt}}\left[\psi (n_N+1)\psi (n_n+1)\right]+{\displaystyle \frac{dn_p}{dt}}\left[\psi (n_N+1)\psi (n_p+1)\right]`$ (36)
$`+`$ $`{\displaystyle \frac{dn_{}}{dt}}\left[\psi (n_\pi +1)\psi (n_{}+1)\right]+{\displaystyle \frac{dn_0}{dt}}\left[\psi (n_\pi +1)\psi (n_0+1)\right]`$
$`+`$ $`{\displaystyle \frac{dn_+}{dt}}\left[\psi (n_\pi +1)\psi (n_++1)\right];`$
where $`n_N=n_n+n_p`$ is the total number of nucleons and $`n_\pi =n_{}+n_0+n_+`$ is the total number of pions. Eq. (36) can be also written in the form:
$`{\displaystyle \frac{d\mathrm{log}𝒜^{(𝒩)}}{dt}}`$ $`=`$ $`{\displaystyle \frac{dn_n}{dt}}{\displaystyle \underset{k=1+n_n}{\overset{n_N}{}}}{\displaystyle \frac{1}{k}}+{\displaystyle \frac{dn_p}{dt}}{\displaystyle \underset{k=1+n_p}{\overset{n_N}{}}}{\displaystyle \frac{1}{k}}`$ (37)
$`+`$ $`{\displaystyle \frac{dn_{}}{dt}}{\displaystyle \underset{k=1+n_{}}{\overset{n_\pi }{}}}{\displaystyle \frac{1}{k}}+{\displaystyle \frac{dn_0}{dt}}{\displaystyle \underset{k=1+n_0}{\overset{n_\pi }{}}}{\displaystyle \frac{1}{k}}+{\displaystyle \frac{dn_+}{dt}}{\displaystyle \underset{k=1+n_+}{\overset{n_\pi }{}}}{\displaystyle \frac{1}{k}};`$
Eqs. (3337) offer the final result for the $`SU(2)`$ case . Notably, they imply a relation for the number of neutral pions in a system. We thus see, that when the case of the non-abelian symmetry is carefully considered and not ignored, one can get relations determining also the ‘neutral’ members of multiplets. In the standard approach the multiplicity of neutral particles is obtained by introducing a subsidiary chemical potential which is related to the lack of chemical equilibrium of a system or to the residual interaction with the environment .
## 5 Conclusions and Outlook
We have shown how constraints due to the preservation of the symmetry properties of a multiparticle macroscopic state define the path of evolution of the system. One should notice that results we presented are general and do not depend on the particular choice of the representation of the symmetry group. For different initial symmetry group representations one gets different paths, but they are all of a similar “shape”, considering the hypothesis that the global behaviour of a macrosystem (in a sense of statistical physics) should not be altered if a number of particles is changed by a very small (“microscopic”) amount.
We have explicitly presented the example how our constraint works in the simplest non - trivial case of $`SU(2)`$ symmetry group.
Although we have studied and implemented the discreet symmetry using quantum states, whenever we referred here to a dynamical equation we considered the limit of incoherent state evolution described by the Boltzmann equation. The dynamical evolution we consider thus is described in terms of diagonal density matrix. This is the appropriate approach given that our main objective is to arrive at a dynamical derivation of symmetry deformed statistical distribution.
We recall that quantum correlations (without symmetry) alone are responsible for the deformation of the Boltzmann distribution into Bose/Fermi distributions, and that the Boltzmann equation yields this result when we allow for Fermi blocking/Bose enhancement in the collision term. In that line of thought, the next step would be to show that it is possible to obtain now within a dynamical Boltzmann equation calculation the evolution of a many body system into symmetry-deformed statistical equilibrium distribution. We are also exploring the possibility that the methods here presented allow the formulation of a microscopic transport theory which would obey the long range correlations introduced by the macroscopic quantum and symmetry constraints.
## Acknowledgments
Work supported in part by a grant from the U.S. Department of Energy, DE-FG03-95ER40937 , and in part by the Polish Committee for Scientific Research under contract KBN - 2 P03B 030 18 . |
warning/0003/math0003163.html | ar5iv | text | # Untitled Document
1200
A Partition Theorem
Saharon Shelah
Institute of Math, The Hebrew University, Jerusalem Israel
Department of Math., Rutgers University, New Brunswick NJ USA
footnote Presented in the third Turan lecture of the author in Hungary, Feb 1998; Publ. Number 679. Done 1-2/98. Partially supported by the Partially supported by the United States Israel Binational Science Foundation.
Abstract
We deal with some relatives of the Hales Jewett theorem with primitive recursive bounds.
Anotated Content
0 Introduction
1 Basic Definitions
2 Proof of the Partition Theorem with a bound
3 : Higher Dimension Theorems
4 The Main Theorem
Key words: Ramsey theory, Hales Jewett theorem, finite combinatorics
Classification: Primary 05A99, Secondary 15A03
0: Introduction
We prove the following: there is a primitive recursive function $`f_{}^{}(,)`$ , in the three variables, such that: for every natural numbers $`t,n>0`$, and $`c`$, for any natural number $`kf_t^{}(n,c)`$ the following holds. Assume $`\mathrm{\Lambda }`$ is an alphabet with $`n>0`$ letters, $`M`$ is the family of non empty subsets of $`\{1,\mathrm{},k\}`$ with $`t`$ members and $`V`$ is the set of functions from $`M`$ to $`\mathrm{\Lambda }`$ and lastly $`d`$ is a $`c`$colouring of $`V`$ (i.e. a function with domain $`V`$ and range with at most $`c`$ members). Then there is a $`d`$monochromatic $`V`$line, which means that there are $`w\{1,\mathrm{},k\}`$, with at least $`t`$ members and function $`\rho `$ from $`\{uM:u`$ not a subset of $`w\}`$ to $`\mathrm{\Lambda }`$ such that letting $`L=\{\eta V:\eta `$ extend $`\rho `$ and for each $`s=1,\mathrm{},t`$ it is constant on $`\{uM:uw`$ has $`s`$ members $`\}\}`$, we have $`dL`$ is constant (for $`t=1`$ those are the Hales Jewett numbers).
A second theorem relates to the first just as the affine Ramsey theorem of Graham, Leob and Rothschild (which continue the n-parameter Ramsey theorem of Graham and Rothschild), relates to the Hales Jewett theorem. We also note an infinitary related theorem parallel to the Galvin Prikry theorem and the Carlson Sympson theorem.
Let us review history and background, not repeating \[GRS 80\]. In the late seventies, Furstenberg and Sarakozy independently prove that if $`𝐩(𝐱)`$ is a polynomial in $`Z[𝐱]`$ satisfying $`p(0)=0`$ and $`AN`$ is a set of positive density then for some $`a,bA`$ and $`nN`$ we have $`ab=𝐩(n)`$. Bergelson and Leibman \[BL96\] continuing Furstenberg \[Fu\] prove (this is a special of a density theorem like Szemeredi ): if $`r,k,t,m`$ are natural numbers, $`𝐩_{\mathrm{},s}(x)`$ for $`\mathrm{}=1,\mathrm{},k`$ and $`s=1,\mathrm{},t`$ are polynomials with rational coefficients, taking integer values at integers, and vectors $`\overline{v}_1,\mathrm{},\overline{v}_t{}_{}{}^{m}Z`$ and any $`r`$colouring of $`{}_{}{}^{m}Z`$ there are $`\overline{a}mZ`$ and $`nZ(n0)`$ such that the set $`S(\overline{a},n)=\{\overline{a}+\mathrm{\Sigma }_{j=1,t}𝐩_{i,j}(n)\overline{v}_j:i=1,\mathrm{},k\}`$ is monochromatic.
Bergelson and Leibman \[ BL 9x\] prove a theorem, set polynomial extension, which is, in different formulation, like the first theorem describe above but without a bound (i.e. the primitive recursiveness). Their method is infinitary so does not seem to give even the weak bound in 2.5 (one with triple induction), and certainly does not give primitive recursive bounds.
Naturally our proofs continue \[Sh 329\]. We thank the referee for telling us on \[BL 9x\] and other helpful comments.See a discussion of related problems in \[Sh 702\].
0.1 Notation:
(a) We use $`\mathrm{\Lambda }`$ for a finite alphabet, always non empty, members of which are denoted by $`\alpha ,\beta ,\gamma .`$
(b) We use $`M,N`$ to denote structures which serve as index sets, so we call them index models. We use $`\tau `$ to denote vocabularies, , (see Definition 1.1), $`F`$ to denote function symbols.
(c) We use $`n,m,k,\mathrm{},i,j,c,r,s,t`$ to denote natural numbers, but usually $`n`$ is the number of letters, i.e. the number of members in an alphabet; $`k`$ the dimension of the index models and $`c1`$ the number of colours.
(d) $`|X|`$ and also card$`(X)`$ denote the number of elements of the set $`X.`$
(e) We use $`\eta ,\nu ,\rho `$ to denote members of spaces, we use $`V,U`$ to denote spaces and $`a,b`$ to denote elements of $`M,N`$ and $`d`$ to denote colourings, $`p`$ to denote the ‘type‘ of a point in a line and p to denote type of a line or a space (see Definition 1.7(3)). We use $`L`$ to denote (combinatorial) lines, $`S`$ to denote (combinatorial) subspaces.
(f) $`A`$ bar on a symbol, say $`\overline{x}`$ denote a finite sequence of such objects, of length lg$`(\overline{x})`$ the ???$`i`$th object being $`x_i`$ (and of $`\overline{x}_m`$ or $`\overline{x}^m`$ it is $`x_i^m).`$
0.2 DEFINITION: (1) For $`m1`$, let $`E`$<sub>m</sub> be the minimal class of functions from natural numbers to natural numbers (with any number of places) closed under composition, which for $`m=1`$ contains $`0,1,x+1`$ and the projection functions, and for $`m>1`$ contains any function which we get by inductive definition on functions from $`E`$<sub>m-1</sub> (see \[Ro84\], so $`E`$<sub>3</sub> is the family of polynomials, $`E`$<sub>4</sub> contains the tower function and $`E`$<sub>5</sub> contains the waw function and $`_{m1}E_m`$ is the family of primitive recursive functions, and the ‘simplest‘ function not there is the Akerman function.) We allow an object like $`\overline{\mathrm{\Lambda }}`$ to be one of the arguments meaning a natural number coding of it (in the cases used this does not matter). Abusing notation, we may say $`f`$ is in $`E`$<sub>n</sub> instead of $`f`$ is bounded by a function from $`E`$<sub>n</sub>, also writing $`f_{\overline{\mathrm{\Lambda }}}(,\mathrm{})`$ we count $`\overline{\mathrm{\Lambda }}`$ as one of the arguments.
(2) We can define the Akerman function $`A_n(m)`$ by double induction (in as sense it is the simplest, smallest function which is not primitive recursive).
0.3 DEFINITION: (1) Let RAM$`(t,\mathrm{},c)`$ be the Ramsey number, i.e. the first $`k`$ such that $`k(t)_c^{\mathrm{}}`$ which mean that if $`A`$ is a set with $`k`$ elements, and $`d`$ is a $`c`$colouring of $`[A]^{\mathrm{}}=^{df}\{B:B`$ is a subset of $`A`$ with $`\mathrm{}`$ elements$`\}`$, that is a function with this domain and range of cardinality $`c`$, then for some $`A_1[A]^t`$ we have $`d[A_1]^{\mathrm{}}`$ is constant.
(2) Let HJ$`(n,m,c)`$ be the Hales Jewett number for getting a monochromatic subspace of dimension $`m`$, when the colouring has $`c`$ colours and for an alphabet with $`n`$ members (this is, by our subsequent definitions, $`f^1(\overline{\mathrm{\Lambda }},m,c)`$ when $`\tau (\overline{\mathrm{\Lambda }})=\{`$id$`\}`$, and $`\mathrm{\Lambda }_{\mathrm{id}}`$ has $`n`$ members, see Definition 1.9).
Section 1 : Basic definitions
We can look at Hales Jewett theorem in geometric terms: $`R`$ is replaced by $`\mathrm{\Lambda }`$; a finite alphabet, the $`k`$dimensional euclidean space $`R^k`$ is replaced by $`{}_{}{}^{[1,k]}\mathrm{\Lambda }`$ (or $`{}_{}{}^{[0,k)}\mathrm{\Lambda }`$), essentially the set of sequences of length $`k`$ of members of the alphabet $`\mathrm{\Lambda }`$; a subspace is replaced by the set of solutions $`(x_1,\mathrm{},x_k){}_{}{}^{[1,k]}\mathrm{\Lambda }`$ of a family of linear equations, which here means just $`x_i=\alpha `$ (where $`\alpha \mathrm{\Lambda },1ik)`$ or $`x_i=x_j`$. Here the basic set $`[1,k]`$ is replaced by a structure $`M`$, a $`\tau `$fim. Such basic definitions are given in this section.
We define a ‘space over an index model of dimension $`k`$, over an alphabet $`\mathrm{\Lambda }`$ of size $`n\mathrm{`}`$, lines and more. We then define the function $`f^1`$, such that for every $`n`$, if $`k`$ is $`f_\tau ^1(n,c)`$ then for every colouring of the space by $`c`$ colours, there is a monochromatic line (in the appropriate interpretation.) Of course the use of id as a special function symbol is not really needed, also we can waive the linear order on $`P^M`$, and the set of automorphisms of the resulting structure are natural for our purpose, but not for the structures from 1.10(3); but at present those decisions does not matter.
1.1 DEFINITION: (1) We call $`M`$ a full index model \[fim or $`\tau `$fim or fim for $`\tau ]`$ if:
(a) the vocabulary $`\tau =\tau _M=\tau (M)`$ of $`M`$ includes a unary predicates $`P`$, a binary predicate $`<`$, and finitely many function symbols $`F`$, $`F`$ being arity$`(F)`$place and no other symbols (so $`F`$ vary over such function symbols). We may write arity$`{}_{}{}^{\tau }(F)`$ for arity$`(F).`$ We usually treat $`\tau `$ as the set of function symbols in $`\tau .`$
(b) the universe of $`M`$ is finite (non empty of course).
$`(c)<^M`$ is a linear order of $`P^M`$, so $`x<^My`$ implies $`x,yP^M.`$
$`(d)F^M`$ is a partial function such that if $`F^M(a_1,\mathrm{},a_r)`$ is well defined (so $`r=`$ arity$`(F))`$ then $`a_1,\mathrm{},a_rP^M`$ and the function is symmetric, i.e. does not depend on the order of the arguments, so if not said otherwise we assume $`a_1^Ma_2^M\mathrm{}^Ma_r.`$
(e) if $`F_1^M(a_1,\mathrm{},a_r)=F_2^M(b_1,\mathrm{},b_t)`$ then $`F_1=F_2`$ (hence $`r=t)`$ and $`a_{\mathrm{}}=b_{\mathrm{}}`$ for $`\mathrm{}=1,\mathrm{},t`$ (under the convention from clause $`(d))`$ and every $`bMP^M`$ has this form. So we let base$`{}_{M}{}^{}(b)=^{df}\{a_1,\mathrm{},a_r\}`$ and let base$`{}_{\mathrm{}}{}^{}(b)=`$ base$`{}_{M,\mathrm{}}{}^{}=_{}^{df}a_{\mathrm{}}`$ where $`b=F^M(a_1,\mathrm{},a_r)`$ (and $`a_1^Ma_2^M\mathrm{}^Ma_r`$ of course) and $`F_{M,b}=^{df}F;`$ those are well defined by the demand above.
$`(f)`$ $`P^M`$ is non empty and we call its cardinality dim$`(M)`$, the dimension of $`M.`$
(g) id<sup>M</sup> is the identity function on $`P^M`$, so id is a unary function symbol of $`\tau `$.
(h) each $`F^M(a_1,\mathrm{},a_{\mathrm{arity}(F)})`$ is well defined iff $`a_1,\mathrm{},a_{\mathrm{arity}(F)}`$ are from $`P^M`$ ( and the value does not depend on the order, of course)
(2) For $`\tau `$ as in part (1), let arity$`(\tau )`$ be Max$`\{`$arity$`(F):F\tau \}`$, so it is at least 1 and let $`\overline{m}[\tau ]=^{df}m_t^\tau :t=1,\mathrm{}`$,arity$`(\tau )`$ where $`m_t^\tau `$ is the number of $`F\tau `$ with arity $`t;`$ and we call $`\overline{m}^\tau `$ the signature of $`\tau `$, of course when saying “the signature of $`M`$” we mean “of $`\tau (M)`$”.
(3) For $`M`$ a fim we call $`BM`$ closed in $`M`$ (or $`M`$closed ) if for $`b=F^M(a_1,\mathrm{},a_s)`$ we have $`bB`$ iff $`a_1,\mathrm{},a_sM.`$ Let the closure of $`A`$ in $`B`$ or cl$`{}_{M}{}^{}(A)`$ for $`AM`$, be the minimal $`M`$closed set $`B`$ $`M`$ which include $`A`$. A close (non empty) subset of $`M`$ is actually a submodel. We do not strictly distinguish between a closed subset $`B`$ of $`M`$ and the model $`MB`$ (which are fims with the same vocabulary).
(4) For $`\tau `$index models $`M,N`$ let PHom$`(M,N)`$ be the set of functions $`f`$ from $`P^M`$ into $`P^N`$ such that $`x^Myf(x)^Nf(y).`$ Let Hom$`(M,N)`$ be the set of functions $`f`$ from $`M`$ into $`N`$ such that $`fP^M`$ PHom$`(M,N)`$ and $`b=F^M(a_1,\mathrm{},a_t)`$ implies $`f(b)=F^N(f(a_1),\mathrm{},f(a_t)).`$ Let PHm$`(M,N)`$ be the set of functions $`f`$ from $`P^M`$ into $`P^N`$, and let Hm$`(M,N)`$ be the set of functions $`f`$ from $`M`$ into $`N`$ such that $`fP^MPHm(M,N)`$ and $`b=F^M(a_1,\mathrm{},a_t)`$ implies $`f(b)=F^N(f(a_1),\mathrm{},f(a_t))`$.
(5) Let Sort$`{}_{}{}^{M}(F)`$ be the range of $`F^M.`$
1.2 Fact: (1) For any $`f`$ PHom$`(M,N)`$ there is a unique $`\widehat{f}`$ Hom$`(M,N)`$ which extend $`f.`$
(2) For any $`f`$ PHm$`(M,N)`$ there is a unique $`\widehat{f}`$ Hm$`(M,N)`$ which extend $`f.`$
1.3 CLAIM/DEFINITION: (1) For any fim $`M`$ there is a polynomial p$`(x)`$, with rational coefficients but positive integers as values for $`x`$ a positive integer (really sum of binomial coefficients binom$`(x,m_1,\mathrm{},m_n)`$ for $`m_i=1,\mathrm{}`$, arity$`(\tau ))`$ such that for $`uP^M`$, the set cl$`{}_{M}{}^{}(u)`$ has exactly p$`(|u|)`$ members. Now p$`(x)`$ depend on the signature of $`\tau `$ only and so we shall denote it by p$`{}_{\tau }{}^{}(x)`$ or p$`{}_{M}{}^{}(x).`$ Note that p$`{}_{\tau }{}^{}(0)=0.`$
1.4 DEFINITION: (1) We say that $`\tau `$ is canonical vocabulary for $`t`$ (or $`t`$canonical) and write $`\tau =\tau _t`$ if $`\tau =\{F_1,\mathrm{},F_t,P,<\}`$ where arity$`(F_s)`$ is $`s.`$
(2) We say that $`M`$ is a $`(J,t)`$canonical fim if:
$`(a)`$ $`J`$ is a finite linear order
$`(b)`$ $`M`$ is a fim with the $`t`$canonical vocabulary
$`(c)`$ $`(P^M,<^M)`$ is $`J`$
$`(d)`$ $`F_1^M`$ is the identity on $`P^M`$
$`(e)`$ for $`r=2,\mathrm{},t`$ the function $`F_r^M`$ is $`F_r^M(a_1,\mathrm{},a_r)=^{df}\{a_1,\mathrm{},a_r\}.`$
1.5 DEFINITION: (1) Let $`M`$ be a fim with vocabulary $`\tau =\tau _M`$ and let $`\{A_1,A_2\}`$ be a partition of $`P^M`$ to convex sets such that $`A_1<^MA_2`$ which means that $`(a_1A_1)(a_2A_2)[a_1<^Ma_2]`$. We define a vocabulary $`\tau _{M,A_1,A_2}.`$ It contains, in addition to the symbols $`P`$, $`<`$, for each function symbol $`F`$ of $`\tau `$ and a $`^M`$increasing sequence $`\overline{a}_1`$ from $`A_1`$ and a $`^M`$increasing sequence $`\overline{a}_2`$ from $`A_2`$ such that lg$`(\overline{a}_1)+`$lg$`(\overline{a}_2)<`$ arity$`{}_{}{}^{\tau }(F)`$ a function symbol called $`F_{\overline{a}_1,\overline{a}_2}`$ with arity arity$`{}_{}{}^{\tau }(F)`$ lg$`(\overline{a}_1)`$ lg$`(\overline{a}_2).`$
We identify $`F\tau `$ with $`F_,`$ and so consider $`\tau _{M,\overline{a}_1,\overline{a}_2}`$ an extension of $`\tau .`$
(2) Let $`\overline{m}=\overline{m}[\tau ,k_0,k_1]`$ be $`\overline{m}[\tau _{M,A_0,A_1}]`$, the signature of $`\tau _{M,A_0,A_2}`$ whenever $`M`$ is a $`\tau `$fim of dimension $`k_0+k_1`$ and $`A_0`$ is the set of $`k_0`$ first members of $`P^M`$ and $`A_1`$ is the set of $`k_1`$ last members of $`P^M.`$
(3) Let $`M_k^\tau `$ be a fim of vocabulary $`\tau `$ and dimension $`k`$, say $`P^M=\{1,\mathrm{},k\}.`$ Let $`\tau ^{[k,\mathrm{}]}`$ be $`\tau _{M_{k+\mathrm{}}^\tau ,A_0,A_1}`$ where $`A_0`$ is the set of the first $`k`$ members of $`P^M`$ and $`A_1`$ is the set of the last $`\mathrm{}`$ members of $`P^M.`$
1.6 DEFINITION: (1) Let $`\overline{\mathrm{\Lambda }}`$ denote a sequence $`\mathrm{\Lambda }_F:F\tau `$ where $`\mathrm{\Lambda }_F`$ is a finite alphabeth, and we let $`\tau [\overline{\mathrm{\Lambda }}]=\tau `$, as $`\overline{\mathrm{\Lambda }}`$ determine $`\tau .`$ We call $`\overline{\mathrm{\Lambda }}`$ an alphabet sequence (for $`\tau )`$ or a $`\tau `$alphabet sequence. We may write $`(\tau ,\mathrm{\Lambda })`$ instead $`\overline{\mathrm{\Lambda }}`$ if $`\tau =\tau [\overline{\mathrm{\Lambda }}]`$ and $`\mathrm{\Lambda }_F=\mathrm{\Lambda }`$ for every $`F\tau .`$
(2) We say $`p`$ is a $`\overline{\mathrm{\Lambda }}`$-type if $`p`$ is a function with domain $`\tau `$ such that $`p(F)\mathrm{\Lambda }_F`$; let p, q denote non empty sets of $`\overline{\mathrm{\Lambda }}`$-types; we identify them with their characteristic functions that they define, so we assume that from p we can reconstruct $`\overline{\mathrm{\Lambda }}`$ hence $`\tau [\overline{\mathrm{\Lambda }}].`$ Let p$`_{\overline{\mathrm{\Lambda }}}`$ be the set of $`\overline{\mathrm{\Lambda }}`$-types. We may write $`\mathrm{\Lambda }`$ instead $`\overline{\mathrm{\Lambda }}`$ is $`\mathrm{\Lambda }_F=\mathrm{\Lambda }`$ for every $`F\tau `$ and then let p<sub>τ,Λ</sub> be the set of constant $`(\tau ,\mathrm{\Lambda })`$types.
1.7 DEFINITION: (1) For $`\overline{\mathrm{\Lambda }}`$ a $`\tau `$alphabet sequence, let $`V=`$ Space$`{}_{\overline{\mathrm{\Lambda }}}{}^{}(M)`$ be defined as follows:
its set of elements is the set of functions $`\eta `$ with domain $`M`$, such that $`b`$ Sort$`{}_{}{}^{M}(F)\eta (b)\mathrm{\Lambda }_F;`$ we assume that from $`V`$ we can reconstruct $`M`$ and $`\mathrm{\Lambda }.`$
(2) We say $`d`$ is a $`C`$colouring of $`V`$, if $`d`$ is a function form $`V`$ into $`C`$, we say $`c`$colouring if $`C`$ has $`c`$ members and the default value of $`C`$ is \[0,$`c)=\{0,1,\mathrm{},c1\}.`$
(3) We say $`L`$ is a $`V`$line or a line of $`V`$ if for q $`=𝔭_{\overline{\mathrm{\Lambda }}}`$ we have : $`L`$ is a $`(V,𝔮)`$line or a q -line of $`V;`$ this is defined for q a (non empty) subset of p$`_{\overline{\mathrm{\Lambda }}}`$ and it means:
$`L`$ is a subset of $`V`$ such that for some subset supp$`(L)=`$ supp$`{}_{M}{}^{}(L)`$ of $`M`$ we have:
(a) supp$`(L)P^M`$ is non empty and we call it supp$`{}_{}{}^{P}(L)`$
(b) supp$`(L)`$ is the $`M`$th closure of supp$`{}_{}{}^{P}(L)`$
(c) for any $`\eta ,\nu L`$ we have $`\eta (M`$ supp$`{}_{M}{}^{}(L))=\nu (M`$ supp$`{}_{M}{}^{}(L))`$
(d) for any $`\eta L`$ for some $`p𝔮`$ we have : if $`b`$ supp$`{}_{M}{}^{}((L)`$ then $`\eta (b)=p(F_{M,b})`$
(e) For any $`p𝔮`$ there is $`\eta L`$ as in clause $`(d).`$
(5) For $`L`$ as above and $`p𝔭`$ let pt$`{}_{L}{}^{}(p)`$ be the unique $`\nu L`$ such that for every $`a`$ $``$ supp$`{}_{M}{}^{}(L)`$ we have $`\nu (a)=p(F_{M,b}).`$ For q$`{}_{}{}^{}𝔮`$, the q$`{}_{}{}^{}`$ subline of a q-line $`L`$ is $`\{`$pt$`{}_{L}{}^{}(p):p𝔮^{}\}.`$
(6) For a colouring $`d`$ of $`V`$, we say a $`V`$line (or $`(V,𝔮)`$line) $`L`$ is $`d`$monochromatic if $`d`$ is constant on $`L.`$
(7) When we are given $`M,\tau ,\overline{\mathrm{\Lambda }},V`$ as in part (4) and in addition we are given $`m`$, we define when $`S`$ is an $`m`$dimensional $`V`$subspace , or $`m`$dimensional subspace for $`V.`$ It means that for some sequence $`M_{\mathrm{}}:\mathrm{}<m`$ we have
(a) each $`M_{\mathrm{}}`$ is a submodel of $`M`$,
(b) if $`\mathrm{}_1<\mathrm{}_2<m`$ then $`M_\mathrm{}_1,M_\mathrm{}_2`$ are disjoint,
(c) for some $`\rho `$, a function with domain ($`M`$ cl$`(\{M_{\mathrm{}}:\mathrm{}<m)\}))`$ such that $`\rho (b)\mathrm{\Lambda }_{F_{M,b}}`$ for every $`b`$ Dom$`(\rho )`$, and some $`m`$dimensional $`\tau `$fim $`K`$ say $`K=M_{[0,m)}^\tau `$ and letting $`N`$ be the submodel of $`M`$ with universe cl$`{}_{M}{}^{}(_{\mathrm{}<m}M_{\mathrm{}})`$ there is $`f`$ Hm$`(N,K)`$ which is onto $`K`$ such that $`fP^M_{\mathrm{}}`$ is constant for each $`\mathrm{}`$ and: $`\nu S`$ iff $`\nu `$ extend $`\rho `$ and for some $`\varrho `$ Space$`{}_{\overline{\mathrm{\Lambda }}}{}^{}(K)`$ we have $`bN\nu (b)=\varrho (\widehat{f}(b))`$
(8) We call $`S`$ convex if
(a) for $`\mathrm{}_1<\mathrm{}_2<m`$ and $`a_1M_\mathrm{}_1`$ and $`a_2M_\mathrm{}_2`$ we have $`a_1<^Ma_2`$ and
(b) $`f`$ Hom$`(N,K)`$.
(9) For $`S`$ as above and (see Definition 1.5(3)) $`\varrho `$ Space$`{}_{\mathrm{\Lambda }}{}^{}(M_{[0,m)}^\tau )`$ we define pt$`{}_{S}{}^{}(\varrho )`$ as the unique $`\nu S`$ as above in part (8).
We may define now a natural function, which is our main concern here :
1.8 DEFINITION: (1) Let $`f^1(𝔭,c)`$ where p $``$ p$`_{\overline{\mathrm{\Lambda }}}`$ ( and $`\overline{\mathrm{\Lambda }}`$ an alphabet sequence) be the minimal $`k`$ such that for any $`\tau _{\overline{\mathrm{\Lambda }}}`$fim $`M`$ of dimension $`k`$, we have:
for any $`c`$colouring $`d`$ of $`V=`$ Space$`{}_{\overline{\mathrm{\Lambda }}}{}^{}(M)`$ there is a p-line $`L`$ of $`V`$ which is $`d`$monochromatic, i.e. such that $`p,q𝔭`$ implies that pt$`{}_{L}{}^{}(p)`$ , pt$`{}_{L}{}^{}(q)`$ have the same colour (by $`d).`$
If $`k`$ does not exist we may say it is $`\omega `$ or is $`\mathrm{}.`$ We may write $`f_\tau ^1(𝔭,c)`$ or $`f^1(𝔭,c;\tau )`$ to stress the role of $`\tau `$.
(2) If p $`=𝔭_{\overline{\mathrm{\Lambda }}}`$ we may write $`f^1(\overline{\mathrm{\Lambda }},c).`$ If $`\mathrm{\Lambda }_F=\mathrm{\Lambda }`$ for every $`F\tau =\tau [\overline{\mathrm{\Lambda }}]`$ then we may write $`f_\tau ^1(\mathrm{\Lambda },c);`$ in this case we can replace $`\mathrm{\Lambda }`$ by $`|\mathrm{\Lambda }|.`$ Clearly only $`\overline{m}^\tau `$ is important so we may write only it. Also we may write $`f_\tau ^1(\overline{n},c)`$ for $`f_\tau ^1(\overline{\mathrm{\Lambda }},c)`$ whenever $`\overline{n}=n_F:F\tau `$ and $`|\mathrm{\Lambda }_F|=n_F.`$
We can of course use the multidimensional versions of those definitions
1.9 DEFINITION: Let $`f^1(\overline{\mathrm{\Lambda }},m,c)`$ be the minimal $`k`$ such that for any $`\tau `$fim $`M`$ of dimension $`k`$ we have: for any $`c`$colouring $`d`$ of Space$`{}_{\overline{\mathrm{\Lambda }}}{}^{}(M)`$ there is a convex subspace $`S`$ of $`V`$ of dimension $`m`$ which is $`d`$monochromatic, i.e. such that all the points in $`S`$ have the same colour (by $`d)`$, if $`k`$ does not exist we say it is $`\omega `$ or is $`\mathrm{}.`$ We may write $`f_\tau ^1(\overline{\mathrm{\Lambda }},m,c)`$ or $`f_\tau ^1(\mathrm{\Lambda },m,c)`$ etc. as before. Clearly only $`\overline{m}^\tau `$ is important (rather than $`\tau `$), so we may write only it. We may replace $`\mathrm{\Lambda }`$ by $`|\mathrm{\Lambda }|.`$ We may replace $`\overline{\mathrm{\Lambda }}`$ by $`n_F:F\tau `$ when $`n_F=|\mathrm{\Lambda }_F|.`$
At present, it does not really matter if we omit the demand convex above.
The function has some obvious monotonicity properties, we mention those we shall actually use.
1.10 Claim: (1) For $`\mathrm{}=1,2`$ assume $`\overline{\mathrm{\Lambda }}`$ is an alphabet sequence for the vocabulary $`\tau ^{\mathrm{}}`$ and arity$`(\tau ^1)`$ arity$`(\tau ^2)`$ and for each $`m=1,\mathrm{}`$, arity$`(\tau ^1)`$ we have
$`\mathrm{\Pi }\{|\mathrm{\Lambda }_F^1|:F\tau ^1`$ has arity $`m\}`$ $`\mathrm{\Pi }\{|\mathrm{\Lambda }_F^2|:F\tau ^2`$ has arity $`m\}.`$
Then $`f^1(\overline{\mathrm{\Lambda }}^1,c)f^1(\overline{\mathrm{\Lambda }}^2,c).`$
(2) For $`\mathrm{}=1,2`$ assume $`\overline{\mathrm{\Lambda }}`$ is an alphabet sequence for the vocabulary $`\tau ^{\mathrm{}}`$ and $`\tau ^1\tau ^2`$ and $`\overline{\mathrm{\Lambda }}`$$`{}_{}{}^{1}=\overline{\mathrm{\Lambda }}^2\tau ^1`$ and $`F\tau ^2\tau ^1|\mathrm{\Lambda }_F^2|=1.`$
Then $`f^1(\overline{\mathrm{\Lambda }}^1,c)=f^1(\overline{\mathrm{\Lambda }}^2,c).`$
Proof: Straightforward.
$`\text{}_{1.10}`$
1.11 DEFINITION: We define, for $`\mathrm{}=1,2,3`$ what is a fim, we just replace in Def.$`1.1`$ clauses $`(d),(e)`$ by
$`(d)_{\mathrm{}}`$ $`F^M`$ is a partial function such that if $`F^M(a_1,\mathrm{},a_r)`$ is well defined (so $`r=`$ arity$`(F))`$ then $`a_1,\mathrm{},a_mP^M`$ and $`\mathrm{}=1`$ implies the function is symmetric, i.e. does not depend on the order of the variables, so if not said otherwise we assume $`a_1^Ma_2^M\mathrm{}^Ma_r.`$
$`(e)_{\mathrm{}}`$ if $`F_1^M(a_1,\mathrm{},a_r)=F_2^M(b_1,\mathrm{},b_t)`$ and $`\mathrm{}\{1,2\}`$ then $`F_1=F_2`$ (hence $`r=t)`$ and $`\mathrm{}=2_{s=1,\mathrm{},r}a_s=b_s`$
and
$`\mathrm{}=1_{s=1,\mathrm{},r1}a_s^Ma_{s+1}`$ $`_{s=1,\mathrm{},r1}b_s^Mb_{s+1}`$ $`_{s=1,\mathrm{},r}a_s=b_s.`$ So we let base$`{}_{M}{}^{}(b)=^{df}\{a_1,\mathrm{},a_r\}`$ and when $`\mathrm{}=1,2`$ let base$`{}_{s}{}^{}(b)=`$ base$`{}_{M,s}{}^{}(b)=^{df}a_s`$ where $`b=F^M(a_1,\mathrm{},a_r)`$ (and if $`\mathrm{}=1`$ then $`a_1^Ma_2^M\mathrm{}^Ma_r`$ , of course) and $`F_{M,b}=^{df}F;`$ those are well defined by the demand above.
$`(e)_{\mathrm{}}^{^{}}`$ if $`\mathrm{}\{1,2,3\}`$ and $`bMP^M`$ then for some $`F\tau `$ and $`a_1,\mathrm{},a_{arity[F]}P^M`$ we have $`b=F^M(a_1,\mathrm{},a_{arity[F]})`$. So $`\mathrm{}=1`$ is the old notion and for $`\mathrm{}=3`$ we require very little. We define $`f_\lambda ^{\mathrm{}}(\overline{\mathrm{\Lambda }},c)`$ as in Definition 1.9 for fim (so again $`\mathrm{}=1`$ is our standard case.)
1.12 Claim: Let $`\tau `$ be a vocabulary and $`\tau _{}=\{G_{F,\pi }:`$ $`F\tau `$ and $`\pi `$ is a permutation of $`\{1,\mathrm{}`$,arity$`(F)\}\}`$
with arity$`(G_{F,\pi })=`$ arity$`(F).`$
Then
$`(\alpha )`$ If $`\overline{\mathrm{\Lambda }}`$ is a $`\tau `$alphabet sequence and $`\overline{\mathrm{\Lambda }}`$$`{}_{}{}^{}=\mathrm{\Lambda }_G^{}:G\tau _{}`$ where $`\mathrm{\Lambda }_{G_{F,\pi }}^{}=\mathrm{\Lambda }_F`$ then $`f^2`$$`{}_{\tau }{}^{}(\overline{\mathrm{\Lambda }},c)f_\tau _{}^1(\overline{\mathrm{\Lambda }}^{},c)`$
$`(\beta )`$ For $`\overline{\mathrm{\Lambda }}`$ a $`\tau `$alphabet sequence we have: $`f_\tau ^3(\overline{\mathrm{\Lambda }},c)`$ is at most RAM$`(f_\tau ^2(\overline{\mathrm{\Lambda }}^{},c)`$, arity$`(\tau ),c^{})`$ where e.g. $`c^{}`$ depend on $`\tau `$ only (and RAM stand for Ramsey number).
$`(\gamma )f_\tau ^1(\overline{\mathrm{\Lambda }},c)f_\tau ^2(\overline{\mathrm{\Lambda }},c)`$
$`(\delta )f_\tau ^2(\overline{\mathrm{\Lambda }},c)f_\tau ^3(\overline{\mathrm{\Lambda }},c)`$
Proof: Straightforward.
$`\text{}_{1.12}`$
Section 2 : Proof of the partition Theorem with a bound
Except Def 2.1,2.2 this section is for the reader convenience only, as it give a proof of a weaker version of the first theorem (with a bound which we get by triple induction). Later in 4.1-4.10 we give a complete proof with the primitive recursive bound, formally not depending on the proofs here. The strategy is to make the $`bM`$ with $`|\mathrm{base}_M(b)|`$ maximal immaterial. We first define some help functions.
2.1 DEFINITION: (1) We call a vocabulary $`\tau `$ monic if there is a unique function symbol of maximal arity, we then denote it by $`F_\tau ^{\mathrm{max}}.`$
(2) For a $`P^M`$ let $`M_a`$ be cl$`{}_{M}{}^{}(\{P^M\{a\})`$
(3) For $`V=`$ Space$`{}_{\overline{\mathrm{\Lambda }}}{}^{}(M)`$ and $`N`$ a closed subset of $`M`$ and $`H\tau `$, we say that a colouring $`d`$ of $`V`$ is $`(N,\alpha ,H)`$invariant if : $`\alpha \mathrm{\Lambda }_H`$, and the following holds, for any $`aP^N:`$
(\*) if $`\nu ,\eta V`$ and $`\nu M_a=\eta M_a`$ and $`\left[bM\mathrm{base}(b)=\{a\}F_{M,b}=H\nu (b)=\alpha =\eta (b)\right]`$ then $`d(\nu )=d(\eta ).`$
(4) In part (3) we write $`(\mathrm{},\alpha ,H)`$monochromatic if above $`N`$ is such that $`P^N`$ is the set of the last $`\mathrm{}`$ members of $`P^M.`$ We write $`(M,\alpha ,H)`$monochromatic if in part (3) we have $`M=N`$.
(5) In parts (3) and (4) we may omit $`H`$ when $`\tau `$ is monic and $`H=F_\tau ^{\mathrm{max}}.`$ Replacing $`\alpha `$ by $`\mathrm{\Lambda }^{}`$ mean that $`\mathrm{\Lambda }^{}`$ is a subset of $`\mathrm{\Lambda }_H`$ and the demand holds for every $`\alpha \mathrm{\Lambda }^{}.`$
2.2 DEFINITION: Let $`f^0`$ be defined as follows. First, $`f_{\overline{\mathrm{\Lambda }}}^0(n,\mathrm{},c)=`$ $`f_{\tau ,\overline{\mathrm{\Lambda }}}^0(n,\mathrm{},c)`$ is defined iff $`\tau `$ is monic with $`H=F_\tau ^{\mathrm{max}}`$ and $`\overline{\mathrm{\Lambda }}`$ an alphabet sequence for $`\tau `$ and $`n|\mathrm{\Lambda }_H|`$ and $`n<|\mathrm{\Lambda }_H|(n=|\mathrm{\Lambda }_H|\mathrm{}=0)`$. Second, $`f_{\overline{\mathrm{\Lambda }}}^0(n,\mathrm{},c)`$ is the first $`k`$ (natural number, if not defined we can understand it as $`\mathrm{}`$ or $`\omega `$ or ‘does not exist‘ ) such that (\*)<sub>k</sub> below holds,
where:
(\*)<sub>k</sub> If clauses (a)-(f) below hold then there is a $`d`$monochromatic line of $`V`$, where :
(a) $`M`$ is a fim of vocabulary $`\tau `$
(b) the dimension of $`M`$ is $`k`$
(c) $`V=`$ Space$`{}_{\overline{\mathrm{\Lambda }}}{}^{}(M)`$
(d) $`\mathrm{\Lambda }^{}`$ is a subset of $`\mathrm{\Lambda }_H`$ with exactly $`n`$ members
(e) $`d`$ is an $`(M,\mathrm{\Lambda }^{},H)`$invariant colouring of $`V`$
(f) if $`\mathrm{}0`$, then there is an $`\alpha `$ such that $`\alpha \mathrm{\Lambda }_H\mathrm{\Lambda }^{}`$ and $`d`$ is $`(\mathrm{},\alpha ,H)`$invariant.
Immediate connections are:
2.3 Observation: (1) The function $`f_{\tau ,\overline{\mathrm{\Lambda }}}^0(n,\mathrm{},c)`$ increases with $`c`$ and decreases with $`\mathrm{}`$ and $`n.`$
(2) The function $`f_{\tau ,\overline{\mathrm{\Lambda }}}^0(n,\mathrm{},c)`$ depends just on $`n,\mathrm{},c`$ and the set $`\{`$(arity$`(F),|\mathrm{\Lambda }_F|):F\tau \}`$ (possibly with multiple membership), so we may replace $`\tau `$ by its $`\overline{m}^\tau `$ (similarly for other such functions).
(3) In definitions 1.8,1.9,2.2 the demand holds for any larger $`k.`$
(4) $`f_{\overline{\mathrm{\Lambda }}}^0(0,0,c)=f^1(\overline{\mathrm{\Lambda }},c).`$
(5) If $`\tau `$ is monic and $`H=F_\tau ^{\mathrm{max}}`$ and $`\tau ^{}=\tau \{H\}`$ then $`f_{\overline{\mathrm{\Lambda }}}^0(|\mathrm{\Lambda }_H|,0,c)=f^1(\overline{\mathrm{\Lambda }}\tau ^{},c).`$
(6) If $`\mathrm{}^{}=f_{\overline{\mathrm{\Lambda }}}^0(n+1,0,c)`$ then $`f_{\overline{\mathrm{\Lambda }}}^0(n,\mathrm{}^{},c)=\mathrm{}^{}.`$
Proof: Trivial.
2.4 MAIN Claim: Assume
(a) $`\overline{\mathrm{\Lambda }}`$ is an alphabet sequence for a vocabulary $`\tau =\tau [\overline{\mathrm{\Lambda }}]`$, and $`n<|\mathrm{\Lambda }_H|`$
(b) $`\tau `$ is a monic vocabulary with $`H=F_\tau ^{\mathrm{max}}`$
(c) $`k_0f_{\overline{\mathrm{\Lambda }}}^0(n,\mathrm{}+1,c)`$ and $`k_0>\mathrm{}`$
(d) $`K`$ is a $`\tau `$fim of dimension $`k_01`$ with $`A_2`$ the last $`\mathrm{}`$ elements and $`A_1`$ the first $`(k_0\mathrm{}1)`$elements (this $`K`$ serve just for notation)
(e) $`\tau ^{}`$ is the vocabulary $`(\tau _{K,A_1,A_2})\{H\}`$ see Definition 1.5(3); so
(i) arity$`(\tau ^{})<`$ arity$`(\tau ))`$,
(ii) proj is the following function from $`\tau ^{}`$ to $`\tau :`$ it map $`F_{K,\overline{a}_1,\overline{a}_2}`$ to $`F`$ so proj$`\tau `$ is the identity, and
(iii) $`\overline{\mathrm{\Lambda }}^{}=^{df}\mathrm{\Lambda }_F^{}:F\tau ^{}`$ where $`\mathrm{\Lambda }_F^{}=\mathrm{\Lambda }_{\mathrm{proj}(F)}`$.
(f) $`c^{}=^{df}c^{\mathrm{card}(\mathrm{Space}_{\overline{\mathrm{\Lambda }}}(K))}.`$
Then
$`f_{\overline{\mathrm{\Lambda }}}^0(n,\mathrm{},c)k_0+`$ $`f^1(\overline{\mathrm{\Lambda }}^{},c^{})1`$
Proof: Let $`k_1=f^1(\overline{\mathrm{\Lambda }}^{},c^{})`$ and let $`k=k_0+k_11`$, so it suffice to prove that $`kf_{\overline{\mathrm{\Lambda }}}^0(n,\mathrm{},c).`$ For this it is enough to check (\*)<sub>k</sub> from Definition 2.1(1), so let $`\mathrm{\Lambda }^{}`$ be a subset of $`\mathrm{\Lambda }_H`$ with $`n`$ elements and $`\alpha ^{}\mathrm{\Lambda }_H\mathrm{\Lambda }^{}`$, also let $`M`$ be a fim of vocabulary $`\tau `$ and dimension $`k`$ (i.e. $`P^M`$ is with $`k`$ members), $`V=`$ Space$`{}_{\overline{\mathrm{\Lambda }}}{}^{}(M)`$, and $`d`$ an $`(\mathrm{},\alpha ^{},H)`$invariant and $`(M,\mathrm{\Lambda }^{},H)`$invariant $`C`$colouring of $`V`$ such that $`C`$ has $`c`$ members. So we just have to prove that the conclusion of Definition 2.2 holds, which means there is a $`d`$monochromatic line of $`V`$.
Let $`w_1=^{df}\{a:aP^M`$ and the number of $`b<^M`$ $`a`$ is $`k_0\mathrm{}1`$ but is $`<k_0\mathrm{}1+k_1\}`$ hence in $`w_1`$ there are $`k_1`$ members, and let $`w_0`$ be the set of first $`k_0\mathrm{}1`$ members of $`P^M`$ by $`<^M`$, and lastly let $`w_2`$ be the set of the $`\mathrm{}`$ last members of $`M`$ by $`<^M.`$ So $`w_0,w_1,w_2`$ form a convex partition of $`P^M.`$
Now we let $`K`$ be $`M`$ restricted to cl$`{}_{M}{}^{}(w_0w_2)`$, (note that this gives no contradiction to the assumption on $`K`$ i.e. clause (d) of the assumptions, as concerning $`K`$ there, only its vocabulary and dimension are important and they fit). Let $`K^+`$ be a fim with vocabulary $`\tau `$ and dimension $`k_0`$, let $`g_0`$ PHom$`(M,K^+)`$ be the following function from $`P^M`$ onto $`P^{K^+}:`$ it maps all the members of $`w_1`$ to one member of $`P^{K^+}`$ which we call $`b^{}`$, it is a one to one order preserving function from $`w_2`$ onto $`\{bP^{K^+}:b^{}<^{K^+}b\}`$ and it is a one to one order preserving function from $`w_0`$ onto $`\{bP^{K^+}:b<^{K^+}b^{}\}.`$ Let $`g`$ Hom$`(M,K^+)`$ be the unique extension of $`g_0;`$ without loss of generality $`g_0`$ is the identity on $`w_0`$ and on $`w_2`$ hence without loss of generality $`g`$ is the identity on $`K`$, it exist by 1.2.
Next recall that the vocabulary $`\tau ^{}=\tau _{K,w_o,w_2}\{H\}`$ is a well defined vocabulary ( see Definition 1.5(1) and remember that $`\tau \tau _{K,w_0,w_2}`$ so $`H\tau _{K,w_0,w_2}`$). Next we shall define a $`\tau ^{}`$model $`N.`$ Its universe is $`MKA^{}`$ where $`A^{}=^{df}\{bM:`$ base$`{}_{M}{}^{}(b)w_1`$ and $`F_{M,b}=H\}`$, we let $`P^N`$ be $`w_1`$ and $`<^N`$ be $`<^MP^N.`$ Now we have to define each function $`F_{K,\overline{a}_1,\overline{a}_2}^N`$, say of arity $`r`$, where $`F\tau ,\overline{a}_1`$ a non decreasing sequence form $`w_0`$ and $`\overline{a}_2`$ a non decreasing sequence from $`w_2`$, and lg$`(\overline{a}_1)+`$ lg$`(\overline{a}_2)<`$ arity(F) and arity$`(F_{K,\overline{a}_1,\overline{a}_2})`$ $`<`$ arity$`(\tau ).`$ Note that the last condition is equivalent to : if $`F=H`$ then at least one of the sequences $`\overline{a}_1,\overline{a}_2`$ is not empty.
For $`b_1^N\mathrm{}^Nb_rP^N`$ we let $`F_{\overline{a}_1,\overline{a}_2}^N(b_1,\mathrm{},b_t)`$ be equal to
$`b=F^M(\overline{a}_1,b_1,\mathrm{},b_t,\overline{a}_2)`$ $`=F^M(a_1^1,a_2^1,\mathrm{},a_{\mathrm{lg}(\overline{a}_1)}^1`$, $`b_1,\mathrm{},b_t,a_1^2,\mathrm{},a_{\mathrm{lg}(\overline{a}_2)}^2).`$
It is easy to check that the number of arguments is right and also the sequence they form is $`^M`$increasing, so this is well defined and belongs to $`M`$, but still we have to check that it belongs to $`N.`$ First note that it does not belong to $`K`$, as if $`bK`$ then base$`{}_{\mathrm{lg}(\overline{a}_1)+1}{}^{}(b)K`$ and it is just $`b_1`$ which belongs to $`w_1`$, contradiction. Second note that it does not belongs to $`A^{}`$, this holds as we have substructed $`H`$ when we have defined $`\tau ^{}.`$ Lastly it is also trivial to note that every member of $`N`$ has this form. It is easy to check that $`N`$ is really a $`\tau ^{}`$fim.
We next let $`V^{}=`$ Space$`{}_{\overline{\mathrm{\Lambda }}^{}}{}^{}(N)`$ and let $`C^{}=\{g:g`$ is a function from Space$`{}_{\overline{\mathrm{\Lambda }}}{}^{}(K)`$ to $`C\}`$ and define a $`C^{}`$-colouring $`d^{}`$ of $`V^{}.`$ For $`\eta V^{}`$ let $`d^{}(\eta )`$ be the following function from Space$`{}_{\overline{\mathrm{\Lambda }}}{}^{}(K)`$ to $`C`$, letting $`\varrho `$ be the function with domain $`A^{}`$ which is constantly $`\alpha ^{}:`$ for $`\nu K`$ we let $`\left(d^{}(\eta )\right)(\nu )=d(\eta \nu \varrho ).`$
Clearly the function $`d^{}(\eta )`$ is a $`C^{}`$colouring of Space$`{}_{\overline{\mathrm{\Lambda }}}{}^{}(K)`$. How many such functions there are? The domain has clearly card(Space$`{}_{\overline{\mathrm{\Lambda }}}{}^{}(K))`$ members, (we can get slightly less if $`\mathrm{}>0`$, but with no real influence). The range has at most $`c`$ members, so the number of such functions is at most $`c^{\mathrm{card}(\mathrm{Space}_{\overline{\mathrm{\Lambda }}}(K))}`$, a number which we have called $`c^{}`$.
So $`d^{}`$ is a $`c^{}`$colouring.
Now as we have chosen $`k_1=f^1(\overline{\mathrm{\Lambda }}^{},c^{})`$ we can apply Definition 2.2 to $`V^{}=`$ Space$`{}_{\overline{\mathrm{\Lambda }}^{}}{}^{}(N)`$ and $`d^{};`$ so we can find a $`d^{}`$monochromatic $`V^{}`$line and we call it $`L^{}.`$ Let $`h`$ be the function from $`U=^{df}`$ Space$`{}_{\overline{\mathrm{\Lambda }}}{}^{}(K^+)`$ to $`V`$ defined as follows:
(\*) $`h(\rho )=\nu `$ iff :
(a) $`\nu V,\rho U`$,
(b) $`\nu K=\rho K`$
(c) if $`bN`$ supp$`{}_{N}{}^{}(L^{})`$ (see Def 1.7(3)) then $`\nu (b)=\eta (b)`$ for every $`\eta L^{}.`$
(d) if a $`A^{}`$ cl<sub>M</sub>( supp$`{}_{N}{}^{}(L^{}))`$ then $`\rho (a)=\alpha ^{}.`$
(e) if a $``$ supp$`{}_{N}{}^{}(L^{})`$, (so a $`N`$, $`F_{N,a}=F_{K,\overline{a}_1,\overline{a}_2}`$, base$`{}_{N}{}^{}(a)`$ supp$`{}_{N}{}^{P}(L^{}))`$, and $`bK^+`$, $`F_{K^+,b}=F,b=F(\overline{a}_1,b^{},\mathrm{},b^{},\overline{a}_2)`$ (with the number of cases of $`b^{}`$ being arity( $`F_{K,\overline{a}_1,\overline{a}_2}))`$ then $`\rho (b)=\nu (a).`$
(f) if $`aA^{}`$cl<sub>M</sub>(sup$`{}_{N}{}^{}(L^{}))`$ and $`b`$ $`K^+`$ is $`H(b^{},\mathrm{},b^{})`$ then $`\rho (b)=\nu (a).`$
Let the range of $`h`$ be called $`S.`$ Now clearly
$`_1`$ $`(\alpha )`$ $`h`$ is a one to one function from $`U`$ to $`SV.`$
$`(\beta )`$ $`S`$ has $`|`$Space$`{}_{\overline{\mathrm{\Lambda }}}{}^{}(K^+)|`$ members
$`(\gamma )`$ $`S`$ is a subspace of $`V`$ of dimension $`k_0`$, such that $`h(\rho )=`$ pt$`{}_{S}{}^{}(\rho )`$, see 1.7(7).
Now clearly
$`_2`$ there is a $`C`$colouring $`d^{}`$ of $`U`$ such that:
$`d^{}(\nu )=d(h(\nu ))`$ for $`\nu U.`$
and
$`_3`$ (a) $`d^{}`$ is $`(K^+,\mathrm{\Lambda }^{})`$invariant
(b) $`d^{}`$ is $`(\mathrm{}+1,\alpha ^{},H))`$invariant
\[WHY? Reflect\]
Applying the definition of $`k_0f_{\tau ,\overline{\mathrm{\Lambda }}}^0(n,\mathrm{}+1,c)`$ , that is Definition 2.2 to $`\overline{\mathrm{\Lambda }}`$, $`\alpha ^{},U,d^{}`$ we can conclude that there is a $`d^{}`$monochromatic $`U`$line $`L^{}.`$ Let $`L=^{df}\{h(\rho ):\rho L^{}\}.`$ It is easy to check that $`L`$ is as required.
$`\text{}_{2.4}`$
As a warm up for the later bounds we prove:
2.5 Theorem: (1) The function $`f_\tau ^1(\overline{\mathrm{\Lambda }},c)`$ is well defined, i.e. always get value, a natural number.
Moreover has a bound which we have got by triple induction.
(2) Similarly the function $`f^0.`$
Proof: (1) The proof follows by induction, the main induction is on $`t=`$ arity$`(\tau _{\overline{\mathrm{\Lambda }}}).`$ Now by observation 1.10(1) without loss of generality $`\tau `$ is monic, i.e. has a unique function symbol of arity $`t`$, called $`H=^{df}F_\tau ^{\mathrm{max}}.`$ Fixing $`t`$, we prove by induction on $`s=|\mathrm{\Lambda }_H|.`$
CASE 0: $`t=1`$
This is Hales-Jewett theorem (on a bound see \[Sh:329\] and \[GRS80\])
CASE 1: $`t>1,s=1`$
By claim 1.10(2) we can decrease $`t.`$
CASE 2: $`t>1,s2`$
We note that $`f^1(\overline{\mathrm{\Lambda }},c)=f_{\overline{\mathrm{\Lambda }}}^0(0,0,c)`$ by 2.3(4) so it is enough to bound the later one. But by 2.3(5) we know $`f_{\overline{\mathrm{\Lambda }}}^0(|\mathrm{\Lambda }_H|,0,c)=f^1(\overline{\mathrm{\Lambda }}\tau ^{},c)`$ where $`\tau ^{}=^{df}\tau \{H\}`$, but for the later one we have a bound by the induction hypothesis on $`t`$ as arity$`(\tau ^{})t`$, so we have a bound on $`f_{\overline{\mathrm{\Lambda }}}^0(|\mathrm{\Lambda }_H|,0,c).`$ By the last two sentences together, it is enough to find a bound to $`f_{\overline{\mathrm{\Lambda }}}^0(n,0,c)`$ by downward induction on $`n|\mathrm{\Lambda }_H|`$, and we have the starting case : $`n=|\mathrm{\Lambda }_H|`$ and the case $`n=0`$ gives the desired conclusion. So assume we know for $`n+1`$ and we shall do it for $`n.`$ Let $`\mathrm{}^{}=^{df}f_{\overline{\mathrm{\Lambda }}}^0(n+1,0,c)`$, so we know that $`\mathrm{}^{}=f_{\overline{\mathrm{\Lambda }}}^0(n,\mathrm{}^{},c)`$ by 2.3(6), so we by downward induction on $`\mathrm{}\mathrm{}^{}`$ give a bound to $`f_{\overline{\mathrm{\Lambda }}}^0(n,\mathrm{},c).`$ So we are left with bounding $`f_{\overline{\mathrm{\Lambda }}}^0(n,\mathrm{},c)`$ given bound for $`f_{\overline{\mathrm{\Lambda }}}^0(n,\mathrm{}+1,c)`$ ( and also $`f_\tau _{}^1(\overline{\mathrm{\Lambda }}^{},c^{})`$ whenever arity$`(\tau _{})<t).`$ For this 2.4 was designed, it says
$`f_{\overline{\mathrm{\Lambda }}}^0(n,\mathrm{},c)f_{\overline{\mathrm{\Lambda }}}^0(n,\mathrm{}+1,c)+f_\tau ^{}^1(\overline{\mathrm{\Lambda }}^{},c)+1`$
where $`\tau ^{},\overline{\mathrm{\Lambda }}^{}`$ were defined there and arity$`(\tau ^{})<`$ arity$`(\tau )`$, (well, we have to assume that $`\mathrm{}<f_{\overline{\mathrm{\Lambda }}}^0(n,\mathrm{}+1,c)`$, but otherwise use $`\mathrm{}+1+f_\tau ^{}^1(\overline{\mathrm{\Lambda }}^{},c)+1`$
$`\text{}_{2.5}`$
Section 3 : Higher Dimension Theorems
Concerning the multidimensional case (see Def 1.9):
3.1 Conclusion: (1) For any $`\overline{\mathrm{\Lambda }}`$, $`m`$ and $`c`$, we have $`f^1(\overline{\mathrm{\Lambda }},m,c)`$ is well defined (with bound as in the proof, actually using one further induction using only $`f_{\tau _i}^1(\overline{\mathrm{\Lambda }},c)`$ for suitable $`\tau _i`$-s in teh $`i`$step.)
(2) We can naturally defined $`\tau `$fim of dimension $`\mathrm{}_0`$ and convex subspaces, and prove that for any $`\tau `$fim $`M`$ of dimension $`\mathrm{}_0`$ and alphabet sequence $`\overline{\mathrm{\Lambda }}`$, if Space$`{}_{M}{}^{}(\overline{\mathrm{\Lambda }})`$ is the union of finitely many Borel subsets, then some convex subspace $`S`$ of dimension $`\mathrm{}_0`$ is included in one of those Borel subsets.
Proof: (1) For simplicity (and without loss of generality by 1.10(1)) we have $`\overline{\mathrm{\Lambda }}`$ is constantly $`\mathrm{\Lambda }`$, so each $`\mathrm{\Lambda }_F`$ is $`\mathrm{\Lambda }`$, a fixed alphabet. We choose by induction on $`i=0,\mathrm{},m`$ the objects $`M_i,\tau _i,k_i`$ and $`c_i`$ such that
(a) $`k_0=0`$ and $`k_i<k_{i+1}`$
(b) $`M_i`$ is a fim for $`\tau `$ of dimension $`k_i`$ (we allow empty fim, if you do not like it start with $`k_0=1`$)
(c) $`M_{i+1}`$ is an end extension of $`M_i`$
(d) $`\tau _i=\tau _{M_i,P^{M_i},\mathrm{}}`$ (see Definition 1.5(1) )
(e) $`c_0`$ is $`c`$ and $`c_{i+1}`$ is $`c^{|\mathrm{Space}_\mathrm{\Lambda }(k_i+mi)|}`$
(f) $`k_{i+1}=k_i+f_{\tau _i}^1(\mathrm{\Lambda },c_i).`$
There is no problem to carry the definition and we can prove that $`k_mf_\tau ^1(\mathrm{\Lambda },m,c).`$
The proof is straight.
2) Such theorems are closed relatives to theorems on appropriate forcing notions, as anyhow it is a set theoretical theorem we use such approach. Specifically we use the general treatment of creature forcing of \[RoSh 470\]. For any finite non empty $`u\omega `$ let $`M_u^\tau `$ be a $`\tau `$model with $`(P^{M_u^\tau },^{M_u^\tau })=(u,)`$, and without loss of generality $`u_1u_2M_{u_1}^\tau M_{u_2}^\tau `$. So for infinite $`u\omega `$ we have $`M_u^\tau =\{M_{u_1}^\tau :u_1u`$ finite $`\}`$ is well defined.
A $`\overline{\mathrm{\Lambda }}`$-creature $`𝔠`$ consist of a convex subspace $`S^𝔠=S[𝔠]`$ of some $`M_u^\tau `$ for some finite non empty $`u=u[𝔠]`$ of the form $`[n,m]=[n_𝔠,m_𝔠]`$ .
Fot creatures $`𝔠_1,\mathrm{},𝔠_k`$ we let $`\mathrm{\Sigma }(𝔠_1,\mathrm{},𝔠_k)`$ be well defined iff $`m_𝔠_{\mathrm{}}=n_{𝔠_{\mathrm{}+1}}`$ for $`\mathrm{}[1,k)`$ and it is the set of $`\overline{\mathrm{\Lambda }}`$creatures $`𝔠`$ such that $`n_𝔠=n_{𝔠_1},m_𝔠=m_{𝔠_k}`$ and $`\eta S^𝔠\mathrm{}[1,k)\eta u[𝔠_{\mathrm{}}]S^𝔠_{\mathrm{}}`$.
So the forcing notion $`Q`$ is well defined by \[RoSh 470\] for the case the lim-sup of the norms is infinity”. So a condition $`p`$ has the form $`\eta ,𝔠_1,𝔠_2,\mathrm{}=\eta ^p,𝔠_1^p,𝔠_2^p,\mathrm{}`$ where for $`t=1,2,\mathrm{}`$, $`𝔠_t^p`$ is a $`\overline{\mathrm{\Lambda }}`$creature , $`m_{𝔠_{t+1}^p}=n_{𝔠_t^p}`$. Let $`B=^{df}\mathrm{Space}_{\overline{\mathrm{\Lambda }}}(M_\omega ^\tau )=\{\rho :\rho \text{ is a function with domain }M_\omega ^\tau \text{ satisfying }f(b)\mathrm{\Lambda }_{F(b)}\}`$ where $`F(b)=F_{M_\omega ^\tau ,b}`$. We say that $`\rho B`$ obeys $`pQ`$ if $`\eta ^p\rho `$ and for $`t=1,2,\mathrm{}`$ we have $`\rho u^{𝔠_t^p}S^{𝔠_t^p}`$. It is proved there that such forcing notions has many good properties. In particular letting cont$`(p)=\{\rho :\rho B\text{ obeys }p\}`$ is a function with domain and defining the $`Q`$-name $`\widehat{f}=\{f^p:p\widehat{G}_Q\}`$. Now note that
(a) $`p_Q\text{}\widehat{f}`$ cont$`(p)`$
(b) if $`N((\chi ),)`$ is countable, the definition of those countably many Borel sets belongs to $`N`$, and $`pQN`$, then we can find $`q`$ such that
(i) $`pq`$
(ii) every $`f\mathrm{cont}(q)`$ is a generic for $`Q`$ over $`N`$
(iii) for some $`p^{},n^{}`$ we have $`pp^{}NQ,p^{}q`$ and $`p^{}_Q\text{}\widehat{f}A_n^{}\text{}`$
Together we conclude that $`\mathrm{cont}(q)A_n^{}`$ and we are done.
$`\text{}_{3.1}`$
We turn to relating the old results from Bergelson Leibman \[BL96\]
3.2 Conclusion: (1) Assume that
(a) $`\tau `$ is a $`t`$canonical vocabulary (see 1.4)
(b) $`k=f_\tau ^1(\mathrm{\Lambda },c)`$, $`\mathrm{\Lambda }`$ a (finite) alphabet
(c) $`R`$ is a ring, and $`r_1,\mathrm{},r_kR`$
(d) for $`\alpha \mathrm{\Lambda },𝐩_\alpha (x)`$ is a polynomial over $`R`$ (i.e. with parameters in $`R`$).
(e) $`d`$ is a $`c`$colouring of $`R`$ ( actually enough to consider a finite subset, the range of $`g`$ in the proof below)
Then we can find $`y,z`$ and $`w\{1,\mathrm{},k\}`$ such that
$`(\alpha )`$ $`yR`$ and $`z=\mathrm{\Sigma }_\mathrm{}wr_{\mathrm{}}R`$
$`(\beta )`$ the set $`\{y+𝐩_\alpha (z):\alpha \mathrm{\Lambda }\}`$ is $`d`$monochromatic
(2) Assume that
(a) $`\tau `$ is a vocabulary of arity $`t`$, such that for each $`s=1,\mathrm{},t`$ in $`\tau `$ there are exactly $`m^{}`$ function symbols of arity $`s`$
(b) $`k=f_\tau ^1(\mathrm{\Lambda },c)`$, $`\mathrm{\Lambda }`$ a (finite) alphabet
(c) $`R`$ is a ring, and $`r_1,\mathrm{},r_kR`$
(d) for $`\alpha \mathrm{\Lambda }`$ and $`m<m^{},𝐩_{\alpha ,m}(x)`$ is a polynomial over $`R`$ (i.e. with coefficients in $`R`$).
(e) $`d`$ is a $`c`$colouring of $`R^m^{}=\{y_m:m<m^{}:y_0,\mathrm{},y_{m^{}1}R\}`$ (actually enough to consider a finite subset, the range of $`g`$ in the proof below).
Then we can find $`y,z`$ and $`w\{1,\mathrm{},k\}`$ such that
$`(\alpha )`$ $`yR`$ and $`z=\mathrm{\Sigma }_\mathrm{}wr_{\mathrm{}}R`$
$`(\beta )`$ the set $`\{y+𝐩_{\alpha ,m}(z):m<m^{}:\alpha \mathrm{\Lambda }\}`$ is $`d`$monochromatic
Proof: (1) Let $`M`$ be a fim for $`\tau `$ of dimension $`k`$ and let $`h`$ be a one to one order preserving function from $`P^M`$ to $`\{1,\mathrm{},k\}.`$ We define a function $`g`$ from $`V=`$ Space$`{}_{\mathrm{\Lambda }}{}^{}(M)`$ to $`R.`$ For $`\eta V`$ we let $`g(\eta )=\mathrm{\Sigma }_{bM}g_b(\eta (b))`$ where $`g_b`$ is the following function from $`\mathrm{\Lambda }`$ to $`R`$. For $`b=F(b_1,\mathrm{},b_t)M`$ and $`\alpha \mathrm{\Lambda }`$ we let $`g_b(\alpha )`$ be zero if $`b_1,b_2,\mathrm{},b_t`$ is with repetitions and otherwise we consider $`𝐩_\alpha (\mathrm{\Sigma }_{i=1,t}r_{h(b_i)})`$, expand it as sum of monoms in $`r_1,\mathrm{},r_k`$ , and let $`g_b(\alpha )`$ be the sum of those monoms for which $`\{r_j:j\{1,\mathrm{},k\}`$ and $`r_j`$ appear in the monom $`\}`$ $`=`$ $`\{h(b_1),\mathrm{},h(b_t)\}`$. Now we define a $`c`$colouring $`d^{}`$ of $`V`$ by $`d^{}(\eta )=d(g(\eta )).`$ Let $`L`$ be a $`d^{}`$monochromatic line of $`V`$ , let supp$`{}_{M}{}^{}(L)=N.`$ Now let $`y=^{df}\mathrm{\Sigma }_{bMN}g_b`$(pt$`{}_{L}{}^{}(\alpha ))`$, note that all the $`\alpha \mathrm{\Lambda }`$ gives the same value. Let $`w=^{df}\{h(b):b`$ supp$`{}_{M}{}^{P}(L)\}`$, recalling Def 1.7(5) and so $`z=\mathrm{\Sigma }_\mathrm{}wr_{\mathrm{}}`$, now check.
Note that algebraically it is more natural to defined $`g`$ differently, working by the rank of the monom rather that by the set of variables appearing.
(2) Similarly, left to the reader.
$`\text{}_{3.2}`$
3.3 Discussion: It is natural to ask:
(1) Can we generalize the Graham Rothschild theorem? (see \[GR 71\], \[GRS 80\])
(2) Can we get here primitive recursive bounds?
(3) Can we prove the density version of the theorem (2.11)?
Below we answer positively questions (1),(2), we believe that the answer to question (3) is positive too but probably it require methods of dynamical systems, see the book Furstenberg \[Fu81\].
3.4 DEFINITION: We define $`f^4(\overline{\mathrm{\Lambda }},t,\mathrm{},c)=f_\tau ^4(\overline{\mathrm{\Lambda }},t,\mathrm{},c)`$ where 0 $`\mathrm{}<t`$ as follows. It is the minimal $`k`$ such that: if $`M`$ is fim for $`\tau ,V=`$ Space$`{}_{\overline{\mathrm{\Lambda }}}{}^{}(M)`$ and $`d`$ is a $`c`$colouring of $`\{S:S`$ is an $`\mathrm{}`$subspace of $`V\}`$ then for some subspace $`U`$ of $`V`$ of dimension $`t`$, all the $`\mathrm{}`$subspaces of $`U`$ (equivalently, $`\mathrm{}`$subspaces of $`V`$ which are contained in $`U`$) have the same colour by $`d.`$
3.5 Theorem : (1) For any $`\overline{\mathrm{\Lambda }}`$, $`t,\mathrm{},c`$ as in Definition 3.3, the function $`f^4(\overline{\mathrm{\Lambda }},t,\mathrm{},c)`$ is well defined, i.e. is finite.
(2) Let $`m=`$ RAM$`(t,\mathrm{},c)`$, see Definition 0.3(1), where $`\tau `$ is a vocabulary and $`\overline{\mathrm{\Lambda }}`$ is a $`\tau `$alphabet sequence, and define $`k_i`$ for $`i=0,\mathrm{},m`$ by induction on $`i`$ as follows (on $`\tau ^{[k,r]}`$ see 1.5(3)):
$`k_0=0,\overline{\mathrm{\Lambda }}^0=\overline{\mathrm{\Lambda }}`$ and $`k_{i+1}=k_i+f_{\tau _i}^1(\overline{\mathrm{\Lambda }}^i,c_i)`$ where $`\tau _i=^{df}\tau ^{[k_i,mi]}`$ and $`\overline{\mathrm{\Lambda }}^i`$ is a $`\tau ^{[k,mi]}`$alphabet sequence, and $`\mathrm{\Lambda }_{F_{M_{k_i+mi}^\tau ,\overline{a}_1,\overline{a}_2}}^i`$ has $`|\mathrm{\Lambda }_F^i|+\mathrm{}+|M_{k_i+mi}^\tau |`$ members and $`c_i=c^{\mathrm{card}(\mathrm{Space}_{\overline{\mathrm{\Lambda }}^i}(M_{k_i+mi}^\tau ))}.`$
Then $`f_\tau ^4(\mathrm{\Lambda },t,\mathrm{},c)k_m.`$
Proof:(1) Follows from (2).
(2) Let $`N=M_{\mathrm{}}^\tau `$ (see notation in 1.5(3), recall that $`\mathrm{}`$ is the dimension of the subspaces we are colouring) and let $`\{\gamma _a:aN\}`$ list a set disjoint to $`\mathrm{\Lambda }`$ without repetitions.
We choose for $`i=0,\mathrm{},m`$ the objects $`k_i,\tau _i,\overline{\mathrm{\Lambda }}^i`$ (consistently with what is said in the statement of the theorem) and $`M_i,M_i^+`$, by induction on $`i`$ as follows:
$`_1`$ (a) $`k_0=0`$ and $`k_i<k_{i+1}`$
(b) $`M_i`$ is a fim for $`\tau `$ of dimension $`k_i`$ (we allow empty fim, the space is the a singleton, if you do not like it start with 1)
(c) $`M_{i+1}`$ an end extension of $`M_i`$ and $`M_i^+`$ is an end extension of $`M_i`$ (so both have vocabulary $`\tau )`$ and has dimension $`k_i+mi`$
(d) $`\tau _i=\tau _{M_i^+,P^{M_i},P^{M_i^+}P^{M_i}}`$ (see Definition 1.5(3) )
(e) $`\overline{\mathrm{\Lambda }}^0=\overline{\mathrm{\Lambda }}`$ and $`\mathrm{\Lambda }_{F_{\overline{a}_1,\overline{a}_2}}^i`$ is the disjoint union of $`\mathrm{\Lambda }_F,\mathrm{\Lambda }_F^{}=^{df}\{\gamma _b:bN`$ and $`F_{N,b}=F\}`$ and $`\{\beta _b:bM_i^+`$ such that $`F_{M_i^+,b}=F\}`$ (and no two letter are incidentally equal, of course).
(f) $`c_0`$ is $`c`$ and $`c_{i+1}`$ is $`c^{\mathrm{card}(\mathrm{Space}_{\overline{\mathrm{\Lambda }}_i}(M_{k_i+mi}^\tau ))}`$
(g) $`k_{i+1}=k_i+f_{\tau _i}^1(\overline{\mathrm{\Lambda }}^i,c_i).`$
Let $`k=k_m,M=M_k`$ and let $`V_i=`$ Space$`{}_{\overline{\mathrm{\Lambda }}}{}^{}(M_i^\tau )`$ and $`V=V_m.`$ We shall regard an $`\mathrm{}`$subspace $`\mathrm{\Phi }`$ of $`V`$ as a function from $`M`$ to $`\mathrm{\Lambda }^{}=\{\gamma _b:bN\}\mathrm{\Lambda }`$, such that (and where):
$`_2`$ (a) $`\mathrm{\Lambda }=_{F\tau }\mathrm{\Lambda }_F`$,
(b) $`\mathrm{\Phi }(b)\mathrm{\Lambda }_{F_{M,b}}^{}\mathrm{\Lambda }_{F_{M,b}}`$, see clause (e) of $`_1`$
(c) if $`bM,\alpha \mathrm{\Lambda }`$ and $`(\nu )[\nu \mathrm{\Phi }\nu (b)=\alpha ]`$ then $`\mathrm{\Phi }(b)=\alpha `$
(d) if $`bM,aN`$ and for every $`\rho `$ Space$`{}_{\overline{\mathrm{\Lambda }}}{}^{}(N)`$ we have (pt$`{}_{\mathrm{\Phi }}{}^{}(\rho ))(b)=\rho (a)`$ then $`\mathrm{\Phi }(b)=\gamma _a.`$
(Reflect on the meaning of $`\mathrm{}`$subspace of $`M`$, i.e. Definition 1.7(7) and it should be clear.) Let $`d`$ be a $`c`$colouring of the set of $`\mathrm{}`$subspaces of $`V.`$ We shall define by downward induction on $`i<m`$ a (non empty) subset $`A_i`$ of $`P^{M_{i+1}}`$ disjoint to $`M_i`$ and a function $`\varrho _i`$ from $`B_i=^{df}M`$ cl$`{}_{M}{}^{}(M_i_{j=i,\mathrm{},m1}A_j)_{j=i+1,\mathrm{},m1}B_j`$ into $`\mathrm{\Lambda }\{\beta _a:aM_i^+\}.`$ We let $`𝐑_i`$ denote the family of $`\mathrm{}`$subspaces $`\mathrm{\Phi }`$ of $`V`$ which satisfies:
(\*)<sub>1</sub> (a) if $`j`$ satisfies $`ij<m`$ and $`bB_j`$ and $`\varrho _j(b)\mathrm{\Lambda }`$ then $`\mathrm{\Phi }(b)=\varrho _j(b)`$
(b) if $`j`$ satisfies $`ij<m`$ and $`bB_j`$ and $`\varrho _j(b)=\beta _a`$ where a $`M_j`$ then $`\mathrm{\Phi }(b)=\mathrm{\Phi }(a)`$
(c) if $`b_1,b_2`$ satisfies the following then $`\mathrm{\Phi }(b_1)=\mathrm{\Phi }(b_2)`$ where the demand is:
(i) $`b_1,b_2`$ cl$`{}_{M}{}^{}(M_i_{j=i,\mathrm{},m1}A_j)`$ and
(ii) $`F_{M,b_1}=F_{M,b_2}`$ and for every $`r\{1,\mathrm{}`$, arity($`F_{M,b_1})`$ $`\}`$ we have: base$`{}_{M,r}{}^{}(b_1)=`$ base$`{}_{M,r}{}^{}(b_2)`$ or they both belongs to the same $`A_j`$ for some $`j\{i,\mathrm{},m1\}`$.
Now $`A_i,B_i,\varrho _i`$ will be chosen such that the following condition holds
(\*)<sub>2</sub> If $`\mathrm{\Phi },\mathrm{\Psi }𝐑_i`$ satisfy the clauses (a),(b) below then $`d(\mathrm{\Phi })=d(\mathrm{\Psi })`$ where
(a) $`\mathrm{\Phi }`$ cl$`{}_{M}{}^{}(M_i_{j=i+1,\mathrm{},m1}A_j)=`$ $`\mathrm{\Psi }`$ cl$`{}_{M}{}^{}(M_i_{j=i+1,\mathrm{},m1}A_j)`$
(b) if $`bN`$ and $`\gamma _b`$ Rang$`(\mathrm{\Phi }M_{i+1})`$ then $`\gamma _b`$ Rang$`(\mathrm{\Phi }M_i)`$.
Suppose now that we have carried this induction, and we shall show that this suffice. Let $`S`$ be the following subset of $`V`$:
(\*)<sub>3</sub> $`\eta S`$ iff
(a) if $`i<m`$ and $`bB_i`$ and $`\varrho _i(b)\mathrm{\Lambda }`$ then $`\eta (b)=\varrho _i(b)`$
(b) if $`i<m`$ and $`bB_i`$ and $`\varrho _i(b)=\beta _a`$ and $`aM_i`$ then $`\eta (b)=\eta (a)`$.
Clearly $`S`$ is an $`m`$subspace of $`V`$, and we may by (\*)<sub>2</sub> above show that:
(\*)<sub>4</sub> if $`\mathrm{\Phi }`$ is an $`\mathrm{}`$subspace of $`S`$, then $`d(\mathrm{\Phi })`$ can be computed from $`J[\mathrm{\Phi }]=^{df}\{`$ Min $`\{i:\mathrm{\Phi }A_i`$ is constantly the $`r`$th member of $`P^N\}:r<\mathrm{}\}.`$
So for some function $`e`$, with domain the family of subsets of $`\{0,\mathrm{},m1\}`$ with $`\mathrm{}`$ elements, we have : if $`\mathrm{\Phi }`$ is an $`\mathrm{}`$subspace of $`S`$ then $`d(\mathrm{\Phi })=e(J[\mathrm{\Phi }])`$. Clearly the set $`\mathrm{Rang}(e)`$ has $`|\mathrm{Rang}(\mathrm{d})|`$ elements.
By Ramsey theorem and the choice of $`m`$, there is a subset $`w`$ of $`\{0,\mathrm{},m1\}`$ with $`t`$ members such that the function $`e`$ is constant on the family of subsets of $`w`$ with $`\mathrm{}`$ elements. Let $`U`$ be a subspace of $`S`$ of dimension $`t`$ such that if $`bM`$, base$`(b)`$ not a subset of $`_{iw}A_i`$ then $`\nu (b):bU`$ is constant (and the constant value belongs to $`\mathrm{\Lambda }_{M,b}`$.
Clearly $`U`$ is as required. The construction, i.e. the inductive choice of $`A_i,\varrho _i`$ is straight.
$`\text{}_{3.5}`$
Section 4: The main Theorem
Now we turn to the obtainment of primitive recursive bounds. The idea is that we decrease the dependency from below, dealing with the unary functions each time (rather than dealing with $`H\tau `$ of maximal arity).
In the definition below, we shall use the case $`r=1.`$
4.1 DEFINITION: (1) Recall that for $`a`$ $`P^M`$ we let $`M_a`$ be cl$`{}_{M}{}^{}(\{P^M\{a\})`$, that is $`M`$ restricted to this set.
(2) For $`V=`$ Space$`{}_{\overline{\mathrm{\Lambda }}}{}^{}(M)`$ and $`N`$ a closed subset of $`M`$ we say that a colouring $`d`$ of $`V`$ is $`(N,r)`$base-invariant if the following holds, for any $`a`$ $`P^N:`$
(\*) if $`\nu ,\eta V`$ and $`\nu M_a=\eta M_a`$ and $`[bMr<|\{i:i=1,\mathrm{}`$, arity$`(F_{M,b})`$ and base$`{}_{M,i}{}^{}(b)=a\}|\nu (b)=\eta (b)]`$ then $`d(\nu )=d(\eta ).`$
(3) We write $`(\mathrm{},r)`$base-invariant if above $`N`$ is such that $`P^N`$ is the set of the last $`\mathrm{}`$ members of $`P^M.`$
4.2 DEFINITION: Let $`f^6`$ be defined as follows. First, $`f_{\overline{\mathrm{\Lambda }}}^6(\mathrm{},c)=`$ $`f^6(\overline{\mathrm{\Lambda }},\mathrm{},c)=f_\tau ^6(\overline{\mathrm{\Lambda }},\mathrm{},c)`$ is defined iff $`\overline{\mathrm{\Lambda }}`$ is an alphabet sequence for a vocabulary $`\tau `$. Second, let $`f_{\overline{\mathrm{\Lambda }}}^6(\mathrm{},c)`$ be the first $`k`$ (natural number, if not defined we can understand it as $`\mathrm{}`$ or $`\omega `$ or ‘does not exist‘ ) such that (\*)<sub>k</sub> below holds, where:
(\*)<sub>k</sub> If clauses (a)-(d) below hold then there is a $`d`$monochromatic line of $`V`$, where :
(a) $`M`$ is a fim of vocabulary $`\tau `$
(b) the dimension of $`M`$ is $`k`$
(c) $`V=`$ Space$`{}_{\overline{\mathrm{\Lambda }}}{}^{}(M)`$
(d) $`d`$ is an $`(\mathrm{},1)`$base-invariant colouring of $`V`$.
Immediate connections are:
4.3 Observation: (1) The function $`f_{\overline{\mathrm{\Lambda }}}^6(\mathrm{},c)`$ increases with $`c`$ and decreases with $`\mathrm{}.`$
(2) We have $`f_{\tau _1}^6(\overline{\mathrm{\Lambda }}^1,\mathrm{}_1,c_1)`$ $`f_{\tau _2}^6(\overline{\mathrm{\Lambda }}^2`$, $`\mathrm{}_2,c_2)`$ if:
(a) $`c_1c_2,\mathrm{}_1\mathrm{}_2`$, and
(b) $`s`$ arity$`(\tau _1)\mathrm{\Pi }\{|\mathrm{\Lambda }_F^1|:F\tau _1`$, arity$`(F)=s\}`$ $`\mathrm{\Pi }\{|\mathrm{\Lambda }_F^2|:F\tau _2`$, arity$`(F)=s\}`$ and
(c) arity$`(\tau _1)<s`$ arity $`(\tau _2)F\tau _2`$ arity$`(F)=s`$ $`|\mathrm{\Lambda }_F^2|=1`$
(3) In definition 4.2 the demand holds for any larger $`k.`$
(4) $`f_{\overline{\mathrm{\Lambda }}}^6(0,c)=f^1(\overline{\mathrm{\Lambda }},c).`$
Proof: Trivial.
4.4 Claim: Assume
(a) $`\tau `$ is a vocabulary of arity $`>1`$ and $`\overline{\mathrm{\Lambda }}`$ is a $`\tau `$alphabet sequence
$`(b)`$ $`\tau ^{}`$ is the following vocabulary: $`\{G_{F,e}:F\tau `$, arity$`(F)>1`$ and $`e`$ is a convex equivalence relation on $`\{1,\mathrm{}`$, arity$`(F)\}`$ such that each $`e`$equivalence class has at least two elements $`\}`$
with arity$`(G_{F,e})=`$ the number of $`e`$equivalence classes and for some $`H\tau `$ of maximal arity, letting $`e=^{df}\{(i,j):i,j[1`$,arity$`(H)]\}`$ we identify $`G_{H,e}`$ with id
$`(c)`$ $`\overline{\mathrm{\Lambda }}^{}`$ is the following $`\tau ^{}`$alphabet sequence: $`\mathrm{\Lambda }_{G_{F,e}}^{}=\mathrm{\Lambda }_F.`$
$`(d)`$ $`\mathrm{}^{}=f_\tau ^{}^1(\overline{\mathrm{\Lambda }}^{},c)`$.
Then $`f_{\overline{\mathrm{\Lambda }}}^6(\mathrm{}^{},c)\mathrm{}^{}`$.
Proof: Let $`M`$ be a fim of vocabulary $`\tau `$ and dimension $`\mathrm{}^{}`$ and $`V=`$ Space$`{}_{\overline{\mathrm{\Lambda }}}{}^{}(M)`$ and $`d`$ is a $`c`$colouring of $`V`$ which is $`(\mathrm{}^{},1)`$base-invariant; it suffice to find a monochromatic $`V`$line $`L..`$
Let $`M^{}`$ be a fim of vocabulary $`\tau ^{}`$ and dimension $`\mathrm{}^{}`$ and $`V^{}=`$ Space$`{}_{\overline{\mathrm{\Lambda }}^{}}{}^{}(M^{}).`$ Let $`g_0`$ be an isomorphism from $`(P^M,<^M)`$ onto $`(P^M^{},<^M^{}).`$ We define a partial function $`g`$ from $`M`$ into $`M^{}`$ as follows; if $`b=F^M(b_1,\mathrm{},b_t)`$ so $`t=`$ arity$`{}_{\tau }{}^{}(F)`$ and $`b_1^Mb_2^M\mathrm{}^Mb_t`$ and $`e=\{(i,j):b_i=b_j\}`$ and $`G_{F,e}`$ $`\tau ^{}`$ is well defined (i.e. every $`e`$equivalence class has at least two elements) and the $`e`$equivalence classes are \[$`s_i,s_{i+1})`$ for $`i=1,\mathrm{}`$,arity$`(G_{F,e})1`$ and 1$`=s_1<s_2<\mathrm{}`$ $`<s_{\mathrm{arity}(G_{F,e})}=t+1`$ then $`g(b)=G_{F,e}^M^{}(g_0(b_{s_1}),\mathrm{}`$, $`g_0(b_{s_{\mathrm{arity}(G_{F,e})1}})).`$
Note:
$`()_1`$ $`g`$ is really a partial function from $`M`$ to $`M^{}`$
$`()_2`$ if $`\eta ,\nu V`$ and $`\eta \mathrm{Dom}(g)=\nu \mathrm{Dom}(g)`$ then $`d(\eta )=d(\nu )`$
\[Why? By the transitivity of equality, it is enough to consider the case that for some $`a^{}M\mathrm{Dom}(g)`$ we have $`\{a^{}\}=\{aM:\eta (a)\nu (a)\}`$. Now by the definition of $`g`$ for some $`aP^M`$ we have $`(!i)[\mathrm{base}_{M,i}(b)=a]`$. Now read the Definition 4.1(2),(3) of the base invariant\]
(\*)<sub>3</sub> we can define a $`c`$colouring $`d^{}`$ of $`V^{}`$ such that $`\eta V,\nu V^{}`$, and $`[b`$ Dom$`(g)\eta (b)=\nu (g(b))]`$ then $`d(\eta )=d^{}(\nu )`$
\[why? by (\*)$`{}_{2}{}^{}]`$
(\*)<sub>4</sub> for any $`V^{}`$line $`L^{}`$ there is a $`V`$line $`L`$ such that for every $`\eta L`$ for some $`\nu L^{}`$ we have $`d(\eta )=d^{}(\nu )`$
\[Why? Reflect. In details, let $`w^{}=`$ supp$`{}_{}{}^{P}(L^{})`$ and $`N^{}=`$ supp$`(L^{})`$ and $`\nu ^{}`$ is the function with domain $`M^{}N^{}`$ such that for every $`b`$ from this set and $`\nu L^{}`$ we have $`\nu (b)=\nu ^{}(b).`$ Let $`w=^{df}\{bP^M:g_0(b)w^{}\}`$ and let $`N=^{df}`$ cl$`{}_{M}{}^{}(w)`$ and choose a function $`\eta ^{}`$ with domain $`MN`$ such that for every $`bMN`$ we have $`\eta ^{}(b)=\nu ^{}(g(b))`$ if $`b\mathrm{Dom}(g)`$ and is any member of $`\mathrm{\Lambda }_{F_{M,b}}`$ otherwise. Let $`L`$ be the $`V`$line such that supp$`(L)=N`$ and for every $`\eta L`$ we have $`\eta `$ extend $`\eta ^{}.`$ Clearly $`L`$ is a $`V`$line and let $`\eta L`$ and we should check the desired conclusion. So there is $`p𝔭_{\overline{\mathrm{\Lambda }}}`$ such that $`\eta =`$ pt$`{}_{L}{}^{}(p);`$ now we define $`q𝔭_{\overline{\mathrm{\Lambda }}^{}}`$ as follows: $`q(G_{F,e})=p(F)`$, the later belongs to $`\mathrm{\Lambda }_F`$ which is equal to $`\mathrm{\Lambda }_{G_{F,e}}^{}.`$ Let $`\nu =`$pt$`{}_{L^{}}{}^{}(q)`$ and we should just check that $`\eta ,\nu `$ are as in (\*)<sub>3</sub> above so we are done.\]
By the assumption $`\mathrm{}^{}=f^1(\overline{\mathrm{\Lambda }}^{},c))`$ (see clause $`(d)`$ in the assumption), hence there is a $`d^{}`$monochromatic $`V^{}`$line $`L^{}.`$ Apply (\*)<sub>4</sub> to it, so there is a $`d`$monochromatic $`V`$line and so we are done.
$`\text{}_{4.4}`$
4.5 DEFINITION: (1) Assume the following:
(i) $`\overline{\mathrm{\Lambda }}`$ is an alphabet sequence for the vocabulary $`\tau `$
(ii) $`P\{(p,q):p,q`$ are $`\overline{\mathrm{\Lambda }}`$-types $`\}`$, see Def 1.6
(iii) $`m,c>0`$.
We define $`f_{\overline{\mathrm{\Lambda }}}^7(P,m,c)`$ as the first $`k`$ (if there is no such $`k`$ it is $`\omega `$ or $`\mathrm{}`$ or undefined) such that (\*)<sub>k</sub> stated below holds, where
(\*)<sub>k</sub> if clauses (a)-(e) below hold then there is a subspace $`S`$ of $`V`$ of dimension $`m`$, satisfying:
if $`L`$ is a $`V`$line $`S`$, and $`(p,q)P`$ then $`d`$(pt$`{}_{L}{}^{}(p))=d`$(pt$`{}_{L}{}^{}(q))`$
where
(a) $`M`$ is a fim of vocabulary $`\tau `$
(b) $`M`$ has dimension $`k`$
(c) $`V=`$ Space$`{}_{\overline{\mathrm{\Lambda }}}{}^{}(M)`$
(d) $`P`$ is a subset of $`\{(p,q):p,q𝔭_{\overline{\mathrm{\Lambda }}}`$ and $`[F\tau `$ arity$`(F)>1p(F)=q(F)]\}`$
(e) $`d`$ is a $`c`$colouring of $`V`$
(2) Let $`P`$$`{}_{\overline{\mathrm{\Lambda }}}{}^{}=`$ $`\{(p,q):p,q𝔭_{\overline{\mathrm{\Lambda }}}`$ and $`[F\tau `$ arity$`(F)>1p(F)=q(F)]\}`$
4.6 MAIN Claim: Assume
(a) $`\overline{\mathrm{\Lambda }}`$ is an alphabet sequence for a vocabulary $`\tau =\tau [\overline{\mathrm{\Lambda }}].`$
(b) $`k_0f_{\overline{\mathrm{\Lambda }}}^6(\mathrm{}+1,c)`$ and $`k_0>\mathrm{}.`$
(c) $`K`$ is a $`\tau `$fim of dimension $`k_01`$ with $`A_2`$ the last $`\mathrm{}`$ elements and $`A_1`$ the first $`(k_0\mathrm{}1)`$elements (this $`K`$ surve just for notation).
(d) $`\tau ^{}`$ is the vocabulary $`\tau _{K,A_1,A_2}`$, see Definition 1.5(1) and proj is the following function from $`\tau ^{}`$ to $`\tau :`$ it map $`F_{K,\overline{a}_1,\overline{a}_2}`$ to $`F`$ and $`\overline{\mathrm{\Lambda }}`$$`{}_{}{}^{}=_{}^{df}\mathrm{\Lambda }_F^{}:F\tau ^{}`$ where $`\mathrm{\Lambda }_F^{}=\mathrm{\Lambda }_{\mathrm{proj}(F)}`$, so proj $`\tau `$ is the identity.
(e) $`c^{}=^{df}c^{\mathrm{card}(\mathrm{Space}_{\overline{\mathrm{\Lambda }}}(K))}.`$
Then
$`f_{\overline{\mathrm{\Lambda }}}^6(\mathrm{},c)k_0+`$ $`f_{\overline{\mathrm{\Lambda }}^{}}^7(P_{\overline{\mathrm{\Lambda }}^{}},1,c^{})1`$
REMARK: This is similar to the proof of 2.4, but for completeness we do it in full.
Proof: Let $`k_1=f_{\overline{\mathrm{\Lambda }}^{}}^7(1,P_{\overline{\mathrm{\Lambda }}^{}},c^{})`$ and let $`k=k_0+k_11`$, so it suffice to prove that $`kf_{\overline{\mathrm{\Lambda }}}^0(n,\mathrm{},c).`$ For this it is enough to check (\*)<sub>k</sub> from Definition 4.2, also let $`M`$ be a fim of vocabulary $`\tau `$ and dimension $`k`$ (that is $`P^M`$ is with $`k`$ members), $`V=`$ Space$`{}_{\overline{\mathrm{\Lambda }}}{}^{}(M)`$, and $`d`$ an $`(\mathrm{},1)`$base-invariant $`C`$colouring of $`V`$ such that $`C`$ has $`c`$ members. So we just have to prove that the conclusion of Definition 4.2 holds, that is there is a monochromatic $`V`$line.
Let $`w_1=^{df}\{a:aP^M`$ and the number of $`b<^M`$ a is $`k_0\mathrm{}1`$ but is $`<k_0\mathrm{}1+k_1\}`$ hence in $`w_1`$ there are $`k_1`$ members, and let $`w_0`$ be the set of first $`k_0\mathrm{}1`$ members of $`P^M`$ by $`<^M`$, lastly let $`w_2`$ be the set of the $`\mathrm{}`$ last members of $`M`$ by $`<^M.`$ So $`w_0,w_1,w_2`$ form a convex partition of $`P^M.`$
Now we let $`K`$ be $`M`$ restricted to cl$`{}_{M}{}^{}(w_0w_2)`$, (note that this gives no contradiction to the assumption on $`K`$, as concerning $`K`$ there, only its vocabulary and dimension are important and they fit). Let $`K^+`$ be a fim with vocabulary $`\tau `$ and dimension $`k_0`$, let $`g_0`$ PHom$`(M,K^+)`$ be the following function from $`P^M`$ onto $`P^{K^+}:`$ it maps all the members of $`w_1`$ to one member of $`P^{K^+}`$ which we call $`b^{}`$, it is a one to one order preserving function from $`w_2`$ onto $`\{bP^{K^+}:b^{}<^{K^+}b\}`$ and it is a one to one order preserving function from $`w_0`$ onto $`\{bP^{K^+}:b<^{K^+}b^{}\}.`$ Let $`g`$ Hom$`(M,K^+)`$ be the unique extension of $`g_0;`$ without loss of generality $`g_0`$ is the identity on $`w_0`$ and on $`w_2`$ hence without loss of generality $`g`$ is the identity on $`K`$, it exist by 1.2.
Next recall that the vocabulary $`\tau ^{}=\tau _{K,w_o,w_2}`$ is a well defined vocabulary (see Definition 1.5(1)). Next we shall define a $`\tau ^{}`$fim $`N.`$ Its universe is $`MK;`$ we let $`P^N`$ be $`w_1`$ and $`<^N`$ be $`<^MP^N.`$ Now we have to define the function $`F_{K,\overline{a}_1,\overline{a}_2}^N`$, say of arity $`r`$, where $`F\tau ,\overline{a}_1`$ a non decreasing sequence from $`w_0`$ and $`\overline{a}_2`$ a non decreasing sequence from $`w_2`$, and lg($`\overline{a}_1)+`$ lg$`(\overline{a}_2)<`$ arity$`(F)`$. So $`r=\mathrm{arity}(F)\mathrm{lg}(\overline{a}_1)\mathrm{lg}(a_2)`$.
For $`b_1^N\mathrm{}^Nb_rP^N`$ we let $`F_{\overline{a}_1,\overline{a}_2}^N(b_1,\mathrm{},b_r)`$ be equal to
$`b=F^M(\overline{a}_1,b_1,\mathrm{},b_r,\overline{a}_2)`$ $`=F^M(a_1^1,a_2^1,\mathrm{},a_{\mathrm{lg}(\overline{a}_1)}^1`$, $`b_1,\mathrm{},b_r,a_1^2,\mathrm{},a_{\mathrm{lg}(\overline{a}_2)}^2).`$
It is easy to check that the number of arguments is right and also the sequence they form is $`^M`$increasing, so this is well defined and belongs to $`M`$, but still we have to check that it belongs to $`N.`$ But $`N=MK`$ and if $`bK`$ then base$`{}_{\mathrm{lg}(\overline{a}_1)+1}{}^{}(b)K`$ and it is just $`b_1`$ which belongs to $`w_1`$, contradiction. Lastly it is also trivial to note that every member of $`N`$ has this form. It is easy to check that $`N`$ is really a $`\tau ^{}`$fim.
We next let $`V^{}=`$ Space$`{}_{\overline{\mathrm{\Lambda }}^{}}{}^{}(N)`$ let $`C^{}=\{g:g`$ is a function from $`\mathrm{Space}_{\overline{\mathrm{\Lambda }}}(K)`$ to $`C\}`$ and define a $`C^{}`$ colouring $`d^{}`$ of $`V^{}.`$ For $`\eta V^{}`$ let $`d^{}(\eta )`$ be the following function from Space$`{}_{\overline{\mathrm{\Lambda }}}{}^{}(K)`$ to $`C:`$ for $`\nu K`$ we let $`\left(d^{}(\eta )\right)(\nu )=d(\eta \nu ).`$
Clearly the function $`d^{}(\eta )`$ is a $`C^{}`$colouring of $`K.`$ How many such functions, that is members of $`C^{}`$ there are? The domain has clearly card(Space$`{}_{\overline{\mathrm{\Lambda }}}{}^{}(K))`$ members, (we can get slightly less if $`\mathrm{}>0`$, but with no real influence). The range has at most $`c`$ members, so the number of such functions is at most $`c^{\mathrm{card}(\mathrm{Space}_{\overline{\mathrm{\Lambda }}}(K))}`$, a number which we have called $`c^{}`$ in the claim’s statement.
Hence $`d^{}`$ is a $`c^{}`$colouring.
So as we have chosen $`k_1=`$ $`f_{\overline{\mathrm{\Lambda }}^{}}^7(P_{\overline{\mathrm{\Lambda }}^{}},1,c^{})`$ we can apply Definition 4.5 to $`V^{}=`$ Space$`{}_{\overline{\mathrm{\Lambda }}^{}}{}^{}(N)`$ and $`d^{};`$ so we can find a $`d^{}`$monochromatic $`V^{}`$line $`L^{}.`$ Let $`h`$ be the function from $`U=^{df}`$ Space$`{}_{\overline{\mathrm{\Lambda }}}{}^{}(K^+)`$ to $`V`$ defined as follows:
(\*) $`h(\rho )=\nu `$ iff :
(a) $`\nu V,\rho U`$,
(b) $`\nu K=\rho K`$
(c) if $`bN`$ supp$`{}_{N}{}^{}(L^{})`$ then $`\nu (b)=\eta (b)`$ for every $`\eta L^{}.`$
(d) if a $``$ supp$`{}_{N}{}^{}(L^{})`$, (so a $`N`$, $`F_{N,a}=F_{K,\overline{a}_1,\overline{a}_2}`$, base$`{}_{N}{}^{}(a)`$ supp$`{}_{N}{}^{P}(L^{}))`$, and $`bK^+`$, $`F_{K^+,b}=F,b=F(\overline{a}_1,b^{},\mathrm{},b^{},\overline{a}_2)`$ (with the number of cases of $`b^{}`$ being arity$`(F_{K,\overline{a}_1,\overline{a}_2}))`$ then $`\rho (b)=\nu (a).`$
Let the range of $`h`$ be called $`S.`$ Now clearly
$`_1(\alpha )h`$ is a one to one function from $`U`$ to $`SV.`$
$`(\beta )`$ $`S`$ has $`|`$Space$`{}_{\overline{\mathrm{\Lambda }}}{}^{}(K^+)|`$ members
$`(\gamma )`$ $`S`$ is a subspace of $`V`$ of dimension $`k_0,h(\rho )=`$ pt$`{}_{S}{}^{}(\rho )`$, see Definition 1.7(7).
Now clearly
$`_2`$ there is a $`C`$colouring $`d^{}`$ of $`U`$ such that:
$`d^{}(\nu )=d(h(\nu ))`$ for $`\nu U.`$
and
$`_3`$ $`d^{}`$ is $`(\mathrm{}+1,1)`$base -invariant
\[WHY? Reflect\]
Applying the definition of $`k_0=f_{\overline{\mathrm{\Lambda }}}^6(\mathrm{}+1,c)`$ , that is Definition 4.2 to $`\overline{\mathrm{\Lambda }},U,d^{}`$ we can conclude that there is a $`d^{}`$monochromatic $`U`$line $`L^{}.`$ Let $`L=^{df}\{h(\rho ):\rho L^{}\}.`$ It is easy to check that $`L`$ is as required.
$`\text{}_{4.6}`$
4.7 Claim: (1) Assume that $`\overline{\mathrm{\Lambda }}`$ is a $`\tau `$alphabet sequence , and $`p^{}𝔭_{\overline{\mathrm{\Lambda }}}`$ and $`P`$$`{}_{}{}^{+}=`$ $`P`$ $`\{(p^{},q):q𝔭_{\overline{\mathrm{\Lambda }}}`$ and
$`[F\tau `$ $``$ arity $`(F)>1q(F)=p^{}(F)\}`$ $`P`$$`_{\overline{\mathrm{\Lambda }}}`$ (see Definition 4.5) and $`n=\mathrm{\Pi }_{F\tau ,\mathrm{arity}(F)=1}|\mathrm{\Lambda }_F|.`$ Then
$`f_{\overline{\mathrm{\Lambda }}}^7(P^+,m,c)`$ HJ$`(n,f_{\overline{\mathrm{\Lambda }}}^7(P,m,c),c)`$
(on HJ see 0.3(2)).
(2) $`f^7(\overline{\mathrm{\Lambda }},P_{\overline{\mathrm{\Lambda }}},m,c)`$ is in $`E`$$`_6.`$
Proof: (1) Straight. Let $`M`$ be a $`\tau `$fim of dimension $`k=^{df}`$ HJ$`(n,f_{\overline{\mathrm{\Lambda }}}^7(P,m,c),c)`$ and let $`V=`$ Space$`{}_{\overline{\mathrm{\Lambda }}}{}^{}(M)`$ and let $`d`$ be a $`c`$colouring of $`V.`$ Let $`\tau ^{}=\{F\tau :`$arity$`(F)=1\}`$ and let $`M^{}`$ be a $`\tau ^{}`$fim of dimension $`k`$, without loss of generality $`M^{}`$ is $`M`$ restricted to $`\tau ^{}`$ and the universe of $`M^{}.`$ Let $`\overline{\mathrm{\Lambda }}`$$`{}_{}{}^{}=\overline{\mathrm{\Lambda }}\tau ^{}`$ and let $`V^{}=`$ Space$`{}_{\overline{\mathrm{\Lambda }}^{}}{}^{}(M^{})`$ and let $`h`$ be the function from $`V^{}`$ to $`V`$ defined as follows: let $`\eta V^{}`$, for $`bM`$ let $`(h(\eta ))(b)`$ be $`\eta (b)`$ is $`bM^{}`$ and be $`p^{}(F_{M,b})`$ if $`bMM^{}.`$ So $`h`$ is a function as required and we define a $`c`$colouring $`d^{}`$ of $`V^{}`$ by $`d^{}(\nu )=d(h(\nu ))`$ for $`\nu V^{}.`$
Now we apply the definition of $`k=`$ HJ$`(n,f_{\overline{\mathrm{\Lambda }}}^7(P,m,c),c)`$ to the space $`V^{}`$ and the colouring $`d^{}`$ and we get a subspace $`S^{}`$ of $`V^{}`$ on which $`d^{}`$ is constant and has dimension $`f_{\overline{\mathrm{\Lambda }}}^7(P,m,c).`$ There is a unique subspace $`S^{}`$ of $`V`$ of dimension $`f_{\overline{\mathrm{\Lambda }}}^7(P,m,c)`$ such that $`\eta S^{}\eta M^{}S^{}.`$ Clearly:
$`()_1`$ if $`L`$ is a $`V`$line which is $`S`$ and $`(p,q)P^+P`$ then $`pt_L(p)=d(pt_L(q))`$
Now, letting $`k^{}=f_{\overline{\mathrm{\Lambda }}}^7(P,m,c)`$ and $`d^{}=dS^{}`$, we can apply the definition of $`f_{\overline{\mathrm{\Lambda }}}^7(P,m,c)`$ and get a subspace $`S`$ of $`S^{}`$ of dimension $`m`$ such that
$`()_2`$ if $`L`$ is a $`V`$line which is included in $`S`$ and $`(p,q)P`$ then $`d^{}(pt_L(p))=d^{}(pt_L(q))`$ which means that $`d(pt_L(p))=d(pt_L(q))`$
By $`()_1+()_2`$, clearly $`S`$ is as required.
(2) Let $`\{p_i^{}:i<i()\}`$ be maximal subset of $`P`$$`_{\overline{\mathrm{\Lambda }}}`$ such that $`i<j<i()p_i^{}\{F\tau :`$ arity$`(F)>1\}`$ $`p_j^{}\{F\tau :`$ arity$`(F)>1\}`$ and let $`P`$$`{}_{j}{}^{}=\{(p_i^{},q):i<j`$ and $`q𝔭_{\overline{\mathrm{\Lambda }}}`$ and $`q\{F\tau :`$ arity$`(F)>1\}=`$ $`p_i^{}\{F\tau :`$ arity$`(F)>1\}\}.`$ By part (1) we have a recursion formula (we use 1.10 freely):
$`f_{\overline{\mathrm{\Lambda }}}^7(P_{i+1},m,c)`$ HJ$`(|\mathrm{\Pi }_{F\tau ,\mathrm{arity}(F)=1}|\mathrm{\Lambda }_F|,f_{\overline{\mathrm{\Lambda }}}^7(P_i,m,c),c)`$
As HJ belongs to $`E`$<sub>5</sub> ( by \[Sh 329, 1.8(2),p.691), we are done.
$`\text{}_{4.7}`$
4.8 DEFINITION: Let $`f^{6,}(\overline{\mathrm{\Lambda }},\mathrm{},t,c)`$ is defined by induction on $`\mathrm{}`$ as follows:
$`f^{6,}(\overline{\mathrm{\Lambda }},0,t,c)=t`$
$`f^{6,}(\overline{\mathrm{\Lambda }},\mathrm{}+1,t,c)`$ is equal to $`k_0+f_{\overline{\mathrm{\Lambda }}[k_0]}^7(P_{\overline{\mathrm{\Lambda }}[k_0]},1,c^{\mathrm{card}(\mathrm{Space}_{\overline{\mathrm{\Lambda }}}(M_{k_0}^\tau )})1`$
where $`k_0=\mathrm{Max}\{,\mathrm{}+1,f^{6,}(\overline{\mathrm{\Lambda }},\mathrm{},t,c)`$ and $`\overline{\mathrm{\Lambda }}`$$`[k_0]`$ is defined from $`\overline{\mathrm{\Lambda }}`$ as in the main claim 4.6.
4.9 Claim: $`f^{6,}`$ beongs to $`E`$<sub>7</sub>
Proof: Straight.
4.10 Theorem: (1) The function $`f^1(\overline{\mathrm{\Lambda }},c)`$ is well defined, i.e. always get value, a natural number and is primitive recursive, in fact belongs to $`E`$$`_8.`$
(2) Similarly the function $`f^6(\overline{\mathrm{\Lambda }},\mathrm{},c).`$
(3) $`f^4`$ is primitive recursive, in fact belongs to $`E`$$`_9.`$
Proof: (1),(2) Let $`\tau =\tau [\overline{\mathrm{\Lambda }}].`$ The proof follows by induction, the main induction is on $`t=`$ arity$`(\tau _{\overline{\mathrm{\Lambda }}})`$ (or, if you prefer $`\mathrm{\Pi }_{F\tau [\overline{\mathrm{\Lambda }}]}(|\mathrm{\Lambda }_F|+1)).`$
CASE 0: arity$`(\tau )=1`$
This is Hales-Jewett theorem (on a bound see \[Sh:329\] or \[GRS80\])
CASE 1: arity$`(\tau )>1`$
Let $`\tau ^{},\overline{\mathrm{\Lambda }}^{}`$ be as in claim 4.4, so arity$`(\tau ^{})`$ $``$ arity$`(\tau )/2`$ and $`|\tau ^{}||\tau |\times 2^{\mathrm{arity}(\tau )}.`$
Let $`\mathrm{}^{}=f^1(\overline{\mathrm{\Lambda }}^{},c)`$ so (by 4.4) clearly $`f_{\overline{\mathrm{\Lambda }}}^6(\mathrm{}^{},c)\mathrm{}^{}`$ hence (by Definition 4.8) clearly $`f^{6,}(\overline{\mathrm{\Lambda }},0,\mathrm{}^{},c)=\mathrm{}^{}=f^1(\overline{\mathrm{\Lambda }}^{},c)`$ ; together we get $`f_{\overline{\mathrm{\Lambda }}}^6(\mathrm{}^{},c)f^{6,}(\overline{\mathrm{\Lambda }},0,\mathrm{}^{},c)`$. Hence (by 4.6 $`+`$ Definition 4.8, we shall prove by induction on $`\mathrm{}\mathrm{}^{})`$ that $`f_{\overline{\mathrm{\Lambda }}}^6(\mathrm{}^{}\mathrm{},c)f^{6,}(\overline{\mathrm{\Lambda }},\mathrm{},\mathrm{}^{},c)`$; for $`\mathrm{}=0`$ this holds by the previous sentnece; for the induction step, i.e. the proof for $`\mathrm{}+1`$ we apply Theorem 4.6 with $`\mathrm{}^{}\mathrm{},\mathrm{}^{}(\mathrm{}+1)`$ here standing for $`\mathrm{}+1,\mathrm{}`$ there and letting $`k_0=\mathrm{Max}\{\mathrm{}^{}\mathrm{},f_{\overline{\mathrm{\Lambda }}}^6(\mathrm{}^{},c)\}`$ and $`\tau ^{},\overline{\mathrm{\Lambda }}^{},c^{}`$ defined as there, and we get that $`f_{\overline{\mathrm{\Lambda }}}^6(\mathrm{}^{}(\mathrm{}+1),c)k_0+f_{\overline{\mathrm{\Lambda }}^{}}^7(P_{\overline{\mathrm{\Lambda }}^{}},1,c^{})1\mathrm{Max}\{\mathrm{}^{}\mathrm{},f_{\overline{\mathrm{\Lambda }}}^6(\mathrm{}^{},c)\}+f_{\overline{\mathrm{\Lambda }}^{}}^7(P_{\overline{\mathrm{\Lambda }}^{}},1,c^{})1`$
but the last expression is exactly $`f^{6,}(\overline{\mathrm{\Lambda }},\mathrm{}+1,\mathrm{}^{},c)`$
So (using $`\mathrm{}=\mathrm{}^{})`$ clearly $`f_{\overline{\mathrm{\Lambda }}}^6(0,c)f^{6,}(\overline{\mathrm{\Lambda }},\mathrm{}^{},\mathrm{}^{},c).`$
Now
$`f^1(\overline{\mathrm{\Lambda }},c)=f_{\overline{\mathrm{\Lambda }}}^6(0,c)`$ $`f^{6,}(\overline{\mathrm{\Lambda }},\mathrm{}^{},\mathrm{}^{},c)`$ $`f^{6,}(\overline{\mathrm{\Lambda }},f^1(\overline{\mathrm{\Lambda }}^{},c),\phi ^1(\overline{\mathrm{\Lambda }}^{},c),c).`$
As $`f^{6,}`$ is from $`E_7`$ by 3.14, this clearly give the desired conclusion.
(3) Should be clear from the proof of 3.5 and the previous parts.
$`\text{}_{4.10}`$
REFERENCE
\[BL96\] V. Bergelson and A. Leibman, Polynomial extensions of van der Waerder’$`s`$ and Szemeredi theorems, JAMS 9(1996)725-753
\[BL 9x\] V. Bergelson and A. Leibman, Set polynomial and polynomial extensions of the Hales Jewett theorem, to appear
\[Fu81\] H. Furstenberg, Recurrence in Ergodic Theory and Combinatorial Number Theory, Princeton University Press 1981
\[GR71\] R.L. Graham, B.L. Rothschild, Ramsey’s theorem for $`n`$parameter sets, TAMS 159(1971)257-292
\[GRS80\] R.L. Graham, B.L. Rothschild and H.J. Spencer, Ramsey Theory Wiley-Interscience Ser. in Discrete Math. New York 1980
\[Ro84\] H.E. Rose, Subrecursion: functions and heirarchies, Oxford Logic Guide 9, Oxford University Press, Oxford 1984
\[Sh:329\] Shelah, Saharon, Primitive recursive bounds for van der Waerden numbers, Journal of the American Mathematical Society, 1 (1988) 683–697 |
warning/0003/hep-th0003291.html | ar5iv | text | # I INTRODUCTION
## I INTRODUCTION
In its original formulation, the AdS/CFT correspondence allows us to study Yang-Mills conformal field theories on $`S^4`$ in terms of an $`S^5`$ compactification of string theory on the hyperbolic space $`H^5`$. \[We use the Euclidean approach throughout; $`S^n`$ is the $`n`$-sphere with the standard Riemmanian metric and conformal structure, and $`H^n`$ is the open ball $`B^n`$ endowed with the metric of constant sectional curvature equal to -1.\] The transition from the gravitational theory in the “bulk” to a non-gravitational theory is effected by a geometric scheme in which $`S^4`$ appears as “infinity” for $`H^5`$. The obvious way $`{\displaystyle \frac{}{}}`$ but, we shall argue, not the only way $`{\displaystyle \frac{}{}}`$ to formulate this concretely is to regard $`S^4`$ as the boundary of the closed ball $`\overline{B}^5`$, after the manner of Penrose compactifications in general relativity.
It is generally agreed that AdS/CFT reflects some very deep property of these conformal field theories. If this is so, then surely the correspondence must work for manifolds other than $`S^4`$. This idea has led to some remarkable insights. For example, replacing $`S^4`$ by $`S^1\times S^3`$, we immediately note that there are at least two candidates for the bulk, $`B^2\times S^3`$ and $`S^1\times B^4`$. The conformal field theory partition function is then naturally defined by a sum of contributions from $`B^2\times S^3`$ and $`S^1\times B^4`$, which leads to a large $`N`$ phase transition related to black hole thermodynamics . Thus AdS/CFT can indeed be pushed beyond $`S^4`$. The obvious question now is : how far can it be pushed ? Are there manifolds for which AdS/CFT definitely does not work? If so, can one classify the manifolds for which it does? In short, what are the “outer limits” for AdS/CFT?
The ideal answer to the first of these questions would be : it works for every compact conformal manifold of dimension less than 11 on which a physically sensible CFT can be defined. Our objective in this work is not, of course, to “prove” such a grandiose assertion; our purpose, instead, is to formulate “physically sensible” in a precise way, and to answer the most obvious objection to this claim : how can it be true of a CFT defined on a compact manifold which is not the boundary of any manifold-with-boundary?
This question was already raised in , where it was suggested that an answer would involve introducing “branes or stringy impurities of some kind” into the bulk. We will argue here that this is indeed precisely the correct answer, and that AdS/CFT does have a chance of working even in this extreme case. Clearly, this will involve a more general formulation of the relationship between the bulk and “infinity” than is usually considered. We will see that recent important results on the geometry of “infinity” , can be interpreted physically as results on the nature of the matter content of the bulk.
We begin in section II by discussing a criterion, hinted at in and stated explicitly in , for a Yang-Mills CFT to be physically reasonable. This leads us, with the aid of the Kazdan-Warner classification , to a precise proposal as to the kinds of manifolds for which AdS/CFT should be expected to work. Among these are some manifolds which cannot be represented as boundaries, and so cannot be Penrose conformal infinity for any bulk. This leads us, in section III, to generalise the concept of conformal infinity in such a way that “infinity” is a hypersurface in a compact manifold, instead of a boundary of a manifold-with-boundary. Finally, in section IV, we use geometric techniques to prove that the bulk must, if “infinity” is not a boundary, contain some kind of “brane or stringy impurity”, and to investigate the nature of these “impurities”.
## II THE STABILITY CONDITION FOR THE CFT
It is pointed out in that the convergence of the path integral for the Yang-Mills CFT is non-trivial $`{\displaystyle \frac{}{}}`$ it depends on the geometry of the underlying manifold. Convergence is not a problem for $`S^4`$ with its standard conformal structure, but, as one moves away from this simplest case, one expects the good behaviour of the CFT to become increasingly questionable. We therefore need to know which properties of $`S^4`$ are essential and which are not. (Throughout this section, all manifolds have dimension $`3`$.)
As a conformal manifold (that is, a manifold on which one is given an equivalence class of conformally related metrics, as is appropriate for the study of conformal field theories), $`S^4`$ is distinguished by being conformally flat, and also by having a conformal structure represented by an Einstein metric. We wish to argue that neither of these properties is essential. Against conformal flatness we adduce the following evidence (apart from the fact that it is obviously an extremely severe restriction and would eliminate too many interesting manifolds). We shall see later that there is a compact manifold of the form $`^4/\mathrm{\Gamma }`$, where $`\mathrm{\Gamma }`$ is of course infinite but discrete, which is a boundary, but on which it is impossible to define a physically reasonable Yang-Mills CFT. This is true for every conformal structure on $`^4/\mathrm{\Gamma }`$. And yet there is a conformally flat conformal structure on this manifold. Evidently the conformal flatness condition has little or no physical significance. Against the Einstein condition we have still stronger evidence, as follows. For any four-dimensional Einstein manifold $`M`$, the Euler characteristic is given by (. page 161)
$$\chi (M)=\frac{1}{8\pi ^2}_M(U^2+W^2)𝑑V,$$
(1)
where $`U`$ is the irreducible component of the curvature tensor determined by the scalar curvature, and $`W`$ is the Weyl tensor. Since $`\chi (S^1\times S^3)=0`$, an Einstein metric on $`S^1\times S^3`$ would necessarily be flat, an impossibility since $`S^1\times S^3`$ is not covered by $`^4`$. So there is no Einstein metric of any kind on $`S^1\times S^3`$, including non-product metrics. (This is just the simplest of several non-existence results of this kind; see .) Thus the Einstein condition would rule out $`S^1\times S^3`$ $`{\displaystyle \frac{}{}}`$ which, as we saw in the Introduction, is one of the most important examples of a manifold to which AdS/CFT can be extended.
Having abandoned constraints on the components of the curvature tensor determined by the Weyl and Ricci tensors, we turn naturally to its last remaining component, $`U`$. The proper criterion , is as follows. Let $`N^n`$ be a compact manifold with a conformal structure $`[g^N]`$. The conformal Laplacian, obviously the physically relevant operator here, is
$$L_{g^N}=\mathrm{\Delta }_{g^N}+\frac{n2}{4(n1)}R(g^N),$$
(2)
where $`\mathrm{\Delta }_{g^N}`$ is the usual Laplacian, $`n=dim(N^n)`$, and $`R(g^N)`$ is the scalar curvature. It is well known that $`L_{g^N}`$ is an elliptic operator with discrete real spectrum bounded from below. Let $`\mu _1(g^N)`$ be the first eigenvalue. Then the Yang-Mills CFT will be stable if $`\mu _1(g^N)>0`$, and sometimes when $`\mu _1(g^N)=0`$, but it is unstable if $`\mu _1(g^N)<0`$. This criterion can be stated more usefully as follows. By Schoen’s solution of the Yamabe problem, $`[g^N]`$ contains a Yamabe metric $`g_Y^N`$ which, by definition, is such that $`R(g_Y^N)`$ is constant on $`N^n`$. This constant is given by
$$R(g_Y^N)=\frac{4(n1)}{n2}\mu _1(g_Y^N).$$
(3)
The stability condition can therefore be expressed in terms of the sign of the scalar curvature. For example, $`S^4`$ with its usual metric has constant positive scalar curvature, so one can construct a stable CFT using this metric. However, $`S^4`$ also has another metric which $`{\displaystyle \frac{}{}}`$ surprisingly $`{\displaystyle \frac{}{}}`$ has constant negative scalar curvature. The CFT will of course be unstable if this metric is used.
Now in fact Kazdan and Warner have given a classification of manifolds according to the behaviour of the scalar curvature. The following theorem is basic.
###### Theorem 1 (Kazdan-Warner).
Every compact connected manifold $`M^n`$, $`n3`$, falls into precisely one of the following three classes.
1. $`P`$ : Every smooth function on $`M^n`$ is the scalar curvature of some metric on $`M^n`$.
2. $`Z`$ : A smooth function on $`M^n`$ is the scalar curvature of some metric if and only if it is either negative at some point or it is identically zero.
3. $`N`$ : A smooth function is the scalar curvature of some metric on $`M^n`$ if and only if it is negative at some point.
To see how to use this theorem, observe for example that since $`S^4`$ admits a metric of positive scalar curvature, it cannot be in $`Z`$ or $`N`$; hence it is in $`P`$; hence every smooth function on $`S^4`$ is the scalar curvature of some metric; hence indeed there is a metric on $`S^4`$ with constant negative scalar curvature, another with identically zero scalar curvature, and so on. Again, the torus $`T^4`$ accepts a flat metric, so it is not in $`N`$; since it is “enlargeable” (, page 306), it admits no metric of positive scalar curvature, so it cannot be in $`P`$; hence it is in $`Z`$. Finally, consider a compact 4-dimensional manifold which accepts a metric of constant negative sectional curvature. Such a manifold has the structure $`^4/\mathrm{\Gamma }`$, for some discrete freely acting group $`\mathrm{\Gamma }`$ with no subgroup of the form $``$. ($`^4/\mathrm{\Gamma }`$ is said to be homotopically atoroidal.) This manifold, too, is enlargeable, so it admits no metric of positive scalar curvature. Nor, however, does it admit a metric of zero scalar curvature, for such a metric on an enlargeable manifold must be flat, but $`^4/\mathrm{\Gamma }`$ is not covered by $`T^4`$. Hence it is in $`N`$. This is the example mentioned earlier : its metric of constant negative sectional curvature is conformally flat (and so its Hirzebruch signature is zero; the signature being, in four dimensions, an isomorphism ( page 92) from the oriented cobordism group to $``$, the manifold is a boundary) and yet every conformal class on $`^4/\mathrm{\Gamma }`$ is represented by a Yamabe metric of negative scalar curvature, so conformal flatness certainly does not ensure satisfactory physical behaviour.
It is important to realise that the Kazdan-Warner classification is a classification of manifolds $`{\displaystyle \frac{}{}}`$ that is, the class to which a space belongs depends only on its topology and differentiable structure. (For this second point, note that $`S^9`$ with its usual differentiable structure belongs to $`P`$, but with a certain exotic differentiable structure , it belongs to $`N`$. The stability of these conformal field theories in nine (and ten) dimensions can therefore depend on the choice of differentiable structure.) Clearly it is not possible to define a stable CFT on a manifold in class $`N`$ no matter which conformal structure we use. The instability is not a geometric phenomenon for manifolds in class $`N`$; it is due to their differential topology. For manifolds in class $`P`$, by contrast, there is always a metric which makes the CFT stable, but there is also another metric which makes it unstable. (Notice that the Kazdan-Warner theorem implies that every compact manifold of dimension greater than two admits a metric of constant negative scalar curvature.) Hence, once it is known that a manifold is in $`P`$, the question of stability becomes a geometric question. (The reader should be aware that deciding the Kazdan-Warner class of a manifold can be non-trivial : for example, simply connected six-dimensional Calabi-Yau manifolds are all in $`P`$, not $`Z`$.)
What of class $`Z`$? In all known examples, these manifolds behave like manifolds in class $`P`$. For example, $`T^4`$ with its flat metric admits a CFT which is perfectly well-behaved . The same appears to be true of manifolds and orbifolds of the form $`T^4/\mathrm{\Delta }`$, where $`\mathrm{\Delta }`$ is a finite group of isometries of $`T^4`$. (All such manifolds/orbifolds are in $`Z`$, since a metric of positive scalar curvature on $`T^4/\mathrm{\Delta }`$ would pull back to a metric of positive scalar curvature on $`T^4`$.) Since $`K3`$ can be regarded as a resolution of a $`T^4`$ orbifold, we expect that the same is true in this case also. (That $`K3`$ is in $`Z`$ follows from the theorem of Lichnerowicz; see .) It is reasonable to conjecture that Kazdan-Warner class $`Z`$ is indeed like class $`P`$ : that is, for each manifold in $`Z`$, there is a conformal structure such that the CFT is stable, while there are of course other conformal structures such that the CFT is unstable.
Throughout this discussion, we have not assumed that “infinity”, the compact manifold on which our CFT is defined, is connected. There is of course a well-defined stable CFT on $`T^4+S^4`$, the disjoint union of $`T^4`$ and $`S^4`$, though the two separate conformal field theories are very different and are decoupled. Now there exists a non-compact five-dimensional spin manifold $`M^5`$ such that $`T^4+S^4`$ is the boundary of a connected manifold-with-boundary $`\overline{M}^5`$ having $`M^5`$ as interior. How can AdS/CFT work in this case? How can one interior be “dual” to two different, decoupled conformal field theories? This question was raised in . The simplest response to this paradox is just to declare that AdS/CFT should not be expected to work for disconnected boundaries (with an exception to be discussed below). In practical terms, this is of course a very minor limitation, since $`T^4`$ and $`S^4`$ are each boundaries, of $`\overline{B}^2\times T^3`$ and $`\overline{B}^5`$ respectively, so the separate conformal field theories can be studied by two applications of ordinary AdS/CFT.
We can now state what we hope to be the full range of AdS/CFT. The claim is that AdS/CFT should work for every compact connected infinity manifold (or perhaps orbifold, etc) of suitable dimension which does not belong to Kazdan-Warner class $`N`$. By “suitable dimension” we mean simply that the bulk should be of dimension 10 or 11, or lower if there is a compactification.
So far, we have concentrated on the conditions to be satisfied by the differential topology and geometry of “infinity”, without concerning ourselves with the details of the physical fields there. This is justifiable, in that the theory at “infinity” is avowedly non-gravitational. For precisely this reason, care should be exercised before imposing geometric conditions on the bulk. Often one assumes that the bulk is an Einstein manifold of Ricci curvature $`n`$, but, while this is legitimate if “infinity” is $`S^n`$ and the conformal structure is not too far from the standard one, it should only be regarded as an approximation in other cases. In more general investigations the metric is only required to be asymptotically Einstein; no particular fall-off rate is assumed, and the metric is certainly not required to be complete $`{\displaystyle \frac{}{}}`$ indeed, in many applications it is definitely incomplete. Our attitude is that conditions on the metric in a gravitational theory should be dictated by the theory itself, not imposed externally. In short, we shall impose no requirements on the bulk metric.
As is observed already in , there is an obvious, strong objection to the idea that AdS/CFT works whenever the CFT is stable (that is, for connected compact manifolds of Kazdan-Warner classes $`P`$ and $`Z`$) : many manifolds are not boundaries. For example, no compact connected four-dimensional manifold is a boundary if its signature is not zero . We now deal with this objection.
## III INFINITY IS JUST ANOTHER BRANE
Let $`\widehat{M}^{n+1}`$ be a compact $`(n+1)`$-dimensional manifold containing a smooth compact boundaryless hypersurface $`N^n`$. Fix a Riemannian metric $`g^M`$ on $`M^{n+1}=\widehat{M}^{n+1}N^n`$ and assume that there exists an “infinity function” f on $`\widehat{M}^{n+1}`$. This is a smooth function which is positive on $`M^{n+1}`$ and vanishes to first order on $`N^n`$, such that $`f^2g^M`$ extends continuously to $`N^n`$. If such a function exists, then $`N^n`$ is “infinitely far from” points in $`M^{n+1}`$, and $`g^M`$ induces a conformal structure (not a Riemannian structure) on $`N^n`$. (Note that $`M^{n+1}`$ need not be complete.) In such a case, we shall say that $`N^n`$ is an infinity hypersurface for $`\widehat{M}^{n+1}`$ with respect to $`g^M`$.
This definition is of course motivated by the formal definition of a Penrose conformal boundary, which is the more usual arena for AdS/CFT. Indeed, any compact manifold-with-boundary with the boundary “at infinity” can be re-interpreted in the above way : simply take two copies, and (adjusting the boundary orientation suitably) identify them along the boundary. The result will be a compact manifold with an infinity hypersurface at the former location of the boundary. One can also do this by beginning with distinct manifolds-with-boundary having diffeomorphic boundaries. For example, $`\overline{B}^2\times S^3`$ can be joined to $`S^1\times \overline{B}^4`$ along their common $`S^1\times S^3`$ boundary, and so the process of summing over distinct interiors can be implemented in a concrete way. Heuristically, there may well be advantages in dethroning “infinity” from its privileged position at the boundary, and thinking of it as “just another brane”, one which happens to be infinitely far away; and certainly compact manifolds are preferable to manifolds-with-boundaries.
Clearly, AdS/CFT can be formulated in this language. Notice, however, that the infinity hypersurfaces obtained in this way have a special property : $`N^n`$ separates $`M^{n+1}`$ into disconnected pieces. By considering infinity hypersurfaces which do not have this effect, we obtain something new. The following family of examples is particularly enlightening.
Let $`P^n`$ be connected, compact, $`n`$-dimensional manifold with a Riemannian metric $`g^P=g_{ij}^Pdx^idx^j`$ and Ricci curvature $`Ric(g^P)=R_{ij}^Pdx^idx^j`$. Let $`\widehat{M}^{n+1}=S^1\times P^n`$ with $`S^1`$ parametrised by $`\theta `$ running from $`0`$ to $`2\pi `$, and let $`M^{n+1}`$ be obtained from $`\widehat{M}^{n+1}`$ by deleting all points with $`\theta =0`$. Define a metric $`g^M`$ on $`M^{n+1}`$ by
$$g^M=\mathrm{cosec}^2(\frac{\theta }{2})[\frac{1}{4}d\theta d\theta +g_{ij}^Pdx^idx^j].$$
(4)
Here the function $`f`$ is $`\mathrm{sin}({\displaystyle \frac{\theta }{2}})`$, which is positive in $`(0,2\pi )`$ and vanishes to first order at $`\theta =0`$, where there is a single copy of $`P^n`$. The infinity hypersurface does not separate $`M^{n+1}`$ into disconnected components. The Ricci tensor of this metric is, in an obvious notation,
$`(R^M)_\theta ^\theta `$ $`=n,`$ (5)
$`(R^M)_j^i`$ $`=n\delta _j^i+\mathrm{sin}^2({\displaystyle \frac{\theta }{2}})[(R^P)_j^i+(n1)\delta _j^i],`$ (6)
all other components being zero. We have expressed the Ricci tensor in $`(1,1)`$ form in order to be able to discuss invariant quantities, namely the eigenvalue functions of the Ricci curvature. (The $`(0,2)`$ components diverge near $`\theta =0`$, but this is merely a coordinate effect.) As $`P^n`$ is compact, the eigenvalue functions of $`Ric(g^P)`$ are bounded, and hence so are those of $`Ric(g^M)`$. For example, if $`P^n`$ is Ricci-flat, the eigenvalue functions of $`Ric(g^M)`$ are bounded above by $`1`$ and below by their asymptotic value, $`n`$.
The structure of this space (Figure 1) is clear : the infinity hypersurface is really one, connected copy of $`P^n`$, and the bulk $`M^{n+1}`$ is just an open submanifold. However, we can (“perversely”) re-interpret the structure of $`M^{n+1}`$as follows. Instead of the compact manifold $`S^1\times P^n`$, let us consider the compact manifold-with-boundary $`[0,2\pi ]\times P^n`$ obtained by artificially distinguishing $`\theta =2\pi `$ from $`\theta =0`$. Now formula (4) still defines a Riemannian metric on the interior, $`(0,2\pi )\times P^n`$. If we insist on setting asunder what belongs together, we can regard $`[0,2\pi ]\times P^n`$ as the Penrose compactification of $`M^{n+1}`$; the boundary now consists of two copies of $`P^n`$, one each at $`\theta =0`$ and $`\theta =2\pi `$. We can lend colour to this imposture by changing coordinates. Let $`x`$ be defined by
$$\mathrm{cosh}(x)=\mathrm{cosec}(\frac{\theta }{2})$$
with $`x0`$ for $`\theta \pi `$, and $`x0`$ for $`\theta \pi `$. A short calculation reveals
$$g^M=dxdx+\mathrm{cosh}^2(x)g_{ij}^Pdx^idx^j,$$
(7)
again suggesting two distinct “infinities”, one at “$`x=\mathrm{}`$”, the other at “$`x=+\mathrm{}`$”. (See page 268 and note that $`g^M`$ is Einstein if $`Ric(g^P)=(n1)g^P`$.)
In short, $`M^{n+1}`$ does not “know” whether it is an open submanifold of a compact manifold, or the interior of a compact manifold-with-boundary. However, the difference, from a physical point of view, is substantial. For the manifold-with-boundary interpretation leads us to a disconnected boundary, and so to the Witten-Yau paradox discussed previously. The “infinity hypersurface” interpretation is both more natural and more physically acceptable.
The metric (4) has no particular physical significance; it was chosen to make the above point. In general, we can take a compact manifold $`\widehat{M}^{n+1}`$ with an infinity hypersurface $`N^n`$, and “split” it along $`N^n`$ to obtain a compact manifold-with-boundary having boundary components $`N_1^n`$, $`N_2^n`$, and so on, with each component diffeomorphic to $`N^n`$. We can then deform a conformal structure on $`N_2^n`$ through physically acceptable conformal structures to obtain a boundary component which is of course still diffeomorphic, but perhaps no longer isometric, to $`N^n`$. It would clearly be absurd to disallow the original manifold on the grounds that it can so artificially be brought into conflict with the Witten-Yau paradox. On the contrary, when we are presented with a metric like (7), representing the interior of a manifold-with-boundary with essentially identical boundary components, our first move should be to make the appropriate identifications and return to the more natural infinity hypersurface interpretation (corresponding to (4)). Of course, if the boundary components are not mutually diffeomorphic (or if they are, but their conformal structures can only be deformed to each other by passing through “unstable” conformal structures) then the paradox will lead to genuine difficulties.
These remarks bring us to our main application.
## IV MANIFOLDS WHICH ARE NOT BOUNDARIES
Suppose that we wish to study a Yang-Mills CFT on some compact manifold $`N^n`$ (such as a four-manifold with non-zero signature) which simply cannot be expressed as the boundary of some manifold-with-boundary. We now have a strategy : represent two copies of $`N^n`$ as a disconnected conformal boundary, and identify them to realise it as an infinity hypersurface in some compact manifold. The CFT should then be dual to some string theory in this “bulk”. As we shall see, this can always be done; the only point at issue is whether the bulk admits a physically reasonable geometry.
When we say that a certain compact manifold $`N^n`$ is not a boundary, we mean “not a boundary by itself”. It is always possible to find a compact manifold $`Q^n`$ such that $`N^n+(Q^n)`$ is a boundary of a connected $`(n+1)`$-dimensional space. (We consider only oriented manifolds; $`Q^n`$ results from reversing orientation.) One says that $`N^n`$ and $`Q^n`$ are cobordant . In general, this will not solve our problem, because it will lead to the Witten-Yau paradox. However, if we choose $`Q^n=N^n`$, then we can represent $`N^n`$ as a connected infinity hypersurface by performing a topological identification. Unfortunately, this seems a somewhat arbitrary proceeding, because there are many other choices for $`Q^n`$ $`{\displaystyle \frac{}{}}`$ cobordant manifolds can be very different. We should ask whether the choice $`Q^n=N^n`$ can be motivated on physical grounds.
Before discussing this, let us consider the concept of spin cobordism. Let $`\overline{M}^{n+1}`$ be a compact manifold-with-boundary with interior $`M^{n+1}`$. If $`M^{n+1}`$ is a spin manifold, a given spin structure induces a spin structure on the boundary in a canonical way (, page 90). Now spin manifolds $`N^n`$ and $`Q^n`$ are spin cobordant if there exists a compact manifold-with-boundary $`\overline{M}^{n+1}`$ having $`N^n+(Q^n)`$ as boundary, and having an interior with a spin structure that induces the given spin structures on $`N^n`$ and $`Q^n`$. Clearly spin cobordism is the appropriate cobordism theory for physical applications. The spin cobordism equivalence classes in a given dimension form an abelian group, $`\mathrm{\Omega }_n^{spin}`$. In low dimensions they are (, page 92):
$`\mathrm{\Omega }_1^{spin}`$ $`=_2,`$ $`\mathrm{\Omega }_2^{spin}`$ $`=_2,`$ (8)
$`\mathrm{\Omega }_3^{spin}`$ $`=0,`$ $`\mathrm{\Omega }_4^{spin}`$ $`=,`$
$`\mathrm{\Omega }_5^{spin}`$ $`=0,`$ $`\mathrm{\Omega }_6^{spin}`$ $`=0,`$
$`\mathrm{\Omega }_7^{spin}`$ $`=0,`$ $`\mathrm{\Omega }_8^{spin}`$ $`=.`$
Let us apply these results. The two-dimensional case is instructive, and we begin with it. The group $`\mathrm{\Omega }_2^{spin}`$ is generated by the torus $`T^2`$. This space has several spin structures, and one of them is not induced by any spin structure on the interior, $`B^2\times S^1`$; so $`T^2`$ with this spin structure is not a spin boundary (though of course it is a boundary as an oriented surface.) Now let $`Q^2`$ be spin cobordant to $`T^2`$, and set up a conformal field theory on $`T^2+(Q^2)`$. As usual, we require stability, meaning that the scalar curvature $`{\displaystyle \frac{}{}}`$ the Gaussian curvature $`{\displaystyle \frac{}{}}`$ of $`Q^2`$ must be positive or zero. But if it is positive, the Gauss-Bonnet theorem identifies $`Q^2`$ as $`S^2`$, which is not spin cobordant to $`T^2`$ with this spin structure. Thus $`Q^2`$ is forced, by the stability condition, to be just another copy of $`T^2`$. We are thus invited to perform the appropriate identification and to study the CFT on $`T^2`$ by regarding it as an infinity hypersurface in a compact three-dimensional manifold.
This example gives us hope that the CFT stability condition might constrain the choice of $`Q^n`$ in higher dimensions, and in fact this is correct for the important case of $`n=4`$. Suppose that $`N^4`$ is a compact connected 4-manifold which admits a stable CFT but which is not a spin boundary. It is known that $`\mathrm{\Omega }_4^{spin}`$ is generated by $`K3`$, so $`N^4`$ is spin cobordant to $`m`$ disjoint copies of $`K3`$, where $`m`$ could be negative but cannot, by hypothesis, be zero. Now the $`\widehat{A}`$ genus is a spin corbordism invariant (. pages 92, 298) and since $`\widehat{A}(K3)=2`$, we see that $`\widehat{A}(N^4)0`$. By the theorem of Lichnerowicz (, page 161) there is no metric of positive scalar curvature on $`N^4`$. Suppose that $`Q^4`$ is spin cobordant to $`N^4`$, so that $`N^4+(Q^4)`$ is a spin boundary; then all of the above applies equally to $`Q^4`$.
The stability of the CFT on $`N^4+(Q^4)`$ now demands that both $`N^4`$ and $`Q^4`$ admit metrics of zero scalar curvature. Now since $`\widehat{A}(N^4)`$ and $`\widehat{A}(Q^4)`$ are non-zero, both admit non-trivial harmonic spinors; but spinors on a compact spin manifold which are harmonic with respect to a scalar-flat metric must in fact be parallel (, page 161). The existence of non-trivial parallel spinors forces the spin holonomy group of a manifold to be special; it also forces the Ricci tensor to vanish . Now the spin holonomy groups of (not necessarily simply connected) compact Ricci-flat Riemannian spin manifolds have been classified and so we can conclude that the metrics on $`N^4`$ and $`Q^4`$ are either flat (which would contradict the non-vanishing of $`\widehat{A}(N^4)`$ and $`\widehat{A}(Q^4)`$) or of spin holonomy precisely $`SU(2)`$. (By “precisely”, we mean that the full, global holonomy group (both linear and spin) is $`SU(2)`$, not just the identity component. There are four-manifolds with disconnected linear holonomy groups having $`SU(2)`$ as identity component, but these are not spin manifolds.) But every compact four-manifold of holonomy precisely $`SU(2)`$ is diffeomorphic to a finite number of copies of $`K3`$ (see , page 365). Since $`\widehat{A}(K3)=2`$, we see that $`m_1`$ copies of $`K3`$ are spin cobordant to $`M_2`$ copies only if $`m_1=m_2`$, and so, since $`N^4`$ is connected, $`Q^4`$ must be diffeomorphic to it. Finally, the moduli space of Einstein metrics on $`K3`$ is connected (, page 366), and since the total scalar curvature is locally constant as a function on moduli space (, page 352), the $`SU(2)`$ metric on $`Q^4`$ can be deformed through scalar-flat metrics so that $`Q^4`$ is isometric to $`N^4`$. Regarding $`N^4+(Q^4)`$ as the boundary of a compact manifold-with-boundary $`\overline{M}^5`$, we can search for an appropriate metric $`g^M`$ on $`M^5`$ such that $`N^4+(Q^4)`$ is the conformal boundary (see below). Performing the identification as usual, we now have $`N^4`$ as an infinity hypersurface in a five-dimensional compact manifold, and we can begin to explore AdS/CFT for $`N^4`$, despite the fact that it is not a boundary.
In dimension 8, the situation is much more complex. The group $`\mathrm{\Omega }_8^{spin}=`$ is generated by the quaternionic projective space $`P^2`$ and by any Joyce manifold $`J_8`$ of holonomy $`Spin(7)`$ (see ). As $`P^2`$ is the symplectic homogeneous space $`{\displaystyle \frac{Sp(3)}{Sp(1)\times Sp(2)}}`$, it admits a metric of positive scalar curvature, and so, therefore, does any simply connected eight-manifold which is spin cobordant to it or to any finite number of copies of it (, page 299). On the other hand, there are scalar-flat metrics on the various candidates for $`J_8`$ and on many topologically distinct manifolds spin cobordant to $`J_8`$ (such as manifolds of linear holonomy $`SU(4)_2`$) and to multiple copies of it (such as Calabi-Yau and hyperKähler manifolds). Thus if $`N^8`$ is a non-boundary eight-manifold, there are many candidates for $`Q^8`$ and there is no good physical justification for selecting $`Q^8=N^8`$. Perhaps this indicates some kind of pathology afflicting AdS/CFT for conformal field theories on manifolds of dimension greater than seven.
## V THE BRANE IN THE BULK
Let $`\widehat{M}^5=S^1\times P^4`$, where $`P^4`$ is a non-boundary four-manifold with a scalar-flat metric. We saw above that in fact $`P^4`$ is Ricci-flat, so if we use equation (4) to define a metric $`g^M`$ on $`M^5`$, then the Ricci tensor of $`g^M`$ satisfies, by equation (5) and (6),
$`(R^M)_\theta ^\theta `$ $`=n,`$ (9)
$`(R^M)_j^i`$ $`=n\delta _j^i+(n1)\mathrm{sin}^2({\displaystyle \frac{\theta }{2}})\delta _j^i.`$ (10)
Thus $`M^5`$ is not an Einstein manifold except “near infinity”.
Now although we should be prepared ultimately to consider non-Einstein metrics on the bulk, we might prefer to begin with Einstein metrics and then to consider perturbations around them. One can of course do this for $`S^4`$ and for $`S^1\times S^3`$; but can one do it for infinities which are not boundaries? The following very remarkable theorem, which follows straightforwardly from results of Witten-Yau and Cai-Galloway , is relevant.
###### Theorem 2 (Witten-Yau-Cai-Galloway).
Let $`\widehat{M}^{n+1}`$ be a compact $`(n+1)`$-dimensional manifold admitting an infinity hypersurface $`N^n`$ with respect to a Riemannian metric $`g^M`$ on the complement (which we assume to be orientable). Suppose that $`N^n`$ corresponds to a non-trivial element of the homology group $`H_n(\widehat{M}^{n+1},)`$ and that the induced conformal structure on $`N^n`$ is represented by a Yamabe metric of positive or zero scalar curvature. Then the equation $`Ric(g^M)=ng^M`$ requires $`g^M`$ to be an incomplete metric.
This result means that, in the case at hand, any attempt to force the bulk to be Einstein everywhere will merely cause the metric to develop some kind of pathology. \[We require all metrics to be differentiable, so incompleteness includes failures of differentiability.\] In physical terms, we can think of $`Ric(g^M)=ng^M`$ as characterising the vacuum, and of the “pathology” as some kind of localised matter, such as a brane. The theorem then simply means that the presence of a non-boundary infinity hypersurface entails the existence of some such object in the bulk.
An example will be helpful. Let $`P^4`$ be a compact non-boundary four-manifold with a scalar-flat (hence, Ricci-flat) metric $`g^P`$. Let $`\widehat{M}^5=S^1\times P^4`$, with $`S^1`$ now parametrised by $`\theta `$ in $`(\pi ,\pi ])`$ (not 0 to $`2\pi `$ as before). Take $`M^5`$ to be $`(S^1\{0\})\times P^4`$, let $`\delta `$ satisfy $`0<\delta <\pi `$, and define a metric $`g_\delta ^M`$ by
$`g_\delta ^M`$ $`=\theta ^2(d\theta d\theta +g^P),`$ $`\theta (0,\delta )(\delta ,0),`$ (11)
$`=h(\theta )(d\theta d\theta +g^P),`$ elsewhere. (12)
Here $`h(\theta )`$ is a function which interpolates continuously and smoothly between the two “$`\theta ^2`$ regions”. Clearly this manifold has a single, connected infinity hypersurface at $`\theta =0`$. Just as for the metric given by equation (4), however, we can make $`\theta =0`$ seem disconnected by changing coordinates :
$$\theta =\{\begin{array}{cc}\pi e^x,\hfill & x0,0<\theta <\delta ;\hfill \\ \pi e^x,\hfill & x<0,\delta <\theta <0,\hfill \end{array}$$
for the metric becomes $`dxdx+\pi ^2e^{2|x|}g^P`$, with infinities apparently at $`x=\pm \mathrm{}`$. Notice the formal similarities to the Randall-Sundrum metric, which has $`e^{2|x|}`$ instead of $`e^{2|x|}`$, and where $`g^P`$ would be flat, not just Ricci-flat. In fact, however, a simple calculation shows that the “pseudo-Randall-Sundrum metric” $`dxdx+\pi ^2e^{2|x|}g^P`$ is an Einstein metric as long as $`g^P`$ is Ricci-flat $`{\displaystyle \frac{}{}}`$ it does not need $`g^P`$ to be flat (see , page 268). So we have
$$Ric(g_\delta ^M)=4g_\delta ^M,\theta (0,\delta )(\delta ,0).$$
(13)
That is, $`g_\delta ^M`$ is precisely Einstein outside the immediate neighbourhood of $`\theta =\pi `$.
The Randall-Sundrum metric has a pathology at $`x=0`$ because of the absolute value function $`|x|`$ in $`e^{2|x|}`$, due to the presence of a brane. Clearly the pseudo-Randall-Sundrum metric has the same property, and the WYCG theorem asserts that setting $`\delta =0`$ will cause $`g_\delta ^M`$ to become incomplete at $`\delta =\pi `$. That is indeed the case, since the connection coefficient $`\mathrm{\Gamma }_{\theta \theta }^\theta `$, for example, is discontinuous there. However, this shows that the pathologies required by the theorem can be rather mild $`{\displaystyle \frac{}{}}`$ one should not think of them as singularities, but rather as branes. If $`\delta `$ is not zero but extremely small, then equation (13) is satisfied everywhere in $`M^5`$ except in an extremely thin slice. Inside that slice, the metric is given by (12), and, turning the WYCG theorem around, we deduce that (13) is certainly not satisfied everywhere in the slice. Again, this is evidence that some kind of localised matter is present.
If we now relax the condition that $`g^M`$ be an Einstein metric, we can hope to learn more about the nature of the matter whose existence is necessitated by this interpretation of “infinity”. Set
$$Ric(g^M)=ng^M+S(g^M),$$
(14)
so that the tensor $`S`$ measures the failure of $`g^M`$ to be Einstein. In the Lorentzian case, one could try to impose sign conditions on $`S`$ by means of the strong energy condition , but that is not appropriate in the Euclidean regime. \[The Euclidean counterpart of a Lorentzian metric which satisfies the strong energy condition need not obey any sign condition; consider for example the Euclidean version of a FRW dust metric, where the Ricci curvature is unbounded both above and below.\] Nevertheless, even in the Euclidean case, $`S`$ does have non-negative eigenvalue functions for certain kinds of matter, such as scalar fields with positive potentials. Thus, the sign of $`S`$ can give general information on the kind of matter which causes a given metric to be non-Einstein. The following theorem is therefore relevant. \[The proof is again a straightforward consequence of results of Cai and Galloway \].
###### Theorem 3 (Cai-Galloway).
Let all conditions be as in the WYCG theorem, except that $`Ric(g^M)=ng^M`$ is weakened to the condition that all of the eigenvalues of $`S(g^M)`$ decay no more slowly than inverse-quartically (see for details) towards infinity. Then either $`g^M`$ is incomplete or some eigenvalue function of $`S(g^M)`$ takes a value strictly less than $`0`$.
\[Notice that the metric (4) escapes the conclusions of this theorem because its Ricci tensor does not tend to $`ng^M`$ quickly enough.\] Thus, for example, if the metric given by (11) and (12) is forced to be complete by an appropriate choice of $`h(\theta )`$, then $`h(\theta )`$ must be such that some eigenvalue of the Ricci tensor falls below $`4`$, which is the value outside the slice.
The above theorem leads us to ask: what kind of matter has an $`S`$-tensor with negative eigenvalues? For example, can p-branes give rise to fields with such an $`S`$-tensor? The answer is yes, as we now show. Consider the following action.
$$[R+\lambda \frac{1}{2}(\varphi )^2\frac{1}{2(p+2)!}e^{a\varphi }F^2]𝑑vol,$$
where $`R`$ is the scalar curvature, $`\lambda `$ is a constant to be chosen, $`\varphi `$ is a scalar field, $`a`$ is a constant, and $`F`$ is a $`(p+2)`$-form derived from a potential in the usual way. Now p-brane solutions of the corresponding Euler-Lagrange equations are studied in . A particularly interesting sub-class of non-singular solutions is obtained when $`\varphi `$ and $`a`$ can consistently be set equal to zero, so that only $`F`$ contributes to $`S`$. This is possible \[see , end of section 3.2\] for
$$(n,p)=(4,1),(5,1),(9,3),$$
that is, for strings in 5 and 6 dimensions, and for 3-branes in 10 dimensions. Selecting $`\lambda `$ so that the coefficient of $`g^M`$ is $`n`$ \[as in (14)\], we obtain in this case, relative to a coordinate basis,
$$S_{\mu \nu }=\frac{1}{2(p+1)!}[F_\mu \mathrm{}F_\nu ^{\mathrm{}}\frac{p+1}{(p+2)(n1)}F^2g_{\mu \nu }],$$
\[see , section 2.1\], and so \[recalling that the space is $`(n+1)`$-dimensional\] we have
$$g^{\mu \nu }S_{\mu \nu }=\frac{(n2p3)F^2}{2(p+2)!(n1)}.$$
Now for a 3-brane in 10 dimensions this vanishes, so $`S`$ is traceless. For a non-trivial solution, this means that some eigenvalue of S must indeed be negative somewhere. For a string in 5 dimensions, the trace is negative, while for a string in 6 dimensions it vanishes; in both cases, $`S`$ is again forced to have eigenvalues which are negative somewhere. For all of these solutions, $`F`$ decays towards infinity as \[,section 4\] $`1/r^{(np1)}`$, so $`S`$ decays at least as rapidly as $`1/r^4`$, as required by the above theorem. In short, the theorem is at least consistent with the hypothesis that some kind of p-brane is present.
We can summarise as follows. The WYCG and Cai-Galloway theorems imply that the presence of a topologically non-trivial infinity hypersurface imposes conditions on the geometry of the bulk. We argue that those conditions suggest that the bulk has been contaminated with “branes or stringy impurities of some kind”, precisely as predicted in .
## VI CONCLUSION
We have argued that the problem of formulating AdS/CFT for manifolds which are not boundaries suggests that we need an alternative framework, which does away with the need to deal with manifolds-with-boundaries. This framework, which uses compact manifolds with infinity hypersurfaces, may find uses even for cases where infinity can be represented as a boundary. For example, it could be interesting to investigate a CFT on the four-torus $`T^4`$ by thinking of it as a submanifold of $`T^5`$. Using the metric (11) above with $`\delta =0`$, we obtain an Einstein (in fact, a locally (Euclidean) AdS) space with a “pseudo-Randall-Sundrum” brane at the antipode to infinity.
In this work, we have followed the usual practice, relating the CFT to a gravitational theory in one more dimension. As string and $`M`$ theories are defined in 10 or 11 dimensions, however, this really just means that we are considering products, like $`AdS_5\times S^5`$, for the bulk. Presumably a generic bulk will not have this special structure, but a 10-dimensional manifold-with-boundary does not have a four-dimensional boundary. As a boundary cannot itself have a boundary, one cannot work stepwise down to four dimensions within the “infinity as a boundary” interpretation. By contrast, one can easily consider a four-dimensional submanifold in a generic compact 10-dimensional manifold endowed with a metric which puts that submanifold “infinitely far” from the points in the bulk. It would be interesting to develop AdS/CFT in this more general setting. |
warning/0003/math0003123.html | ar5iv | text | # Classification des triples de Manin pour les algèbres de Lie réductives complexes
## 0 Introduction
Let $`𝔤`$ be a finite dimensional, complex, reductive Lie algebra. One says that a symmetric, $`𝔤`$-invariant, $``$(resp. $``$)-bilinear form on $`𝔤`$ is a Manin form if and only if its signature is $`(dim_{}𝔤,dim_{}𝔤)`$ (resp. is non degenerate). Recall that a Manin-triple in $`𝔤`$ is a triple $`(B,𝔦,𝔦^{})`$, where $`B`$ is a real (resp. complex) Manin form and where $`𝔦,𝔦^{}`$ are real (resp. complex) Lie subalgebras of $`𝔤`$, isotropic for $`B`$, and such that $`𝔦+𝔦^{}=𝔤`$. Then this is a direct sum and $`𝔦,𝔦^{}`$ are of real dimension equal to the complex dimension of $`𝔤`$. Our goal is to classify all Manin-triples of $`𝔤`$, up to conjugacy under the action on $`𝔤`$ of the connected, simply connected Lie group $`G`$ with Lie algebra $`𝔤`$, by induction on the rank of the derived algebra of $`𝔤`$. One calls af-involution (resp. f-involution) of a complex semi-simple Lie algebra $`𝔪`$, any $``$(resp. $``$)-linear involutive automorphism of $`𝔪`$, $`\sigma `$, such that there exists : 1) an ideal $`\stackrel{~}{𝔪}_0`$ of $`𝔪`$, which is reduced to zero for f-involutions, and a real form $`\stackrel{~}{𝔥}_0`$ of $`\stackrel{~}{𝔪}_0,`$ 2) simple ideals of $`𝔪`$, $`𝔪_j^{}`$, $`𝔪_j^{\prime \prime }`$, $`j=1,\mathrm{},p`$, 3) an isomorphism of complex Lie algebras, $`\tau _j`$, between $`𝔪_j^{}`$ and $`𝔪_j^{\prime \prime }`$, $`j=1,\mathrm{},p`$, such that, denoting by $`𝔥_j:=\{(X,\tau _j(X))|X𝔪_j^{}\}`$, and by $`𝔥`$ the fixed point set of $`\sigma `$, one has :
$$𝔪=\stackrel{~}{𝔪}_0(_{j=1,\mathrm{},p}(𝔪_j^{}𝔪_j^{\prime \prime }))$$
$$𝔥=\stackrel{~}{𝔥}_0(_{j=1,\mathrm{},p}𝔥_j)$$
Notice that an $``$-linear involutive automorphism of $`𝔪`$ is determined by its fixed point set, as the set of antiinvariant points is just the orhogonal of the fixed point set, for the Killing form of $`𝔪`$, viewed as a real Lie algebra. The following Theorem generalizes previous results of E.Karolinsky (cf \[K2\], Theorem 1 (i) and \[K1\] for the proof, see also \[K3\] Proposition 3.1), as we do not make any restriction on the Manin form. If we are dealing with $``$(resp. $``$)-bilinear Manin form, Manin triple will mean real (resp. complex) Manin triple Theorem 1 Let $`B`$ be a Manin form and let $`𝔦`$ be a Lagrangian subalgebra of $`𝔤`$ for $`B`$, i.e. a real (resp. complex) Lie subalgebra of $`𝔤`$ whose real dimension is equal to the complex dimension of $`𝔤`$ and which is isotropic for $`B`$ . Then : a) If we denote by $`𝔭`$ the normalizer in $`𝔤`$ of the nilpotent radical, $`𝔫`$, of $`𝔦`$, then $`𝔭`$ is a parabolic subalgebra of $`𝔤`$, with nilpotent radical equals to $`𝔫`$. b) Let $`𝔩`$ be a Levi subalgebra of $`𝔭`$ (i.e. $`𝔩`$ is a reductive Lie subalgebra of $`𝔭`$ with $`𝔭=𝔩𝔫`$), and denote by $`𝔪`$ its derived ideal and $`𝔞`$ its center, then the intersection, $`𝔥`$, of $`𝔦`$ and $`𝔪`$ is the fixed point set of an af-involution (resp. f-involution) of $`𝔪`$, which is isotropic for $`B`$. c) The intersection, $`𝔦_𝔞`$, of $`𝔞`$ and $`𝔦`$, is isotropic for $`B`$, and its real dimension equals the complex dimension of $`𝔞`$. d) One has $`𝔦=𝔥𝔦_𝔞𝔫`$. Reciprocally, any real (resp. complex) Lie subalgebra, $`𝔦`$, of $`𝔤`$, which is of this form is Lagrangian for $`B`$. Then, one says that $`𝔦`$ is under $`𝔭`$ One chooses a Cartan subalgebra $`𝔧_0`$ of $`𝔤`$, and a Borel subalgebra of $`𝔤`$ containing $`𝔧_0`$, $`𝔟_0`$. Then, from Theorem 1 and the Bruhat decomposition, one sees (cf. Proposition 1) that every Manin triple is conjugated, under $`G`$, to a Manin triple $`(B,𝔦,𝔦^{})`$ such that $`𝔦`$ is under $`𝔭`$ and $`𝔦^{}`$ is under $`𝔭^{}`$, with $`𝔭`$ containing $`𝔟_0`$ and $`𝔭^{}`$ containing the opposite Borel subalgebra to $`𝔟_0`$, with respect to $`𝔧_0`$. A Manin triple satisfying these conditions will be called standard, under $`(𝔭,𝔭^{})`$. If $`𝔯`$ is a real subalgebra of $`𝔤`$, we denote by $`R`$ the analytic subgroup of $`G`$ with with Lie algebra $`𝔯`$ . If $`𝔢`$ is an abelian real subalgebra of $`𝔤`$, $`𝔯`$ is an $`𝔢`$-invariant subspace of $`𝔤`$, let $`\mathrm{\Delta }(𝔯,𝔢)`$ the set of weights of $`𝔢`$ in $`𝔯`$, which is the subset of $`Hom_{}(𝔢,)`$ whose elements are joint eigenvalues of elements in $`𝔢`$ acting on $`𝔯`$. The weight space of $`\alpha `$ in $`\mathrm{\Delta }(𝔯,𝔢)`$ is denoted by $`𝔯^\alpha `$. Let $`𝔭`$, $`𝔭^{}`$ be given as above, and let $`B`$ be a Manin form on $`𝔤`$. Theorem 2 If there exists a standard Manin triple $`(B,𝔦,𝔦^{})`$ under $`(𝔭,𝔭^{})`$, then $`𝔭`$ or $`𝔭^{}`$ is different from $`𝔤`$ . Theorem 3 Let $`(B,𝔦,𝔦^{})`$ be a real (resp. complex) standard Manin triple under $`(𝔭,𝔭^{})`$). Let $`p^𝔫^{}`$ be the projection of $`𝔤`$ on the $`𝔧_0`$-invariant supplementary subspace of the nilpotent radical $`𝔫^{}`$ of $`𝔭^{}`$ , with kernel $`𝔫^{}`$ . Let $`𝔩𝔫`$ the Langlands decomposition of $`𝔭`$, such that $`𝔩`$ contains $`𝔧_0`$. Set $`𝔦_1=p^𝔫^{}(\stackrel{~}{𝔥}𝔭^{})`$, where $`\stackrel{~}{𝔥}=𝔦𝔩`$. One defines similarly $`𝔦_1^{}`$. Then $`𝔦_1,𝔦_1^{}`$ are contained in $`𝔩𝔩^{}`$. Moreover, if $`B_1`$ denotes the restriction of $`B`$ to $`𝔩𝔩^{}`$, $`(B_1,𝔦_1,𝔦_1^{})`$ is a real (resp. complex) Manin triple for $`𝔩𝔩^{}`$. We set $`𝔤_1:=𝔩𝔩^{}`$. We will use freely the notation of Theorem 1 for $`(B_1,𝔦_1,𝔦_1^{})`$, which is called the predecessor of the standard Manin triple $`(B,𝔦,𝔦^{})`$ . Theorem 4 Every real (resp. complex ) Manin triple under $`(𝔭,𝔭^{})`$ is conjugate, by an element of $`PP^{}`$, to a real (resp. complex) Manin triple under $`(𝔭,𝔭^{})`$, $`(B,𝔦,𝔦^{})`$, whose all successive predecessors, $`(B,𝔦_1,𝔦_1^{}),(B,𝔦_2,𝔦_2^{}),\mathrm{}`$, are standard Manin triples in $`𝔤_1=𝔩𝔩^{},𝔤_2,\mathrm{}`$, with respect to the intersection of $`𝔟_0,𝔟_0^{}`$, with $`𝔤_1,𝔤_2,\mathrm{}`$, and such that the intersection $`𝔣_0`$ ( resp. $`𝔣_0^{}`$) of $`𝔧_0`$ with $`𝔦`$ (resp. $`𝔦^{}`$) is a fundamental Cartan subalgebra of $`𝔦`$ (resp. $`𝔦^{}`$), contained in $`𝔦_1,𝔦_2,\mathrm{}`$ (resp. $`𝔦_1^{},𝔦_2^{},\mathrm{})`$ . Such a Manin triple will be called strongly standard. The smallest integer, $`k`$, such that $`𝔤_k=𝔧_0`$, is called the height of the strongly standard Manin triple.
Now, we assume that $`B`$ is $``$-bilinear. One defines : $`^+:=\{\lambda ^{}|Re\lambda <0,orRe\lambda =0etIm\lambda >0\},^{}=^{}^+`$. If $`B`$ is a complex Manin form on $`𝔤`$, one denotes by $`𝔤_+`$ (resp. $`𝔤_{}`$) the sum of the ideals of $`𝔤`$, $`𝔤_i`$ , for which the restriction of $`B`$ to $`𝔤_i`$ is equal to $`\lambda _iK_{𝔤_i}`$, where $`\lambda _i^+`$ (resp. $`^{}`$), and $`K_{𝔤_i}`$ is the Killing form of $`𝔤_i`$. Set $`𝔧_+=𝔧_0𝔤_+`$, $`𝔧_{}=𝔧_0𝔤_{}`$ . The restriction from $`𝔧_0`$ to $`𝔧_+`$ identifies the roots from $`𝔧_0`$ in $`𝔤_+`$ to those from $`𝔧_+`$ in $`𝔤_+`$. Let us denote by $`\stackrel{~}{R}_+`$ the set of these roots and by $`\mathrm{\Sigma }_+`$, the set of simple roots of the set of positive roots, $`\stackrel{~}{R}_+^+`$, of $`\stackrel{~}{R}_+`$, whose elements are the non zero weights of $`𝔧^+`$ in $`𝔟_0𝔤^+`$. One defines similarly $`\stackrel{~}{R}_{}`$. One defines also $`\mathrm{\Sigma }_{}`$, the set of simple roots of the set of positive roots, $`\stackrel{~}{R}_{}^+`$, of $`\stackrel{~}{R}_{}`$, whose elements are the non zero weights of $`𝔧_{}`$ in $`𝔟_0^{}𝔤^{}`$. For $`\alpha \stackrel{~}{R}=\stackrel{~}{R}_+\stackrel{~}{R}_{}`$, let $`H_\alpha 𝔧_0`$ be the coroot of $`\alpha `$. Let $`𝒲=(H_\alpha ,X_\alpha ,Y_\alpha )_{\alpha \mathrm{\Sigma }}`$, be a Weyl system of generator of $`[𝔤,𝔤]`$, where $`\mathrm{\Sigma }=\mathrm{\Sigma }_+\mathrm{\Sigma }_{}`$. More precisely, for all $`\alpha ,\beta \mathrm{\Sigma }`$, one has :
$$[X_\alpha ,Y_\beta ]=\delta _{\alpha \beta }H_\beta $$
$$[H_\alpha ,X_\beta ]=N_{\alpha \beta }X_\beta $$
$$[H_\alpha ,Y_\beta ]=N_{\alpha \beta }Y_\beta $$
where :
$$N_{\alpha \beta }=\beta (H_\alpha )=2K_𝔤(H_\alpha ,H_\beta )/K_𝔤(H_\alpha ,H_\alpha )$$
Definition One calls $`(A,A^{},𝔦_𝔞,𝔦_𝔞^{})`$ generalized Belavin-Drinfeld data with respect to $`B`$, if : 1) $`A`$ is a bijection from a subset $`\mathrm{\Gamma }_+`$ of $`\mathrm{\Sigma }_+`$ on a subset $`\mathrm{\Gamma }_{}`$ of $`\mathrm{\Sigma }_{}`$, such that :
$$B(H_{A\alpha },H_{A\beta )}=B(H_\alpha ,H_\beta ),\alpha ,\beta \mathrm{\Gamma }_+$$
2) $`A^{}`$ is a bijection from a subset $`\mathrm{\Gamma }_+^{}`$ of $`\mathrm{\Sigma }_+`$ on a subset $`\mathrm{\Gamma }_{}^{}`$ of $`\mathrm{\Sigma }_{}`$, such that :
$$B(H_{A^{}\alpha },H_{A^{}\beta })=B(H_\alpha ,H_\beta ),\alpha ,\beta \mathrm{\Gamma }_+^{}$$
3) Let $`C=^{\prime \prime }A^1A^{}^{\prime \prime }`$ be the map defined on $`domC=\{\alpha \mathrm{\Gamma }_+^{}|A^{}\alpha \mathrm{\Gamma }_{}\}`$ by $`C\alpha =A^1A^{}\alpha `$, $`\alpha domC`$. Then $`C`$ satisfies : For all $`\alpha domC`$, there exists $`n^{}`$ such that $`\alpha ,\mathrm{},C^{n1}\alpha domC`$ and $`C^n\alpha domC`$. 4) $`𝔦_𝔞`$ (resp.$`𝔦_𝔞^{}`$) is a complex vector subspace of $`𝔧_0`$, i, included and Lagrangian in the orthogonal, $`𝔞`$ (resp. $`𝔞^{}`$) for the Killing form of $`𝔤`$ (or for $`B`$), to the suspace generated by $`H_\alpha `$, $`\alpha \mathrm{\Gamma }:=\mathrm{\Gamma }_+\mathrm{\Gamma }_{}`$ (resp. $`\mathrm{\Gamma }^{}:=\mathrm{\Gamma }_+^{}\mathrm{\Gamma }_{}^{}`$. 5) Let $`𝔣`$ be the subspace of $`𝔧_0`$ generated by the family $`H_\alpha +H_{A\alpha }`$, $`\alpha \mathrm{\Gamma }_+`$. One defines similarly $`𝔣^{}`$. Then :
$$(𝔣𝔦_𝔞)(𝔣^{}𝔦_𝔞^{})=\{0\}$$
We will denote by $`R_+`$ the sub-system of roots of $`\stackrel{~}{R}`$ whose elements those of $`\stackrel{~}{R}`$ which are linear combination of elements of $`\mathrm{\Gamma }_+`$. One defines similarly $`R_{}`$, $`R_+^{}`$, $`R_{}^{}`$. We will denote also by $`A`$ (resp. $`A^{}`$) the $``$-linear extension of $`A`$ (resp. $`A^{}`$), which defines a bijection from $`R_+`$ on $`R_{}`$ (resp. $`R_+^{}`$ on $`R_{}^{}`$). If $`A`$ satisfies the condition 1) above, there exists a unique isomorphism $`\tau `$ from the subalgebra $`𝔪_+`$ of $`𝔤`$, generated by $`X_\alpha ,H_\alpha ,Y_\alpha `$, $`\alpha \mathrm{\Gamma }_+`$, on the subalgebra $`𝔪_{}`$ of $`𝔤`$, generated by $`X_\alpha ,H_\alpha ,Y_\alpha `$, $`\alpha \mathrm{\Gamma }_{}`$, such that :
$$\tau (H_\alpha )=H_{A\alpha },\tau (X_\alpha )=X_{A\alpha },\tau (Y_\alpha )=Y_{A\alpha },\alpha \mathrm{\Gamma }^+$$
Theorem (cf. Proposition 8 et Théorème 5) (i) Let $`𝒟=(A,A^{},𝔦_𝔞,𝔦_𝔞^{})`$ be generalized Belavin-Drinfeld data, with respect to $`B`$. Let $`𝔭`$ be the parabolic subalgebra of $`𝔤`$, containing $`𝔟_0`$ and $`𝔪`$. Its Langlands decomposition $`𝔭=𝔩𝔫`$, where $`𝔩`$ contains $`𝔧_0`$, satisfies $`𝔩=𝔪𝔞`$. Let $`𝔦`$ be equal to $`𝔥𝔦_𝔞𝔫`$, where $`𝔥:=\{X+\tau (X)|X𝔪_+\}`$. One defines similarly $`𝔦^{}`$. Then $`(B,𝔦,𝔦^{})`$ is a strongly standard Manin triple. (ii) Every Manin triple is conjugate by an element of $`G`$ to a unique Manin triple of this type. If the original triple is moreover strongly standard, the element of $`G`$ can be taken in $`J_0`$, or, in other words, every strongly standard is of the precceeding type, if one allows to change $`𝒲`$. One shows easily that the preceeding Theorem implies the classification of certain $`R`$-matrix given by Belavin and Drinfeld (\[BD\], Theorem 6.1). One proves also results for real Manin triples. One retrieves a result of A. Panov \[P1\] which classifies certain Lie bialgebras structures on a real simple Lie algebra.
Aknowledgment : I thank very much C. Klimcik for suggesting me this work, and for many interesting discussions. I thank J.L. Brylinski for pointing out to me the work of E. Karolinsky. I thank also B. Enriquez and Y. Kosmann-Scwarzbach for telling to me the relations of my earlier work \[De\] with the work of A. Belavin and G. Drinfeld \[BD\], and A. Panov \[P1\].
## 1 Sous-algèbres de Lie Lagrangiennes
Dans tout l’article, algèbre de Lie voudra dire algèbre de Lie de dimension finie. Si $`𝔤`$ est une algèbre de Lie on notera souvent $`𝔤^{der}`$ son idéal dérivé. Soit $`𝔞`$ est une algèbre de Lie abélienne sur $`𝕂=`$ ou $``$, $`V`$ un $`𝔞`$-module complexe. Pour $`\lambda Hom_𝕂(𝔞,)`$, on note $`V^\lambda :=\{vV|Xv=\lambda (X)v,X𝔞\}`$, qui est appelé le sous-espace de poids $`\lambda `$ de $`V`$. On dit que $`\lambda `$ est un poids de $`𝔞`$ dans $`V`$ si $`V^\lambda `$ est non nul et on note $`\mathrm{\Delta }(V,𝔞)`$ l’ensemble des poids non nuls de $`𝔞`$ dans $`V`$. Si $`G`$ est un groupe de Lie, on notera $`G^0`$ sa composante neutre.
###### Lemme 1
(i)Soit $`𝔤`$ une algèbre de Lie semi-simple complexe, $`𝔤_1,\mathrm{},𝔤_n`$ ses idéaux simples. Toute forme $``$ (resp.$``$)-bilinéaire $`𝔤`$-invariante sur $`𝔤`$ est du type $`B_\lambda `$ ou $`B_\lambda ^𝔤`$ (resp. $`K_\lambda `$ ou $`K_\lambda ^𝔤`$), où $`\lambda =(\lambda _1,\mathrm{},\lambda _n)^n`$ et :
$$B_\lambda (X_1+\mathrm{}+X_n,Y_1+\mathrm{}+Y_n)=\underset{i=1,\mathrm{},n}{}Im(\lambda _iK_{𝔤_i}(X_i,Y_i))$$
$$K_\lambda (X_1+\mathrm{}+X_n,Y_1+\mathrm{}+Y_n=\underset{i=1,\mathrm{},n}{}\lambda _iK_{𝔤_i}(X_i,Y_i),$$
si les $`X_i,Y_i`$ sont des éléments de $`𝔤_i`$. Ici $`K_{𝔤_i}`$ désigne la forme de Killing de $`𝔤_i`$. En particulier, une telle forme est symétrique et les idéaux simples sont deux à deux orthogonaux pour une telle forme. (ii) La forme $`B_\lambda `$ (resp. $`K_\lambda `$) est non dégénérée si et seulement si chacun des $`\lambda _i`$ est non nul. La forme $`B_\lambda `$ est alors de signature $`(dim_{}𝔤,dim_{}𝔤)`$. (iii) La restriction de $`B_\lambda ^𝔤`$ à une sous-algèbre complexe et simple, $`𝔰`$, de $`𝔤`$, est de la forme $`B_\mu ^𝔰`$, où $`\mu =_{i=1,\mathrm{},n}q_i\lambda _i`$ et pour chaque $`i`$, $`q_i`$ est un nombre rationnel positif. De plus $`q_i`$ est non nul si et seulement si $`𝔰`$ a un crochet non nul avec $`𝔤_i`$.
Démonstration : On traite le cas des formes $``$-bilinéaires, celui des formes $``$-bilinéaires étant semblable. Traitons d’abord le cas où $`𝔤`$ est simple, auquel cas $`n=1`$. La donnée d’une forme $``$-bilinéaire $`𝔤`$-invariante sur $`𝔤`$, équivaut à celle d’une application $``$-linéaire entre $`𝔤`$ et $`Hom_{}(𝔤,)`$, qui commute à l’action de $`𝔤`$, regardée comme algèbre de Lie réelle. Mais la partie imaginaire de la forme de Killing de $`𝔤`$ (regardée comme complexe) détermine un isomorphisme de $`𝔤`$-modules entre $`𝔤`$ et $`Hom_{}(𝔤,)`$. Finalement, la donnée de notre forme équivaut à la donnée d’un endomorphisme, $`T`$, $``$-linéaire de $`𝔤`$, commutant à l’action de $`𝔤`$. Soit $`𝔨`$ une forme rélle compacte de $`𝔤`$. On écrit :
$$T(X)=ReT(X)+iImT(X),X𝔨,$$
$`ReT(X),ImT(X)𝔨`$. Alors $`ReT`$, $`ImT`$ sont des éléments de $`Hom_𝔨(𝔨,𝔨)`$. Comme $`𝔨`$ est simple, il résulte du lemme de Schur que cet espace est égal à $`Id_𝔨`$. Donc, il existe $`\lambda `$ tel que :
$$T(X)=\lambda X,X𝔨$$
Maintenant, si $`X,Y𝔨`$, on a :
$$T(i[X,Y])=T([iX,Y])=[iX,TY]$$
la dernière égalité résultant du fait que $`T`$ commute à l’action de $`𝔤`$. Joint à ce qui précède, cela donne :
$$T(i[X,Y])=\lambda i[X,Y],X,Y𝔨$$
Comme $`[𝔨,𝔨]=𝔨`$, on conclut que $`T`$ est la multiplication par $`\lambda `$. D’où (i) dans le cas où $`𝔤`$ est simple. Supposons maintenant que $`𝔤`$ soit la somme directe de deux idéaux $`𝔤^{}`$, $`𝔤^{\prime \prime }`$. Soit $`B`$ une forme $``$-bilinéaire $`𝔤`$-invariante sur $`𝔤`$. On a :
$$B([X^{},Y^{}],X^{\prime \prime })=B(Y^{},[X^{},X^{\prime \prime }])=0,X^{},Y^{}𝔤^{},X^{\prime \prime }𝔤^{\prime \prime },$$
la dernière égalité résultant du fait que $`𝔤^{}`$, $`𝔤^{\prime \prime }`$ commutent entre eux. Comme $`[𝔤,𝔤]=𝔤`$, on a aussi $`[𝔤^{},𝔤^{}]=𝔤^{}`$. Finalement $`𝔤^{}`$ et $`𝔤^{\prime \prime }`$ sont orthogonaux. Alors, on déduit (i) pour $`𝔤`$ semi-simple du cas où $`𝔤`$ est simple. (ii) L’assertion sur la non nullité des $`\lambda _i`$ est claire. Pour l’étude de la signature, on se ramène au cas où $`𝔤`$ est simple. Supposons $`\lambda `$, non nul. Soit $`𝔨`$ une forme réelle compacte de $`𝔤`$. On fixe une base $`X_1,\mathrm{},X_l`$ de $`𝔨`$. On choisit une racine carrée $`\mu `$ de $`i\lambda ^1`$.On pose $`Y_i=\mu X_i`$, $`Z_i=i\mu X_i`$. Alors $`B_\lambda (Y_i,Z_j)=0`$, $`B_\lambda (Y_i,Y_j)=\delta _{i,j}`$, $`B_\lambda (Z_i,Z_j)=\delta _{i,j}`$. D’où l’on déduit l’assertion voulue sur la signature. (iii) On utilisera le fait suivant : Si $`\rho `$ est une représentation complexe d’une algèbre de Lie simple complexe $`𝔰`$ dans un espace de dimension finie $`V`$, on a :
$$tr(\rho (X)\rho (Y))=qK_𝔰(X,Y)$$
$`q`$ est un nombre rationel positif. De plus $`q`$ est nul si et seulement si $`\rho `$ est triviale. L’existence d’un coefficient de proportionnalité $`q`$ est claire , car la forme de Killing est, à un scalaire mutiplicatif près, la seule forme $``$-bilinéaire invariante sur $`𝔰`$. On se ramène, pour l’étude de $`q`$, au cas où $`\rho `$ est simple. On considère, sur $`V`$, un produit scalaire invariant par une forme réelle compacte, $`𝔨`$, de $`𝔰`$. Si $`X𝔨`$, $`\rho (X)`$ est antihermitien et :
$$tr(\rho (X)\rho (X))=tr(\rho (X)\rho (X)^{})0$$
cette trace étant nulle seulement si $`\rho (X)`$ est nul. On en déduit que $`q>0`$ si $`\rho `$ est non triviale. Puis on prend un élément non nul d’une sous-algèbre de Cartan $`𝔧`$ de $`𝔰`$, sur lequel tous les poids entiers de $`𝔧`$ sont entiers. On en déduit que $`K_𝔤(X,X)`$ et $`tr(\rho (X)\rho (X))`$ sont des entiers, le premier nombre étant non nul car égal à la somme sur toutes les racines, $`\alpha `$, de $`(\alpha (X))^2`$. La rationalité de $`q`$ en résulte.
###### Définition 1
Si $`𝔤`$ est une algèbre de Lie réductive complexe, une forme $``$(resp. $``$)-bilinéaire symétrique sur $`𝔤`$ et invariante par $`𝔤`$ est dite forme de Manin si et seulement si elle est de signature $`(dim_{}𝔤,dim_{}𝔤)`$ (resp. si et seulement si elle est non dégénérée). Une forme de Manin est dite forme spéciale si sa restriction à toute sous-algèbre de Lie complexe semi-simple est non dégénérée.
###### Lemme 2
(i) Une forme $``$(resp. $``$)-bilinéaire symétrique $`𝔤`$-invariante sur $`𝔤`$ est spéciale si et seulement si sa restriction à $`𝔤^{der}`$ est spéciale et si sa restriction au centre, $`𝔷`$, de $`𝔤`$ est de signature $`(dim_{}𝔷,dim_{}𝔷)`$ (resp. est non dégénérée). (ii) La restriction d’une forme spéciale à une sous-algèbre de Lie semi- simple complexe de $`𝔤`$ est spéciale. (iii) La restriction d’une forme spéciale au centralisateur d’un élément semi-simple de $`𝔤`$, dont l’image par la représentation adjointe de $`𝔤`$ n’a que des valeurs propres réelles, est spéciale. (iv) Si $`𝔤`$ est semi-simple, et $`B=B_\lambda ^𝔤`$ (resp. $`K=K_\lambda ^𝔤`$), où $`\lambda =(\lambda _1,\mathrm{},\lambda _n)^n`$ vérifie :
$$Si\underset{i=1,\mathrm{},n}{}q_i\lambda _i=0avecq_i^+,alorslesq_isonttousnuls$$
(1.1)
alors $`B`$ (resp. $`K`$) est spéciale. On note que (1.1) est satisfait dès que les $`\lambda _i`$ sont indépendants sur $``$, ou bien tous strictement positifs. (v) Si $`𝔤`$ est simple, toute forme $``$ (resp. $``$)-bilinéaire symétrique $`𝔤`$-inva- riante sur $`𝔤`$ est spéciale.
Démonstration : On traite le cas des formes $``$-bilinéaires, celui des formes $``$-bilinéaires étant semblable. (i) résulte du fait que toute sous-algèbre semi-simple de $`𝔤`$ est contenue dans $`𝔤^{der}`$ et que, pour toute forme $``$-bilinéaire symétrique $`𝔤`$-invariante sur $`𝔤`$, le centre de $`𝔤`$ et $`𝔤^{der}`$ sont orthogonaux. (ii) est clair. Montrons (iii). Comme le centralisateur d’un élément de $`𝔤`$ est la somme de son intersection avec $`𝔤^{der}`$ et de celle avec le centre, on se réduit aisément, grâce à (i) au cas où $`𝔤`$ est semi-simple, ce que l’on suppose dans la suite. Soit $`X`$ un élément semi-simple de $`𝔤`$ tel que $`adX`$ n’a que des valeurs propres réelles, soit $`𝔩`$ son centralisateur et $`𝔠`$ le centre de $`𝔩`$. D’après (i) et (ii), il suffit de voir que la restriction d’une forme spéciale à $`𝔠`$ est de signature $`(dim_{}𝔠,dim_{}𝔠)`$. Soit $`𝔧`$ une sous-algèbre de Cartan de $`𝔤`$ contenant $`X`$. Alors $`𝔠`$ est égal à l’intersection des noyaux des racines de $`𝔧`$ dans $`𝔤`$ s’annulant sur $`X`$. Cela montre que $`𝔠`$ est la somme de ses intersections avec les idéaux simples de $`𝔤`$. Il suffit alors de prouver notre assertion sur la signature dans le cas où $`𝔤`$ est simple. Alors $`B=B_\lambda `$,avec $`\lambda `$ non nul. Soit $`𝔧_{}`$ l’espace formé des éléments de $`𝔧`$ sur lesquels toutes les racines de $`𝔧`$ dans $`𝔤`$ sont réelles, qui est une forme réelle de $`𝔧`$. Il est clair que $`𝔠`$ est la somme directe de $`𝔠_{}:=𝔠𝔧_{}`$ avec $`i𝔠_{}`$. On fixe une base orthonormée, $`X_1,\mathrm{},X_l`$, de $`𝔠_{}`$, pour la forme de Killing de $`𝔤`$. Celle-ci existe car la forme de Killing est définie positive sur $`𝔧_{}`$. On choisit une racine carrée $`\mu `$ de $`i\lambda ^1`$. On pose $`Y_i=\mu X_i`$, $`Z_i=i\mu X_i`$. Alors $`B_\lambda (Y_i,Z_j)=0`$, $`B_\lambda (Y_i,Y_j)=\delta _{i,j}`$, $`B_\lambda (Z_i,Z_j)=\delta _{i,j}`$. D’où l’on déduit l’assertion voulue sur la signature, ce qui prouve (iii). (iv) est une conséquence immédiate du Lemme 1 et (v) est un cas particulier de (iv)
###### Corollaire 1
(i) La restriction d’une forme $``$ (resp.$``$)-bilinéaire, symétrique, $`𝔤`$-invariante sur $`𝔤`$, et non dégénérée, $`B`$ , au centralisateur, $`𝔩`$ d’un élément semi-simple de $`𝔤`$, dont l’image par la représentation adjointe de $`𝔤`$ n’a que des valeurs propres réelles, est non dégénérée. Il en va de même de sa restriction à $`𝔩^{der}`$ et au centre $`𝔞`$ de $`𝔩`$ . (ii) Si $`B`$ est une forme de Manin sur $`𝔤`$, sa restriction à $`𝔩`$ (resp. $`𝔩^{der}`$, $`𝔞`$) est une forme de Manin sur $`𝔩`$ (resp. $`𝔩^{der}`$, $`𝔞`$). (iii) Une forme bilinéaire symétrique, $`𝔤`$-invariante est une forme de Manin si et seulement si sa restriction à $`𝔤^{der}`$ et sa restriction au centre de $`𝔤`$ sont des formes de Manin.
Démonstration : On traite le cas des formes $``$-bilinéaires, celui des formes $``$-bilinéaires étant semblable. Montrons (i). L’algèbre de Lie $`𝔩`$ est la somme du centre $`𝔷`$ de $`𝔤`$ avec la somme de ses intersections avec les idéaux simples de $`𝔤`$. Ces sous-algèbres sont deux à deux orthogonales pour $`B`$, d’après le Lemme 1. La restriction de $`B`$ à $`𝔷`$ est non dégénérée, car $`𝔤^{der}`$ et $`𝔷`$ sont orthogonaux. De plus la restriction de $`B`$ à l’intersection de $`𝔩`$ avec un idéal simple de $`𝔤`$ est non dégénérée, d’après le Lemme 2 (iii), appliqué à cet idéal simple. D’où l’on déduit que la restriction de $`B`$ à $`𝔩`$ est non dégénérée, ce qui implique le même fait pour sa restriction à $`𝔩^{der}`$ et $`𝔞`$. Si $`B`$ est de signature $`(dim_{}𝔤,dim_{}𝔤)`$, sa restriction à $`𝔷`$ est de signature $`(dim_{}𝔷,dim_{}𝔷)`$. Alors, les assertions sur la signature se démontre comme ci-dessus, grâce au Lemme 2 (ii), (v), appliqué aux idéaux simples de $`𝔤`$. D’où (ii). La partie si de (iii) est claire. La partie seulement si résulte de (ii) On rappelle que le radical d’une algèbre de Lie, $`𝔤`$, est son plus grand idéal résoluble, et que son radical nilpotent, est le plus grand idéal, dont les éléments sont représentés par des endomorphismes nilpotents dans chaque représentation de dimension finie de $`𝔤`$. Suivant Bourbaki, on appelle sous-algèbre de Levi d’une algèbre de Lie, toute sous-algèbre semi-simple supplé- mentaire du radical. Si $`𝔤`$ est une algèbre de Lie semi-simple complexe on appelle décomposition de Langlands d’une sous-algèbre parabolique $`𝔭`$ de $`𝔤`$ une décomposition de la forme $`𝔭=𝔩𝔫`$, où $`𝔫`$ est le radical nilpotent et $`𝔩`$ est une sous-algèbre de Lie complexe de $`𝔤`$, réductive dans $`𝔤`$. Rassemblons dans un Lemme quelques propriétés des décompositions de Langlands d’une sous-algèbre parabolique.
###### Lemme 3
Soit $`𝔭`$ une sous-algèbre parabolique de $`𝔤`$, $`𝔫`$ son radical nilpotent. (i) Si $`𝔧`$ est une sous-algèbre de Cartan de $`𝔭`$, c’est une sous-algèbre de Cartan de $`𝔤`$ , et il existe une seule décomposition de Langlands de $`𝔭`$, $`𝔭=𝔩𝔫`$, telle que $`𝔧`$ soit contenue dans $`𝔩`$. (ii) Si $`𝔧`$, $`𝔧^{}`$ sont deux sous-algèbres de Cartan de $`𝔤`$, contenues dans $`𝔭`$, elles sont conjuguées par un élément de $`P`$, qui conjugue les algèbres $`𝔩`$ et $`𝔩^{}`$ correspondantes. (iii) Si $`𝔭=𝔩𝔫`$ est une décomposition de Langlands de $`𝔭`$, toute sous-algèbre de Cartan de $`𝔭`$ est une sous-algèbre de Cartan de $`𝔤`$.
Démonstration : Revenant à la définition des sous-algèbres paraboliques (cf. \[Bou\], Ch. VIII, Paragraphe 3.4, Définition 2, par exemple), on voit qu’il existe une sous-algèbre de Cartan de $`𝔤`$, $`𝔧_1`$, et une décomposition de Langlands de $`𝔭`$, $`𝔭=𝔩_1𝔫`$, avec $`𝔧_1`$ contenue dans $`𝔩_1`$, telle que $`𝔩_1`$ soit la somme des sous-espaces poids de $`𝔧_1`$ dans $`𝔭`$, qui ne rencontrent pas $`𝔫`$. En particulier les poids de $`𝔧_1`$ dans $`𝔩_1𝔭/𝔫`$ sont distincts de ceux dans $`𝔫`$. Si $`𝔭=𝔩_1^{}𝔫`$ est une décomposition de Langlands de $`𝔭`$, avec $`𝔧_1`$ contenue dans $`𝔩_1^{}`$, $`𝔩_1^{}`$ est somme des sous-espaces poids de $`𝔧_1`$ dans $`𝔩_1^{}𝔤/𝔫`$. D’où l’égalité de $`𝔩_1`$ et $`𝔩_1^{}`$. Ceci assure l’unicité de $`𝔩`$ pour $`𝔧=𝔧_1`$. Maintenant toutes les sous-algèbres de Cartan de $`𝔭`$ sont conjuguées à $`𝔧_1`$, par un élément de $`P`$ (cf. \[Bour\], Chapitre VII, Paragraphe 3.3, Théorème 1). On en déduit (i) par transport de structure, et (ii) résulte de la preuve de (i). Montrons (iii). Si $`𝔭=𝔩𝔫`$ est une décomposition de Langlands de $`𝔭`$, l’action du centre de $`𝔩`$ sur $`𝔤`$ est semi-simple, puisque $`𝔩`$ est réductive dans $`𝔤`$. Cela implique que, si $`𝔧`$ une une sous-algèbre de Cartan de $`𝔩`$, $`𝔧`$ est abélienne et formée d’éléments semi-simples de $`𝔤`$. Mais $`𝔩`$ est isomorphe à $`𝔭/𝔫`$, qui est une algèbre réductive de même rang que $`𝔤`$. Pour des raisons de dimension, on voit que $`𝔧`$ est une sous-algèbre de Cartan de $`𝔤`$.
###### Définition 2
On appelle af-involution (resp. f-involution), où a vaut pour antilinéaire et f pour flip, d’une algèbre de Lie semi-simple complexe $`𝔪`$, tout automorphisme involutif, $``$-linéaire (resp. $``$-linéaire), $`\sigma `$, de $`𝔪`$, pour lequel il existe : 1) un idéal $`\stackrel{~}{𝔪}_0`$ de $`𝔪`$, qui est en outre est réduit à zéro pour les f-involutions, et une forme réelle, $`\stackrel{~}{𝔥}_0`$, de $`\stackrel{~}{𝔪}_0`$ . 2) des idéaux simples de $`𝔪`$, $`𝔪_j^{}`$, $`𝔪_j^{\prime \prime }`$, $`j=1,\mathrm{},p`$. 3) un isomorphisme d’algèbres de Lie complexes, $`\tau _j`$, entre $`𝔪_j^{}`$ et $`𝔪_j^{\prime \prime }`$, $`j=1,\mathrm{},p`$, tel que, notant $`𝔥_j:=\{(X,\tau _j(X))|X𝔪_j^{}\}`$, et notant $`𝔥`$, l’ensemble des points fixes de $`\sigma `$, on ait :
$$𝔪=\stackrel{~}{𝔪}_0(_{j=1,\mathrm{},p}(𝔪_j^{}𝔪_j^{\prime \prime }))$$
$$𝔥=\stackrel{~}{𝔥}_0(_{j=1,\mathrm{},p}𝔥_j)$$
Il est bon de remarquer qu’un automorphisme involutif de $`𝔪`$ est caractérisé par son espace de points fixes, car l’espace des éléments anti-invariants est juste l’orthogonal de celui-ci, pour la forme de Killing de $`𝔪`$ regardée comme réelle. On remarque qu’une f-involution est en particulier une af-involution. Débutons par quelques propriétés élémentaires.
###### Lemme 4
Soit $`𝔥`$ une forme réelle simple d’une algèbre de Lie semi-simple complexe $`𝔰`$. (i) L’algèbre $`𝔰`$ n’est pas simple, si et seulement si $`𝔥`$ admet une structure complexe (ii) Dans ce cas, $`𝔰`$ est le produit de deux idéaux simples, $`𝔰_1`$, $`𝔰_2`$, isomorphes à $`𝔥`$. (iii) Toujours dans ce cas, il existe un isomorphisme antilinéaire, $`\tau `$, entre les algèbres de Lie $`𝔰_1`$, $`𝔰_2`$, regardées comme réelles, tel que :
$$𝔥=\{(X,\tau (X))|X𝔰_1\}$$
Démonstration : Les points (i) et (ii) sont bien connus. Montons (iii). Comme $`𝔥`$ est une forme réelle de $`𝔰_1𝔰_2`$, la projection de $`𝔥`$ sur chacun des deux facteurs est non nulle, donc induit un isomorphisme $``$-linéaire de $`𝔥`$ avec chacun de ces facteurs. Il en résulte que $`𝔥`$ a la forme indiquée, mais on sait seulement que $`\tau `$ est $``$-linéaire. Mais alors $`𝔥`$ apparaît comme l’ensemble des points fixes de l’automorphisme involutif $``$-linéaire de $`𝔰`$ défini par :
$$(X,Y)(\tau ^1(Y),\tau (X)),X𝔰_1,Y𝔰_2$$
D’après la remarque qui précède le Lemme, cette involution doit être égale à la conjugaison par rapport à $`𝔥`$, donc elle est antilinéaire. Ceci implique l’antilinéarité de $`\tau `$.
###### Lemme 5
On se donne une af-involution, $`\sigma `$, d’une algèbre de Lie semi-simple de $`𝔪`$. Les idéaux simples de $`𝔪`$ sont permutés par $`\sigma `$. On note $`𝔪_j`$, $`j=1,\mathrm{},r`$, les idéaux simples de $`𝔪`$. On définit une involution $`\theta `$ de $`\{1,\mathrm{},r\}`$ caractérisée par : $`\sigma (𝔪_j)=𝔪_{\theta (j)},j=1,\mathrm{},r`$. Pour $`j=1,\mathrm{},r`$, l’une des propriétés suivantes est vérifiée : 1) $`\theta (j)=j`$ et la restriction de $`\sigma `$ à $`𝔪_j`$ est un automorphisme antilinéaire de $`𝔪_j`$. 2) $`\theta (j)j`$ et la restriction de $`\sigma `$ à $`𝔪_j`$ est un isomorphisme antilinéaire de $`𝔪_j`$ sur $`𝔪_{\theta (j)}`$. 3) $`\theta (j)j`$ et la restriction de $`\sigma `$ à $`𝔪_j`$ est un isomorphisme $``$-linéaire de $`𝔪_j`$ sur $`𝔪_{\theta (j)}`$. Si on est dans le cas 1) ou 2), $`𝔪_j`$ est contenu dans l’idéal $`\stackrel{~}{𝔪}_0`$ de la définition des af-involutions. En particulier si $`\sigma `$ est une f-involution, on est toujours dans le cas 3).
Démonstration : En effet, soit $`𝔥_{p+l}`$, $`l=1,\mathrm{},q`$, les idéaux simples de $`\stackrel{~}{𝔥}_0`$. Comme $`\stackrel{~}{𝔥}_0`$ est une forme réelle de $`\stackrel{~}{𝔪}_0`$, $`\stackrel{~}{𝔪}_0`$ est la somme directe des $`𝔥_l+i𝔥_l`$, qui sont en outre des idéaux. Si $`𝔥_{p+l}+i𝔥_{p+l}`$ est simple, c’est un idéal simple de $`𝔪`$ et on est dans le cas 1). Sinon $`𝔥_l+i𝔥_l`$ est le produit de deux idéaux simples et l’on est dans le cas 2), d’après le Lemme précédent. On traite de même le cas où $`𝔪_l`$ est égal à l’un des $`𝔪_j^{}`$, $`𝔪_j^{\prime \prime }`$, $`j=1,\mathrm{},p`$, en remarquant que $`𝔥_j`$ est l’ensemble des points fixes de l’involution $``$-linéaire de $`(𝔪_j^{},𝔪_j^{\prime \prime })`$ donnée par :
$$(X,Y)(\tau ^1(Y),\tau (X)),X𝔪_j^{},Y𝔪_j^{\prime \prime }$$
(1.2)
###### Lemme 6
Tout isomorphisme $``$-linéaire entre deux algèbres de Lie simples complexes est soit $``$-linéaire, soit antilinéaire.
Démonstration : Deux algèbres de Lie semi-simples complexes qui sont isomorphes comme algèbres réelles, sont isomorphes comme algèbres de Lie complexes. En effet, leurs systèmes de racines restreintes sont alors isomorphes. Mais chacun de ceux-ci est isomorphe au système de racine associé à une sous-algèbre de Cartan . D’où l’assertion. Ceci implique que l’on peut se limiter, pour prouver le Lemme, aux automorphismes $``$-linéaire d’une algèbre de Lie simple complexe, $`𝔤`$. Considérant la conjugaison par rapport à une forme réelle de $`𝔤`$, $`X\overline{X}`$, l’algèbre de Lie $`𝔤^{}:=\{(X,\overline{X})|X𝔤\}`$ est une forme réelle de $`𝔤\times 𝔤`$ isomorphe à $`𝔤`$. Alors tout automorphisme, $`\sigma `$, $``$-linéaire de $`𝔤`$, définit, par transport de structure, un automorphisme de $`𝔤^{}`$, $`\sigma ^{}`$, qui possède un unique prolongement $``$-linéaire à $`𝔤\times 𝔤`$, $`\sigma ^{\prime \prime }`$ . Il existe deux automorphismes $``$-linéaires de $`𝔤`$, $`\tau `$ et $`\sigma `$, tels que $`\sigma ^{\prime \prime }`$ vérifie :
$$\sigma ^{\prime \prime }(X,Y)=(\tau (X),\theta (Y))oubien(\tau (Y),\theta (X)),(X,Y)𝔤\times 𝔤$$
Ecrivant la définition de $`\sigma ^{}`$, la stabilité de $`𝔤^{}`$, par $`\sigma ^{\prime \prime }`$ implique que $`\sigma `$ est $``$-linéaire dans le premier cas et antilinéaire dans le second. Corollaire Une involution $``$ (resp. $``$)-linéaire d’une algèbre de Lie semi-simple complexe est une af-involution (resp. f-involution) si et seulement si sa restriction à tout idéal simple qu’elle laisse invariant est antilinéaire (resp . si elle ne laisse aucun idéal simple invariant) Démonstration : En effet, d’après le Lemme 5, il suffit de voir que tout automorphisme du produit de deux algèbres de Lie simples complexes $`𝔰_1`$, $`𝔰_2`$, permutant les facteurs, est soit $``$-linéaire, soit antilinéaire. Mais un tel automorphisme est de la forme :
$$(X,Y)(\tau ^1(Y),\tau (X)),X𝔰_1,Y𝔰_2,$$
$`\tau `$ est un isomorphisme $``$-linéaire entre $`𝔰_1`$, $`𝔰_2`$ . On conclut grâce au Lemme précédent. Soit $`E`$ un espace vectoriel complexe de dimension finie muni d’une forme $``$ (resp. $``$)-bilinéaire symétrique non dégénérée. Tout sous-espace vectoriel réel (resp. complexe) isotrope est de dimension réelle inférieure ou égale à la dimension complexe de $`E`$ Un sous-espace vectoriel réel (resp. complexe) de $`E`$, muni d’une d’une forme $``$ (resp. $``$)-bilinéaire symétrique non dégénérée est dit Lagrangien s’il est isotrope et de dimension réelle égale à la dimension complexe de $`E`$. Un tel espace existe si et seulement la forme est de signature $`(dim_{}E,dim_{}E)`$ (resp. si $`E`$ est de dimension complexe paire). Comme on l’a indiqué dans l’introduction, le Théorème suivant généralise des résultats d’E. Karolinsky (cf. \[K1\], Théorème 3 (i) et \[K3\] Proposition 3.1)
###### Théorème 1
Soit $`𝔤`$ une algèbre de Lie réductive complexe et $`B`$ une forme de Manin $``$ (resp. $``$)-bilinéaire. Soit $`𝔦`$ une sous-algèbre de Lie réelle (resp. complexe) de $`𝔤`$, Lagrangienne pour $`B`$. On a les propriétés suivantes : (i) Si l’on note $`𝔭`$ le normalisateur dans $`𝔤`$ du radical nilpotent, $`𝔫`$, de $`𝔦`$, $`𝔭`$ est une sous-algèbre parabolique de $`𝔤`$, contenant $`𝔦`$, de radical nilpotent $`𝔫`$. (ii) Soit $`𝔭=𝔩𝔫`$ une décomposition de Langlands de $`𝔭`$, $`𝔞`$ le centre de $`𝔩`$ et $`𝔪`$ son idéal dérivé. On note $`𝔥`$ l’intersection de $`𝔦`$ et $`𝔪`$. Elle est isotrope pour $`B`$. De plus $`𝔥`$ est l’espace des points fixes d’une af-involution (resp. f-involution), $`\sigma `$, de $`𝔪`$. Si $`B`$ est réelle et spéciale, celle-ci est antilinéaire et $`𝔥`$ est une forme réelle de $`𝔪`$. (iii) L’intersection $`𝔦_𝔞`$ de $`𝔞`$ et $`𝔦`$ est Lagrangienne pour la restriction de $`B`$ à $`𝔞`$. (iv) On a $`𝔦=𝔥𝔦_𝔞𝔫`$. Réciproquement si une sous-algèbre de Lie réelle, $`𝔦`$, de $`𝔤`$ est de la forme ci-dessus, elle est Lagrangienne pour $`B`$. On dit alors que $`𝔦`$ est sous $`𝔭`$.
Début de la démonstration du Théorème 1 : Soit $`𝔦`$ une sous-algèbre de Lie réelle (resp. complexe) de $`𝔤`$ Lagrangienne pour $`B`$. On note $`𝔯`$ son radical et on pose :
$$𝔫:=\{X𝔯𝔤^{der}|ad_𝔤(X)estnilpotent\}$$
(1.3)
Soit $`𝔥`$ une sous-algèbre de Levi de $`𝔦`$.
###### Lemme 7
L’ensemble $`𝔫`$ est un idéal de $`𝔦`$ et $`[𝔦,𝔯]`$ est contenu dans $`𝔫`$.
Démonstration : Montrons que $`𝔫`$ est un idéal de $`𝔯`$ contenant $`[𝔯,𝔯]`$. En effet, comme $`𝔯`$ est résoluble, dans une base, sur $``$, bien choisie de $`𝔤`$, les $`ad_𝔤(X)`$, $`X𝔯`$ s’écrivent sous forme de matrices triangulaires supérieures. Pour $`X𝔯`$, les entrées de la diagonale de cette matrice sont notées $`\lambda _1(X),\mathrm{},\lambda _p(X)`$, où les $`\lambda _i`$ sont des caractères de $`𝔯`$. Alors $`𝔫`$ est l’intersection des noyaux de ces caractères avec $`𝔤^{der}`$. Donc $`𝔫`$ est un idéal de $`𝔯`$ contenant $`[𝔯,𝔯]`$. Si $`𝔣`$ est une sous-algèbre de Cartan de $`𝔥`$, $`𝔣𝔯`$ est encore une algèbre de Lie résoluble car $`[𝔣,𝔣]=\{0\}`$ et $`[𝔣,𝔯]𝔯`$. Un argument similaire à celui ci-dessus montre que $`[𝔣,𝔯]`$ est contenu dans $`𝔫`$. La réunion de toutes les sous-algèbres de Cartan de $`𝔥`$ est dense dans $`𝔥`$, d’après la densité des éléments réguliers (cf.\[Bou\], Ch. VII, Paragraphe 2.2 et Paragraphe 2.3, Théorème 1). Par continuité et densité, on en déduit que $`[𝔥,𝔯]𝔫`$.
###### Lemme 8
Soit $`k`$ un entier compris entre 0 et la dimension réelle (resp. complexe) de $`𝔯/𝔫`$. Il existe un sous-espace réel (resp. complexe), abélien, $`𝔞_k`$, de $`𝔯`$, de dimension $`k`$, tel que : (i) $`𝔞_k𝔫=\{0\}.`$ (ii) $`𝔞_k`$ est formé d’éléments semi-simples de $`𝔤`$. (iii) $`𝔞_k`$ et $`𝔥`$ commutent.
Démonstration : On procède par récurrence sur $`k`$. Si $`k=0`$, le Lemme est clair. Supposons le démontré pour $`k<dim_𝕂(𝔯/𝔫)`$ (où $`𝕂=`$, resp. $``$) et montrons le au rang $`k+1`$. Alors $`𝔥𝔞_k`$ est une algèbre de Lie réductive dans $`𝔤`$, regardée comme réelle (resp. complexe). D’autre part, comme $`𝔫`$ contient $`[𝔦,𝔯]`$ d’après le Lemme précédent, on voit que $`𝔦`$ et donc $`𝔥𝔞_k`$ agit trivialement sur $`𝔯/𝔫`$. Ceci implique que :
$$[𝔥𝔞_k,𝔞_k𝔫]𝔫$$
Donc, $`𝔞_k𝔫`$ est un $`(𝔥𝔞_k)`$-sous-module de $`𝔯`$, qui admet un supplémentaire dans $`𝔯`$ commutant à $`𝔥𝔞_k`$, puisque $`𝔥𝔞_k`$ est réductive dans $`𝔤`$ et que le quotient $`𝔯/𝔞_k𝔫`$ est un $`𝔥𝔞_k`$-module trivial. On choisit un élément non nul de ce supplémentaire, $`X`$. Alors :
$$X𝔯,X𝔥𝔞_k,et[X,𝔥𝔞_k]=\{0\}$$
(1.4)
On écrit $`X=X_s+X_n`$, où $`X_n`$ est un élément de $`𝔤^{der}`$, $`X_s`$ est un élément de $`𝔤`$ commutant à $`X_n`$ tels que $`ad_𝔤X_s`$ est semi-simple et $`ad_𝔤X_n`$est nilpotent. On sait qu’alors $`ad_𝔤X_s`$, $`ad_𝔤X_n`$ sont des polynômes en $`ad_𝔤X`$. Joint à (1.2), cela implique :
$$[X_s,𝔥𝔞_k]=\{0\},[X_n,𝔥𝔞_k]=\{0\}$$
(1.5)
Montrons que $`X_n`$ appartient à $`𝔦`$. Soit $`𝔧`$ une sous-algèbre de Cartan de $`𝔥`$. Alors $`𝔧𝔯`$ est résoluble. On peut donc choisir une base de $`𝔤`$ dans laquelle les $`ad_𝔤Y`$, $`Y𝔧𝔯`$, sont représentés par des matrices triangulaires supérieures. On peut choisir cette base de sorte qu’elle soit la réunion de bases des idéaux simples de $`𝔤`$ avec une base du centre de $`𝔤`$, ce que l’on fait dans la suite. Comme $`ad_𝔤X_n`$ est un polynôme en $`ad_𝔤X`$, et que $`X𝔯`$, il est représenté dans cette base par une matrice triangulaire supérieure. Comme cet endomorphisme est nilpotent, sa diagonale est nulle. On en déduit que, pour tout $`Y𝔧𝔯`$, les composantes de $`X_n`$ et $`Y`$ dans les idéaux simples de $`𝔤`$ sont deux à deux orthogonales pour la forme de Killing de $`𝔤`$. Alors, il résulte de l’orthogonalité, pour $`B`$, du centre de $`𝔤`$ à $`𝔤^{der}`$ et du Lemme 1 (i), que :
$$B(X_n,Y)=0,Y𝔧𝔯$$
En utilisant la densité dans $`𝔥`$ de la réunion de ses sous-algèbres de Cartan, on en déduit que $`X_n`$ est orthogonal à $`𝔦`$ pour $`B`$. Mais $`𝔦`$ est un sous-espace isotrope pour $`B`$ de dimension maximale. Donc $`X_n`$ est élément de $`𝔦`$ comme désiré. Ecrivons $`X_n=H+R`$ avec $`H𝔥`$, $`R𝔯`$. Comme $`X_n`$ commute à $`𝔥`$, d’après (1.3) et que $`[𝔥,𝔯]𝔯`$, on voit que $`H`$ commute à $`𝔥`$. Donc $`H`$ est nul puisque $`𝔥`$ est semi-simple. Finalement $`X_n𝔯`$, et en fait $`X_n𝔫`$, d’après la définition de $`𝔫`$. Comme $`X`$ appartient à un supplémentaire de $`𝔞_k+𝔫`$ dans $`𝔯`$ et que $`X=X_s+X_n`$, on a :
$$X_s𝔯,X_s𝔞_k+𝔫$$
On pose $`𝔞_{k+1}=𝔞_k+𝕂X_s`$. D’après (1.3) et la semi-simplicité de $`ad_𝔤X_s`$, $`𝔞_{k+1}`$ vérifie les propriétés voulues. Suite de la démonstration du Théorème 1: On pose $`𝔦_𝔞:=𝔞_p`$, avec $`p=dim_𝕂𝔯/𝔫`$, de sorte que $`𝔦=𝔥𝔦_𝔞𝔫`$, où $`𝔦_𝔞`$ est formé d’éléments semi-simples de $`𝔤`$ avec :
$$[𝔥,𝔦_𝔞]=\{0\},𝔯=𝔦_𝔞𝔫$$
Comme $`𝔦_𝔞𝔫`$ est résoluble, il existe une sous-algèbre de Borel, $`𝔟`$, de $`𝔤`$, contenant $`𝔦_𝔞𝔫`$. A noter que $`𝔫`$ est contenue dans le radical nilpotent, $`𝔳`$, de $`𝔟`$, d’après la définitionde $`𝔫`$ et les propriétés du radical nilpotent d’une sous-algèbre de Borel. Montrons que $`𝔦_𝔞`$ est contenue dans une sous-algèbre de Cartan de $`𝔤`$, contenue dans $`𝔟`$. En effet, d’après \[Bor\], Proposition 11.15, la sous-algèbre de Borel $`𝔟`$ contenant $`𝔦_𝔞`$, elle contient une sous-algèbre de Borel du centralisateur $`𝔩`$ de $`𝔦_𝔞`$ dans $`𝔤`$. Celle-ci contient une sous-algèbre de Cartan $`𝔧`$ de $`𝔩`$. Celle-ci est aussi une sous-algèbre de Cartan de $`𝔤`$ contenant $`𝔦_𝔞`$ (cf. \[Bou\], Ch. VII, Paragraphe 2.3, Proposition 10). Soit $`𝔲`$ la somme des sous espaces poids de $`𝔦_𝔞`$, dans $`𝔳`$, , pour des poids non nuls. Alors $`𝔭:=𝔩𝔲`$ est une sous algèbre parabolique de $`𝔤`$, contenant $`𝔟`$. Comme $`𝔩`$ est réductive, $`𝔪:=𝔩^{der}`$ est semi-simple et le radical de $`𝔭`$ est égal à la somme du centre $`𝔞`$ de $`𝔩`$ avec $`𝔲`$. La définition de $`𝔲`$ montre que $`[𝔭,𝔭]=𝔪𝔲`$, donc le radical nilpotent de $`𝔭`$ est égal à $`𝔲`$ (cf. \[Bou\], Ch. I, Paragraphe 5.3, Théorème 1). Comme $`𝔦=𝔥𝔦_𝔞𝔫`$, que $`𝔦_𝔞𝔫`$ est contenu dans $`𝔟`$ et que $`𝔥`$ est contenu dans $`𝔩`$, on a :
$$𝔦𝔭$$
Or $`𝔭`$ (resp. $`𝔲`$) est la somme de ses intersections $`𝔭_i`$ (resp. $`𝔲_i`$) avec les idéaux simples $`𝔤_i`$ de $`𝔤`$. Comme $`𝔭_i`$ est orthogonal à $`𝔲_i`$ pour la forme de Killing de $`𝔤_i`$, on en déduit que $`𝔲`$ est orthogonal à $`𝔭`$ pour $`B`$ (cf. Lemme 1 (i)). Comme $`𝔦`$ est un sous-espace isotrope pour $`B`$, de dimension maximale et contenu dans $`𝔭`$, $`𝔲`$ est inclus dans $`𝔦`$. Il résulte alors de la définition de $`𝔫`$, que $`𝔲`$ est contenu dans $`𝔫`$. Par suite, on a :
$$𝔦=𝔲(𝔦𝔩),𝔫=𝔲(𝔫𝔩)$$
(1.6)
Remarquons que $`𝔞`$ contient $`𝔦_𝔞`$. On a $`𝔳=𝔲(𝔳𝔪)`$. Comme $`𝔫𝔳`$ et $`𝔲𝔫`$, on en déduit que $`𝔫=𝔲(𝔫𝔪)`$. On déduit alors de (1.6) que : $`𝔫𝔩=𝔫𝔪`$. Finalement, on a :
$$𝔦=𝔥𝔦_𝔞(𝔫𝔪)𝔲$$
Alors, posant :
$$𝔦^{}:=𝔦𝔪,$$
on a :
$$𝔦^{}=𝔥(𝔫𝔪)$$
C’ est une sous-algèbre isotrope de $`𝔪`$ pour la restriction de $`B`$ à $`𝔪`$, donc, d’après le Corollaire du Lemme 2 (ii), de dimension réelle inférieure ou égale à la dimension complexe de $`𝔪`$. De même, $`𝔦_𝔞`$ est un sous espace isotrope de $`𝔞`$ pour la restriction de $`B`$ à $`𝔞`$. D’après le Corollaire du Lemme 2, la restriction de $`B`$ à $`𝔞`$ est de signature $`(dim_{}𝔞,dim_{}𝔞)`$ (resp. est non dégénérée). Il en résulte que la dimension réelle de $`𝔦_𝔞`$ est inférieure ou égale à $`dim_{}𝔞`$. Mais $`dim_{}𝔤=dim_{}𝔪+dim_{}𝔞+dim_{}𝔲`$. Comme $`dim_{}𝔦=dim_{}𝔤`$, on déduit de ce qui précède que l’on a :
$$dim_{}𝔦^{}=dim_{}𝔪,dim_{}𝔦_𝔞=dim_{}𝔞$$
(1.7)
###### Lemme 9
L’algèbre de Lie $`𝔫^{}:=𝔫𝔪`$ est réduite à zéro et $`𝔥`$ a la forme indiquée dans le Théorème.
Démonstration : Si $`𝔣`$ est une sous-algèbre de Cartan de $`𝔥`$, $`𝔣+𝔫^{}+𝔦𝔫^{}`$ est une sous-algèbre de Lie réelle et résoluble de $`𝔪`$. On peut donc choisir une base de $`𝔪`$, réunion de bases des idéaux simples de $`𝔪`$ , telle que, pour tout $`X𝔣+𝔫^{}+𝔦𝔫^{}`$, $`ad_𝔪X`$ soit représenté, dans cette base, par une matrice triangulaire supérieure. De plus, si $`X`$ est élément de $`𝔫^{}+𝔦𝔫^{}`$, les éléments diagonaux de cette matrice sont nulles. On voit, grâce au Lemme 1 (i), que $`𝔫^{}+i𝔫^{}`$ est orthogonal à $`𝔣+𝔫^{}`$, pour la restriction, $`B^{}`$, de $`B`$ à $`𝔪`$. Ceci étant vrai pour tout $`𝔣`$, $`𝔫^{}+i𝔫^{}`$ est orthogonal à $`𝔦^{}`$ ($`=𝔥+𝔫^{}`$), pour $`B^{}`$. Mais $`B^{}`$ est non dégénérée, d’après le Corollaire du Lemme 2, donc, d’après le Lemme 1 (ii) et (1.7), $`𝔦^{}`$ est un sous-espace isotrope de $`𝔪`$, pour $`B^{}`$, de dimension maximale. Il en résulte que $`𝔫^{}+i𝔫^{}`$ est contenu dans $`𝔦^{}`$. Mais $`𝔫^{}+i𝔫^{}`$ est aussi contenu dans $`𝔳𝔤^{der}`$. Finalement $`𝔫^{}+i𝔫^{}`$ est contenudans l’intersection de $`𝔦^{}`$ avec $`𝔫`$, d’après la définition de celui-ci. Mais, comme $`𝔦^{}`$ est contenu dans $`𝔪`$, $`𝔦^{}𝔫=𝔫^{}`$. Alors on a : $`𝔫^{}+i𝔫^{}𝔫^{}`$, c’est à dire que $`𝔫^{}`$ est un sous-espace vectoriel complexe de $`𝔤`$, ce qui bien sur évident dans le cas complexe.
Soit $`𝔥_j`$, $`j=1,\mathrm{},r`$, les idéaux simples de $`𝔥`$. Comme $`𝔥_ji𝔥_j`$ est un idéal de l’algèbre de Lie simple réelle $`𝔥_j`$, il y a deux possiblités pour $`𝔥_j`$. Ou bien $`𝔥_ji𝔥_j=\{0\}`$, et alors $`𝔥_j+i𝔥_j`$ est une algèbre de Lie semi-simple complexe dont $`𝔥_j`$ est une forme réelle. Ou bien $`𝔥_ji𝔥_j=𝔥_j`$ et $`𝔥_j`$ est une sous-algèbre simple complexe de $`𝔤`$. On remarquera que cette deuxième possibilité est exclue, si $`B`$ est spéciale, puique $`𝔥_j`$ serait alors semi-simple complexe et isotrope pour $`B`$. On suppose que, pour $`j=1,\mathrm{},p`$, $`𝔥_ji𝔥_j=\{0\}`$, et que pour $`j=p+1,\mathrm{},r`$, $`𝔥_ji𝔥_j=𝔥_j`$. Si $`j=1,\mathrm{},p`$, on note $`𝔨_j=𝔥_ji𝔥_j`$. Dans le cas complexe, i.e. si $`B`$ est $``$-bilinéaire, $`𝔥_j`$ est toujours complexe et $`p=0`$. Si $`j=p+1,\mathrm{},r`$, on note $`𝔨_j`$ la somme des projections de $`𝔥_j`$ dans les idéaux simples de $`𝔪`$. On note aussi $`𝔨_j^{}=𝔥_j+i𝔥_j`$, pour $`j=1,\mathrm{},r`$. On note $`𝔨=_{j=1,\mathrm{},r}𝔨_j`$ et $`𝔨^{}=_{j=1,\mathrm{},r}𝔨_j^{}=𝔥+𝔦𝔥`$ , qui est contenu dans $`𝔨`$. Par ailleurs, deux éléments, $`X`$ et $`Y`$, de $`𝔪`$ commutent si et seulement $`X`$ et $`iY`$ commutent (resp. $`X`$ commute à chacune des projections de $`Y`$ dans les idéaux simples de $`𝔪`$). Il en résulte que, pour $`jl`$, $`𝔨_j`$ commute à $`𝔥_l`$, donc $`𝔨_j(_{lj}𝔨_l)`$ est contenu dans le centre de $`𝔨_j`$ , qui est semi-simple complexe. Il en résulte que :
$$𝔨=_{j=1,\mathrm{},r}𝔨_j,𝔨^{}=_{j=1,\mathrm{},r}𝔨_j^{}$$
(1.8)
Alors $`𝔨`$, $`𝔨^{}`$ sont des sous-algèbres de Lie semi-simples complexes de $`𝔪`$. Montrons que :
$$𝔨𝔫^{}=\{0\}$$
(1.9)
En effet $`𝔫^{}`$ est un idéal dans $`𝔦^{}=𝔥+𝔫^{}`$, puisque $`𝔦^{}=𝔦𝔪`$ et $`𝔫^{}`$ est l’intersection de l’idéal $`𝔫`$ de $`𝔦`$ avec $`𝔪`$. C’est donc un $`𝔥`$-module, et aussi un $`𝔨^{}`$-module puisque $`𝔫^{}`$ est un espace vectoriel complexe. Donc $`𝔨𝔫^{}`$ est un sous-$`𝔨^{}`$-module, et aussi une sous-algèbre résoluble de $`𝔨`$. Il est clair que les $`𝔨_j`$ sont des sous-$`𝔨^{}`$-modules de $`𝔨`$, qui n’ont aucun sous-quotient simple en commun. En effet, d’une part $`𝔨_l^{}`$ agit trivialement sur $`𝔨_j`$, si $`jl`$. D’autre part, d’après les définitions, on voit que les sous-quotients simples du $`𝔨_j^{}`$-module $`𝔨_j`$ sont isomorphes à des sous-quotients de $`𝔨_j^{}`$, dont aucun n’est trivial, puique $`𝔨_j^{}`$ est une algèbre de Lie semi-simple. Donc, si $`𝔨𝔫^{}`$ est non nul, il a une intersection non nulle, $`𝔨^{\prime \prime }`$, avec l’un des $`𝔨_j`$, qui est un $`𝔨_j^{}`$-sous-module. Comme $`𝔨_j`$ est une sous-algèbre de Lie de $`𝔪`$, il en va de même de $`𝔨^{\prime \prime }`$, qui est de plus résoluble, puisque c’est le cas de
$`𝔫^{}`$. Si $`j=1,\mathrm{},p`$, $`𝔨_j=𝔨_j^{}`$ et un $`𝔨_j^{}`$-sous-module de $`𝔨_j`$ est isomorphe un idéal de $`𝔨_j`$. Alors $`𝔨𝔫^{}`$ est à la fois semi-simple et résoluble. Une contradiction qui montre (1.9) dans ce cas. Si $`j=p+1,\mathrm{},r`$, $`𝔨_j^{}=𝔥_j`$ est simple, donc l’une des projections de $`𝔨^{\prime \prime }`$ sur un idéal simple de $`𝔪`$ est isomorphe à $`𝔥_j`$. Cette projection étant un morphisme d’algèbres de Lie, il en résulte que l’ algèbre de Lie résoluble $`𝔨^{\prime \prime }`$, admet un quotient semi-simple. Une contradiction qui achève de prouver (1.9). Pour $`j=p+1,\mathrm{},r`$, $`𝔥_j`$ ne peut être contenu dans un idéal simple de $`𝔪`$. En effet, d’après le Corollaire du Lemme 2, la restriction de $`B`$ à $`𝔪`$ est une forme de Manin. D’après le Lemme 1 (ii) et le Lemme 2 (v), la restriction de $`B`$ à un idéal simple de $`𝔪`$ est spéciale, et notre assertion en résulte, car pour $`j=p+1,\mathrm{},r`$, $`𝔥_j`$ est isotrope et semi-simple complexe . Pour $`j=p+1,\mathrm{},r`$, on notera $`n_j`$, le nombre d’idéaux simples de $`𝔪`$ dans lesquels $`𝔥_j`$ a une projection non nulle, et pour $`j=1,\mathrm{},p`$, on pose $`n_j=1`$. On vient de voir que :
$$n_j2,j=p+1,\mathrm{},r$$
(1.10)
Montrons que $`𝔫^{}=\{0\}`$. On a évidemment :
$$dim_{}𝔦^{}=(\underset{j=1,\mathrm{},r}{}dim_{}𝔥_j)+dim_{}𝔫^{}$$
Alors, en posant $`p_j=1`$, pour $`j=1,\mathrm{},p`$ et $`p_j=2`$, pour $`j=p+1,\mathrm{},r`$, on a :
$$dim_{}𝔦^{}=(\underset{j=1,\mathrm{},r}{}p_jdim_{}𝔨_j^{})+dim_{}𝔫^{}$$
(1.11)
Soit $`𝔧_𝔨`$ une sous-algèbre de Cartan de $`𝔨`$, $`𝔟_𝔨`$ une sous-algèbre de Borel de $`𝔨`$, contenant $`𝔧_𝔨`$, de radical nilpotent $`𝔫_𝔨`$. Alors $`𝔟_𝔨𝔫^{}`$ est une algèbre de Lie résoluble, contenue dans $`𝔪`$, donc contenue dans une sous-algèbre de Borel, $`𝔟_𝔪`$, de $`𝔪`$. Alors $`𝔧_𝔨`$ est contenue dans une sous algèbre de Cartan, $`𝔧_𝔪`$, contenue dans $`𝔟_𝔪`$ (voir avant (1.6)). De plus $`𝔫_k`$ est contenu dans le radical nilpotent de $`𝔟_𝔪`$, $`𝔫_𝔪`$, qui vérifie :
$$𝔫_𝔪=\{X𝔟_𝔪|ad_𝔪Xestnilpotent\}$$
En effet, les éléments de $`𝔫_𝔨`$ sont représentés, dans toute représentation de dimension finie de $`𝔫_𝔨`$, et donc de $`𝔟_𝔪`$, par des opérateurs nilpotents. De même, $`𝔫^{}`$ est contenudans $`𝔪`$, car pour tout $`X𝔫^{}`$, $`ad_𝔤X`$, et donc $`ad_𝔪X`$ est nilpotent. On note $`𝔧^{\prime \prime }`$ (resp. $`𝔫^{\prime \prime }`$) un supplémentaire de $`𝔧_𝔨`$ dans $`𝔧_𝔪`$, resp. $`𝔫_k𝔫^{}`$ dans $`𝔫_𝔪`$. Un calcul immédiat montre :
$$dim_{}𝔪=dim_{}𝔨+2dim_{}𝔫^{}+dim_{}𝔧^{\prime \prime }+2dim_{}𝔫^{\prime \prime }$$
(1.12)
En posant $`n_j=1`$ pour $`j=1,\mathrm{},p`$, on a immédiatement :
$$dim_{}(𝔨)=\underset{j=1,\mathrm{},r}{}n_jdim_{}𝔨_j^{}$$
(1.13)
Alors (1.11), joint à la première égalité de (1.7), et à (1.12) , (1.13), implique :
$$2dim_{}𝔫^{}+\underset{j=1,\mathrm{},r}{}p_jdim_{}𝔨_j^{}=2dim_{}𝔫^{}+dim_{}𝔧^{\prime \prime }+2dim_{}𝔫^{\prime \prime }+\underset{j=1,\mathrm{},r}{}n_jdim_{}𝔨_j^{}$$
Comme $`n_j`$ est supérieur ou égal à $`p_j`$, d’après (1.10) et la définition des $`n_j`$, $`p_j`$, on en déduit :
$$p_j=n_j,j=1,\mathrm{},r,et𝔫^{\prime \prime }=𝔧^{\prime \prime }=\{0\}$$
Alors $`𝔨𝔫^{}`$ contient la sous-algèbre de Borel, $`𝔟_𝔪`$, de $`𝔪`$. C’est une sous-algèbre parabolique dont le radical est égal à $`𝔫^{}`$, donc est nilpotent. Elle est donc égale à $`𝔪`$ et son radical nilpotent $`𝔫^{}`$ est réduit à zéro, comme désiré. En outre $`𝔪=𝔨`$, donc les $`𝔨_j`$ sont des idéaux de $`𝔪`$. On pose $`\stackrel{~}{𝔪}_0=_{j=1,\mathrm{},p}𝔨_j`$, $`\stackrel{~}{𝔥}_0=_{j=1,\mathrm{},p}𝔥_j`$ . On pose $`q=rp`$. Pour $`l=1,\mathrm{},q`$, $`𝔨_{p+l}`$ est somme de deux idéaux simples, $`𝔪_l^{}`$, $`𝔪_l^{\prime \prime }`$, car $`n_{p+l}=2`$. La projection de $`𝔥_{l+p}`$ sur chacun de ces idéaux est bijective, sa surjectivité résultant de la définition de $`𝔨_l`$, son injectivité résultant de la simplicité de $`𝔥_{l+p}`$ et de la non nullité de ce morphisme d’algèbres de Lie. Donc $`𝔥_{p+l}:=\{(X,\tau _l(X))|X𝔪_l^{}\}`$, où $`\tau _l`$ est un isomorphisme $``$-linéaire de l’algèbre de Lie $`𝔪_l^{}`$ sur $`𝔪_l^{\prime \prime }`$. Donc $`𝔥`$ a la forme voulue
###### Lemme 10
Aucun poids non nul de $`𝔞`$ dans $`𝔤`$ n’est nul sur $`𝔦_𝔞`$.
Démonstration : Raisonnons par l’absurde et supposons qu’il existe un poids non nul $`\alpha `$ de $`𝔞`$ dans $`𝔤`$, nul sur $`𝔦_𝔞`$. Soit $`H_\alpha 𝔞`$ tel que :
$$K_𝔤(H_\alpha ,X)=\alpha (X),X𝔞$$
(1.14)
Alors $`H_\alpha `$ appartient à l’un des idéaux simples de $`𝔤`$. En effet, soit $`𝔧`$ une sous-algèbre de Cartan de $`𝔤`$, contenant $`𝔞`$. Alors $`\alpha `$ est la restriction à $`𝔞`$ d’une racine $`\beta `$ de $`𝔧`$ dans $`𝔤`$ et l’on a :
$$K_𝔤(H_\alpha ,H_\alpha )>0$$
(1.15)
Soit $`H_\beta 𝔧`$ tel que :
$$K_𝔤(H_\beta ,X)=\beta (X),X𝔧$$
Alors $`𝔧`$ (resp. $`𝔞`$) est la somme directe de ses intersections avec les idéaux simples de $`𝔤`$, et $`H_\beta `$ (resp. $`H_\alpha `$) appartient à l’une de celles-ci. On déduit alors du Lemme 1 (i), qu’ il existe $`\mu `$ , non nul car $`B`$ est non dégénérée, tel que :
$$B(\lambda H_\alpha ,X)=Im(K_𝔤(\lambda \mu H_\alpha ,X)),\lambda X𝔤$$
(1.16)
Comme $`\alpha `$ est nulle sur $`𝔦_𝔞`$, il résulte de (1.14) et (1.16) que $`H_\alpha `$ est orthogonale à $`𝔦_a`$, pour $`B`$. Comme $`B`$ est une forme de Manin, la restriction de $`B`$ à $`𝔞`$ est de signature $`(dim_{}𝔞,dim_{}𝔞)`$. Tenant compte de (1.7), on voit que $`𝔦_𝔞`$ est un sous-espace de $`𝔞`$, isotrope pour $`B`$, de dimension maximale. Alors, ce qui précède montre que $`H_\alpha `$ est contenu dans $`𝔦_𝔞`$. Par ailleurs, si $`\lambda `$ est une racine carrée de $`i\mu ^1`$, $`B(\mu H_\alpha ,\mu H_\alpha )`$ est non nul d’après (1.15) et (1.16). Une contradiction avec le fait que $`𝔦_𝔞`$ est isotrope qui achève de prouver le Lemme. Fin de la démonstration du Théorème 1: Montrons la propriété suivante :
$`Toutesousalg\stackrel{`}{e}breparabolique𝔮de𝔤est\stackrel{´}{e}galeaunormalisateur`$
$`dans𝔤desonradicalnilpotent𝔴`$ (1.17)
D’abord $`𝔮`$ normalise $`𝔴`$. Donc, le normalisateur $`𝔯`$ de $`𝔴`$ dans $`𝔤`$ contient $`𝔮`$. C’est donc une sous-algèbre parabolique de $`𝔤`$, qui contient $`𝔴`$ comme idéal. Compte tenu de \[Bou\], Ch. I, Paragraphe 5.3, Remarque 2, $`𝔴`$ est contenu dans le radical nilpotent, $`𝔵`$, de $`𝔯`$. On a alors : $`𝔭𝔯,𝔴𝔵`$, d’où l’on déduit facilement que $`𝔵=𝔴`$ et $`𝔯=𝔮`$, comme désiré. Donc $`𝔮`$ est bien le normalisateur de $`𝔫`$. Ceci achève de prouver (1.17). Montrons que $`𝔫`$ est le radical nilpotent de $`𝔦`$. En effet, comme $`𝔫`$ est somme de sous-espaces poids sous $`𝔞`$ et qu’aucun de ces poids n’est nul sur $`𝔦_𝔞`$, d’après le Lemme précédent, on a :
$$[𝔦_𝔞,𝔫]=𝔫$$
Donc $`[𝔦,𝔦]=𝔥𝔫`$ et l’intersection de $`[𝔦,𝔦]`$ avec le radical $`𝔯=𝔦_𝔞+𝔫`$ de $`𝔦`$ est égal à $`𝔫`$. Donc, d’après \[Bou\] Ch. I, Paragraphe 5.3, Théorème 1, $`𝔫`$ est bien le radical nilpotent de $`𝔦`$. On a donc montré que $`𝔦`$ s’écrit de la manière voulue, pour une décomposition de Langlands particulière du normalisateur, $`𝔭`$, de $`𝔫`$ . Si $`𝔭=𝔩^{}𝔫`$ est une autre décomposition de Langlands de $`𝔭`$, $`𝔩`$ et $`𝔩^{}`$ sont isomorphes, puisqu’elles sont toutes les deux isomorphes à $`𝔤/𝔫`$. Comme $`𝔦`$ contient $`𝔫`$, les intersections de $`𝔦`$ avec $`𝔩`$ et $`𝔩^{}`$ se correspondent dans cet isomorphisme, et la décomposition de $`𝔦`$ qu’on en déduit, relativement à cette nouvelle décomposition de Langlands de $`𝔭`$, a les propriétés voulues. Etudions la partie réciproque du Théorème. Une sous-algèbre parabolique de $`𝔤`$ est la somme de ses intersections avec les idéaux simples de $`𝔤`$. En outre, elle est orthogonale à son radical nilpotent, pour la forme de Killing de $`𝔤`$. On conclut que si $`𝔦`$ a une décomposition comme dans l’énoncé, elle est isotrope pour $`B`$, et de dimension réelle égale à la dimension complexe de $`𝔤`$.
###### Définition 3
On rappelle qu’une sous-algèbre de Cartan d’une algèbre de Lie semi-simple réelle est une sous-algèbre de Cartan fondamentale si et seulement si elle contient des éléments réguliers dont l’image par la représen- tation adjointe n’a que des valeurs propres imaginaires pures. Cela équivaut au fait qu’aucune racine de cette sous-algèbre de Cartan n’est réelle. Une sous-algèbre de Cartan d’une algèbre de Lie réelle est dite fondamentale si sa projection dans une sous-algèbre de Levi, parallèlement au radical, est une sous-algèbre de Cartan fondamentale de cette algèbre de Lie semi-simple réelle.
Comme toutes les sous-algèbres de Cartan fondamentales d’une algèbre de Lie réelle semi-simple sont conjuguées entre elles par des automorphismes intérieurs, il en va de même pour les sous-algèbres de Cartan fondamentales d’une algèbre de Lie réelle. En effet il suffit d’adapter la preuve du fait que (ii) implique (i), dans \[Bou\], Ch. VII, Paragraphe 3.5, Proposition 5, en remarquant pour cela que tout automorphisme intérieur d’une algèbre de Levi d’une algèbre de Lie réelle, s’étend en un automorphisme intérieur de l’algèbre de Lie.
###### Lemme 11
On conserve les hypothèses et notations du Théorème 1. (i) Si $`𝔣`$ est une sous-algèbre de Cartan de $`𝔥`$, il existe des éléments réguliers de $`𝔤`$ contenus dans $`𝔣𝔦_𝔞`$. Le centralisateur dans $`𝔤`$, $`𝔧`$, de $`\stackrel{~}{𝔣}:=𝔣𝔦_𝔞`$, est une sous-algèbre de Cartan de $`𝔤`$, contenue dans $`𝔩`$, vérifiant $`𝔧=(𝔧𝔪)𝔞`$. (ii) Si $`𝔣`$ est une sous-algèbre de Cartan de $`𝔥`$, $`𝔣𝔦_𝔞`$ est une sous-algèbre de Cartan de $`𝔦`$. (iii) Toute sous-algèbre de Cartan de $`𝔦`$ (resp. sous-algèbre de Cartan de $`𝔦`$ contenue dans $`𝔥+𝔦_𝔞`$) est conjuguée, par un automorphisme intérieur de $`𝔦`$, (resp. égale) à une algèbre de ce type. (iv) Soit $`\stackrel{~}{𝔣}^{}`$ une sous-algèbre de Cartan de $`𝔦`$. Il existe une unique décomposi- tion de Langlands $`𝔩^{}+𝔫`$ de $`𝔭`$, telle que $`𝔩^{}`$ contienne $`\stackrel{~}{𝔣}^{}`$. Alors, notant $`𝔣^{}:=\stackrel{~}{𝔣}^{}𝔩^{der}`$, on a $`\stackrel{~}{𝔣}^{}=𝔣^{}+(𝔦𝔞^{})`$, où $`𝔞^{}`$ est le centre de $`𝔩^{}`$. De plus $`\stackrel{~}{𝔣}^{}`$ est une sous-algèbre de Cartan fondamentale de $`𝔦`$, si et seulement si $`𝔣^{}`$ est une sous-algèbre de Cartan fondamentale de $`𝔥^{}:=𝔦𝔩^{der}`$
Démonstration : Montrons (i). On raisonne par l’absurde. On note $`𝔧^{}`$ une sous-algèbre de Cartan de $`𝔤`$, qui contient $`𝔣𝔞`$. Celle-ci existe puisque les éléments de $`𝔣𝔞`$ sont semi-simples. Supposons qu’aucun élément de $`\stackrel{~}{𝔣}:=𝔣𝔦_𝔞`$ ne soit régulier dans $`𝔤`$. Alors, pour tout $`X\stackrel{~}{𝔣}`$, il existe une racine $`\alpha _X`$ de $`𝔧^{}`$ dans $`𝔤`$, nulle sur $`X`$. Pour une racine donnée, l’intersection de son noyau avec $`\stackrel{~}{𝔣}`$ est un fermé de $`\stackrel{~}{𝔣}`$. Notre hypothèse montre que $`\stackrel{~}{𝔣}`$ est la réunion de ces fermés. Il en résulte que l’un de ces sous-espaces vectoriels est d’intérieur non vide, donc égal à $`\stackrel{~}{𝔣}`$. Cela signifie qu’une racine de $`𝔧^{}`$ s’annule sur $`\stackrel{~}{𝔣}`$. Alors, d’après le Lemme 7, dont on vérifie aisément qu’il est valable pour toute décomposition de Langlands de $`𝔭`$ , celle-ci doit être nulle sur $`𝔞`$. C’est donc une racine de $`𝔧^{}`$ dans $`𝔪`$, qui ne peut être nulle sur $`𝔣`$. Une contradiction qui prouve la première partie de (i). Le centralisateur $`𝔧`$ de $`\stackrel{~}{𝔣}`$ est donc une sous-algèbre de Cartan. Par ailleurs $`𝔧`$ contient $`𝔞`$. On en déduit la deuxième partie de (i). D’après le Lemme 7, le nilespace de $`\stackrel{~}{𝔣}`$ dans $`𝔫`$ est réduit à zéro. Comme $`𝔣`$ est une sous-algèbre de Cartan de $`𝔥`$, le nilespace de $`\stackrel{~}{𝔣}`$ dans $`𝔥`$ est égal à $`𝔣`$. Finalement le nilespace de $`\stackrel{~}{𝔣}`$ dans $`𝔦`$ est égal à $`\stackrel{~}{𝔣}`$. Alors (ii) résulte de \[Bou\], Ch. VII, Paragraphe 2.1, Proposition 3. Montrons (iii). Soit $`\stackrel{~}{𝔣}^{}`$ une autre sous-algèbre de Cartan de $`𝔦`$. La projection, $`\stackrel{~}{𝔣}^{\prime \prime }`$, de $`\stackrel{~}{𝔣}^{}`$ sur $`\stackrel{~}{𝔥}:=𝔥+𝔦_𝔞`$, parallèlement à $`𝔫`$, est une sous-algèbre de Cartan de $`\stackrel{~}{𝔥}`$ (cf. \[Bou\], Ch. VII, Paragraphe 2.1, Corollaire 2 de la Proposition 4), donc de la forme $`𝔣^{}+𝔦_𝔞`$, où $`𝔣^{}`$ est une sous-algèbre de Cartan de $`𝔥`$. Alors $`\stackrel{~}{𝔣}^{}`$ et $`\stackrel{~}{𝔣}^{\prime \prime }`$ sont deux sous algèbres de Cartan de $`𝔦`$, ayant la même projection sur $`\stackrel{~}{𝔥}`$, parallèlement à $`𝔫`$, donc conjuguées par un automorphisme intérieur de $`𝔦`$, d’après \[Bou\], Ch. VII, Paragraphe 3.5, Proposition 5 (voir aussi après la Définition 3). Si de plus $`\stackrel{~}{𝔣}^{}`$ est contenue dans $`\stackrel{~}{𝔥}`$, le raisonnement ci-dessus montre qu’elle a la forme indiquée. Ce qui prouve (iii). Prouvons (iv). Grâce à (iii), on se ramène, par conjugaison, au cas où $`\stackrel{~}{𝔣}^{}`$ est contenue dans $`𝔩`$ et comme dans (i). Si $`𝔩^{}+𝔫`$ est une décomposition de Langlands de $`𝔭`$, où $`𝔩^{}`$ contient $`\stackrel{~}{𝔣}^{}`$, $`𝔩^{}`$ contient un élément régulier de $`𝔤`$, contenu dans $`\stackrel{~}{𝔣}^{}`$, dont le centralisateur dans $`𝔤`$ est une sous-algèbre de Cartan de $`𝔤`$, contenue dans $`𝔩^{}`$. Celle-ci est égale au centralisateur dans $`𝔤`$ de $`\stackrel{~}{𝔣}^{}`$. D’où l’unicité de $`𝔩^{}`$, grâce aux propriétés des décompositions de Langlands (cf. Lemme 3 (i)). L’assertion sur les sous-algèbre de Cartan fondamentales est claire car $`𝔥^{}`$ est une sous-algèbre de Levi de $`𝔦`$, d’après le Théorème 1.
## 2 Triples de Manin pour une algèbre de Lie réductive complexe : Descente
Dans toute la suite $`𝔤`$ désignera une algèbre de Lie réductive complexe. On fixe, $`𝔧_0`$, une sous algèbre de Cartan de $`𝔤`$, $`𝔟_0`$ une sous-algèbre de Borel de $`𝔤`$, contenant $`𝔧_0`$. On note $`𝔟_0^{}`$ la sous-algèbre de Borel opposée à $`𝔟_0`$, relativement à $`𝔧_0`$.
###### Définition 4
Un triple de Manin pour $`𝔤`$ est un triplet $`(B,𝔦,𝔦^{})`$, où $`B`$ est une forme de Manin sur $`𝔤`$, $`𝔦`$ et $`𝔦^{}`$ sont des sous-algèbres de Lie réelles de $`𝔤`$, isotropes pour $`B`$, telles que $`𝔤=:𝔦𝔦^{}`$. La signature de $`B`$ étant égale à $`(dim_{}𝔤,dim_{}𝔤)`$, $`𝔦`$ et $`𝔦^{}`$ sont Lagrangiennes. Un s-triple est un triple de Manin pour lequel la forme est spéciale. Si $`𝔦`$ est sous $`𝔭`$ et $`𝔦^{}`$ est sous $`𝔭^{}`$, on dit que le triple de Manin est sous $`(𝔭,𝔭^{})`$
###### Remarque 1
D’après le Lemme 2 (v), si $`𝔤`$ est simple, la notion de s-triple et de triple de Manin coincident.
On note $`G`$ le groupe connexe , simplement connexe, d’algèbre de Lie $`𝔤`$. Si $`𝔰`$ est une sous-algèbre de $`𝔤`$, on note $`S`$ le sous-groupe analytique de $`G`$, d’algèbre de Lie $`𝔰`$. Comme $`𝔤`$ est complexe, les sous-groupes paraboliques de $`𝔤`$ sont connexes (cf. \[Bor\], Théoréme 11.16). Donc, si $`𝔭`$ est une sous-algèbre parabolique de $`𝔤`$, $`P`$ est le sous-groupe parabolique de $`G`$, d’algèbre de Lie $`𝔭`$. On remarque que $`G`$ agit sur l’ensemble des triples de Manin , en posant, pour tout triple de Manin $`(B,𝔦,𝔦^{})`$ et tout $`gG`$ :
$$g(B,𝔦,𝔦^{}):=(B,Adg(𝔦),Adg(𝔦^{}))$$
Notre but est construire, par récurrence sur la dimension de $`𝔤^{der}`$, tous les triples de Manin modulo cette action de $`𝔤`$.
###### Proposition 1
Tout triple de Manin est conjugué, sous l’action de $`G`$, à un triple de Manin $`(B,𝔦,𝔦^{})`$ sous $`(𝔭,𝔭^{})`$ , avec $`𝔟_0𝔭`$ et $`𝔟_0^{}𝔭^{}`$ (un tel triple de Manin sera dit standard ). De plus $`𝔭`$ et $`𝔭^{}`$ sont uniques.
Démonstration : On rappelle (cf. \[Bor\], Corollaire 14.13) que :
$$L^{}intersectiondedeuxsousalg\stackrel{`}{e}bresdeBorelde𝔤,\underset{¯}{𝔟},\underset{¯}{𝔟}^{},$$
$$contientunesousalg\stackrel{`}{e}bredeCartande𝔤$$
(2.1)
Soit $`(B,\underset{¯}{𝔦},\underset{¯}{𝔦}^{})`$ un triple de Manin sous $`(\underset{¯}{𝔭},\underset{¯}{𝔭}^{})`$. Soit $`\underset{¯}{𝔟}`$ (resp. $`\underset{¯}{𝔟}^{}`$) une sous-algèbre de Borel de $`𝔤`$, contenue dans $`\underset{¯}{𝔭}`$ (resp. $`\underset{¯}{𝔭}^{})`$. On a $`𝔤=\underset{¯}{𝔦}+\underset{¯}{𝔦}^{}\underset{¯}{𝔭}+\underset{¯}{𝔭}^{}`$. Donc $`\underset{¯}{𝔭}+\underset{¯}{𝔭}^{}`$ est égal à $`𝔤`$ et $`\underset{¯}{P}\underset{¯}{P}^{}`$ est ouvert dans $`G`$. Mais $`\underset{¯}{P}\underset{¯}{P}^{}`$ est réunion de $`(\underset{¯}{B},\underset{¯}{B}^{})`$-doubles classes, qui sont en nombre fini (Bruhat). L’une de ces doubles classes contenues dans $`\underset{¯}{P}\underset{¯}{P}^{}`$ doit donc être ouverte. Soit $`p\underset{¯}{P}`$ et $`p^{}\underset{¯}{P}^{}`$, tels que $`\underset{¯}{B}pp^{}\underset{¯}{B}^{}`$ soit un ouvert de $`G`$. On pose $`B_1=p^1\underset{¯}{B}p`$, $`B_1^{}=p^{}\underset{¯}{B}^{}p^1`$. Alors le sous-groupe de Borel de $`G`$, $`B_1`$ (resp. $`B_1^{}`$), est contenu dans $`P`$ (resp. $`P^{}`$) et $`B_1B_1^{}`$ est ouvert dans $`G`$. Donc, on a $`𝔟_1+𝔟_1^{}=𝔤`$ et l’intersection de $`𝔟_1`$ et $`𝔟_1^{}`$ contient une sous-algèbre de Cartan de $`𝔤`$, $`𝔧_1`$ (cf. (2.1)). Pour des raisons de dimension, cette intersection est réduite à $`𝔧_1`$. Alors $`𝔟_1`$ et $`𝔟_1^{}`$ sont opposées relativement à $`𝔧_1`$ . D’après \[Bor\], Proposition 11.19, il existe $`g^{}G`$ tel que $`Adg^{}(𝔟_1)=𝔟_0`$, $`Adg^{}(𝔧_1)=𝔧_0`$. Alors $`Adg^{}(𝔟_1^{})`$ est égal à $`𝔟_0^{}`$. Alors, notant $`𝔦=Adg^{}(\underset{¯}{𝔦})`$, $`𝔦^{}=Adg^{}(\underset{¯}{𝔦}^{})`$, on voit que $`(B,𝔦,𝔦^{})`$ vérifie les propriétés voulues. L’unicité de $`𝔭`$ résulte du fait que deux sous-algèbres paraboliques de $`𝔤`$, conjuguées par un élément de $`G`$ et contenant une même sous-algèbre de Borel, sont égales (cf. \[Bor\], Corollaire 11.17). On fixe désormais $`𝔭`$ (resp. $`𝔭^{}`$) une sous-algèbre parabolique de $`𝔤`$, contenant $`𝔟_0`$ (resp. $`𝔟_0^{})`$. On note $`𝔭=𝔩𝔫`$ (resp. $`𝔭^{}=𝔩^{}𝔫^{}`$) la décomposition de Langlands de $`𝔭`$ (resp. $`𝔭^{}`$) telle que $`𝔩`$ (resp. $`𝔩^{}`$) contienne $`𝔧_0`$ (cf. Lemme 3). On note $`𝔪=𝔩^{der}`$, $`𝔞`$ le centre de $`𝔩`$. Si $`𝔦`$ est une sous-algèbre de Lie réelle de $`𝔤`$, Lagrangienne pour une forme de Manin, on notera $`𝔥=𝔦𝔪`$, $`𝔦_𝔞=𝔦𝔞`$, $`\stackrel{~}{𝔥}=𝔥𝔦_a`$. On introduit des notations similaires pour $`𝔭^{}`$. Comme $`𝔟_0𝔭`$ (resp. $`𝔟_0^{}𝔭^{}`$), $`𝔫`$ (resp. $`𝔫^{}`$) est contenu dans le radical nilpotent de $`𝔟_0`$ (resp. $`𝔟_0^{}`$). Ces derniers sont d’intersection réduite à zéro, donc :
$$𝔫𝔫^{}=\{0\}$$
(2.2)
Décomposant $`𝔭𝔭^{}`$ en sous-espaces poids sous $`𝔧_0`$, on voit que :
$$𝔭𝔭^{}=(𝔩𝔩^{})(𝔫𝔩^{})(𝔫^{}𝔩).$$
(2.3)
###### Proposition 2
(i) Si un élément de $`G`$ conjugue deux triples de Manin sous $`(𝔭,𝔭^{})`$, c’est un élément de $`PP^{}`$. (ii) Le groupe $`LL^{}`$ est égal au sous-groupe analytique de $`G`$, d’algèbre de Lie $`𝔩𝔩^{}`$. Notons $`N_L^{}`$ (resp. $`N_L^{}`$) , le sous-groupe analytique de $`G`$, d’algèbre de Lie $`𝔫𝔩^{}`$ (resp. $`𝔫^{}𝔩`$). Alors on a :
$$PP^{}=(LL^{})N_L^{}N_L^{}$$
De plus $`N_L^{}`$ et $`N_L^{}`$ commutent entre eux.
Démonstration : Si $`(B,𝔦,𝔦^{})`$ et $`(B,\underset{¯}{𝔦},\underset{¯}{𝔦}^{})`$ sont deux triples de Manin sous $`(𝔭,𝔭^{})`$, conjugués par un élément, $`g`$, de $`G`$, celui-ci conjugue le radical nilpotent de $`𝔦`$ avec celui de $`\underset{¯}{𝔦}`$, donc normalise $`𝔫`$, puisque les deux triples de Manin sont sous $`(𝔭,𝔭^{})`$. Mais un élément du normalisateur, $`Q`$, dans $`G`$ de $`𝔫`$, normalise le normalisateur dans $`𝔤`$ de $`𝔫`$, c’est à dire $`𝔭`$, comme on l’a vu plus haut ( cf. (1. 17)). Comme $`P`$ est connexe, les éléments de $`Q`$ normalisent $`P`$. Donc $`Q`$ est inclus dans $`P`$ et $`gP`$. De même, on a $`gP^{}`$. D’où (i) Montrons (ii). Il est clair que $`PP^{}`$ est un sous-groupe de Lie de $`G`$, d’algèbre de Lie $`𝔭𝔭^{}`$. On a :
$$[𝔫𝔩^{},𝔫^{}𝔩][𝔫,𝔩][𝔫^{},𝔩^{}]𝔫𝔫^{}$$
Donc $`N_L^{}`$ et $`N_L^{}`$ commutent entre eux, d’après (2.1). Alors $`(LL^{})^0N_L^{}N_L^{}`$ est un sous-goupe ouvert et connexe de $`PP^{}`$, donc on a :
$$(PP^{})^0=(LL^{})^0N_L^{}N_L^{}$$
(2.4)
Soit $`gPP^{}`$. Alors $`Adg(𝔧_0)`$ est une sous-algèbre de Cartan de $`𝔤`$, contenue dans $`𝔭𝔭^{}`$, c’est donc une sous-algèbre de Cartan de $`𝔭𝔭^{}`$ (cf \[Bou\], Ch. VIII, Paragraphe 2.1, Exemple 3), donc conjugué, par un élément $`g^{}`$ de $`(PP^{})^0`$, à $`𝔧_0`$, puisqu’il s’agit d’algèbres de Lie complexes. Donc $`Adg^{}g(𝔧_0)=𝔧_0`$ et $`g^{}g`$ est un élément de $`PP^{}`$. En utilisant la décomposition de Bruhat de $`G`$ et $`P`$, pour $`B_0`$, on voit que $`g^{}g`$ centralise le centre $`𝔞`$ de $`𝔩`$. De même on voit que $`g^{}g`$ centralise le centre $`𝔞^{}`$ de $`𝔩^{}`$. Donc $`g^{}g`$ est un élément du centralisateur, $`L^{\prime \prime }`$, de $`𝔞+𝔞^{}`$ dans $`G`$. Mais $`P^{\prime \prime }:=L^{\prime \prime }(N_L^{}N_L^{})`$, est une décomposition de Langlands du sous-groupe parabolique de $`G`$, d’algèbre de Lie :
$$𝔭^{\prime \prime }=(𝔭^{}𝔩)𝔫=(𝔩𝔩^{})(𝔩𝔫^{})𝔫$$
Or $`P^{\prime \prime }`$ est connexe, puisque $`G`$ est complexe. Donc $`L^{\prime \prime }`$ est connexe. Par ailleurs, il contient $`LL^{}`$ et a même algèbre de Lie que $`LL^{}`$. Donc on a :
$$L^{\prime \prime }=LL^{}=(LL^{})^0$$
(2.5)
On conclut alors que $`g^{}g(LL^{})^0`$. Donc $`g`$ est un élément de $`(PP^{})^0`$. Ce qui précède montre que :
$$PP^{}=(PP^{})^0$$
On achève de prouver $`(ii)`$, grâce à (2.4) et (2.5). Le Lemme suivant est une conséquence facile de résultats de Gantmacher (cf. \[G\]).
###### Lemme 12
Si $`\sigma `$ et $`\sigma ^{}`$ sont deux automorphismes involutifs et antilinéaires d’une algè- bre de Lie semi-simple complexe, $`𝔪`$, celle-ci contient au moins un élément non nul et invariant par ces deux involutions.
Démonstration : Avec nos hypothèses $`\sigma \sigma ^{}`$ est un automorphisme $``$-linéaire de $`𝔪`$, dont l’espace des points fixes, $`𝔪^{\sigma \sigma ^{}}`$, est un espace vectoriel complexe, non réduit à zéro d’après \[G\], Théorème 28. Mais $`𝔪^{\sigma \sigma ^{}}`$ , est égal à $`\{X𝔪|\sigma (X)=\sigma ^{}(X)\}`$ donc aussi égal à $`𝔪^{\sigma ^{}\sigma }`$. Si $`X𝔪^{\sigma \sigma ^{}}`$, on a donc $`\sigma ^{}(\sigma (X))=X`$, soit encore $`\sigma ^{}(\sigma (X))=\sigma (\sigma (X))`$. Donc $`\sigma (X)`$ est élément de $`𝔪^{\sigma \sigma ^{}}`$. Par suite $`\sigma `$, restreint à $`𝔪^{\sigma \sigma ^{}}`$ est une involution antilinéaire de $`𝔪^{\sigma \sigma ^{}}`$. L’ensemble de ses points fixes est une forme réelle de $`𝔪^{\sigma \sigma ^{}}`$, donc il est non réduit à zéro. Mais cet ensemble est égal à $`𝔪^\sigma 𝔪^\sigma ^{}`$.
###### Proposition 3
Si $`\sigma `$ et $`\sigma ^{}`$ sont deux af-involutions d’une algèbre de Lie semi-simple complexe, $`𝔪`$, elle contient au moins un élément non nul et invariant par ces deux involutions.
Démonstration : On note $`𝔪_j`$, $`j=1,\mathrm{},r`$, les idéaux simples de $`𝔪`$. On définit une involution $`\theta `$ de $`\{1,\mathrm{},r\}`$ caractérisée par : $`\sigma (𝔪_j)=𝔪_{\theta (j)},j=1,\mathrm{},r`$. Nous allons d’abord étudier le cas suivant :
$$Ilexistejtelque\theta (j)=\theta ^{}(j)=j$$
(2.6)
Dans ce cas, la restriction de $`\sigma `$ et $`\sigma ^{}`$ à $`𝔪_j`$, sont deux automorphismes involutifs et antilinéaires de $`𝔪`$, d’après le Corollaire du Lemme 6, qui ont des points fixes non nuls en commun, d’après le Lemme précédent. La Proposition en résulte, dans ce cas. Supposons maintenant :
$$Ilexistejtelque\theta (j)=\theta ^{}(j)j$$
(2.7)
On note $`j^{}:=\theta (j)`$. Il est clair que : $`(𝔪_j\times 𝔪_j^{})^\sigma =\{(X,\sigma (X))|X𝔪_j\}`$ et de même pour $`(𝔪_j\times 𝔪_j^{})^\sigma ^{}`$. Il existe un élément non nul, $`X`$ de $`𝔪_j`$ tel que $`(\sigma _{}^{}{}_{}{}^{1}\sigma )(X)=X`$, car $`\sigma _{}^{}{}_{}{}^{1}\sigma `$ est un automorphisme $``$-linéaire de $`𝔪_j`$ (cf. \[G\], Théorème 28). Alors $`(X,\sigma (X))`$ est un élément non nul de $`(𝔪_j\times 𝔪_j^{})^\sigma (𝔪_j\times 𝔪_j^{})^\sigma ^{}`$ , ce qui prouve la Proposition dans ce cas. Il nous reste à étudier le cas suivant :
$$Pourtoutj,\theta (j)\theta ^{}(j)$$
(2.8)
On construit, pour tout $`j`$, par récurrence sur $`n`$, une suite $`j_1=j,j_2,\mathrm{},j_n,\mathrm{}`$ , telle que :
$$pourtoutn,j_{n+1}j_n$$
(2.9)
$$pourtoutn,j_{n+1}=\theta (j_n)ou\theta ^{}(𝔧_n)$$
(2.10)
Plus précisément, on pose
$$j_2=\theta (1)si\theta (1)1,j_2=\theta ^{}(1)sinon$$
(2.11)
et, pour $`n2`$, on pose :
$$j_{n+1}=\theta (j_n)sij_n=\theta ^{}(j_{n1})et\theta (j_n)j_n$$
$$j_{n+1}=\theta ^{}(j_n)sij_n=\theta ^{}(j_{n1})et\theta (j_n)=j_n$$
$$j_{n+1}=\theta ^{}(j_n)sij_n=\theta (j_{n1})et\theta ^{}(j_n)j_n$$
(2.12)
$$j_{n+1}=\theta (j_n)sij_n=\theta (j_{n1})et\theta ^{}(j_n)=j_n$$
Ces relations définissent la suite $`(j_n)`$, car, à cause de (2.8), on a nécessaire- ment $`\theta (j_n)\theta ^{}(j_n)`$ et $`\theta (j_{n1})\theta ^{}(j_{n1})`$. Par ailleurs les relations (2.9) et (2.10) sont vérifiées, la première résultant d’une récurrence immédiate. On obtient également les relations suivantes :
$$Pourn2,si\theta (j_n)j_netsi\theta ^{}(j_n)j_nona:$$
$$(j_{n1},j_n,j_{n+1})est\stackrel{´}{e}gal\stackrel{`}{a}(\theta (j_n),j_n,\theta ^{}(j_n))ou\stackrel{`}{a}(\theta ^{}(j_n),j_n,\theta (j_n))$$
(2.13)
$$Pourn2etsi\theta (j_n)=j_nousi\theta ^{}(j_n)=j_nona:j_{n1}=j_{n+1}$$
(2.14)
Après renumérotation des $`𝔤_j`$, on peut supposer que le début de la suite $`(j_n)`$, s’écrit $`j_1=1,j_2=2,\mathrm{},j_p=p,j_{p+1}=k<p`$ On fait d’abord la convention suivante :
$$S^{}ilexistejtelque\theta (j)=jou\theta ^{}(j)=j,onsupposequ^{}on$$
$$l^{}achoisicommepremier\stackrel{´}{e}l\stackrel{´}{e}ment,et,quitte\stackrel{`}{a}\stackrel{´}{e}changerlerole$$
$$de\theta et\theta ^{},qu^{}ilestfix\stackrel{´}{e}par\theta .Onaalors\theta (1)=1,\theta ^{}(1)=2$$
(2.15)
Traitons le cas où $`p=2`$. Alors $`j_1=j_3`$, et (2.8), (2.13) montrent que $`\theta (2)`$ ou $`\theta ^{}(2)`$ est égal à 2. Alors on doit avoir $`\theta (1)=1`$, d’après (2.15), puis $`\theta ^{}(1)=2`$ d’après (2.11). Comme $`\theta ^{}(1)=2`$, on a nécessairement $`\theta (2)=2`$. Dans ce cas, un élément $`(X_1,X_2)𝔪_1𝔪_2`$ est invariant par $`\sigma `$ et $`\sigma ^{}`$ si et seulement si on a :
$$X_1=\sigma (X_1),X_2=\sigma (X_2),X_2=\sigma ^{}(X_1)$$
ce qui équivaut au système :
$$X_1=\sigma (X_1),X_1=(\sigma ^1\sigma \sigma ^{})(X_1),X_2=\sigma ^{}(X_1)$$
Mais la restriction de $`\sigma `$ à $`𝔪_1`$ (resp. $`𝔪_2`$) est un automorphisme involutif antilinéaire, puisque $`\theta (1)=1`$ et $`\theta (2)=2`$ (cf. le Corollaire du Lemme 6). De plus, la restriction de $`\sigma ^{}`$ à $`𝔪_1`$ est soit $``$-linéaire, soit antilinéaire, d’après le Lemme 6. Alors, la restriction de $`\sigma ^1\sigma \sigma ^{}`$ à $`𝔪_1`$ est un automorphisme involutif antilinéaire. Alors, dans le cas $`p=2`$, la Proposition résulte du Lemme 12. On suppose maintenant :
$$p>2$$
(2.16)
On remarque d’abord que :
$$Sij=2,\mathrm{},p1,ona\theta (j)jet\theta ^{}(j)j$$
(2.17)
En effet, si on avait par exemple $`\theta (j)=j`$, (2.14) conduirait à $`j1=j+1`$ une contradiction qui prouve (2.17). Montrons maintenant que :
$$k=1oup1$$
(2.18)
Supposons $`k1`$. Alors, on a $`1<kp1`$. Alors d’après (2.13) et (2.14), on a l’égalité d’ensembles :
$$\{\theta (k),\theta ^{}(k)\}=\{k1,k+1\},$$
(2.19)
ce qui implique :
$$\theta (k)k,\theta ^{}(k)k$$
(2.20)
Comme $`j_{p+1}=k`$, on déduit de (2.20) et (2.13) que la séquence $`(j_p,j_{p+1},j_{p+2})`$ est égale soit à $`(\theta (k),k,\theta ^{}(k))`$, soit à $`(\theta ^{}(k),k,\theta (k))`$, c’est à dire , grâce à (2.19)), soit à $`(k1,k,k+1)`$, soit à $`(k+1,k,k1)`$. Mais $`j_p=p`$, est différent de $`k1`$. Donc $`p=k+1`$, i.e. $`k=p1`$. Ceci achève de prouver (2.18). Traitons d’abord le cas :
$$k=1$$
(2.21)
Comme $`p>2`$, on a $`1p1`$. Donc $`j_{p1}=p1`$, est différent de $`j_{p+1}=1`$. Alors, (2.14), (2.13) impliquent l’égalité d’ensembles :
$$\{\theta (p),\theta ^{}(p)\}=\{p1,1\}$$
(2.22)
Supposons, d’abord que $`\theta ^{}(1)=2`$, ce qui implique, d’après (2.11), que $`\theta (1)=1`$. Comme $`p>2`$, ni $`\theta (p)`$, ni $`\theta ^{}(p)`$ ne peut être égal à 1. Une contradiction avec l’équation précédente qui montre que l’on doit avoir, d’après (2.11) :
$$\theta (1)=2$$
(2.23)
Comme $`p>2`$, la seule possibilité laissée par (2.22) est :
$$\theta (p1)=pet\theta ^{}(p)=1$$
(2.24)
On déduit de (2.23) et (2.17), joints à (2.13), que, pour $`j=1,\mathrm{},p1`$, on a :
$$\theta (j)=j+1,sijestimpair(resp.\theta ^{}(j)=j+1sijestpair)$$
(2.25)
ce qui, joint à (2.24), implique que $`p`$ est pair. Notons $`p=2q`$. On déduit de (2.25) et (2.23) qu’un élément $`(X_1,\mathrm{},X_p)`$ de $`𝔪_1\mathrm{}𝔪_p`$, est invariant à la fois par $`\sigma `$ et $`\sigma ^{}`$ si et seulement si le système suivant est vérifié:
$$\sigma (X_1)=X_2,\sigma ^{}(X_2)=X_3$$
$$\mathrm{},\mathrm{}$$
$$\sigma (X_{2j1})=X_{2j},\sigma ^{}(X_{2j})=X_{2j+1}$$
$$\mathrm{},\mathrm{}$$
$$\sigma (X_{2q1})=X_{2q},\sigma ^{}(X_{2q})=X_1$$
Notons $`\tau `$ la restriction de $`(\sigma ^{}\sigma )^q`$ à $`𝔪_1`$, qui est un automorphisme $``$-linéaire de $`𝔪_1`$. Ce système possède une solution non nulle si et seulement si l’équa- tion :
$$X_1=\tau (X_1),X_1𝔪_1$$
possède une solution non nulle. C’est le cas, d’après \[G\], Théorème 28. Ceci achève de prouver la Proposition dans le cas $`k=1`$. On suppose maintenant :
$$k=p1>1$$
(2.26)
Comme $`(j_{p1,}j_p,j_{p+1})=(p1,p,p1)`$, on déduit de (2.13), (2.14) et (2.8), que l’on a soit :
$$\theta (p)=p,\theta ^{}(p)=p1$$
(2.27)
soit :
$$\theta (p)=p1,\theta ^{}(p)=p$$
(2.28)
Alors, d’après notre convention (2.15), on a $`\theta (1)=1`$. Supposons (2.27) vérifié. Comme ci-dessus, ceci joint à (2.23) et (2.13), montre que $`p`$ est pair et que , pour $`j=1,\mathrm{},p1`$, on a :
$$\theta ^{}(j)=j+1,sijestimpair(resp.\theta (j)=j+1sijestpair)$$
(2.29)
On note $`p=2q`$. On déduit de (2.28) et (2.29) qu’un élément $`(X_1,\mathrm{},X_p)`$ de $`𝔪_1\mathrm{}𝔪_p`$, est invariant à la fois par $`\sigma `$ et $`\sigma ^{}`$ si et seulement si le système suivant est vérifié :
$$\sigma (X_1)=X_1,\sigma ^{}(X_1)=X_2$$
$$\mathrm{},\mathrm{}$$
$$\sigma (X_{2j})=X_{2j+1},\sigma ^{}(X_{2j+1})=X_{2j+2}$$
$$\mathrm{},\mathrm{}$$
$$\sigma (X_{2q2})=X_{2q1},\sigma ^{}(X_{2q1})=X_{2q}$$
$$\sigma (X_{2q})=X_{2q}$$
Notant $`\tau `$ la restriction de $`(\sigma ^{}\sigma )^q`$ à $`𝔪_1`$, qui est un automorphisme $``$-linéai- re de $`𝔪_1`$, ce système possède une solution non nulle si et seulement si le sys- tème :
$$X_1=\sigma (X_1),X_1=(\tau ^1\sigma \tau )(X_1),X_1𝔪_1$$
(2.30)
possède une solution non nulle. La restriction de $`\sigma `$ à $`𝔪_1`$ et $`𝔪_p`$ est antilinéaire. Par ailleurs $`\tau `$ est soit $``$-linéaire, soit antilinéaire, d’après le Lemme 6. Donc la restriction à $`𝔪_1`$ de $`\tau ^1\sigma \tau `$ est antilinéaire. Il résulte alors du Lemme 12, que (2.30) à une solution non nulle. Ce qui achève la preuve de la Proposition dans le cas étudié. Le cas où (2.28) est satisfait se traite de manière similaire, mais alors $`p`$ est impair. Ceci achève notre discussion et la preuve de la Proposition.
###### Théorème 2
Si $`𝔤`$ n’est pas commutative et si $`(B,𝔦,𝔦^{})`$ est un triple de Manin de $`𝔤`$, sous $`(𝔭,𝔭^{})`$, $`𝔩𝔩^{}`$ est différent de $`𝔤`$.
Démonstration : Raisonnons par l’absurde et supposons qu’il existe un triple de Manin , $`(B,𝔦,𝔦^{})`$, sous $`(𝔤,𝔤)`$, et que $`𝔤`$ ne soit pas commutative. Alors $`𝔥:=𝔦𝔤^{der}`$ (resp. $`𝔥^{}:=𝔦^{}𝔤^{der}`$) est l’espace des points fixes d’une af-involutions $`\sigma `$ (resp. $`\sigma ^{}`$) de $`𝔤^{der}`$, d’après le Théorème 1. En appliquant la Proposition précédente, on aboutit à une contradiction avec l’hypothése $`𝔦𝔦^{}=\{0\}`$, ce qui achève de prouver le Théorème . Soit $`V`$ un sous-espace $`𝔧_0`$ invariant de $`𝔤`$. On suppose qu’il est lasomme de sous-espaces poids de $`𝔤`$ pour $`𝔧_0`$, ce qui s’écrit aussi :
$$V=\underset{\{\lambda 𝔧_0^{}|V^\lambda \{0\}\}}{}𝔤^\lambda $$
Alors $`V`$ admet un unique supplémentaire $`𝔧_0`$-invariant, $`V^{}`$, qui est égal à la somme des sous-espaces poids de $`𝔤`$ qui ont une intersection nulle avec $`V`$, soit encore :
$$V^{}=\underset{\{\lambda 𝔧_0^{}|𝔤^\lambda V=\{0\}\}}{}𝔤^\lambda $$
On note $`p_V`$ (resp. $`p^V`$, la projection de $`𝔤`$ sur $`V`$ (resp. $`V^{}`$) parallèlement à $`V^{}`$ (resp. $`V`$). Tout sous-espace $`𝔧_0`$-invariant est stable sous $`p^V`$ et $`p_V`$. Si de plus $`V`$ est $`𝔩`$-invariant, $`V^{}`$ est aussi $`𝔩`$-invariant. En effet, comme $`𝔩`$ est réductive dans $`𝔤`$, $`V`$ admet un supplémentaire $`𝔩`$-invariant qui n’est autre que $`V^{}`$. On voit aussi que dans ce cas, $`V^{}`$ ne dépend pas du choix de la sous-algèbre de Cartan $`𝔧_0`$ de $`𝔤`$, contenue dans $`𝔩`$. On a le même fait pour $`𝔩^{}`$ et $`𝔩𝔩^{}`$.
###### Théorème 3
Soit $`B`$ une forme de Manin réelle (resp. complexe) sur $`𝔤`$ et $`𝔦`$, $`𝔦^{}`$ des sous-algèbres de Lie Lagrangiennes de $`𝔤`$ pour $`B`$, avec $`𝔦`$ sous $`𝔭`$ et $`𝔦^{}`$ sous $`𝔭^{}`$. On a, grâce au Théorème 1, $`𝔦=𝔥𝔦_𝔞𝔫`$, où $`𝔥=𝔦𝔪`$, $`𝔦_𝔞=𝔦𝔞`$. On note $`\stackrel{~}{𝔥}=𝔦𝔩`$. On fait de même pour $`𝔦^{}`$. Les conditions (i) et (ii) suivantes sont équivalentes : (i) $`(B,𝔦,𝔦^{})`$ est un triple de Manin (ii) Notant $`𝔦_1=p^𝔫^{}(\stackrel{~}{𝔥}𝔭^{})`$, $`𝔦_1^{}=p^𝔫(\stackrel{~}{𝔥}^{}𝔭)`$, on a : a) $`𝔦_1`$ et $`𝔦_1^{}`$ sont contenues dans $`𝔩𝔩^{}`$, et $`(B_1,𝔦_1,𝔦_1^{})`$ est un triple de Manin dans $`𝔩𝔩^{}`$, où $`B_1`$ désigne la restriction de $`B`$ à $`𝔩𝔩^{}`$. b) $`𝔫𝔥^{}`$ et $`𝔫^{}𝔥`$ sont réduits à zéro. Si l’une de ces conditions est vérifiée, on appellera $`(B_1,𝔦_1,𝔦_1^{})`$ l’antécédent du triple de Manin $`(B,𝔦,𝔦^{})`$.
Démonstration : Montrons que (i) implique (ii). Supposons que $`(B,𝔦,𝔦^{})`$ soit un triple de Manin dans $`𝔤`$. Pour des raisons de dimension, ceci équivaut à $`𝔦𝔦^{}=\{0\}`$. Ceci implique immédiatement la propriété b) de (ii). Montrons ensuite que $`𝔦_1`$ est une sous-algèbre de Lie réelle (resp. complexe) de $`𝔩𝔩^{}`$, et isotrope pour $`B_1`$. Etudiant les sous-espaces poids sous $`𝔧_0`$, on voit que :
$$𝔩𝔭^{}=(𝔩𝔩^{})(𝔩𝔫^{})$$
(2.31)
Comme $`𝔩𝔩^{}`$ est $`𝔧_0`$-invariant et que $`p^𝔫^{}(𝔩𝔫^{})`$ est réduit à zéro, on a :
$$p^𝔫^{}(𝔩𝔭^{})𝔩𝔩^{}$$
Il en résulte que $`𝔦_1`$ est bien contenu dans $`𝔩𝔩^{}`$. Par ailleurs, la restriction de $`p^𝔫^{}`$ à $`𝔭^{}`$ est la projection sur $`𝔩^{}`$, parallèlement à $`𝔫^{}`$. C’est donc un morphisme d’algèbres de Lie, ce qui implique que $`𝔦_1`$ est une sous-algèbre de Lie rélle de $`𝔩𝔩^{}`$. Soit $`X`$, $`X_1𝔦_1`$. Ce sont des éléments de $`𝔩𝔩^{}`$, et il existe $`N^{}`$ et $`N_1^{}𝔫^{}`$ tels que $`Y`$ et $`Y_1`$ soient éléments de $`𝔦`$, où :
$$Y:=X+N^{},Y_1:=X_1+N_1^{}$$
Par ailleurs $`𝔫^{}`$ et $`𝔭^{}`$ sont orthogonaux pour $`B`$ (cf. la fin de la démonstration du Théorème 1). Un calcul immédiat montre alors que $`B(Y,Y_1)`$ est égal à $`B_1(X,X_1)`$. Comme $`Y,Y_1𝔦`$, $`B(Y,Y_1)`$ est nul. Finalement, $`𝔦_1`$ est isotrope pour $`B_1`$. On montre de même des propriétés similaires pour $`𝔦_1^{^{}}`$. Montrons $`𝔦_1+𝔦_1^{}=𝔩𝔩^{}`$. Soit $`X𝔩𝔩^{}`$. Alors $`X=I+I^{}`$, avec $`I𝔦`$, $`I^{}𝔦^{}`$. Ecrivons $`I=H+N,I^{}=H^{}+N^{}`$$`H\stackrel{~}{𝔥},H^{}\stackrel{~}{𝔥}^{},N𝔫,N^{}𝔫^{}`$. On a donc :
$$X=H+N+H^{}+N^{}$$
(2.32)
ce qui implique : $`H=XH^{}N^{}N`$. On voit ainsi que $`H`$ est élément de $`(𝔭^{}+𝔫)𝔩`$. Décomposant sous l’action de $`𝔧_0`$, on voit que :
$$(𝔭^{}+𝔫)𝔩=𝔭^{}𝔩$$
(2.33)
Finalement $`H`$ est élément de $`\stackrel{~}{𝔥}𝔭^{}`$. de même, on voit que $`H^{}`$ est élément de $`\stackrel{~}{𝔥}^{}𝔭`$. Par ailleurs, $`𝔫`$ et $`𝔫^{}`$ sont des sous-espaces $`𝔧_0`$-invariants et en somme directe avec $`𝔩𝔩^{}`$. Donc, appliquant $`p_{𝔩𝔩^{}}`$ à (2.32), on a :
$$X=p_{𝔩𝔩^{}}(H)+p_{𝔩𝔩^{}}(H^{})$$
De (2.31), on déduit que la restriction de $`p_{𝔩𝔩^{}}`$ à $`𝔩𝔭^{}`$ est égale à la restriction de $`p^𝔫^{}`$ à $`𝔩𝔭^{}`$. Donc, on a :
$$p_{𝔩𝔩^{}}(H)=p^𝔫^{}(H)𝔦_1$$
on obtient de même :
$$p_{𝔩𝔩^{}}(H)𝔦_1^{}$$
et l’on conclut que :
$$X𝔦_1+𝔦_1^{}$$
Ceci achève de prouver que :
$$𝔩𝔩^{}=𝔦_1+𝔦_1^{}$$
(2.34)
Par ailleurs :
$$𝔩𝔩^{}estlecentralisateurd^{}un\stackrel{´}{e}l\stackrel{´}{e}mentsemisimplede𝔤,dont$$
$$l^{}imageparlarepr\stackrel{´}{e}sentationadjointen^{}aquedesvaleurspropres$$
$$r\stackrel{´}{e}elles$$
(2.35)
En effet $`(𝔩𝔩^{})((𝔫^{}𝔩)𝔫)`$ est une décomposition de Langlands d’une sous-algèbre parabolique de $`𝔤`$. Alors la restriction $`B_1`$ de $`B`$ à $`𝔩𝔩^{}`$ est une forme de Manin (cf. Corollaire du Lemme 2 ), et $`𝔦_1`$, $`𝔦_1^{}`$, qui sont isotropes pour $`B_1`$, sont de dimensions réelles inférieures ou égales à la dimension complexe de $`𝔩𝔩^{}`$. La somme dans (2.34) est nécesssairement directe, ce qui achève de prouver que (i) implique (ii). Montrons que (ii) implique (i). Supposons satisfaites les conditions a) et b) de (i). Montrons que $`𝔦𝔦^{}`$ est réduit à zéro. Soit $`X`$ un élément de $`𝔦𝔦^{}`$. Alors :
$$X=H+N=H^{}+N^{},o\stackrel{`}{u}H\stackrel{~}{𝔥},H^{}\stackrel{~}{𝔥}^{},N𝔫,N^{}𝔫^{}$$
(2.36)
On a alors :
$$H=H^{}+N^{}N𝔩(𝔭^{}+𝔫)$$
(2.33) implique que $`H𝔩𝔭^{}`$. De même, on montre que $`H^{}𝔩^{}𝔭`$. Appliquant $`p_{𝔩𝔩^{}}`$ à (2.36), on voit que :
$$p_{𝔩𝔩^{}}(X)=p_{𝔩𝔩^{}}(H)=p_{𝔩𝔩^{}}(H^{})$$
et, grâce à la première partie de la démonstration, cela conduit à :
$$p_{𝔩𝔩^{}}(X)=p^𝔫^{}(H)=p^𝔫(H^{})𝔦_1𝔦_1^{}$$
Donc on a :
$$p^𝔫^{}(H)=p^𝔫(H^{})=0$$
Mais $`p^𝔫^{}`$ est injective sur $`\stackrel{~}{𝔥}𝔭^{}`$, car $`𝔫^{}\stackrel{~}{𝔥}`$ est réduit à zéro. En effet $`𝔫^{}\stackrel{~}{𝔥}`$ est contenu dans $`𝔫^{}𝔩`$. On voit que cette dernière intersection est égal à $`𝔫^{}𝔪`$. Donc $`𝔫^{}\stackrel{~}{𝔥}`$ est égal à $`𝔫^{}𝔥`$, qui est réduit à zéro, d’après b). Donc $`H`$ est nul et il en va de même de $`H^{}`$. Alors $`X`$ est un élément de $`𝔫𝔫^{}`$, qui est réduit à zéro, d’après nos hypothèses sur $`𝔭`$, $`𝔭^{}`$. Donc $`X`$ est nul et $`𝔦𝔦^{}`$ est réduit à zéro. Alors la somme $`𝔦+𝔦^{}`$ est directe, et l’on a $`𝔤=𝔦𝔦^{}`$ pour des raisons de dimension. Ceci achève de prouver le Théorème.
###### Proposition 4
Si $`(B,𝔦,𝔦^{})`$ est un triple de Manin sous $`(𝔭,𝔭^{})`$, d’antécédent $`(B_1,𝔦_1,𝔦_1^{})`$, et si $`g=nn^{}xPP^{}`$, où $`xLL^{}`$, $`nN_L^{}`$, $`n^{}N_L^{}`$, l’antécédent de $`(B,Adg(𝔦),Adg(𝔦^{}))`$ est égal à $`(B_1,Adx(𝔦_1),Adx(𝔦_1^{}))`$.
Démonstration : Ecrivons $`\underset{¯}{𝔦}=Adg(𝔦)`$ et $`\underset{¯}{\overset{~}{𝔥}}=\underset{¯}{𝔦}𝔩`$, etc. On note $`(B_1,\underset{¯}{𝔦}_1,\underset{¯}{𝔦}_1^{})`$, l’antécédent de $`(B,Adg(𝔦),Adg(𝔦^{}))`$. On a, grâce à la Proposition 2 :
$$g=n^{}nx=n^{}x(x^1nx)$$
Donc :
$$Adg(𝔦)=Adn^{}x(𝔦)$$
puisque $`x^1nxNI`$. Comme $`n^{}x(LL^{})N_L^{}L`$, cela implique :
$$\underset{¯}{\overset{~}{𝔥}}=Adn^{}x(\stackrel{~}{𝔥}),$$
$`\stackrel{~}{𝔥}=𝔦𝔩`$. Mais $`n^{}x`$ est aussi élément de $`P^{}`$. Alors, on a :
$$\underset{¯}{\overset{~}{𝔥}}𝔭^{}=Adn^{}x(\stackrel{~}{𝔥}𝔭^{})$$
D’où l’on déduit :
$$\underset{¯}{𝔦}_1=p^𝔫^{}(Adn^{}x(\stackrel{~}{𝔥}𝔭^{}))$$
Mais il est clair que la restriction de $`p^𝔫^{}`$ à $`𝔭^{}`$, n’est autre que la projection sur $`𝔩^{}`$, parallélement à $`𝔫^{}`$. Cette restriction entrelace l’action adjointe de $`P^{}`$ sur $`𝔭^{}`$ avec l’action naturelle de $`P^{}`$ sur $`𝔩^{}`$, identifié au quotient de $`𝔭^{}`$ par $`𝔫^{}`$ ($`N^{}`$ agit trivialement). Il en résulte :
$$\underset{¯}{𝔦}_1=Adx(p^𝔫^{}(\stackrel{~}{𝔥}𝔭^{}))=Adx(𝔦_1)$$
comme désiré. On traite de manière similaire $`\underset{¯}{𝔦}_{}^{}{}_{1}{}^{}`$.
###### Proposition 5
Tout triple de Manin sous $`(𝔭,𝔭^{})`$ est conjugué, par un élé- ment de $`PP^{}`$ à un triple de Manin , $`(B,𝔦,𝔦^{})`$, sous $`(𝔭,𝔭^{})`$, d’antécédent $`(B_1,𝔦_1,𝔦_1^{})`$ pour lequel il existe une sous-algèbre de Cartan fondamentale $`\stackrel{~}{𝔣}`$ (resp. $`\stackrel{~}{𝔣}^{}`$), de $`𝔦`$ (resp. $`𝔦^{}`$), contenue dans $`𝔦_1`$ (resp. $`𝔦_1^{}`$). On dit que le triple $`(B,𝔦,𝔦^{})`$ est lié à son antécédent, avec lien $`(\stackrel{~}{𝔣},\stackrel{~}{𝔣}^{})`$.
Démonstration : Démontrons (i). Soit $`(B,\underset{¯}{𝔦},\underset{¯}{𝔦}^{})`$ un triple de Manin pour $`𝔤`$, sous $`(𝔭,𝔭^{})`$. On note $`\underset{¯}{𝔥}=\underset{¯}{𝔦}𝔪`$, etc.. On note $`\underset{¯}{\sigma }`$ (resp. $`\underset{¯}{\sigma }^{}`$), l’af-involution de $`𝔪`$ (resp. $`𝔪^{}`$) ayant $`\underset{¯}{𝔥}`$ (resp. $`\underset{¯}{𝔥}^{}`$) pour espace de points fixes. On définit de même $`\underset{¯}{𝔦}_𝔞`$ . Comme $`\underset{¯}{𝔦}+\underset{¯}{𝔦}^{}=𝔤`$, on a :
$$\underset{¯}{𝔦}+𝔭^{}=𝔤$$
Appliquant $`p_𝔩`$ à cette égalité, on en déduit :
$$(\underset{¯}{𝔦}𝔩)+(𝔭^{}𝔩)=𝔩$$
On applique encore la projection de $`𝔩`$ sur $`𝔪`$, parallèlement à $`𝔞`$ pour obtenir :
$$\underset{¯}{𝔥}+(𝔭^{}𝔪)=𝔪$$
En conséquence, $`\underset{¯}{H}(P^{}M)^0`$ est ouvert dans $`M`$. Or $`(P^{}M)^0`$ est le sous groupe parabolique de $`M`$, d’algèbre de Lie $`𝔭^{}𝔪`$. Par ailleurs $`\sigma `$ étant une af-involution, $`𝔪`$ est le produit d’idéaux $`𝔪_j`$, invariants par $`\sigma `$ et sur lesquels induit : soit une conjugaison par rapport à une forme réelle, soit ”l’ échange des facteurs ” de deux idéaux isomorphes dont $`𝔪_j`$ est la somme. Il résulte alors de \[M2\], \[M1\], que $`\underset{¯}{𝔥}𝔭^{}`$ contient une sous-algèbre de Cartan fondamentale $`\underset{¯}{𝔣}`$ de $`\underset{¯}{𝔥}`$ et une sous-algèbre de Borel, $`𝔟`$, de $`𝔪`$, contenant $`\underset{¯}{𝔣}`$, contenue dans $`𝔭^{}𝔪`$, et telle que :
$$\underset{¯}{\sigma }(\underset{¯}{𝔟})+\underset{¯}{𝔟}=𝔪$$
(2.37)
D’après le Lemme 11, :
$$\underset{¯}{\overset{~}{𝔣}}:=\underset{¯}{𝔣}+\underset{¯}{𝔦}_𝔞$$
est une sous-algèbre de Cartan fondamentale de $`\underset{¯}{𝔦}`$ et le centralisateur dans $`𝔤`$, $`\underset{¯}{𝔧}`$, de $`\underset{¯}{\overset{~}{𝔣}}`$ est une sous-algèbre de Cartan de $`𝔤`$, contenue dans $`𝔩`$. De la définition des af-involutions, il résulte que toute sous-algèbre de Cartan de $`\underset{¯}{𝔥}`$ contient des éléments réguliers de $`𝔪`$. Il résulte alors de \[Bor\], Proposition 11.15, que $`\underset{¯}{𝔧}𝔪`$ est contenu dans $`𝔟`$. Donc $`\underset{¯}{𝔧}=(\underset{¯}{𝔧}𝔪)𝔞`$ est contenu dans $`𝔭^{}𝔩`$. C’est une sous-algèbre de Cartan de $`𝔭^{}𝔩`$, donc elle est conjuguée à $`𝔧_0`$, par un élément du sous-groupe analytique de $`G`$, d’algèbre de Lie $`𝔭^{}𝔩`$. Mais, d’après (2.31), on a : $`𝔭^{}𝔩=(𝔩𝔩^{})(𝔫^{}𝔩)`$ et ce sous-groupe analytique est égal à $`(LL^{})N_L^{}`$, puisque $`LL^{}`$ est connexe, d’après la Proposition 2. Donc, il existe $`n^{}N_L^{}`$, $`xLL^{}`$, tels que
$$Adxn^{}(\underset{¯}{𝔧})=𝔧_0$$
soit encore :
$$Adn^{}(\underset{¯}{𝔧})=Adx^1(𝔧_0)𝔩𝔩^{}$$
(2.38)
On trouve de même $`\stackrel{~}{\underset{¯}{𝔣}^{}}`$, $`\underset{¯}{𝔟}^{}`$, $`\underset{¯}{𝔧}^{}`$ et $`x^{}LL^{}`$, $`nN_L^{}`$, vérifiant des propriétés similaires. On pose :
$$u=nn^{},𝔦=Adu(\underset{¯}{𝔦}),𝔦^{}=Adu(\underset{¯}{𝔦}^{})$$
Comme $`n`$ et $`n^{}`$ commutent et que $`𝔫`$ est un idéal de $`𝔦`$, on a :
$$Adu(\underset{¯}{𝔦})=Adn^{}(\underset{¯}{𝔦})$$
et de même :
$$Adu(\underset{¯}{𝔦}^{})=Adn(\underset{¯}{𝔦}^{})$$
On pose alors :
$$\stackrel{~}{𝔣}=Adn^{}(\underset{¯}{\overset{~}{𝔣}}),\stackrel{~}{𝔣}^{}=Adn(\underset{¯}{\overset{~}{𝔣}}^{}),𝔟=Adn^{}(\underset{¯}{𝔟}),𝔟^{}=Adn(\underset{¯}{𝔟}^{})$$
(2.39)
On voit alors que $`(B,𝔦,𝔦^{})`$ est un triple de Manin, conjugué par $`u`$ à $`(B,\underset{¯}{𝔦},\underset{¯}{𝔦}^{})`$ et sous $`(𝔭,𝔭^{})`$. On va voir que $`\stackrel{~}{𝔣}`$ a les propriétés voulues. D’abord, comme $`\underset{¯}{\overset{~}{𝔣}}`$ est une sous-algèbre de Cartan fondamentale de $`\underset{¯}{𝔦}`$, par conjugaison, on en déduit que $`\stackrel{~}{𝔣}`$ est une sous-algèbre de Cartan fondamentale de $`𝔦`$. Le centralisateur, $`𝔧`$, de $`\stackrel{~}{𝔣}`$ vérifie
$$𝔧=Adn^{}(\underset{¯}{𝔧})$$
(2.40)
donc est contenu dans $`𝔩𝔩^{}`$, d’après (2.38). Alors, d’après le Lemme 11, on a bien $`\stackrel{~}{𝔣}=𝔣(\stackrel{~}{𝔣}𝔞)`$, où $`𝔣=\stackrel{~}{𝔣}𝔥`$, et $`\stackrel{~}{𝔣}𝔞`$ est égal à $`𝔦_𝔞`$. On a vu que $`\stackrel{~}{𝔣}`$ est contenu dans $`𝔦𝔩𝔩^{}`$, donc dans $`\stackrel{~}{𝔥}𝔭^{}`$. De plus $`p^𝔫^{}`$ est l’identité sur $`𝔩^{}`$. Donc $`\stackrel{~}{𝔣}`$ est contenu dans $`𝔦_1`$. Par ailleurs, comme $`\stackrel{~}{𝔣}`$ est une sous-algèbre de Cartan de $`𝔦`$, contenue dans $`\stackrel{~}{𝔥}𝔭^{}`$, c’est une sous-algèbre de Cartan de $`\stackrel{~}{𝔥}𝔭^{}`$ (cf. \[Bou\], Ch. VII, Paragraphe 2.1, Exemple 3), et par projection , c’est une sous-algèbre de Cartan de $`𝔦_1`$ (cf. l.c., Corollaire 2 de la Proposition 4). Il reste à voir que cette sous-algèbre de Cartan de $`𝔦_1`$ est fondamentale. On suppose que $`𝔦_1`$ est sous $`𝔭_1`$. D’après le Lemme 11 (iv), il existe une unique décompo- sition de Langlands $`𝔭_1=𝔩_1𝔫_1`$, telle que $`𝔩_1`$ contienne $`\stackrel{~}{𝔣}`$. On note $`𝔪_1=𝔩_{1}^{}{}_{}{}^{der}`$. Il suffit de voir que $`\stackrel{~}{𝔣}𝔪_1`$ est une sous-algèbre de Cartan fondamentale de $`𝔥_1:=𝔦_1𝔪_1`$. Pour cela, il suffit de voir qu’aucune racine de $`\stackrel{~}{𝔣}𝔪_1`$ dans $`𝔪_1`$ n’est réelle. D’après le Lemme 11 (iv), $`\stackrel{~}{𝔣}=(\stackrel{~}{𝔣}𝔪_1)(\stackrel{~}{𝔣}𝔞_1)`$, où $`𝔞_1`$ est le centre de $`𝔩_1`$. Alors, une racine $`\alpha `$ de $`\stackrel{~}{𝔣}𝔪_1`$ dans $`𝔪_1`$, prolongée par zéro sur $`\stackrel{~}{𝔣}𝔞_1`$ est une racine de $`\stackrel{~}{𝔣}`$ dans $`𝔪`$. Mais alors, comme $`𝔣`$ est une sous-algèbre de Cartan fondamentale de $`𝔥`$, $`\alpha `$ n’est pas réelle sur $`𝔣`$. Ceci prouve que $`\stackrel{~}{𝔣}`$ est une sous-algèbre de Cartan fondamentale de $`𝔦_1`$. On montre de même que $`\stackrel{~}{𝔣}^{}`$ est une sous-algèbre de Cartan fondamentale de $`𝔦^{}`$ et $`𝔦_1^{}`$. Ceci achève la preuve de la Proposition.
###### Théorème 4
Tout triple de Manin réel (resp. complexe) sous $`(𝔭,𝔭^{})`$ est conjugué, par un élément de $`PP^{}`$, à un triple de Manin réel (resp. complexe) sous $`(𝔭,𝔭^{})`$, $`(B,𝔦,𝔦^{})`$, dont tous les antécédents successifs, $`(B,𝔦_1,𝔦_1^{}),(B,𝔦_2,𝔦_2^{}),`$ $`\mathrm{}`$, sont des triples de Manin standard dans $`𝔤_1=𝔩𝔩^{},𝔤_2,\mathrm{}`$, relativement à l’intersection de $`𝔟_0,𝔟_0^{}`$, avec $`𝔤_1,𝔤_2,\mathrm{}`$, et tel que l’intersection $`𝔣_0`$ ( resp. $`𝔣_0^{}`$) de $`𝔧_0`$ avec $`𝔦`$ (resp. $`𝔦^{}`$) soit une sous-algèbre de Cartan fondamentale de $`𝔦`$ (resp. $`𝔦^{}`$), contenue dans $`𝔦_1,𝔦_2,\mathrm{}`$ (resp. $`𝔦_1^{},𝔦_2^{},\mathrm{}`$) . Un triple satisfaisant ces propriétés sera appelé triple fortement standard. Le plus petit entier, $`k`$, tel que $`𝔤_k=𝔧_0`$, est appelé la hauteur du triple fortement standard.
Démonstration : On procède par récurrence sur la dimension de $`𝔤^{der}`$. Si celle-ci est nulle, le Théorème est clair. Supposons l’assertion démontrée pour les algèbres réductives dont l’idéal dérivé est de dimension strictement inférieure à celle de $`𝔤^{der}`$. D’après la Proposition 5, le triple donné est conjugué, par un élément de $`PP^{}`$, à un triple de Manin $`𝒯^{}`$, sous $`(𝔭,𝔭^{})`$, lié à son antécédent $`𝒯_1^{}`$. D’après l’hypothèse de récurrence, ce dernier est conjugué par un élément, $`g_1`$, de $`LL^{}`$, a un triple de Manin fortement standard : $`\underset{¯}{𝒯}_1:=g_1𝒯_1^{}`$. Par transport de structure, $`\underset{¯}{𝒯}=g_1𝒯^{}`$ est lié à son antécédent $`\underset{¯}{𝒯}_1`$. On note $`\underset{¯}{𝒯}=(B,\underset{¯}{𝔦},\underset{¯}{𝔦}^{})`$, $`\underset{¯}{𝒯}_1=(B,\underset{¯}{𝔦}_1,\underset{¯}{𝔦}_{}^{}{}_{1}{}^{})`$, et $`(𝔣,𝔣^{})`$ un lien entre ces triples. Comme $`\underset{¯}{𝒯}_1`$ est fortement standard, $`𝔣_1:=𝔧_0\underset{¯}{𝔦}_1`$ (resp. $`𝔣_1^{}:=𝔧_0\underset{¯}{𝔦}_1^{}`$) est une sous-algèbre de Cartan fondamentale de $`\underset{¯}{𝔦}_1`$ (resp. $`\underset{¯}{𝔦}_{}^{}{}_{1}{}^{}`$). Alors $`𝔣`$ et $`𝔣_1`$ (resp. $`𝔣^{}`$ et $`𝔣_1^{}`$ ) sont des sous-algèbre de Cartan fondamentales de $`\underset{¯}{𝔦}_1`$ (resp. $`\underset{¯}{𝔦}_{}^{}{}_{1}{}^{}`$), donc conjuguées par un élément $`i_1`$ de $`\underset{¯}{I}_1`$ (resp. $`i_1^{}`$ de $`\underset{¯}{I}_{}^{}{}_{1}{}^{}`$), i.e. :
$$𝔣=i_1(𝔣_1),𝔣^{}=i_1^{}(𝔣_1^{})$$
On dispose d’une suite exacte de groupes :
$$0N^{}P^{}L^{}0$$
où la flèche de $`P^{}`$ dans $`L^{}`$ est le morphisme dont la différentielle est la restriction de $`p^𝔫^{}`$ à $`𝔭`$. La définition de $`\underset{¯}{𝔦}_1`$ montre que la restriction de ce morphisme à $`(\underset{¯}{\overset{~}{H}}P^{})^0`$ est un morphisme surjectif sur $`\underset{¯}{I}_1`$, de noyau contenu dans $`N_L^{}`$. Comme $`\underset{¯}{\overset{~}{H}}`$ est contenu dans $`\underset{¯}{I}`$, on en déduit qu’il existe $`n^{}N_L^{}`$ et $`i\underset{¯}{l}`$ tels que :
$$i_1=n^{}i$$
De même on trouve $`nN_L^{}`$, $`i^{}\underset{¯}{I}^{}`$ tels que :
$$i_1^{}=ni^{}$$
Montrons que le triple de Manin $`𝒯=(B,𝔦,𝔦^{})`$, défini par $`𝒯=n^1n^1\underset{¯}{𝒯}`$, convient. D’abord c’est un triple sous $`(𝔭,𝔭^{})`$, puisque $`n,n^{}PP^{}`$, conjugué du triple initial. Par ailleurs $`n`$ et $`n^{}`$ commutent (cf. Proposition 2) et $`N`$ est contenu dans $`\underset{¯}{I}`$. Donc, on a :
$$𝔦=n^1(\underset{¯}{𝔦})=n^1i^1(\underset{¯}{𝔦})=i_1^1(\underset{¯}{𝔦})$$
Comme $`𝔣=i_1(𝔣_1)`$ est une sous-algèbre de Cartan fondamentale de $`\underset{¯}{𝔦}`$, $`i_1^1(𝔣)=𝔣_1`$ est une sous-algèbre de Cartan fondamentale de $`𝔦`$, par transport de structure. De plus $`𝔣_1`$ est contenue dans $`\underset{¯}{𝔦}_1`$. De même $`𝔣_1^{}`$ est une sous-algèbre de Cartan fondamentale de $`𝔦^{}`$, contenue dans $`\underset{¯}{𝔦}_1^{}`$. Par ailleurs, comme $`\underset{¯}{𝒯}_1`$ est fortement standard, $`𝔧_0`$ est la somme directe de $`𝔣_1`$ et $`𝔣_1^{}`$. Comme $`𝔦`$ et $`𝔦^{}`$ ont une intersection réduite à zéro, il en résulte que $`𝔣_1`$ (resp. $`𝔣_1^{}`$) est égal à l’intersection de $`𝔧_0`$ avec $`𝔦`$ (resp. $`𝔦^{}`$). Par ailleurs, d’après la Proposition 4, l’antécédent de $`(B,𝔦,𝔦^{})`$ est égal à $`\underset{¯}{𝒯}_1`$. Comme ce dernier est fortement standard, ce qui précède suffit à prouver que $`(B,𝔦,𝔦^{})`$ l’est aussi.
###### Remarque 2
Le Théorème 4 réduit la classification des triples de Manin à celle des triples fortement standard.
###### Proposition 6
Si deux triples de Manin fortement standard sont conjugués par un élément de $`G`$, ils sont de même hauteur. Ceci permet de définir la hauteur d’un triple de Manin comme la hauteur d’un triple fortement standard auquel il est conjugué.
Démonstration : On procède par récurrence sur la dimension de $`𝔤^{der}`$. Si celle-ci est nulle, la Proposition est vraie. Sinon, d’après la Proposition 2 et la Proposition 4, si deux triples de Manin fortement standard sont conjugués par un élément, leurs antécédents sont conjugués par un élément de $`LL^{}`$. La Proposition résulte alors immédiatement de l’application de l’hypothèse de récurrence. Etablissons quelques propriétés des triples fortement standard. On utilisera les deux Lemmes suivant.
###### Lemme 13
Soit $`B`$ une forme de Manin $``$-bilinéaire. Soit $`𝔦`$ une sous-algèbre Lagrangienne complexe sous $`𝔭`$. Soit $`\alpha `$ un poids non nul de $`𝔧_0`$ dans $`𝔤`$, tel que $`𝔤^\alpha `$ soit contenu dans $`𝔦`$. Alors $`𝔤^\alpha `$ est contenu dans $`𝔫`$.
Démonstration : On emploie les notations du Théorème 1. On note $`𝔣=𝔧_0𝔥`$. On sait que si $`\beta `$ est une racine de $`𝔧_0`$ dans $`𝔪`$, $`\sigma (𝔪^\beta )=𝔪^\beta ^{}`$, avec $`\beta ^{}\beta `$ (cf. Lemme 5, pour les f-involutions). De plus si $`\beta _{|𝔣}=\gamma _{|𝔣}`$, pour un autre poids de $`𝔧_0`$ dans $`𝔪`$, on a $`\beta `$ égal à $`\gamma `$ ou $`\gamma ^{}`$. On note :
$$𝔥_\beta :=\{X+\sigma (X)|X𝔪^\beta \}$$
Soit $`R_{}`$ un sous-ensemble de l’ensemble $`R`$, des poids non nuls de $`𝔧_0`$ dans $`𝔪`$, tel que $`R_{}`$ et $`\{\beta ^{}|\beta R_{}\}`$ forme une partition de $`R`$. Alors, on a :
$$𝔥𝔦_𝔞=(𝔣+𝔦_𝔞)(_{\beta R_{}}𝔥_\beta )$$
(2.41)
qui est une décomposition en somme directe de représentations de $`𝔣`$ qui sont deux à deux sans sous-quotients simples isomorphes, toutes étant irréduc- tibles, sauf peut-être la première. Si $`\alpha `$ n’est pas un poids de $`𝔧_0`$ dans $`𝔫`$, comme il est non nul, c’est un poids de $`𝔧_0`$ dans $`𝔪`$. En étudiant l’action de $`𝔣`$, on est conduit à :
$$𝔤^\alpha 𝔥_\alpha $$
Comme $`\alpha \alpha ^{}`$, c’est impossible. Une contradiction qui achève de prouver le Lemme.
###### Lemme 14
Soit $`(B,𝔦,𝔦^{})`$ un triple de Manin complexe fortement standard, et soit $`(B_1,𝔦_1,𝔦_1^{})`$ son antécédent, que l’on suppose sous $`(𝔭_1,𝔭_1^{})`$. On note $`𝔭_1=𝔩_1𝔫_1`$ la décomposition de Langlands de $`𝔭_1`$ telle que $`𝔩_1`$ contienne $`𝔧_0`$ et on note $`𝔪_1`$ l’idéal dérivé de $`𝔩_1`$. Alors $`𝔪_1`$ est égal à l’idéal dérivé de $`(𝔪𝔪^{})\sigma (𝔪𝔪^{})`$ et l’involution, $`\sigma _1`$, de $`𝔪_1`$, dont $`𝔥_1:=𝔦_1𝔪_1`$ est l’ensemble des points fixes, est égale à la restriction de $`\sigma `$ à $`𝔪_1`$.
Démonstration : On réutilise les notations du Lemme précédent. Etudiant la décomposition en représen- tations irréductibles sous $`𝔣`$, de $`\stackrel{~}{𝔥}𝔭^{}`$, et utilisant (2.41), on voit que :
$$\stackrel{~}{𝔥}𝔭^{}=(𝔣+𝔦_𝔞)(_{\beta R_{},𝔥_\beta 𝔭^{}}𝔥_\beta )$$
Mais $`𝔥_\beta 𝔭^{}`$ si et seulement si $`𝔤^\beta `$ et $`𝔤^\beta ^{}`$ sont contenus dans $`𝔭^{}`$. Si $`\beta `$ satisfait cette condition on a :
$$p^𝔫^{}(𝔥_\beta )=𝔤^\beta si𝔤^\beta 𝔪^{}et𝔤^\beta 𝔫^{}$$
$$p^𝔫^{}(𝔥_\beta )=𝔥_\beta si𝔤^\beta ,𝔤^\beta ^{}𝔪^{}$$
Notons, pour $`V`$ sous-espace vectoriel complexe de $`𝔤`$, invariant sous $`𝔧_0`$, $`\mathrm{\Delta }(V,𝔧_0)`$, l’ensemble des poids non nuls de $`𝔧_0`$ dans $`V`$. Alors on a :
$$𝔦_1=𝔲_1𝔳_1𝔦_𝔞$$
(2.42)
où :
$$𝔲_1=𝔣_{\beta R_{}\mathrm{\Delta }(𝔩𝔩^{},𝔧_0),\beta ^{}R_{}\mathrm{\Delta }(𝔩𝔩^{},𝔧_0)}𝔥_\beta $$
(2.43)
et
$$𝔳_1=_{\beta \mathrm{\Delta }(𝔩𝔩^{},𝔧_0),\beta ^{}\mathrm{\Delta }(𝔩𝔫^{},𝔧_0)}𝔤^\beta $$
(2.44)
On remarque que :
$$𝔲_1=((𝔪𝔩^{})\sigma (𝔪𝔩^{}))^\sigma $$
(2.45)
Donc $`𝔲_1`$ est une sous-algèbre de Lie de $`𝔪𝔩^{}`$, réductive, dont le centre est contenu dans $`𝔣`$. Notons $`𝔴_1=(𝔪𝔪^{})\sigma (𝔪𝔪^{})+𝔧_0`$. On voit que $`𝔮_1:=𝔴_1𝔳_1`$ est une une sous-algèbre parabolique, de $`𝔩𝔩^{}`$, dont le radical nilpotent est $`𝔳_1`$, et la décomposition ci-dessus est une décomposition de Langlands de $`𝔮_1`$ avec $`𝔧_0`$ contenu dans $`𝔴_1`$. Par ailleurs $`𝔦_1`$ est contenue dans $`𝔮_1`$, d’après (2.42), (2.43) et (2.44), Tenant compte du fait que, pour $`\beta R`$, $`[𝔣,𝔪^\beta ]=𝔪^\beta `$, d’après la définition de $`𝔣`$ et des f-involutions, on déduit de (1.3) que le radical nilpotent de $`𝔦_1`$ contient $`𝔳_1`$. Par ailleurs, comme le centre de $`𝔲_1`$ est contenu dans $`𝔧_0`$, le radical et, a fortiori, le radical nilpotent de $`𝔦_1`$ est contenu dans $`𝔧_0𝔳_1`$. Mais ce radical nilpotent ne rencontre pas $`𝔧_0`$ (cf. Théorème 1), donc il est égal à $`𝔳_1`$. Alors $`𝔭_1=𝔮_1`$, $`𝔩_1=𝔴_1`$. Donc $`𝔪_1`$ a la forme annoncée. Comme $`𝔥_1=𝔲_1𝔪_1`$, on déduit de (2.45) que $`\sigma _1`$ a la forme annoncée.
## 3 Classification des triples de Manin complexes
Dans toute cette partie triple de Manin voudra dire triple de Manin complexe. On rappelle qu’on a fixé $`𝔧_0`$ et $`𝔟_0`$, $`𝔭`$, $`𝔭^{}`$. On définit :
$$^+:=\{\lambda ^{}|Re\lambda <0,ouRe\lambda =0etIm\lambda >0\},^{}=^{}^+$$
Si $`B`$ est une forme de Manin complexe sur $`𝔤`$, on note $`𝔤_+`$ (resp. $`𝔤_{}`$) la somme de ses idéaux simples, $`𝔤_i`$, pour lesquels la restriction de $`B`$ à $`𝔤_i`$ est égal à $`K_{\lambda _i}^{𝔤_i}`$, avec $`\lambda _i^+`$ (resp. $`^{}`$).
###### Lemme 15
Aucune sous-algèbre de Lie semi-simple complexe de $`𝔤_+`$ (resp. $`𝔤_{}`$) n’est isotrope pour $`B`$.)
On fait la démontration pour $`𝔤_+`$, celle pour $`𝔤_{}`$ étant identique. Soit $`𝔰`$ une sous algèbre de Lie semi-simple complexe de $`𝔤_+`$. On note $`𝔨_𝔰`$ une forme réelle compacte de $`𝔰`$. Alors $`𝔨_𝔰`$ est contenu dans une forme réele compacte de $`𝔤_+`$. En effet, le sous-groupe analytique de $`G_+`$ d’algèbre de Lie $`𝔨_𝔰`$ est compact, comme groupe de Lie connexe d’algèbre de Lie semi-simple compacte. Il est donc contenu dans un sous-groupe compact maximal $`K`$ de $`G_+`$, et son algèbre de Lie, $`𝔨`$, convient. On note $`𝔤_i`$, $`i=1,\mathrm{},p`$, les idéaux simples de $`𝔤_+`$. Alors $`𝔨=_{i=1,\mathrm{},p}𝔨_i`$, où $`𝔨_i=𝔨𝔤_i`$. Le Lemme résultera de la preuve de :
$$B(X,X)0,X𝔨\{0\}$$
(3.1)
Pour cela on remarque que :
$$\underset{i=1,\mathrm{},p}{}\lambda _ix_i0,si,pouri=1,\mathrm{},p,\lambda _i^+etx_i0,nontousnuls$$
(3.2)
Maintenant, si $`X𝔨\{0\}`$ et $`X=_{i=1,\mathrm{},p}X_i`$, où $`X_i𝔨_i`$, $`K_{𝔤_i}(X_i,X_i)`$ est négatif où nul, et non nul pour au moins un indice $`i`$. Alors (3.1), et donc le Lemme, résulte de (3.2) et de la définition de $`𝔤_+`$.
###### Lemme 16
Soit $`𝔦`$ une sous-algèbre Lagrangienne sous $`𝔭`$. On note $`𝔥=𝔦𝔪`$, $`𝔪_+=𝔪𝔤_+`$, $`𝔪_{}=𝔪𝔤_{}`$. On a $`𝔪=𝔪_+𝔪_{}`$. Par ailleurs la f-involution, $`\sigma `$, dont $`𝔥`$ est l’espace des points fixes, induit un morphisme bijectif, $`\tau `$, entre $`𝔪_+`$ et $`𝔪_{}`$ tel que :
$$B(\tau (X),\tau (X))=B(X,X),X𝔪_+$$
Démonstration : L’involution $`\sigma `$ permute, sans point fixe, d’après le Lemme 5, les idéaux simples de $`𝔪`$. Ceux ci sont contenus dans des idéaux simples de $`𝔤`$, donc contenus soit dans $`𝔪_+`$, soit dans $`𝔪_{}`$. Donc $`𝔪=𝔪_+𝔪_{}`$. Si $`\sigma `$ envoyait un idéal simple de $`𝔪_+`$ dans un autre idéal de $`𝔪_+`$, l’algèbre de Lie $`𝔪_+`$, donc aussi $`𝔤_+`$, contiendrait une sous-algèbre semi-simple complexe isotrope. C’est impossible, d’après le Lemme précédent. Donc $`\sigma `$ envoie tout idéal simple de $`𝔪_+`$ dans $`𝔪_{}`$. Le Lemme en résulte immédiatement. Notations On notera $`𝔧_+=𝔧_0𝔤_+`$, $`𝔞_+=𝔞𝔧_+`$. on définit de même $`𝔧_{}`$ et $`𝔞_{}`$. La restriction de $`𝔧_0`$ à $`𝔧_+`$ identifie les racines de $`𝔧_0`$ dans $`𝔤_+`$ à celles de $`𝔧_+`$. On note $`\stackrel{~}{R}_+`$ l’ensemble de celles-ci , $`\mathrm{\Sigma }_+`$, l’ensemble des racines simples de l’ensemble de racines positives, $`\stackrel{~}{R}_+^+`$, de $`\stackrel{~}{R}_+`$, formé des éléments de $`\stackrel{~}{R}_+`$ qui sont des poids de $`𝔧_+`$ dans $`𝔟_0𝔤_+`$. On définit de même $`\stackrel{~}{R}_{}`$, relativement à $`𝔤_{}`$ et $`𝔧_{}`$. On définit aussi $`\mathrm{\Sigma }_{}`$, l’ensemble des racines simples de l’ensemble de racines positives, $`\stackrel{~}{R}_{}^+`$, de $`\stackrel{~}{R}_{}`$, formé des éléments de $`\stackrel{~}{R}_{}`$ qui sont des poids de $`𝔧_{}`$ dans $`𝔟_0^{}𝔤_{}`$. On définit, pour $`𝔦`$ comme dans le Lemme précédent, $`R_+`$, l’ensemble des racines de $`𝔧_+`$ dans $`𝔪_+`$, $`\mathrm{\Gamma }_+=\mathrm{\Sigma }_+R_+`$. Puis on définit comme ci-dessus $`R_{}`$ et $`\mathrm{\Gamma }_{}`$. La restriction de $`\tau `$ à $`𝔞^+:=𝔪_+𝔧_0`$ définit une bijection entre $`𝔞^+`$ et $`𝔞^{}:=𝔪_{}𝔧_0`$, dont l’inverse de la transposée induit une bijection, notée $`A`$, ente $`R_+`$ et $`R_{}`$. On notera, pour $`\alpha R_+`$, $`A\alpha `$, au lieu de $`A(\alpha )`$. Soit $`\alpha \stackrel{~}{R}`$. On note $`H_\alpha `$ l’élément de $`𝔧_0`$, tel que $`\alpha (H_\alpha )=2`$ et qui est orthogonal au noyau de $`\alpha `$ pour la forme de Killing de $`𝔤`$. Soit $`\alpha R_+`$ et $`\beta =A\alpha `$. Alors, on voit facilement que $`H_\alpha `$ (resp. $`H_\beta `$) est élément de $`𝔞^+`$ (resp. $`𝔞^{}`$). Comme $`\tau `$ transporte la forme de Killing de $`𝔪_+`$ sur celle de $`𝔪_{}`$, et que celles-ci sont proportionnelles à la restriction de la forme de Killing de $`𝔤`$, on a :
$$\tau (H_\alpha )=H_\beta $$
Le Lemme 15 implique donc :
$$B(H_{A\alpha },H_{A\beta })=B(H_\alpha ,H_\beta ),\alpha ,\beta \mathrm{\Gamma }_+$$
(3.3)
Soit $`𝔦^{}`$ une autre sous-algèbre Lagrangienne de $`𝔤`$, pour laquelle on introduit des objets similaires, notés avec des . On notera $`C`$, ou parfois $`{}_{}{}^{\prime \prime }A_{}^{1}A^{}^{\prime \prime }`$, l’application définie sur la partie, éventu- ellement vide :
$$domC:=\{\alpha R_+^{}|A^{}\alpha R_{}\}$$
(3.4)
par :
$$C\alpha =A^1A^{}\alpha ,\alpha domC$$
(3.5)
l’image de $`C`$ étant égale à :
$$ImC=\{\alpha R_+|A\alpha R_{}^{}\}$$
(3.6)
###### Lemme 17
Soit $`(B,𝔦,𝔦^{})`$, un triple fortement standard. Avec les notations précédentes, pour tout $`\alpha domC`$, il existe $`n^{}`$ tel que :
$$\alpha ,\mathrm{},C^{n1}\alpha domCetC^n\alpha domC$$
Démonstration : Soit $`\alpha domC`$. Supposons que pour tout $`n^{}`$, $`C^n\alpha `$ soit défini et élément de $`domC`$. Comme $`R_+^{}`$ est un ensemble fini, il existe $`n_1`$, $`n_1^{}^{}`$, distincts, tels que $`C^{n_1}\alpha =C^{n^{}1}\alpha `$. D’où l’on déduit l’existence de $`n^{}`$ tel que $`C^n\alpha =\alpha `$. On note :
$$\alpha _1=\alpha ,\alpha _2=C\alpha ,\mathrm{},\alpha _n=C^{n1}\alpha $$
Montrons que :
$$\{\alpha _1,\mathrm{},\alpha _n\}R_+R_+^{},A(\{\alpha _1,\mathrm{},\alpha _n\})=A^{}(\{\alpha _1,\mathrm{},\alpha _n\})$$
(3.7)
En effet, pour tout $`i`$, $`\alpha _i`$ est élément de $`domCImC`$, ce qui implique la première inclusion. Par ailleurs, pour $`i=1,\mathrm{},n1`$, comme $`{}_{}{}^{\prime \prime }A_{}^{1}A^{}{}_{}{}^{\prime \prime }\alpha _{i}^{}=\alpha _{i+1}`$, on a :
$$A^{}\alpha _i=A\alpha _{i+1},i=1,\mathrm{},n1$$
Maintenant, comme $`{}_{}{}^{\prime \prime }A_{}^{1}A^{}{}_{}{}^{\prime \prime }\alpha _{n}^{}=\alpha _1`$, on a aussi
$$A_{\alpha _n}^{}=A_{\alpha _1}$$
Ceci achève de prouver (3.7). On note $`R_+^0`$ (resp. $`R_{}^0`$), l’intersection de $`R_+`$ (resp. $`R_{}`$) avec le sous-espace vectoriel réel, $`V_0`$, de $`𝔧_0^{}`$, engendré par les $`\alpha _i`$ (resp. $`A\alpha _i`$), $`i=1,\mathrm{},n`$. L’identification de $`𝔧_0`$ à $`𝔧_0^{}`$, à l’aide de la forme de Killing de $`𝔤`$, fait apparaître $`(V_0)_{}`$ comme (le dual d’) un sous-espace vectoriel complexe $`𝔞_0^+`$ de $`𝔞^+`$. On définit de même $`𝔞_0^{}`$. On note $`𝔪_+^0`$ (resp. $`𝔪_{}^0`$) la sous-algèbre de Lie de $`𝔤`$ engendrée par les espaces radiciels $`𝔪_+^\alpha `$, $`\alpha R_+^0`$ (resp. $`𝔪_{}^\alpha `$, $`\alpha R_{}^0`$). On voit immédiatement que $`𝔪_+^0`$ est semi-simple et que :
$$𝔪_+^0=(_{\alpha R_0^+}𝔪_+^\alpha )𝔞_0^+$$
car $`R_+^0`$ est un système de racines dans $`𝔞_0^+`$ (\[Bou\], Ch. 7, Par. 1, Proposition 4). Il en va de même pour $`𝔪_{}^0`$. Alors $`\tau `$ et $`\tau ^{}`$ induisent un isomorphisme entre $`𝔪_+^0`$ et $`𝔪_{}^0`$. Donc $`\tau ^1\tau `$ induit un automorphisme de $`𝔪_0^+`$, qui a donc un point fixe non nul $`X`$. Alors, on a :
$$\tau (X)=\tau ^{}(X)etX+\tau (X)𝔥𝔥^{}𝔦𝔦^{}$$
Une contradiction qui achève de prouver la propriété voulue pour $`C`$.
###### Proposition 7
Soit $`(B,𝔦,𝔦^{})`$, un triple fortement standard. Alors, pour tout $`\alpha R_+`$ (resp. $`R_+^{}`$), $`\alpha `$ et $`A\alpha `$ (resp. $`A^{}\alpha `$), sont de même signes (relativement aux ensembles de racines positives $`\stackrel{~}{R}_+^+`$, $`\stackrel{~}{R}_{}^+`$, définis plus haut)
Début de la démonstration : On raisonne par récurrence sur la dimension de $`𝔤^{der}`$. Si celle-ci est nulle, le résultat est clair. On suppose maintenant que celle-ci n’est pas nulle, et que la Proposition est vraie pour les algèbres réductives dont l’idéal dérivé est de dimension strictement inférieure à celle de $`𝔤^{der}`$. Nous allons commencer par établir plusieurs Lemmes.
###### Lemme 18
Avec les notations précédentes, on a : (i) Si $`\alpha R_+`$ et $`\alpha ImC`$, $`\alpha `$ et $`A\alpha `$ sont de mêmes signes. (ii) Si $`\alpha R_+^{}`$, et $`\alpha DomC`$, $`\alpha `$ et $`A^{}\alpha `$ sont de même signes.
Démonstration : Montrons (i). Raisonnons par l’absurde, et supposons que $`\alpha `$ et $`\beta :=A\alpha `$ soient de signes opposés. L’hypothèse sur $`\alpha `$ équivaut à :
$$\alpha R_+,A\alpha R_{}^{}$$
(3.8)
Quitte à changer $`\alpha `$ en $`\alpha `$, on peut supposer $`\alpha `$ positive. Soit $`X`$ un élément non nul de $`𝔪_+^\alpha 𝔟_0^{}`$. Alors $`\tau (X)𝔪_{}^\beta 𝔟_0^{}`$, d’après notre hypothèse sur $`\alpha `$, et la définition des ensembles de racines positives, $`\stackrel{~}{R}_+^+`$, $`\stackrel{~}{R}_{}^+`$. Enfin $`Y:=X+\tau (X)`$ est un élément non nul de $`𝔥`$. Comme $`𝔪_{}^\beta 𝔟_0^{}`$ et que $`\beta =A\alpha R_{}^{}`$, on a $`𝔪_{}^\beta 𝔫^{}`$. Supposons d’abord $`\alpha R_+^{}`$. Comme $`\alpha `$ est positive, $`𝔪_+^\alpha `$ est contenu dans $`𝔟_0^{}`$ et $`\alpha R_+^{}`$, implique, comme ci-dessus, que $`𝔪_+^\alpha `$ est contenu dans $`𝔫^{}`$. Alors $`X+\tau (X)`$ est un élément non nul de de $`𝔥𝔫^{}𝔦𝔦^{}`$. Une contradiction qui montre qu’on doit avoir $`\alpha R_+^{}`$. Alors $`X𝔪^{}`$, $`\tau (X)𝔫^{}`$, $`Y𝔥𝔭^{}`$ et $`p^𝔫^{}(Y)=X`$. Donc $`𝔪_+^\alpha =𝔤^\alpha `$ est contenu dans $`𝔦_1`$, où $`(B_1,𝔦_1,𝔦_1^{})`$ est le triple antécédent de $`(B,𝔦,𝔦^{})`$. Appliquant le Lemme 13 à $`𝔦_1`$, on voit qu’alors $`𝔤^\alpha `$ est contenu dans $`𝔫_1`$. Mais comme le triple $`(B,𝔦,𝔦^{})`$ est fortement standard les poids de $`𝔧_0`$ dans $`𝔫_1`$ doivent être des poids de $`𝔧_0`$ dans $`𝔟_0`$, ce qui n’est pas le cas de $`\alpha `$. Ceci achève de prouver (i). Pour (ii), l’hypothèse se traduit par une condition analogue à (3.8), en échangeant le rôle de $`𝔦`$ et $`𝔦^{}`$. On déduit donc (ii) de (i).
###### Lemme 19
Si $`\xi R_+ImC`$, il existe $`\alpha domC`$, $`\alpha ImC`$, et $`n^{}`$ tels que :
$$\alpha ,C\alpha ,\mathrm{}C^{n1}\alpha domC,etC^ndomC$$
et vérifiant :
$$\xi =C^i\alpha ,pouruni\{1,\mathrm{},n\}$$
Démonstration : Comme $`C`$ est une bijection de $`domC`$ sur $`ImC`$, car $`A`$ et $`A^{}`$ sont injectives, on note $`C^1`$ la bijection réciproque. Echangeant le rôle de $`𝔦`$ et $`𝔦^{}`$ dans le Lemme 17 , on voit qu’il existe $`n^{}^{}`$ tel que :
$$\xi ,C^1\xi ,\mathrm{},(C^1)^{n^{}1}\xi ImC,(C^1)^n^{}\xi ImC$$
On pose $`\alpha =(C^1)^n^{}\xi domC`$. On a aussi $`\alpha ImC`$. On choisit, grâce au Lemme 17, $`n^{}`$ tel que :
$$\alpha ,C\alpha ,\mathrm{}C^{n1}\alpha domC,etC^n\alpha domC$$
Alors $`\alpha `$ et $`n`$ ont clairement les propriétés voulues.
###### Lemme 20
Soit $`\alpha R_+R_{}`$ tel que $`A\alpha R_{}R_{}^{}`$ (resp. $`A^{}\alpha R_{}R_{}^{}`$). Alors $`\alpha `$ et $`A\alpha `$ (resp. $`A^{}\alpha `$) sont de même signe.
Démonstration : Prouvons l’assertion sur $`A\alpha `$ et soit $`\alpha `$ comme dans l’énoncé. Soit $`X𝔤^\alpha `$. Les hypothèses impliquent que $`𝔤^\alpha `$ et $`𝔤^{A\alpha }`$ sont contenus dans $`𝔪𝔪^{}`$. Donc, on a :
$$𝔤^\alpha (𝔪𝔪^{})\sigma (𝔪𝔪^{})$$
et d’après le Lemme 14, on a $`𝔤^\alpha 𝔪_1`$. De plus, d’après ce même Lemme, l’involution $`\sigma _1`$ est la restriction de $`\sigma `$ à $`𝔪_1`$. On va appliquer l’hypothèse de récurrence du début de la démonstration de la Proposition 7. Pour cela on remarque $`(𝔩𝔩^{})_+=(𝔩𝔩^{})𝔤_+`$ et de même pour $`(𝔩𝔩^{})_{}`$. Donc les racines de $`𝔧_0`$ dans $`𝔩𝔩^{}`$ qui sont positives dans $`𝔩𝔩^{}`$ sont positives dans $`𝔤`$. L’application de l’hypothèse de récurrence montre alors que $`\alpha `$ et $`A\alpha `$ sont de même signe. Ceci achève de prouver (i). Alors l’assertion sur $`A^{}\alpha `$ résulte de celle sur $`A\alpha `$, par échange du rôle de $`𝔦`$ et $`𝔦^{}`$
###### Lemme 21
(i) Soit $`\alpha R_+ImC`$, $`\alpha R_+^{}`$. Alors $`\alpha `$ et $`A\alpha `$ sont de même signe. (ii) Soit $`\alpha R_+^{}domC`$, $`\alpha R_+`$. Alors $`\alpha `$ et $`A^{}\alpha `$ sont de même signe.
Démonstration : Démontrons (i). Soit $`\alpha `$ comme dans l’énoncé. Quitte à changer $`\alpha `$ en $`\alpha `$, on peut supposer que $`\alpha `$ est négative. Raisonnons par l’absurde et supposons $`A\alpha `$ positive. Soit $`X𝔤^\alpha `$. Nos hypothèses montrent que $`𝔤^\alpha `$ est contenu dans $`𝔫^{}`$. Par ailleurs, écrivant $`\alpha =^{\prime \prime }A^1A^{}{}_{}{}^{\prime \prime }\beta `$, où $`\beta domC`$, on a $`A\alpha =A^{}\beta R_{}^{}`$. Donc $`𝔤^{A\alpha }`$ est contenu dans $`𝔪_{}^{}`$. Alors :
$$X+\sigma (X)𝔥𝔭^{},p^𝔫^{}(X+\sigma (X))=\sigma (X)$$
Il en résulte que $`𝔤^{A\alpha }`$ est contenu dans $`𝔦_1`$, donc dans $`𝔫_1`$, d’après le Lemme 13. Alors $`𝔤^{A\alpha }`$ doit être contenu dans $`𝔟_0`$, puisque le triple $`(B_1,𝔦_1,𝔦_1^{})`$ est standard. Mais, comme $`A\alpha R_{}`$ et est positive, ce n’est pas le cas. Une contradiction qui achève de prouver (i). (ii) se déduit de (i) par l’échange de $`𝔦`$ et $`𝔦^{}`$. Fin de la démonstration de la Proposition 7: Soit $`\alpha R_+`$ et montrons que $`A\alpha `$ est de même signe que $`\alpha `$. Distinguons 3 cas. 1) Si $`\alpha ImC`$, cela résulte du Lemme 18. 2) Si $`\alpha ImC`$ et $`\alpha R_+^{}`$, cela résulte du Lemme précédent. 3) Si $`\alpha ImC`$ et $`\alpha R_+^{}`$, on écrit $`\alpha =^{\prime \prime }A^1A^{}{}_{}{}^{\prime \prime }\beta `$, où $`\beta domC`$. Alors on a $`A\alpha =A^{}\beta R_{}R_{}^{}`$. On conclut grâce au Lemme 20. On vient donc de montrer que pour tout $`\alpha R_+`$, $`\alpha `$ et $`A\alpha `$ sont de même signe. On démontre un énoncé similaire pour $`A^{}`$ en échangeant le rôle de $`𝔦`$ et $`𝔦^{}`$. Ceci achève la démonstration de la Proposition. On rappelle que $`H_\alpha `$, $`\alpha \stackrel{~}{R}`$ a été défini avant (3.3). C’est la coracine correspondant à $`\alpha `$. On rappelle qu’un système de générateurs de Weyl de $`𝔤^{der}`$ est une famille $`𝒲=(H_\alpha ,X_\alpha ,Y_\alpha )_{\alpha \mathrm{\Sigma }}`$, où $`\mathrm{\Sigma }=\mathrm{\Sigma }_+\mathrm{\Sigma }_{}`$, telle que, pour tout $`\alpha ,\beta \mathrm{\Sigma }`$, on ait :
$$[X_\alpha ,Y_\beta ]=\delta _{\alpha \beta }H_\beta $$
(3.9)
$$[H_\alpha ,X_\beta ]=N_{\alpha \beta }X_\beta $$
(3.10)
$$[H_\alpha ,Y_\beta ]=N_{\alpha \beta }Y_\beta $$
(3.11)
où:
$$N_{\alpha \beta }=\beta (H_\alpha )=2K_𝔤(H_\alpha ,H_\beta )/K_𝔤(H_\alpha ,H_\alpha )$$
(3.12)
On a alors :
$$N_{\alpha \beta },si\alpha \beta $$
et :
$$adX_\alpha ^{1N_{\alpha \beta }}X_\beta =adY_\alpha ^{1N_{\alpha \beta }}Y_\beta =0,si\alpha \beta $$
(3.13)
Un tel système existe et on obtient les autres par conjugaison par les éléments de $`J_0`$.
###### Définition 5
On appelle donnée de Belavin-Drinfeld généralisée, relativement à $`B`$, la donnée de $`(A,A^{},𝔦_𝔞,𝔦_𝔞^{})`$, où : 1) $`A`$ est une bijection d’une partie $`\mathrm{\Gamma }_+`$ de $`\mathrm{\Sigma }_+`$ sur une partie $`\mathrm{\Gamma }_{}`$ de $`\mathrm{\Sigma }_{}`$, telle que :
$$B(H_{A\alpha },H_{A\beta })=B(H_\alpha ,H_\beta ),\alpha ,\beta \mathrm{\Gamma }_+$$
(3.14)
2) $`A^{}`$ est une bijection d’une partie $`\mathrm{\Gamma }_+^{}`$ de $`\mathrm{\Sigma }_+`$ sur une partie $`\mathrm{\Gamma }_{}^{}`$ de $`\mathrm{\Sigma }_{}`$, telle que :
$$B(H_{A^{}\alpha },H_{A^{}\beta })=B(H_\alpha ,H_\beta ),\alpha ,\beta \mathrm{\Gamma }_+^{}$$
(3.15)
3) On définit $`C=^{\prime \prime }A^1A^{}^{\prime \prime }`$ comme dans (3.4), (3.5). Alors $`C`$ satisfait la ”condition de sortie ” : Pour tout $`\alpha domC`$, il existe $`n^{}`$ tel que $`\alpha ,\mathrm{},C^{n1}\alpha domCetC^n\alpha domC`$. 4) $`𝔦_𝔞`$ (resp. $`𝔦_𝔞^{}`$) est un sous-espace vectoriel complexe de $`𝔧_0`$, contenu et Lagrangien dans l’ orthogonal, $`𝔞`$ (resp. $`𝔞^{}`$) pour la forme de Killing de $`𝔤`$ (ou pour $`B`$), à l’espace engendré par les $`H_\alpha `$, $`\alpha \mathrm{\Gamma }:=\mathrm{\Gamma }_+\mathrm{\Gamma }_{}`$ (resp. $`\mathrm{\Gamma }^{}:=\mathrm{\Gamma }_+^{}\mathrm{\Gamma }_{}^{}`$). 5) Notons $`𝔣`$ le sous-espace de $`𝔧_0`$ engendré par la famille $`H_\alpha +H_{A\alpha }`$, $`\alpha \mathrm{\Gamma }_+`$. On définit de même $`𝔣^{}`$. Alors :
$$(𝔣𝔦_𝔞)(𝔣^{}𝔦_𝔞^{})=\{0\}$$
(3.16)
On notera alors $`R_+`$ le sous-système de racines de $`\stackrel{~}{R}`$ formé des éléments de $`\stackrel{~}{R}`$ qui sont combinaison linéaire d’éléments de $`\mathrm{\Gamma }_+`$. On définit de même $`R_{}`$, $`R_+^{}`$, $`R_{}^{}`$. On notera encore $`A`$ (resp. $`A^{}`$) le prolongement par $``$-linéarité de $`A`$ (resp. $`A^{}`$), qui définit une bijection de $`R_+`$ sur $`R_{}`$ (resp. $`R_+^{}`$ sur $`R_{}^{}`$).
###### Lemme 22
Si $`A`$ vérifie la condition 1) ci dessus, il existe un unique isomorphisme, $`\tau `$, de la sous-algèbre $`𝔪_+`$ de $`𝔤`$, engendrée par les $`X_\alpha ,H_\alpha ,Y_\alpha `$, $`\alpha \mathrm{\Gamma }_+`$, sur la sous-algèbre $`𝔪_{}`$ de $`𝔤`$, engendrée par les $`X_\alpha ,H_\alpha ,Y_\alpha `$, $`\alpha \mathrm{\Gamma }_{}`$, tel que :
$$\tau (H_\alpha )=H_{A\alpha },\tau (X_\alpha )=X_{A\alpha },\tau (Y_\alpha )=Y_{A\alpha },\alpha \mathrm{\Gamma }^+$$
En effet $`𝔪_+`$ et $`𝔪_{}`$sont semi-simples et les familles données sont des systèmes de généra- teurs de Weyl de ces algèbres. Alors, d’après \[Bou\], Chapitre VIII, Paragraphe 4.3, Théorème 1, il suffit, pour montrer le Lemme, de voir que les relations, du type (3.9) à (3.13), satisfaites par ces générateurs se correspondent, c’est à dire qu’il faut montrer :
$$\alpha (H_\beta )=A\alpha (H_{A\beta }),\alpha ,\beta \mathrm{\Gamma }_+$$
(3.17)
Montrons d’abord que, la deuxième définition de $`N_{\alpha \beta }`$ (deuxième égalité de (3.12)), on peut remplacer $`K_𝔤`$ par n’importe qu’elle autre forme $`𝔤`$-invariante non dégénérée. En effet si $`𝔤^\alpha `$ n’est pas dans le même idéal simple que $`𝔤^\beta `$, $`B(H_\alpha ,H_\beta )=K_𝔤(H_\alpha ,H_\beta )=0`$. Sinon $`B`$ et $`K_𝔤`$ sont proportionnelles sur l’idéal simple contenant $`𝔤^\alpha `$ et $`𝔤^\beta `$. D’où notre assertion. Alors (3.17) résulte immédiatement de la condition 1).
###### Proposition 8
Soit $`𝒟=(A,A^{},𝔦_𝔞,𝔦_𝔞^{})`$ une donnée de Belavin-Drinfeld généralisée, relative à $`B`$. On note $`𝔭`$ la sous-algèbre parabolique de $`𝔤`$, contenant $`𝔟_0`$ et $`𝔪`$. Sa décomposition de Langlands $`𝔭=𝔩𝔫`$, où $`𝔩`$ contient $`𝔧_0`$, vérifie $`𝔩=𝔪𝔞`$. On utilise les notations du Lemme précédent. On note $`𝔦:=𝔥𝔦_𝔞𝔫`$, où $`𝔥:=\{X+\tau (X)|X𝔪_+\}`$. On définit de même $`𝔦^{}`$. Alors $`(B,𝔦,𝔦^{})`$ est un triple de Manin fortement standard. On dira que ce triple de Manin est associé à la donnée de Belavin-Drinfeld généralisée, $`𝒟`$, et au système de générateurs de Weyl, $`𝒲`$. On le notera $`𝒯_{𝒟,𝒲}`$. On note $`𝒲_1=(H_\alpha ,X_\alpha ,Y_\alpha )_{\alpha \mathrm{\Gamma }\mathrm{\Gamma }^{}}`$. C’est un système de générateurs de Weyl de $`(𝔩𝔩^{})^{der}`$. Notons $`\mathrm{\Gamma }_{1+}=\mathrm{\Gamma }_+\mathrm{\Gamma }_+^{}A^1(\mathrm{\Gamma }_{}\mathrm{\Gamma }_{}^{})`$. On note $`A_1`$ la restriction de $`A`$ à $`\mathrm{\Gamma }_{1+}`$ et $`\mathrm{\Gamma }_1`$ son image. On définit de même $`A_1^{}`$. On note $`𝔞_1`$ l’intersection des noyaux des éléments de $`\mathrm{\Gamma }_{1+}\mathrm{\Gamma }_1`$ et on note $`𝔦_{𝔞_1}=𝔦_𝔞𝔱_1`$, où $`𝔱_1`$ est l’intersection de $`𝔣`$ avec $`𝔞_1`$. On définit de même $`𝔦_{𝔞_1^{}}`$. Alors $`(A_1,A_1^{},𝔦_{𝔞_1},𝔦_{𝔞_1^{}})`$ est une donnée de Belavin-Drinfeld généralisée, $`𝒟_1`$, pour $`𝔩𝔩^{}`$ et l’antécédent du triple $`𝒯_{𝒟,𝒲}`$ est égal à $`𝒯_{𝒟_1,𝒲_1}`$.
Démonstration : On procède par récurrence sur la dimension de $`𝔤^{der}`$. Si celle-ci est nulle, le résultat est clair. On suppose le résultat est vrai pour les algèbres réductives dont l’idéal dérivé est de dimension strictement inférieure à celle de $`𝔤^{der}`$. Montrons que $`𝔥`$ est isotrope pour $`B`$. Etudions la forme bilinéaire, $`B^{}`$, sur $`𝔪_+`$, définie par :
$$B^{}(X,Y)=B(\tau (X),\tau (Y)),X,Y𝔪_+$$
C’est clairement une forme bilinéaire invariante sur $`𝔪_+`$, qui coincide sur $`𝔞^+=𝔧_0𝔪_+`$ avec la restriction de $`B`$ à $`𝔪_+`$, d’après (3.14) et le Lemme 22. Comme $`𝔞^+`$ est une sous-algèbre de Cartan de $`𝔪_+`$, le Lemme 1 permet de voir que $`B^{}`$ coincide sur $`𝔪_+`$ avec cette restriction. Comme $`𝔤_+`$ et $`𝔤_{}`$ sont orthogonaux pour $`B`$, il en résulte que $`𝔥`$ est isotrope. Par ailleurs, la définition de $`𝔥`$ montre que c’est l’espace des points fixes d’une f-involution, $`\sigma `$, de $`𝔪`$, dont la restriction à $`𝔪_+`$ est égale à $`\tau `$. En particulier $`\sigma `$ permute $`𝔪_+`$ et $`𝔪_{}`$. L’application du Théorème 1 montre que $`𝔦`$ est une sous-algèbre de Lie de $`𝔤`$, Lagrangienne pour $`B`$. On montre de même que $`𝔦^{}`$ est Lagrangienne pour $`B`$. Pour montrer que $`(B,𝔦,𝔦^{})`$ est un triple de Manin de $`𝔤`$, on se propose d’appli- quer le Théorème 3. Montrons que $`𝔥𝔫^{}=\{0\}`$. Soit $`X𝔥𝔫^{}`$. Alors, comme $`𝔫^{}`$ est orthogonal à $`𝔧_0`$, pour la forme de Killing de $`𝔤`$, il en est ainsi de $`X`$. Notons $`R_+`$ l’ensemble des poids non nuls de $`𝔧_0`$ dans $`𝔪_+`$. Ceux-ci sont en particulier nuls sur $`𝔧_{}`$. Alors on a :
$$X=\underset{\alpha R_+}{}(X(\alpha )+\tau (X(\alpha ))),o\stackrel{`}{u}X(\alpha )𝔪_+^\alpha $$
(3.18)
On note que $`\tau (X(\alpha ))`$ est de poids $`\beta `$ sous $`𝔧_0`$, où $`\beta `$ est un poids de $`𝔧_0`$ dans $`𝔪_{}`$, donc nul sur $`𝔧_+`$. La décomposition (3.18) apparait alors comme une décomposition de $`X`$ en vecteurs poids sous $`𝔧_0`$, pour des poids deux à deux distincts. Comme $`X𝔫^{}𝔟_0^{}`$, et que $`𝔪_+𝔤_+`$, il faut que, pour tout $`\alpha R_+`$, avec $`\alpha `$ positive, on ait $`X(\alpha )=0`$. D’autre part, si $`\alpha R_+`$, avec $`\alpha `$ négative, $`\tau (X(\alpha ))`$ est un élément de $`𝔪_{}^\beta `$, avec $`\beta `$ négative. En effet $`\beta `$ est l’image de $`\alpha `$ par le prolongement $``$-linéaire de $`A`$ à l’espace vectoriel réel engendré par $`\mathrm{\Gamma }_+`$. Alors $`\tau (X(\alpha ))`$ appartient à $`𝔟_0`$. Comme $`X𝔫^{}`$, on doit avoir $`\tau (X(\alpha ))=0`$, donc aussi $`X(\alpha )=0`$. Finalement $`X=0`$, comme désiré. Donc $`𝔥𝔫^{}=\{0\}`$. On montre de même $`𝔥^{}𝔫=\{0\}`$. On note
$$\stackrel{~}{𝔥}=𝔥𝔦_𝔞,𝔣=𝔧_0𝔥,\stackrel{~}{𝔣}=𝔣𝔦_𝔞$$
Si $`\alpha R_+`$, on note :
$$𝔥_\alpha =\{X+\tau (X)|X𝔪_+^\alpha \}$$
Alors, d’après (2.41) où l’on prend $`R_{}=R_+`$, on a :
$$\stackrel{~}{𝔥}=\stackrel{~}{𝔣}(_{\alpha _{R_+}}𝔥_\alpha )$$
(3.19)
Cette décomposition apparait comme une décomposition de $`\stackrel{~}{𝔥}`$ en représenta- tions de $`𝔣`$ sans sous-quotients simples équivalents, les $`𝔥_\alpha `$ étant de plus irréduc- tibles. Comme $`𝔭^{}`$contient $`\stackrel{~}{𝔣}`$, on en déduit :
$$\stackrel{~}{𝔥}𝔭^{}=\stackrel{~}{𝔣}(_{\alpha R_+,𝔥_\alpha 𝔭^{}}𝔥_\alpha )$$
(3.20)
On note encore $`A`$ le prolongement $``$-linéaire de $`A`$ au sous-espace vectoriel réel de $`𝔧_0^{}`$ engendré par $`\mathrm{\Gamma }_+`$. On a alors, pour tout $`\alpha R_+`$, $`\sigma (𝔤^\alpha )=𝔤^{A\alpha }`$ . Comme $`𝔭^{}`$ est somme de sous-espaces poids sous $`𝔧_0`$, on a $`𝔥_\alpha 𝔭^{}`$ si et seulement si $`\alpha `$ et $`A\alpha `$ sont des poids de $`𝔧_0`$ dans $`𝔭^{}`$ c’est à dire :
$$\alpha R_+(R_+^{}\stackrel{~}{R}^{})etA\alpha R_{}(R_{}^{}\stackrel{~}{R}^+)$$
(3.21)
Calculons $`p^𝔫^{}(𝔥_\alpha )`$ pour $`\alpha `$ satisfaisant (3.21). On distingue quatre cas. 1)
$$\alpha R_+R_+^{},A\alpha R_{}R_{}^{}$$
Dans ce cas , si $`X𝔤^\alpha `$, on a $`X𝔪_+^{}`$, $`\tau (X)𝔪_{}^{}`$. Alors $`p^𝔫^{}(X+\tau (X))=X+\tau (X)`$. Donc :
$$p^𝔫^{}(𝔥_\alpha )=𝔥_\alpha $$
2)
$$\alpha R_+R_+^{},A\alpha \stackrel{~}{R}^+R_{}^{}$$
On trouve comme ci-dessus que si $`X𝔤^\alpha `$, on a $`X𝔪_+^{}`$, $`\tau (X)𝔫^{}`$. Alors $`p^𝔫^{}(X+\tau (X))=X`$. D’où l’on déduit que :
$$p^𝔫^{}(𝔥_\alpha )=𝔪_+^\alpha $$
3)
$$\alpha \stackrel{~}{R}^{}R_+^{},A\alpha R_{}R_{}^{}$$
Si $`X𝔤^\alpha `$, on a $`X𝔫^{}`$, $`\tau (X)𝔪_{}^{}`$. Alors $`p^𝔫^{}(X+\tau (X))=\tau (X)`$. Donc :
$$p^𝔫^{}(𝔥_\alpha )=𝔪_{}^{A\alpha }$$
4)
$$\alpha \stackrel{~}{R}^{}R_+^{},A\alpha \stackrel{~}{R}^{}R_{}^{}$$
Alors on aurait $`𝔥_\alpha 𝔫^{}`$ et cette possibilité est exclue, d’après ce qu’on a vu plus haut. Notons :
$$𝔫_1:=(_{\alpha R_+^+R_+^{},A\alpha R_{}^{}}𝔪_+^\alpha )(_{\alpha R_+^{}R_+^{},A\alpha R_{}^{}}𝔪_{}^{A\alpha })$$
(3.22)
et
$$𝔪_1=((𝔪𝔪^{}\sigma (𝔪𝔪^{}))^{der},𝔩_1=𝔪_1+𝔧_0$$
(3.23)
L’analyse des poids sous $`𝔧_0`$ montre que
$$𝔩_1=𝔧_0(_{\alpha _{R_1}}𝔤^\alpha )$$
où :
$$R_1=(R_+R_+^{}A^1(R_{}R_{}^{}))(R_{}R_{}^{}A(R_+R_+^{}))$$
En particulier $`𝔩_1`$ et $`𝔫_1`$ ont une intersection réduite à zéro. On note $`\sigma _1`$ la restriction de $`\sigma `$ à $`𝔪_1`$, qui est clairement une f-involution. En effet $`𝔪_1`$ est la somme directe de son intersection avec $`𝔤_+`$ et $`𝔤_{}`$, et $`\sigma _1`$ permute ces deux idéaux. On note $`𝔥_1`$ l’ensemble des points fixes de $`\sigma _1`$. On déduit de la définition de $`𝔪_1`$ ( cf. (3.23)), et de $`\sigma _1`$ que :
$$𝔥_1=(_{\alpha R_+R_+^{},A\alpha R_{}R_{}^{}}𝔥_\alpha )(𝔣𝔪_1)$$
(3.24)
Notons :
$$𝔦_1=p^𝔫^{}(\stackrel{~}{𝔥}𝔭^{})$$
Grâce à (3.20) et ce qui précéde on voit que $`𝔦_1`$ est la somme de son intersection, $`\stackrel{~}{𝔣}`$, avec $`𝔧_0`$, avec son intersection avec l’orthogonal de $`𝔧_0`$ pour la forme de Killing de $`𝔤`$. Tenant compte du fait que $`A`$ préserve le signe des racines, on voit, grâce à la discussion ci-dessus et à (3.20), que l’intersection de $`𝔦_1`$ avec l’orthogonal de $`𝔧_0`$ pour la forme de Killing de $`𝔤`$, coincide avec celle de $`𝔥_1𝔫_1`$. On en déduit facilement l’égalité :
$$𝔦_1=(𝔥_1+\stackrel{~}{𝔣})𝔫_1$$
(3.25)
Montrons que $`𝔭_1:=𝔩_1𝔫_1`$ est une sous-algèbre parabolique de $`𝔩𝔩^{}`$. D’après la définition de $`𝔭_1`$, et de celle de $`𝔪_1`$, $`𝔫_1`$ (cf (3.22), (3.23)), $`𝔭_1`$ est la somme de son intersection $`𝔭_{1+}`$ avec $`𝔤_+`$ avec celle avec $`𝔤_{}`$, $`𝔭_1`$. Il suffit donc d’étudier séparément ces intersections. On ne traite que $`𝔭_{1+}`$, $`𝔭_1`$ se traitant de la même manière. Soit :
$$E=\{\alpha R_+^+R_+^{}|A\alpha R_{}^{}\},F=R_+R_+^{}A^1(R_{}R_{}^{})$$
Alors $`E`$, $`F`$, $`EF`$ sont des parties closes du système de racines de $`𝔧_0`$ dans $`𝔩𝔩^{}𝔤_+`$, $`R_+R_+^{}`$, $`EF`$ en étant une partie parabolique, contenant $`R_+^+R_+^+`$ et dont $`E`$ est un idéal (cf \[War\], 1.1.2.9, 1.1.2.13, pour la terminologie). Cela résulte du fait que $`R_+^+`$, $`R_+^{}`$ et $`R_{}^{}`$ sont des parties closes, d’après leur définition (cf. fin de la Définition 5). Le seul point non immédiat est le fait que si $`\alpha E`$, $`\beta F`$ et $`\alpha +\beta \stackrel{~}{R}`$, alors $`\alpha +\beta `$ appartient à $`E`$. D’après la définition de $`R_{}^{}`$, on voit que nos hypothèses impliquent que $`A(\alpha +\beta )R_{}^{}`$. Il reste à voir que $`\alpha +\beta `$ est positive. L’étude de ses composantes dans la base $`\mathrm{\Sigma }`$, montre que l’une de celles-ci, correspondant à une racine $`\gamma \mathrm{\Gamma }_+A^1\mathrm{\Gamma }_{}^{}`$, est strictement positive, car pour l’une au moins de ces racines , la composante de $`\alpha `$ est strictement positive tandis que celle de $`\beta `$ est nulle. Cela implique que $`\alpha +\beta `$ est positive et finalement dans $`E`$. Joint à la définition de $`𝔭_{1+}`$, et celles de $`𝔪_1`$ et $`𝔫_1`$ , cela implique que $`𝔭_{1+}`$ est une sous-algèbre parabolique de $`𝔩𝔩^{}𝔤_+`$. On conclut que $`𝔭_1`$ est une sous-algèbre parabolique de $`𝔩𝔩^{}`$, qui contient $`𝔟_0𝔩𝔩^{}`$. La démonstration montre aussi que $`𝔫_1`$ en est le radical nilpotent. Notons $`𝔣_1=𝔣𝔪_1`$ et $`𝔱_1`$ l’intersection de $`𝔣`$ avec le centre $`𝔞_1`$ de $`𝔩_1`$. Montrons que :
$$𝔣=𝔣_1𝔱_1$$
(3.26)
En effet tout élément de $`𝔣`$ est de la forme $`X+\tau (X)`$, où $`X`$ est un élément de $`𝔧_0𝔪_+`$. On note que $`𝔪`$ et $`𝔪_1`$ sont la somme directe de leurs intersection avec $`𝔤_+`$ et $`𝔤_{}`$. Il en va de même de $`𝔧_0`$, $`𝔞`$, $`𝔞_1`$. Alors $`X`$ est la somme d’un élément de $`𝔧_0𝔪_+`$ avec un élément de $`𝔞_1𝔤_+`$. La décomposition ci-dessus en résulte aussitot. On pose :
$$𝔦_{𝔞_1}=𝔱_1𝔦_𝔞𝔞_1$$
(3.27)
de sorte que (3.25) se réécrit
$$𝔦_1=𝔥_1𝔦_{𝔞_1}𝔫_1$$
On a :
$$dim_{}𝔧_0=dim_{}(𝔧_0𝔪_+)+dim_{}(𝔧_0𝔪_{})+dim_{}𝔞=dim_{}𝔣+dim_{}𝔞$$
Comme $`𝔦_𝔞`$ est Lagrangienne dans $`𝔞`$, tenant compte de ce qui précède, on a :
$$dim_{}𝔧_0=dim_{}𝔣_1+dim_{}𝔱_1+dim_{}𝔦_𝔞$$
(3.28)
d’où l’on déduit :
$$dim_{}𝔧_0=dim_{}(𝔧_0𝔪_1)+dim_{}𝔞_1$$
(3.29)
Montrons :
$$dim_{}(𝔧_0𝔪_1)=dim_{}𝔣_1$$
(3.30)
En effet, $`𝔣_1=𝔣𝔪_1`$, vérifie :
$$𝔣_1=\{X+\tau (X)|X𝔧_0𝔪_{1+}\}$$
et
$$𝔧_0𝔪_1=(𝔧_0𝔪_{1+})(𝔧_0𝔪_1)$$
les deux facteurs étant échangés par $`\tau `$. (3.30) en résulte. On déduit de (3.28) à (3.30) que :
$$dim_{}𝔦_{𝔞_1}=dim_{}𝔞_1$$
Comme $`𝔱_1𝔥`$ et $`𝔦_𝔞`$ sont isotropes pour $`B`$, et orthogonales pour $`B`$, puisque le centre d’une algèbre réductive est orthogonal à son idéal dérivé pour toute forme invariante, ce qui précède montre que :
$$𝔦_{𝔞_1}estLagrangiennedans𝔞_1$$
(3.31)
Cela implique que $`𝔦_1`$ est Lagrangienne pour $`B`$. On introduit de la même manière $`𝔦_1^{}`$. On prend $`A_1`$ égal à la restriction de $`A`$ à $`\mathrm{\Gamma }_{1+}=\mathrm{\Gamma }_+\mathrm{\Gamma }_+^{}A^1(\mathrm{\Gamma }_{}\mathrm{\Gamma }_{}^{})`$. On définit de même $`A_1^{}`$. Montrons que $`(A_1,A_1^{},𝔦_{a_1},𝔦_{𝔞_1^{}})`$ est une donnée de Belavin-Drinfeld généralisée. Les conditions 1), 2), 3) résultent immédiatement des conditions satisfaites par $`(A,A^{},𝔦_𝔞,𝔦_𝔞^{})`$. La condition 4) résulte de (3.31). Enfin 5) résulte de la condition 5) pour $`𝒟`$, joint à (3.26) et (3.27). L’application de l’hypothèse de récurrence montre que $`(B_1,𝔦_1,𝔦_1^{})`$ un triple de Manin associé à $`𝒟_1`$ et $`𝒲_1`$. Le Théorème 3 permet de conclure que $`(B,𝔦,𝔦^{})`$ est un triple de Manin, d’antécédent $`(B_1,𝔦_1,𝔦_1^{})`$. Ceci achève la preuve de la Proposition.
###### Théorème 5
(i) Tout triple de la forme $`𝒯_{𝒟,𝒲}`$, où $`𝒟`$ est une donnée de Belavin-Drinfeld généra- lisées, et $`𝒲`$ un ensemble de générateurs de Weyl de $`𝔤`$, est un triple de Manin fortement standard. (ii) Réciproquement tout triple fortement standard est de cette forme. (iii) Soit $`𝒯_{𝒟,𝒲}`$ (resp $`𝒯_{\underset{¯}{𝒟},\underset{¯}{𝒲}}`$), où $`𝒟`$, $`\underset{¯}{𝒟}`$ sont des données de Belavin-Drinfeld généralisées, et $`𝒲`$, $`\underset{¯}{𝒲}`$ des ensembles de générateurs de Weyl de $`𝔤^{der}`$. Ces triples de Manin sont conjugués sous $`G`$, si et seulement si $`𝒟=\underset{¯}{𝒟}`$. Alors ils sont conjugués par l’élément de $`J_0`$ qui conjugue $`𝒲`$ et $`\underset{¯}{𝒲}`$.
Démonstration : Le point (i) a été vu à la Proposition précédente. Prouvons (ii). Soit $`(B,𝔦,𝔦^{})`$ un triple fortement standard sous $`(𝔭,𝔭^{})`$. On utilise les notations qui suivent le Lemme 16. D’après la Proposition 7, l’application $`A`$ est une bijection de $`R_+`$ sur $`R_{}`$ qui préserve les signes des racines. Donc elle induit une bijection de $`\mathrm{\Gamma }_+`$ sur $`\mathrm{\Gamma }_{}`$, notée encore $`A`$. On a des propriétés analogues pour $`A^{}`$. On note $`𝔦_𝔞=𝔦𝔞`$, $`𝔦_𝔞^{}^{}=𝔦^{}𝔞^{}`$. On voit facilement que $`𝔣`$ est l’espace vectoriel engendré par les $`H_\alpha +H_{A\alpha }`$, $`\alpha \mathrm{\Gamma }_+`$ et que $`𝔣𝔦_𝔞=𝔦𝔧_0`$. On a des propriétés similaires pour $`𝔦^{}`$. Alors (3.16) résulte du fait que $`𝔦`$ et $`𝔦^{}`$ ont une intersection réduite à zéro. D’après le Théorème 1, le Lemme 17 et (3.3), $`(A,A^{},𝔦_𝔞,𝔦_𝔞^{}^{})`$ est une donnée de Belavin-Drinfeld généralisée, notée $`𝒟`$. Il reste à trouver un système de générateurs de Weyl de $`𝔤`$, vérifiant les relations du Lemme 22, avec $`\tau `$ comme dans le Lemme 16. Si $`\alpha \mathrm{\Gamma }_+^{}`$ est comme dans le Lemme 18, i.e. $`\alpha domC`$, $`\alpha ImC`$, on choisit $`X_\alpha 𝔤^\alpha `$, $`Y_\alpha 𝔤^\alpha `$, tels que $`[X_\alpha ,Y_\alpha ]=H_\alpha `$. Puis notant $`\alpha _i=C^i\alpha `$, $`i=0,\mathrm{},n`$, on définit par récurrence sur $`i`$, pour $`Z=X`$ ou $`Z=Y`$ :
$$Z_{\alpha _{i+1}}:=\tau ^1\tau ^{}(Z_{\alpha _i})$$
(3.32)
l’expression du membre de droite étant bien définie, pour $`i=0,\mathrm{},n1`$, car alors $`\alpha _idomC`$. Puis on pose, pour $`Z=X`$ ou $`Z=Y`$ :
$$Z_{A\alpha _i}=\tau (Z_{\alpha _i}),si\alpha _i\mathrm{\Gamma }_+,Z_{A^{}\alpha _i}=\tau ^{}(Z_{\alpha _i}),si\alpha _i\mathrm{\Gamma }_+^{}$$
(3.33)
Cette définition est cohérente, car si $`\beta =A\alpha _i=A^{}\alpha _i^{}`$, on a $`\alpha _i^{}domC`$ et $`C\alpha _i^{}=\alpha _i`$. Mais alors on a $`i=i^{}+1`$. Il résulte de (3.32) que les deux définitions de $`Z_\beta `$ de (3.33) coincident. Maintenant, soit un élément $`\beta `$ de $`\mathrm{\Gamma }_+\mathrm{\Gamma }_{}`$ qui n’est pas de la forme $`\alpha _i`$, pour un $`\alpha `$ comme ci-dessus. En particulier on a $`\beta domC`$, $`\beta ImC`$. On choisit $`X_\beta 𝔤^\beta `$, $`Y_\beta 𝔤^\beta `$, tels que $`[X_\beta ,Y_\beta ]=H_\beta `$, puis on pose, pour $`Z=X`$ ou $`Z=Y`$ :
$$Z_{A\beta }=\tau (Z_\beta )si\beta \mathrm{\Gamma }_+,Z_{A^{}\beta }=\tau ^{}(Z_\beta )si\beta \mathrm{\Gamma }_+^{}$$
(3.34)
Montrons que cette définition est cohérente. Supposons que $`A\beta `$ soit défini et égal à l’ un des $`A^{}\alpha _i`$ ci dessus. On aurait alors $`\beta ImC`$, ce qui est impossible. Comme $`\beta domC`$, on voit de même que $`A^{}\beta `$, s’il est défini ne peut-être égal à l’un des $`A\alpha _i`$. Enfin si $`A\beta =A^{}\beta ^{}`$, pour deux éléments $`\beta `$, $`\beta ^{}`$ comme ci-dessus,on aurait $`\beta ImC`$, ce qui n’est pas. Les autres égalités à envisager pour voir la cohérence de (3.33) et (3.34) étant exclues, d’après la bijectivité de $`A`$ et $`A^{}`$, cette cohérence est donc prouvée. Enfin si $`\alpha \mathrm{\Sigma }`$, $`\alpha \mathrm{\Gamma }_+\mathrm{\Gamma }_{}`$, on choisit $`X_\alpha 𝔤^\alpha `$, $`Y_\alpha 𝔤^\alpha `$, tels que $`[X_\alpha ,Y_\alpha ]=H_\alpha `$. Il est alors facile de voir que la famille $`(H_\alpha ,X_\alpha ,Y_\alpha )_{\alpha \mathrm{\Sigma }}`$ est un système de générateurs de Weyl, $`𝒲`$. De plus la définition de $`\tau `$, $`\tau ^{}`$ et (3.33) (3.34) montrent que le triple $`(B,𝔦,𝔦^{})`$ est égal à $`𝒯_{𝒟,𝒲}`$. Ceci achève la preuve de (ii). Prouvons (iii). Supposons les deux triples conjugués par un élément de $`G`$. On note le premier triple $`(B,𝔦,𝔦^{})`$ qu’on suppose sous $`(𝔭,𝔭^{})`$, on note $`𝒟=(A,A^{},𝔦_𝔞,𝔦_𝔞^{})`$ et on introduit des notations similaires pour le deuxième triple, en soulignant. Montrons que :
$$𝔦𝔧_0=\underset{¯}{𝔦}𝔧_0,𝔦^{}𝔧_0=\underset{¯}{𝔦}^{}𝔧_0$$
(3.35)
On procède par récurrence sur la dimension de $`𝔤^{der}`$, le résultat étant clair si celle-ci est nulle. Par ailleurs, comme deux sous-algèbres paraboliques standard conjuguées sous $`G`$ sont égales, on a
$$(𝔭,𝔭^{})=(\underset{¯}{𝔭},\underset{¯}{𝔭}^{})$$
(3.36)
et les deux triples sont conjugués par un élément de $`PP^{}`$. De la Proposition 4, on déduit que les antécédents des deux triples, qui sont fortement standard, sont conjugués par un élément de $`LL^{}`$. L’application de l’hypothèse de récurrence conduit au résultat voulu, car si $`(B_1,𝔦_1,𝔦_1^{})`$ est l’antécédent de $`(B,𝔦,𝔦^{})`$, la définition des triples fortement standard montre que $`𝔦𝔧_0=𝔦_1𝔧_0`$, etc. L’égalité (3.36) montre que $`𝔪=\underset{¯}{𝔪}`$, donc $`\mathrm{\Gamma }=\underset{¯}{\mathrm{\Gamma }}`$. De (3.35), on déduit l’égalité de $`𝔣:=𝔪𝔦𝔧_0`$ avec $`\underset{¯}{𝔣}:=𝔪\underset{¯}{𝔦}𝔧_0`$. Comme la définition de $`𝒯_{𝒟,𝒲}`$ montre que $`𝔣`$ est engendré par $`(H_\alpha +H_{A\alpha })_{\alpha \mathrm{\Gamma }_+}`$ et de même pour $`\underset{¯}{𝔣}`$, l’égalité de $`A`$ et $`\underset{¯}{A}`$ en résulte immédiatement. Il en va de même de l’égalité de $`A^{}`$ et $`\underset{¯}{A}^{}`$. Comme $`𝔦_𝔞=𝔦𝔞`$ et de même pour $`\underset{¯}{𝔦}_{\underset{¯}{𝔞}}`$, on déduit l’égalité de ces espaces de (3.35) , car $`𝔞=\underset{¯}{𝔞}`$ et $`𝔦_𝔞=𝔦𝔧_0𝔞`$ et de même pour $`\underset{¯}{𝔦}_{\underset{¯}{𝔞}}`$. On procède de même pour $`𝔦_𝔞^{}^{}`$, $`\underset{¯}{𝔦}_{}^{}{}_{\underset{¯}{𝔞}^{}}{}^{}`$. Donc $`𝒟`$ est égal à $`\underset{¯}{𝒟}`$, comme désiré. Maintenant, il est clair qu’un élément de $`J_0`$ qui conjugue $`𝒲`$ et $`\underset{¯}{𝒲}`$, conjugue les deux triples.
###### Remarque 3
Soit $`𝔤_1`$ une algèbre de Lie simple complexe, $`𝔧_1`$, une sous-algèbre de Cartan de $`𝔤_1`$, $`𝔤=𝔤_1\times 𝔤_1`$, $`𝔧_0=𝔧_1\times 𝔧_1`$ et $`B`$ la forme $``$-bilinéaire sur $`𝔤`$, $`𝔤`$-invariante, égale à $`K_{𝔤_1}`$ sur le premier facteur et à $`K_{𝔤_1}`$ sur le deuxième facteur. La classification de Belavin et Drinfeld de certaines $`R`$-matrices (cf. \[BD\], Théorème 6.1) se réduit, d’après l.c. équations 6.1 à 6.5 , et \[S\], Propositions 1 et 2, à la classification des triples de Manin pour $`𝔤`$, $`(B,𝔦,𝔦^{})`$, où $`𝔦`$ est égal à la diagonale $`diag(𝔤_1)`$ de $`𝔤_1\times 𝔤_1`$, modulo la conjugaison par la diagnale de $`G_1\times G_1`$ . Ceci se fait simplement à l’aide du Théorème précédent.
On remarque que $`𝔤_+=𝔤_1\times \{0\}`$, $`𝔤_{}=\{0\}\times 𝔤_1`$. On fixe pour cela une sous-algèbre de Borel $`𝔟_1`$ de $`𝔤_1`$. On note $`𝔟_1^{}`$ la sous-algèbre de Borel opposée, relativement à $`𝔧_1`$. On pose $`𝔟_0=𝔟_1\times 𝔟_1^{}`$. Soit $`𝒲_1`$ un système de générateurs de Weyl de $`𝔤_1`$, relativement à l’ensemble, $`\mathrm{\Sigma }_1`$, des racines simples de l’ensemble des racines de $`𝔧_1`$ dans $`𝔟_1`$. On note $`𝒲`$ le système de générateurs de Weyl de $`𝔤`$ égal à $`(𝒲_1\times \{0\})(\{0\}𝒲_1`$. D’après le Théorème précédent, il existe une unique donnée de Belavin-Drinfeld généralisée, $`𝒟=(A,A^{},𝔦_a,𝔦_𝔞^{})`$ telle que $`(B,𝔦,𝔦^{})`$ soit conjugué à $`𝒯_{𝒟,𝒲}`$. Il est alors facile de voir que $`𝔦`$ est sous $`𝔤`$. Alors $`𝔞`$ et $`𝔦_𝔞`$ sont réduits à zéro et $`\mathrm{\Gamma }_+=\mathrm{\Sigma }_1\times \{0\}`$. Par ailleurs, la conjugaison des triples se traduit par le fait que l’isomorphisme $`\tau `$ du Lemme 22, de $`𝔤_1`$ sur $`𝔤_1`$, est un automorphisme intérieur de $`𝔤_1`$. Par ailleurs, il préserve $`𝔧_1`$ et induit une permutation, $`A`$, de l’ensemble $`\mathrm{\Sigma }_1`$. Cette permutation doit donc être triviale, i.e. $`A`$ est l’identité (cf. \[Bou\], Chapitre VIII, Paragraphe 5.2). Alors $`\tau `$ est l’identité et la première sous-algèbre isotrope de $`𝒯_{𝒟,𝒲}`$ est la diagonale. L’élément de $`G`$ qui conjugue les deux triples stabilise donc la diagonale. C’est donc un élément de la diagonale. Il résulte de la discussion précédente et du Théorème 5, que l’ensemble $`𝒯_{𝒟,𝒲}`$, où $`𝒟`$ décrit l’ensemble des données de Belavin-Drinfeld généralisées telles que $`A`$ est l’identité de $`\mathrm{\Sigma }_1`$, $`𝔦_𝔞=\{0\}`$, classifie les triples de Manin pour $`𝔤`$, $`(B,𝔦,𝔦^{})`$, où $`𝔦`$ est égal à la diagonale $`diag(𝔤_1)`$ de $`𝔤_1\times 𝔤_1`$, modulo la conjugaison par la diagnale de $`G_1\times G_1`$. Ceci redonne le résultat de Belavin-Drinfeld \[BD\], Théorème 6.1.
## 4 Triples de Manin réels pour une algèbre semi-simple complexe
### 4.1
Soit $`𝔤_1`$ une algèbre de Lie semi-simple complexe, soit $`𝔤_1`$ une forme réelle déployée de $`𝔤_1`$ et $`𝔟_1`$ une sous-algèbre de Borel de $`𝔤_1`$, complexifiée d’une sous-algèbre de Borel, $`𝔟_1`$, de $`𝔤_1`$. Soit $`𝔧_1`$ une sous-algèbre de Cartan déployée de $`𝔤_1`$, contenue dans $`𝔟_1`$ et $`𝔧_1`$ sa complexifiée. On note $`X\overline{X}`$ la conjugaison de $`𝔤_1`$ par rapport à sa forme réelle $`𝔤_1`$. On note $`\eta `$ l’application de $`𝔤_1`$ dans $`𝔤:=𝔤_1\times 𝔤_1`$, définie par :
$$\eta (X)=(X,\overline{X}),X𝔤_1$$
Alors $`\eta (𝔤_1)`$ est une forme réelle de $`𝔤`$ et la conjugaison, $`j`$, par rapport à cette forme réelle vérifie :
$$j(X,Y)=(\overline{Y},\overline{X}),X,Y𝔤_1$$
On note $`𝔟_0:=𝔟_1\times 𝔟_1`$ et $`𝔧_0:=𝔧_1\times 𝔧_1`$ Si $`V`$ est un sous-espace vectoriel réel de $`𝔤_1`$, on noptera $`V_{}=\eta (V)+i\eta (V)𝔤`$. On utilisera les Notations qui suivent le Lemme 16 pour $`𝔤`$. Si $`B`$ est une forme $``$-bilinéaire invariante sur $`𝔤_1`$, on note $`B_{}`$, l’unique forme $``$-bilinéaire invariante sur $`𝔤`$, telle que :
$$B_{}(\eta (X),\eta (X^{}))=B(X,X^{}),X,X^{}𝔤_1$$
(4.1)
On voit aisément que, si $`𝔰`$, est un idéal simple de $`𝔤_1`$, et si la restriction de $`B`$ à $`𝔰`$ est égal à $`Im\lambda K_𝔰`$, pour $`\lambda `$, on a :
$$B_{}((X,Y),(X^{},Y^{}))=\lambda K_𝔰(X,X^{})\overline{\lambda }K_𝔰(Y,Y^{}),X,X^{},Y,Y^{}𝔰$$
(4.2)
La démonstration de la Proposition suivante est immédiate.
###### Proposition 9
On fixe une une forme de Manin réelle sur $`𝔤_1`$. L’application qui à un un triple de Manin réel dans $`𝔤_1`$, $`(B,𝔦,𝔦^{})`$, associe $`(B_{},𝔦_{},𝔦_{}^{})`$, est une bijection entre l’ensemble des triples de Manin réels de $`𝔤_1`$, associés à $`B`$, et l’ensemble des triples de Manin complexes de $`𝔤`$, associés à $`B_{}`$, et pour lesquels chacune des sous-algèbres Lagrangiennes est stable par $`j`$. Cette bijection transforme triples fortement standard, relativement à $`𝔟_1`$, $`𝔧_1`$, en triples fortement standard relativement à $`𝔟_0`$, $`𝔧_0`$. En outre elle transforme les triples conjugués par $`G_1`$ en triples conjugués par $`\eta (G_1)`$
Hypothèse On suppose que, pour tout déal simple de $`𝔤_1`$, la restriction de $`B`$ à $`𝔰`$ est égal à $`Im\lambda K_𝔰`$, pour un $`\lambda `$ réel. On note $`𝔤_{1+}`$ la somme des idéaux simples $`𝔰`$ de $`𝔤_1`$, tels que la restriction de $`B`$ à $`𝔰`$ est égal à $`Im\lambda K_𝔰`$, pour $`\lambda ^+`$. On définit de même $`𝔤_1`$. Alors, au vu de (4.2), on a, pour la forme $`B_{}`$ :
$$𝔤_+=𝔤_{1+}\times 𝔤_1,𝔤_{}=𝔤_1\times 𝔤_{1+}$$
(4.3)
On note $`\stackrel{~}{R}_1`$ l’ensemble des racines de $`𝔧_1`$ dans $`𝔤_1`$ , où $``$ vaut + ou -. Alors, avec les Notations qui suivent le Lemme 16, on a :
$$\stackrel{~}{R}_+=(\stackrel{~}{R}_{1+}\times \{0\})(\{0\}\times \stackrel{~}{R}_1)$$
On note $`\mathrm{\Sigma }_{1+}`$ (resp. $`\mathrm{\Sigma }_1`$) les racines simples de $`𝔧_{1+}:=𝔧_1𝔤_{1+}`$ dans $`𝔟_1𝔤_{1+}`$ (resp. $`𝔧_1:=𝔧_1𝔤_1`$ dans $`𝔟_1^{}𝔤_1`$, où $`𝔟_1^{}`$ est la sous-algèbre de Borel de $`𝔤_1`$ opposée à $`𝔟_1`$, relativement à $`𝔧_1`$)), qu’on identifie à des racines de $`𝔧_1`$ dans $`𝔤_1`$. Alors on a :
$$\mathrm{\Sigma }_+=(\mathrm{\Sigma }_{1+}\times \{0\})(\{0\}\times \mathrm{\Sigma }_1),\mathrm{\Sigma }_{}=((\mathrm{\Sigma }_1)\times \{0\})(\{0\}\times (\mathrm{\Sigma }_{1+}))$$
(4.4)
Soit $`𝒲_1`$ un système de générateurs de Weyl de $`𝔤_1`$, relativement à $`\mathrm{\Sigma }_1`$, dont tous les éléments sont dans $`𝔤_1`$. Soit $`𝒲`$ le système de générateurs de Weyl de $`𝔤`$, relativement à $`\mathrm{\Sigma }:=\mathrm{\Sigma }_+\mathrm{\Sigma }_{}`$, et défini comme suit, en tenant compte de (4.4) . Si $`\alpha \mathrm{\Sigma }_{1+}`$ on pose :
$$H_{(\alpha ,0)}=(H_\alpha ,0),X_{(\alpha ,0)}=(X_\alpha ,0),Y_{(\alpha ,0)}=(Y_\alpha ,0)$$
(4.5)
$$H_{(0,\alpha )}=(0,H_\alpha ),X_{(0,\alpha )}=(0,Y_\alpha ),Y_{(0,\alpha )}=(0,X_\alpha )$$
(4.6)
Si $`\alpha \mathrm{\Sigma }_1`$, on pose :
$$H_{(0,\alpha )}=(0,H_\alpha ),X_{(0,\alpha )}=(0,X_\alpha ),Y_{(0,\alpha )}=(0,Y_\alpha )$$
(4.7)
$$H_{(\alpha ,0)}=(H_\alpha ,0),X_{(\alpha ,0)}=(Y_\alpha ,0),Y_{(\alpha ,0)}=(X_\alpha ,0)$$
(4.8)
On notera $`\alpha \alpha ^f`$ l’échange des facteurs dans $`𝔧_1^{}\times 𝔧_1^{}=𝔧_0^{}`$. On fait de même dans $`J_0`$.
###### Proposition 10
On fixe une forme de Manin réelle, $`B`$, sur $`𝔤_1`$. Soit $`𝒟=(A,A^{},𝔦_𝔞,𝔦_𝔞^{})`$ une donnée de Belavin-Drinfeld généralisée pour $`𝔤`$ et $`B_{}`$. (i) Pour $`tJ_0`$, le triple de Manin $`(B_{},𝔦,𝔦^{})=𝒯_{𝒟,t𝒲}`$ est le complexifié d’un triple réel de $`𝔤_1`$, relativement à $`B`$, si et seulement si : 1) l’application $`\alpha \alpha ^f`$ induit une bijection de $`\mathrm{\Gamma }_+`$ sur $`\mathrm{\Gamma }_{}`$ et l’on a, en prolongeant $`A`$ par $``$-linéarité :
$$(A\alpha )^f=A^1(\alpha ^f),\alpha \mathrm{\Gamma }_+$$
2) A’ vérifie des conditions similaires. 3) $`𝔦_𝔞`$ et $`𝔦_𝔞^{}`$ sont stables par $`j`$ 4)
$$Posantu:=\overline{t}(t^f)^1,onau^\alpha =u^{A\alpha },\alpha \mathrm{\Gamma }_+$$
5) L’élément $`t`$ de $`J_0`$ vérifie des conditions similaires relativement à $`A^{}`$. (ii) On fixe $`𝒟`$ et $`t`$ vérifiant les conditions ci-dessus. Soit $`t_1J_1`$ et $`t^{}=(t_1,\overline{t}_1)t`$. Alors $`𝒟`$ et $`t^{}`$ vérifient les conditions ci-dessus, et les triples complexes $`𝒯_{𝒟,t𝒲}`$, $`𝒯_{𝒟,t^{}𝒲}`$ sont les complexifiés de triples réels de $`𝔤_1`$, conjugués par $`t_1`$.
Démonstration : Si l’automorphisme $``$-linéaire de $`𝔤`$, $`j`$, laisse $`𝔦`$ invariant, il laisse invariant le radical nilpotent $`𝔫`$ de $`𝔦`$, donc aussi $`𝔭`$, qui est le normalisateur de $`𝔫`$. Par ailleurs $`𝔧_0`$ est contenu dans $`𝔭`$ et est invariant par $`j`$. Donc $`j(𝔩)=𝔩`$, d’où l’on déduit $`j(𝔪)=𝔪`$. Alors $`𝔦=𝔥𝔦_𝔞𝔫`$ est invariant par $`j`$ si et seulement si :
$$j(𝔥)=𝔥,j(𝔦_𝔞)=𝔦_𝔞,j(𝔫)=𝔫$$
(4.9)
La seconde égalité de l’équation précédente conduit à 3). L’égalité (cf (3.19)) :
$$𝔥=𝔣(_{\alpha R_+}𝔥_\alpha )$$
(4.10)
et la stabilité de $`𝔧_0`$, et de son orthogonal pour la forme de Killing par $`j`$, montre que la première égalité de (4.9) implique :
$$j(𝔣)=𝔣$$
Mais $`𝔣`$ est engendré par les $`H_\alpha +H_{A\alpha }`$, $`\alpha \mathrm{\Gamma }_+`$. De plus $`j(H_\alpha +H_{A\alpha })`$ est égal à $`H_{\alpha ^f}+H_{(A\alpha )^f}`$. Ce dernier doit être une combinaison linéaire de $`H_\beta +H_{A\beta }`$, $`\beta \mathrm{\Gamma }_+`$. Mais on a :
$$H_\beta 𝔤_+,etH_{A\beta }𝔤_{},H_{\alpha ^f}𝔤_{},H_{(A\alpha )^f}𝔤_+$$
car $`f`$ échange $`𝔤_+`$ et $`𝔤_{}`$ (voir (4.3)). Comme $`f`$ envoie chaque élément de $`\mathrm{\Sigma }`$ sur l’opposé d’un élément de $`\mathrm{\Sigma }`$ (cf. (4.4)), la projection, sur $`𝔤_+`$ et $`𝔤_{}`$, de l’écriture de $`H_{\alpha ^f}+H_{(A\alpha )^f}`$ dans la base de $`𝔣`$, montre qu’il existe $`\beta \mathrm{\Gamma }_+`$ telle que :
$$(A\alpha )^f=\beta ,(\alpha )^f=A\beta $$
(4.11)
Ceci implique immédiatement la condition 1). Comme :
$$j(𝔤^\alpha )=𝔤^{\alpha ^f},\alpha \stackrel{~}{R}$$
au vu de (4.10) et (4.11),la stabilité de $`𝔥`$ par $`j`$ implique alors :
$$j(𝔥_\alpha )=𝔥_\beta $$
$`\alpha \mathrm{\Gamma }_+`$ et $`\beta `$ est comme ci-dessus. Mais, la définition de $`𝔦`$ (cf. Proposition 8 et Lemme 22) montre que $`𝔥_\alpha `$ a pour base :
$$U_\alpha :=t^\alpha X_\alpha +t^{A\alpha }X_{A_\alpha }$$
et $`𝔥_\beta `$ a pour base:
$$V_\beta :=t^\beta Y_\beta +t^{A\beta }Y_{A_\beta }$$
Par ailleurs la définition de $`j`$ et celle de $`𝒲`$ montrent que :
$$j(X_\alpha )=Y_{\alpha ^f},j(X_{A\alpha })=Y_{(A\alpha )^f}$$
Donc, on a :
$$j(U_\alpha ):=\overline{t}^\alpha Y_{\alpha ^f}+\overline{t}^{A\alpha }Y_{(A_\alpha )^f}$$
L’écriture de la proportionnalité de $`j(U_\alpha )`$ à $`V_\beta `$, conduit à :
$$\overline{t}^\alpha t^{A\beta }=\overline{t}^{A\alpha }t^\beta ,\alpha \mathrm{\Gamma }_+$$
En tenant compte de (4.11), ceci implique la condition 4). On procède de même pour $`𝔦^{}`$. Ce qui précède montre que les conditions 1) à 5) sont nécessaires pour que le triple donné soit le complexifié d’un triple réel. Réciproquement, montrons que si ces conditions sont satisfaites le triple donné est bien le complexifié d’un triple réel. En effet, la deuxième égalité de (4.9) est alors satisfaite. Comme les racines de $`𝔧_0`$ dans $`𝔫`$ sont celles de $`𝔧_0`$ dans $`𝔟_1\times 𝔟_1`$ qui ne sont pas combinaison linéaires d’éléments de $`\mathrm{\Gamma }_+\mathrm{\Gamma }_{}`$, on déduit la troisième égalité de (4.9) de la condition 1) et (4.11). Il reste à vérifier la première égalité de (4.9). D’abord, il résulte de 1) et de la discussion ci-dessus que :
$$𝔣eststableparj$$
(4.12)
Par ailleurs les conditions 1) et 4), et la discussion ci-dessus montre que :
$$j(𝔥_\alpha )=𝔥_{(A\alpha )^f},\alpha \mathrm{\Gamma }_+$$
(4.13)
On montre de même que :
$$j(𝔥_\alpha )=𝔥_{(A\alpha )^f},\alpha \mathrm{\Gamma }_+$$
(4.14)
Mais $`𝔥`$, qui est isomorphe à sa projection dans $`𝔤_+`$, est engendrée par les $`𝔥_\alpha `$, $`𝔥_\alpha `$, $`\alpha \mathrm{\Gamma }_+`$. Alors (4.12) à (4.14), joints à 1) montrent que $`𝔥`$ est stable par $`𝔦`$, ce qui achève de prouver que $`𝔦`$ est stable par $`j`$. On procède de même pour $`𝔦^{}`$. Ceci achève de prouver (i). (ii) est une conséquence immédiate de la Proposition précédente
On se fixe une donnée de Belavin-Drinfeld généralisée, qui vérifie les propriétés 1) à 4) de la Proposition pécédente. On introduit $`C=^{\prime \prime }A^1A^{}^{\prime \prime }`$, comme dans la Définition 5. On note $`\mathrm{\Gamma }_0:=domCImC`$. Si $`\alpha \mathrm{\Gamma }_0`$, on note $`𝒞(\alpha )`$, l’ensemble des $`\beta \mathrm{\Gamma }_0`$ tels qu’ il existe $`n`$ tel que :
$$\alpha ,C\alpha ,\mathrm{},C^{n1}\alpha \mathrm{\Gamma }_0,et\beta =C^n\alpha $$
ou :
$$\beta ,C\beta ,\mathrm{},C^{n1}\beta \mathrm{\Gamma }_0,et\alpha =C^n\beta $$
Si $`\alpha \mathrm{\Sigma }_+`$ n’appartient pas à $`\mathrm{\Gamma }_0`$, on pose $`𝒞(\alpha )=\{\alpha \}`$. Suivant Panov \[P1\], on appelle les $`𝒞(\alpha )`$ des chaines. Il est clair que les chaines sont disjointes ou confondues, et forment une partition de $`\mathrm{\Sigma }_+`$. On appellera $`C`$-équivalence la relation d’équivalence sur $`\mathrm{\Sigma }_+`$, dont les classes d’équivalence sont les chaines. On dira que, pour deux éléments distincts $`\alpha `$ et $`\beta `$, de $`\mathrm{\Sigma }_+`$, $`\alpha `$ est $`C`$-lié à $`\beta `$ si $`\alpha domC`$ et $`\beta =C\alpha `$. On définit, pour $`\alpha \mathrm{\Gamma }_0`$ :
$$\stackrel{ˇ}{𝒞}(\alpha ):=\{A^1(\beta ^f)|\beta 𝒞(\alpha )\mathrm{\Gamma }_+\}\{A^1(\beta ^f)|\beta 𝒞(\alpha )\mathrm{\Gamma }_+^{}\}$$
Si $`\alpha \mathrm{\Gamma }_0`$, on pose $`\stackrel{ˇ}{𝒞}(\alpha ):=𝒞(\alpha )`$.
###### Lemme 23
(i) Pour tout $`\alpha \mathrm{\Sigma }_+`$, $`\stackrel{ˇ}{𝒞}(\alpha )`$ est de la forme $`𝒞(\beta )`$, et l’opération $`\stackrel{ˇ}{}`$ est une involution de l’ensemble des chaines. (ii) Si $`t`$ vérifie les conditions 4) et 5) de la Proposition précédente, on a, avec les notations de celle-ci, pour tout $`\alpha \mathrm{\Gamma }_0`$:
$$u^\beta =u^\beta ^{},\beta ,\beta ^{}𝒞(\alpha )$$
$$u^\beta =\overline{u^\alpha },\beta \stackrel{ˇ}{𝒞}(\alpha )$$
Démonstration : Montrons (i). Il suffit d’étudier le cas $`\alpha \mathrm{\Gamma }_0`$. Soit $`\alpha \mathrm{\Sigma }_+`$. On note $`𝒜=𝒞(\alpha )\mathrm{\Gamma }_+`$, $`𝒜^{}=𝒞(\alpha )\mathrm{\Gamma }_+^{}`$. De la définition des chaines, il résulte que, pour toute paire d’éléments distincts $`\beta `$, $`\beta ^{}`$, de $`𝒜`$, il existe $`n^{}`$, $`\beta _1,\mathrm{},\beta _{n+1}𝒜`$, où, pour $`i=1,\mathrm{},n`$, $`\beta _i`$ est $`C`$-lié à $`\beta _{i+1}`$, avec $`\{\beta ,\beta ^{}\}=\{\beta _1,\beta _{n+1}\}`$. On remarque que les conditions 1) et 2) impliquent, par un calcul immédiat, que :
$$Si\beta ,\beta ^{}𝒜sontCli\stackrel{´}{e}s,A^1(\beta ^f)etA^1(\beta ^f)sontCli\stackrel{´}{e}s$$
On a le même énoncé pour $`𝒜^{}`$. De ce qui précède, il résulte que les éléments de $`A^1𝒜`$ (resp. $`A^1𝒜^{}`$) sont $`C`$-équivalents entre eux. Pour achever de prouver (i), il suffit de traiter le cas où $`𝒜`$ et $`𝒜^{}`$ sont non vides et d’intersection vide. De la définition des chaines , et du fait que $`domC`$ (resp. $`ImC`$) est un sous-ensemble de $`\mathrm{\Gamma }_+^{}`$ (resp. $`\mathrm{\Gamma }_+`$), cela n’est possible que si $`𝒞(\alpha )`$ n’a que deux éléments et $`𝒜`$ (resp. $`𝒜^{}`$) un élément $`\beta `$ (resp. $`\beta ^{}`$), où $`\beta ^{}`$ et $`\beta `$ sont $`C`$-liés. En outre :
$$\beta \mathrm{\Gamma }_+\mathrm{\Gamma }_+^{},\beta ^{}\mathrm{\Gamma }_+^{}\mathrm{\Gamma }_+$$
(4.15)
Mais alors $`A^{}\beta ^{}`$ et $`A\beta `$ sont égaux et on note $`\gamma `$ leur valeur commune . Utilisant les conditions 1) et 2) de la Proposition 10, on en déduit que :
$$A^1(\beta ^f)=A^1(\beta ^f)=\gamma ^f,\stackrel{ˇ}{𝒞}(\alpha )=\{\gamma ^f\}$$
Montrons que la classe de $`C`$-équivalence de $`\gamma ^f`$ est réduite à un élément. Pour cela il suffit de voir que $`\gamma ^f`$ n’est $`C`$-lié à aucun élément et qu’aucun élément ne lui est $`C`$-lié. Raisonnons par l’absurde et supposons par exemple que $`\delta `$ soit $`C`$-lié à $`\gamma ^f`$. On a alors :
$$\delta \mathrm{\Gamma }_+^{}etA(\gamma ^f)=A^{}\delta $$
(4.16)
On déduit des conditions 1), 2) de la Proposition 10 :
$$A^1\gamma =A^1(\delta ^f)$$
(4.17)
D’après (4.16) et la condition 1), $`A^1(\delta ^f`$ est un élément de $`\mathrm{\Gamma }_+^{}`$. Par ailleurs, d’après la définition de $`\gamma `$, $`A^1\gamma `$ est égal à $`\alpha `$. Joint à (4.17), cela montre que $`\alpha `$ est élément de $`\mathrm{\Gamma }_+^{}`$, ce qui contredit (4.15). Donc aucun élément n’est $`C`$-lié à $`\gamma ^f`$. On montre de même que $`\gamma ^f`$ n’est lié à aucun élément. Cecci achève de prouver (i). Si $`\beta `$ est $`C`$-lié à $`\beta ^{}`$, on a $`A\beta ^{}=A^{}\beta `$. Par ailleurs, d’après les conditions 4) et 5) de la proposition 10, on a :
$$u^\beta =u^{A^{}\beta },u^\beta ^{}=u^{A\beta ^{}}$$
La première égalité de (ii) en résulte. Pour la deuxième égalité, on remarque, que si $`\alpha \mathrm{\Gamma }_+`$, on a, d’après la condition 4) de la Proposition 10 :
$$u^\alpha =u^{A\alpha }$$
Mais, il résulte de la définition de $`u`$ que :
$$u^f=\overline{u}^1$$
(4.18)
On en déduit le résultat voulu. Le théorème suivant a été suggéré par un résultat de A. Panov (cf. \[P1\], Théorème 6.13) et la preuve que nous en donnons plus loin.
###### Théorème 6
Soit $`B`$ une forme de Manin réelle sur $`𝔤_1`$. Tout triple de Manin réel dans $`𝔤_1`$, relativement à $`B`$, est conjugué par un élément de $`G_1`$ à un triple fortement standard dont le complexifié est de la forme $`𝒯_{𝒟,t𝒲}`$, où $`t`$ est un élément de $`J_0`$ et $`𝒟=(A,A^{},𝔦_𝔞,𝔦_𝔞^{})`$ une donnée de Belavin-Drinfeld généralisée pour $`𝔤`$ et $`B_{}`$, qui vérifie, outre les conditions 1) à 5) de la Proposition 10 : 1)
$$u^2=1$$
2) Pour tout $`\alpha Gamma^+\mathrm{\Gamma }_0`$ : $`u_1^\alpha =1`$. 3) $`t=(v_1,1)`$, où $`u_1J_1`$.
Alors $`v_1^2=1`$
###### Remarque 4
Pour une classe de conjugaison sous $`G_1`$ de triples réels de $`𝔤_1`$, relativement à $`B`$, la donnée $`𝒟`$ est uniquement déterminée, d’après le Théorème 5 et il n’y a qu’un nombre fini de choix possibles pour $`t`$, car les éléments de carré 1 de $`J_1`$ sont en nombre fini.
Démonstration : Tenant compte du fait que $`u`$ n’est pas changé par la mutiplication de $`t`$ par un élément de la forme $`(t_1,\overline{t}_1)`$, on voit qu ’il suffit de trouver $`t`$ satisfaisant toutes les conditions à l’exception de 3). D’après la Proposition 10, il existe $`\underset{¯}{t}`$ et $`𝒟`$ satisfaisant les conditions 1) à 5) de celle-ci (en y changeant $`t`$ en $`\underset{¯}{t}`$). On vérifie aisément que si $`t^{}J_0`$ satisfait:
$$t^\alpha =t^{A\alpha },\alpha \mathrm{\Gamma }_+$$
(4.19)
$$t^\alpha =t^{A^{}\alpha },\alpha \mathrm{\Gamma }_+^{}$$
(4.20)
on a :
$$𝒯_{𝒟,\underset{¯}{t}𝒲}=𝒯_{𝒟,\underset{¯}{t}t^{}𝒲}$$
Il reste à choisir $`t^{}`$ vérifiant (4.19) et (4.20) de telle sorte que $`t=\underset{¯}{t}t^{}`$ vérifie les propriétés voulues. On choisit un sous-ensemble $`\mathrm{\Theta }`$ de $`\mathrm{\Sigma }_+`$, tel que toute classe de $`C`$-équivalence soit de la forme $`𝒞(\alpha )`$ ou $`\stackrel{ˇ}{𝒞}(\alpha )`$ pour un unique $`\alpha \mathrm{\Theta }`$. On remarque que si $`A\alpha =A^{}\beta `$, $`\alpha `$ et $`\beta `$ sont $`C`$-équivalents. On caractérise alors $`t^{}`$ par $`t^(\alpha )`$, $`\alpha \mathrm{\Sigma }`$. On choisit pour tout $`\alpha \mathrm{\Sigma }_+`$, une racine carrée $`z_\alpha `$ de $`\overline{(\underset{¯}{u}^\alpha )}^1`$, où $`\underset{¯}{u}=\overline{\underset{¯}{t}}(\underset{¯}{t}^f)^1`$. On pose alors, pour $`\alpha \mathrm{\Theta }`$ :
$$t^\beta =z_\alpha ,si\beta 𝒞(\alpha ),etsi\alpha \mathrm{\Gamma }_0$$
(4.21)
$$t^\beta =z_\alpha (resp.\overline{z_\alpha }),si\beta 𝒞(\alpha )(resp.\beta \stackrel{ˇ}{𝒞}(\alpha ))etsi𝒞(\alpha )\stackrel{ˇ}{𝒞}(\alpha )$$
(4.22)
$$t^\beta =|u_\alpha |^{1/2}si\alpha \mathrm{\Gamma }_0esttelque𝒞(\alpha )=\stackrel{ˇ}{𝒞}(\alpha ),et\beta 𝒞(\alpha )$$
(4.23)
Ceci détermine $`t^\beta `$ pour $`\beta \mathrm{\Sigma }_+`$, et l’on pose :
$$t^{\beta ^f}=\overline{(t^\beta )^1},\beta \mathrm{\Sigma }_+$$
(4.24)
Ainsi $`t^{}`$ est entièrement caractérisé par les relations (4.21) à (4.24). Il faut voir que $`t^{}`$ vérifie (4.19) et(4.20). Soit $`\beta \mathrm{\Gamma }_+`$. On suppose $`\beta 𝒞(\alpha )`$, $`\alpha \mathrm{\Theta }`$. D’après (4.24), on a :
$$t^{A\beta }=\overline{t}^{(A\beta )^f}$$
Mais, d’après la condition 1) de la Proposition 10, et la définition de $`\stackrel{ˇ}{𝒞}(\alpha )`$, on a : $`(A\beta )^f=A^1(\beta ^f)\stackrel{ˇ}{𝒞}(\alpha )`$, donc , d’après (4.22) et (4.23), on a :
$$\overline{t^{(A\beta )^f}}=t^\beta $$
On traite de même le cas où $`\beta \stackrel{ˇ}{𝒞}(\alpha )`$, et alors $`\alpha \mathrm{\Theta }`$. Ceci prouve que (4.19) est vérifié. On prouve (4.20) de la même manière. Calculons $`v_\beta :=\overline{t}^\beta (t^1)^{\beta ^f}`$, $`\beta \mathrm{\Sigma }_+`$ A l’aide de (4.22) à (4.24) et du Lemme 23, on voit que pour $`\alpha \mathrm{\Theta }`$ :
$$v_\beta =u^\beta ,si\beta 𝒞(\alpha )\stackrel{ˇ}{𝒞}(\alpha ),etsi\alpha \mathrm{\Gamma }_0ousi𝒞(\alpha )\stackrel{ˇ}{𝒞}(\alpha )$$
$$v_\beta =|u|^\beta ,si\alpha \mathrm{\Gamma }_0esttelque𝒞(\alpha )=\stackrel{ˇ}{𝒞}(\alpha ),et\beta 𝒞(\alpha )$$
Par ailleurs il résulte du Lemme 23 :
$$u^\beta estr\stackrel{´}{e}elsi\alpha \mathrm{\Gamma }_0esttelque𝒞(\alpha )=\stackrel{ˇ}{𝒞}(\alpha ),etsi\beta 𝒞(\alpha )$$
Alors il est clair que $`t:=t^{}\underset{¯}{t}`$ vérifie la condition 2) du Théorème et l’élément $`u`$ correspondant vérifie :
$$u^\alpha =1ou1,\alpha \mathrm{\Sigma }_+$$
$$u^f=\overline{u}^1$$
Il en résulte que $`u^2=1`$ comme désiré.
### 4.2 Une autre démonstration d’un résultat d’A. Panov
On suppose maintenant que $`𝔤_1`$ est une algèbre de Lie complexe simple et on pose $`B_1=ImK_{𝔤_1}`$. On fixe une forme réelle $`𝔥_1`$ de $`𝔤_1`$ et $`\sigma _1`$ la conjugaison de $`𝔤_1`$ par rapport à $`𝔥_1`$. On note $`G_1`$ le groupe adjoint de $`𝔤_1`$ (et non son recouvrement universel) et, si $`𝔢`$ est une sous-algèbre de Lie réelle de $`𝔤_1`$, on note $`E`$ le sous-groupe analytique de $`G_1`$, d’algèbre de Lie $`𝔢`$. On s’intéresse aux triples de Manin réels de $`𝔤_1`$, $`(B_1,𝔦_1,𝔦_1^{})`$ tels que $`𝔦_1`$ soit égal à $`𝔥_1`$, qu’on appelle ”$`𝔥_1`$-triple” , à conjugaison près par les élément de $`G_1`$, ou, ce qui revient au même, par ceux de $`G_1^{\sigma _1}`$, qui est l’ ensemble des éléments de $`G_1`$ commutant à $`\sigma _1`$. Ces triples décrivent les structures de bigèbres de Lie sur $`𝔥_1`$, dont le double est isomorphe à $`𝔤_1`$ (cf. \[P1\]). Soit $`𝔣_1`$ une sous-algèbre de Cartan fondamentale de $`𝔥_1`$ et soit $`𝔧_1`$ la complexifiée de $`𝔣_1`$ dans $`𝔤_1`$. On choisit une sous-algèbre de Borel, $`𝔟_1`$, de $`𝔤_1`$, contenant $`𝔧_1`$ et telle que $`\sigma _1(𝔟_1)`$ soit égal à la sous-algèbre de Borel de $`𝔤_1`$, $`𝔟_1^{}`$, opposée à $`𝔟_1`$, relativement à $`𝔧_1`$. On note $`\mathrm{\Sigma }_1`$ l’ensemble des racines simples de $`𝔧_1`$ dans $`𝔟_1`$. On note $`\theta `$ l’involution de $`\mathrm{\Sigma }_1`$, caractérisée par :
$$\sigma _1(𝔤_1^\alpha )=𝔤_1^{\theta (\alpha )}$$
(4.25)
En calulant de deux manières différentes $`t(\sigma _1(X))`$, pour $`X𝔤^\alpha `$, on trouve :
$$\overline{(t^{\sigma _1})^\alpha }=\overline{t^{\theta (\alpha )}}$$
(4.26)
Il est facile de voir qu’on peut choisir un système de générateurs de Weyl de $`𝔤_1`$, $`𝒲_1`$, relativement à $`\mathrm{\Sigma }_1`$, tel que :
$$\sigma _1(H_\alpha )=H_{\theta (\alpha )},\sigma _1(X_\alpha )=\epsilon _\alpha Y_{\theta (\alpha )},\sigma _1(Y_\alpha )=\epsilon _\alpha X_{\theta (\alpha )}$$
(4.27)
où :
$$\epsilon _\alpha =1ou1,et\epsilon _\alpha =1si\theta (\alpha )\alpha $$
(4.28)
On note $`𝔤:=𝔤_1\times 𝔤_1`$, $`𝔧_0=𝔧_1\times 𝔧_1`$, $`𝔟_0=𝔟_1\times 𝔟_1^{}`$. On note $`\eta _1`$ l’aplication de $`𝔤_1`$ dans $`𝔤`$ définie par :
$$\eta _1(X)=(X,\sigma _1(X)),X𝔤_1$$
dont l’image par $`\eta _1`$ est une forme réelle de $`𝔤`$. On note $`j_1`$ la conjugaison de $`𝔤`$ par rapport à cette forme réelle, qui vérifie :
$$j_1(X,Y)=(\sigma _1(Y),\sigma _1(X)),X,Y𝔤_1$$
On note $`𝒲=(𝒲_1\times \{0\})(\{0\}\times 𝒲_1`$). On note $`B`$ la forme de Manin sur $`𝔤`$, égale à $`K_{𝔤_1}`$ sur le premier facteur et à $`K_{𝔤_1}`$ sur le deuxième facteur facteur. Alors $`𝔤_+=𝔤_1\times \{0\}`$, $`𝔤_{}=\{0\}\times 𝔤_1`$ et $`\mathrm{\Sigma }_+=\mathrm{\Sigma }_1`$, $`\mathrm{\Sigma }_{}=\mathrm{\Sigma }_1`$ Si $`V`$ est un sous-espace vectoriel réel de $`𝔤_1`$, on notera :
$$V_{}:=\eta _1(V)+i\eta _1(V)$$
On remarque que $`𝔥_1`$ est égal à la diagonale, $`diag(𝔤_1)`$, dans $`𝔤_1\times 𝔤_1`$. Si $`𝒯_1=(B_1,𝔦_1,𝔦_1^{})`$ est un triple de Manin réel dans $`𝔤_1`$, on appelle triple complexifié de $`𝒯_1`$, le triple de Manin complexe dans $`𝔤`$, $`𝒯_1=(B,𝔦_1,𝔦_1^{})`$. Alors $`𝔦_1`$ et $`𝔦_1^{}`$ sont stables par $`j_1`$. De plus si on a deux triples réels dans $`𝔤_1`$, comme ci-dessus, ils sont conjugués par un élément de $`G_1`$ si et seulement si leurs complexifiés sont conjugués par un élément de $`\eta _1(G_1)=\{(g_1,g_1^{\sigma _1})|g_1G_1\}`$. On remarque que si deux sous-espaces vectoriels complexes de $`𝔤`$ sont conjugués par un élément de $`\eta _1(G_1)`$, si l’un est stable par $`j_1`$, l’autre l’est aussi. De plus le complexifié d’un triple fortement standard relativement à $`𝔟_1`$, $`𝔧_1`$, est fortement standard relativement à $`𝔟_0`$, $`𝔧_0`$, car il est facile de voir que le complexifié d’un antécédent est l’antécédent du complexifié. Comme tout triple de Manin réel dans $`𝔤_1`$ est conjugué par un élément de $`G_1`$, à un triple fortement standard, d’après le Théorème 4, on a immédiatement :
###### Lemme 24
(i) Si $`𝒯_1`$ est un ”$`𝔥_1`$-triple”, son complexifié est conjugué par un élément de $`\eta _1(G_1)`$ à un triple fortement standard $`𝒯=(B,𝔦,𝔦^{})`$ tel que : 1) $`𝔦`$ est conjugué par un élément de $`\eta _1(G_1)`$ à $`diag(𝔤_1)`$. 2) $`𝔦`$ et $`𝔦^{}`$ sont stables par $`j_1`$. (ii) Si $`𝒯=(B,𝔦,𝔦^{})`$ est un triple de Manin fortement standard vérifiant 1), 2) , il existe un ”$`𝔥_1`$-triple”, unique à conjugaison sous $`G_1`$ près (où plutôt $`G_1^{\sigma _1}`$ près) tel que son complexifié soit conjugué à $`𝒯`$ par un élément de $`\eta _1(G_1)`$. Deux triples fortement standard, associés à $`B`$, vérifant 1), 2), conjugués sous $`\eta _1(G_1)`$, conduisent à la même classe de conjugaison sous $`G_1^{\sigma _1}`$ de ”$`𝔥_1`$-triple”.
###### Lemme 25
(i)Tout triple de Manin complexe, fortement standard dans $`𝔤`$ pour $`𝔟_0`$, $`𝔧_0`$, associé à $`B`$, vérifiant les conditions 1), 2), du Lemme précédent est conjugué par un élément de $`\eta _1(G_1)`$ à un triple $`𝒯_{𝒟,(1,t)(𝒲)}`$, vérifiant les mêmes conditions, où $`𝒟`$ est une donnée de Belavin-Drinfeld $`(A,𝔦_𝔞,A^{},𝔦_𝔞^{}^{})`$, relativement à $`B`$ et $`t`$ est un élément de $`J_1`$. (ii) Le fait que $`𝒯_{𝒟,(1,t)(𝒲)}`$ vérifie les conditions 1), 2) du Lemme précédent équivaut à: 1) Il existe un élément $`g_1G_1`$, tel que $`t=g_1^{\sigma _1}g_1^1`$. En particulier on a $`t^{\sigma _1}=t^1`$ 2) $`\mathrm{\Gamma }_+=\mathrm{\Sigma }_1`$, et $`A`$ est l’identité. 3) $`𝔦_𝔞=\{0\}`$ 4)
$$\theta (\mathrm{\Gamma }_+^{})=\mathrm{\Gamma }_{}^{},et\theta (A^{}\alpha )=A^1(\theta (\alpha )),\epsilon _{A^{}\alpha }t^{A^{}\alpha }=\epsilon _\alpha t^\alpha ,\alpha \mathrm{\Gamma }_+^{}$$
5) L’espace $`𝔦_𝔞^{}^{}`$ est stable par $`j_1`$.
Démonstration : Soit $`\underset{¯}{𝒯}=𝒯_{𝒟,(t_1,t_2)(𝒲)}`$ un triple de Manin fortement standard, vérifiant les conditions 1), 2) du Lemme précédent, où $`t_1,t_2J_1`$. Alors $`\underset{¯}{𝒯}`$ est conjugué par $`(t_1,t_1^{\sigma _1})\eta _1(G_1)`$ à $`𝒯:=𝒯_{𝒟,(1,t)(𝒲)}`$, où $`t=t_1^{\sigma _1}t_2`$, qui vérifie les conditions 1) et 2) du Lemme précédent. Montrons que cela implique les propriétés 1) à 5) ci-dessus. Ecrivons $`𝒯=(B,𝔦,𝔦^{})`$. On obtient 1) en écrivant que $`𝔦`$ est conjugué par un élément de $`\eta _1(G_1)`$ à $`diag(𝔤_1)`$. Alors $`𝔦`$ est semi-simple, donc sous $`𝔤`$. Par suite, avec les notations du Théorè me 1, on a $`𝔪=𝔤`$ et $`\mathrm{\Gamma }_+=\mathrm{\Sigma }_1`$. Alors, utilisant les notations du Lemme 22, $`𝔦=\{(X,\tau (X))|X𝔤_1\}`$. Comme $`𝔦`$ est conjugué par un élément de $`\eta _1(G_1)`$ à $`diag(𝔤_1)`$, cela implique que l’automorphisme $`\tau `$ de $`𝔤_1`$ est intérieur. Or, par définition, $`\tau `$ préserve $`𝔧_1`$, et l’inverse du transposé de sa restriction à $`𝔧_1`$ induit $`A`$ sur $`\mathrm{\Sigma }_1`$. Comme $`A`$ préserve $`\mathrm{\Sigma }_1`$, $`A`$ doit être l’identité. On a donc prouvé 2). 3) résulte du fait que $`𝔪=𝔤`$, donc $`𝔞=\{0\}`$. On traduit maintenant le fait que $`𝔦^{}`$ est stable par $`j_1`$. Cela implique que $`𝔭^{}`$ est stable par $`j_1`$. Comme $`𝔧_0`$ est stable par $`j_1`$, joint à (4.25), cela implique aussi que $`𝔪^{}`$ est stable par $`j_1`$. Cela conduit immédiatement à la première égalité de 4). Les deux autres sont obtenues en traduisant la stabilité de $`𝔥^{}`$ par $`j_1`$. Plus précisément on pose :
$$U=(X_\alpha ,t^{A^{}\alpha }X_{A^{}\alpha })𝔥^{},\alpha \mathrm{\Sigma }_1$$
On a, en tenant compte de (4.27) et de l’antilinéarité de $`\sigma _1`$ :
$$j_1(U)=(\overline{t^{A^{}\alpha }}\epsilon _{A^{}\alpha }Y_{\theta (A^{}\alpha )},\epsilon _\alpha Y_{\theta (\alpha )})$$
La stabilité de $`𝔥^{}`$ par $`j_1`$, implique que $`\theta (A^{}\alpha )`$ est élément de $`\mathrm{\Gamma }_+^{}`$ et $`j_1(U)`$ doit être un multiple scalaire de :
$$V=(Y_{\theta (A^{}\alpha )},t^{A^{}\theta (A^{}\alpha )}Y_{A^{}\theta (A^{}\alpha )})$$
Les deux premières égalités de 4) en résulte immédiatement et l’on obtient en outre :
$$\epsilon _{A^{}\alpha }t^{A^{}\alpha }=\epsilon _\alpha t^{\theta (\alpha )}$$
En utilisant (4.26) et le fait que $`t^{\sigma _1}=t^1`$, on aboutit à la troisième égalité de 4). La condition 5) est immédiate. On procède de même pour la réciproque.
###### Lemme 26
Soit $`t`$, $`𝒟`$ vérifiant les conditions 1) à 5) du Lemme précédent. Soit $`t^{}J_1`$ tel que :
$$t^{A^{}\alpha }=t^\alpha ,\alpha \mathrm{\Gamma }_+^{}$$
Alors, on a : $`𝒯_{𝒟,(1,t)(𝒲)}=𝒯_{𝒟,(t^{},t^{}t)(𝒲)}`$, et $`𝒯_{𝒟,(1,t)(𝒲)}`$ est conjugué par $`(t^{},t^{\sigma _1})\eta _1(G_1)`$ à $`𝒯_{𝒟,(1,(t^{}t^{\sigma _1})t)(𝒲)}`$.
Démonstration : Notons $`𝒯_{𝒟,(1,t)𝒲}=(B,𝔦,𝔦^{})`$ et utilisons les notations du Théorème 1. La stabilité de $`𝔦`$ par $`(t^{},t^{})`$ est claire. Par ailleurs, $`(t^{},t^{})J_0`$ laisse stable $`𝔭^{}`$, donc $`𝔫^{}`$ et laisse fixe point par point les éléments de $`𝔞^{}`$. Il reste à voir que $`𝔥^{}`$ est invariant. Mais cette algèbre est engendrée par :
$$(X_\alpha ,t^{A^{}\alpha }X_{A^{}\alpha }),(Y_\alpha ,t^{A^{}\alpha }Y_{A^{}\alpha }),\alpha \mathrm{\Gamma }_+^{}$$
Le Lemme en résulte immédiatement. Si $`\alpha \mathrm{\Gamma }_0:=\mathrm{\Gamma }_+^{}\mathrm{\Gamma }_{}`$, on note $`𝒞(\alpha )`$, l’ensemble des $`\beta \mathrm{\Gamma }_0`$ tels qu’ il existe $`n`$ tel que :
$$\alpha ,A^{}\alpha ,\mathrm{},A^{n1}\alpha \mathrm{\Gamma }_+^{},et\beta =A^n\alpha $$
ou :
$$\beta ,A^{}\beta ,\mathrm{},A^{n1}\beta \mathrm{\Gamma }_+^{},et\alpha =A^n\beta $$
Les $`𝒞(\alpha )`$ sont soit distincts soit confondus. De plus la deuxième égalité de la condition 4) du Lemme 25, montre facilement :
$$\theta (𝒞(\alpha )=𝒞(\theta \alpha )$$
(4.29)
(cf \[P1\], Lemme 4.11, pour un Lemme analogue) Par ailleurs, grâce à la définition de de $`\epsilon _\alpha `$, on a :
$$\epsilon _{\theta (\alpha )}=\epsilon _\alpha ,\alpha \mathrm{\Sigma }_1$$
(4.30)
Le Théorème suivant est du à A. Panov (cf. \[P\], Théorèmes 4.5, 4.13).
###### Théorème 7
(i) Le complexifié d’un $`𝔥_1`$-triple est conjugué par un élément de $`\eta _1(G_1)`$ à un triple $`𝒯_{𝒟,(1,u)(𝒲)}`$, où $`u`$, $`𝒟`$ vérifient les conditions 1) à 5) du Lemme 25 (avec $`t`$ remplacé par $`u`$) et $`u`$ vérifie de plus : 1) $`u^{\sigma _1}=u=u^1`$ 2) $`u^\alpha =1ou1,si\alpha \mathrm{\Sigma }_1`$ 3) $`u^\alpha =1`$, si $`\alpha \mathrm{\Sigma }_1\mathrm{\Gamma }_0`$ et $`\theta (\alpha )\alpha `$ , ou si $`\alpha \mathrm{\Gamma }_0`$ et $`\theta (𝒞(\alpha )𝒞(\alpha )`$ (ii) Réciproquement si $`u`$ et $`𝒟`$ vérifient les propriétés ci-dessus (i.e. 1) et 2) de (i) et les conditions 1) à 5) du Lemme 25), $`𝒯_{𝒟,(1,u)𝒲}`$ est conjugué au complexifié d’un $`𝔥_1`$-triple, unique modulo la conjugaison de $`G_1^{\sigma _1}`$. (iii) Dans le cas où $`𝔣_1`$ est l’algèbre de Lie d’un tore maximal compact dans $`G_1`$, $`\theta `$ est triviale, $`\mathrm{\Gamma }_+`$ est vide et les conditions sur $`u`$ ci-dessus se réduisent aux deux premières, outre celles du Lemme 25.
Démonstration : On sait que le complexifié d’un $`𝔥_1`$-triple est conjugué par un élément de $`\eta _1(G_1)`$ à un triple $`𝒯_{𝒟,(1,t)𝒲}`$, où $`t`$, $`𝒟`$ vérifient les conditions 1) à 5) du Lemme 26 (avec $`t`$ remplacé par $`u`$. L’idée est d’appliquer le Lemme 25, avec $`t^{}`$ bien choisi.Il suffit de définir $`t^\alpha ,\alpha \mathrm{\Sigma }_1`$. On établit d’abord quelques résultats auxiliaires. On rappelle que d’après (4.26) et la condition 1) du Lemme 25, on a :
$$\overline{t^{\theta (\alpha )}}=t^\alpha ,\alpha \mathrm{\Sigma }_1$$
(4.31)
Par ailleurs, d’après la troisième égalité de la condition 4) du Lemme 25 , on a :
$$t^{A^{}\alpha }=\epsilon _{A^{}\alpha }\epsilon _\alpha t^\alpha ,\alpha \mathrm{\Gamma }^{}$$
d’où l’on déduit :
$$t^\beta =\epsilon _\beta \epsilon _\beta ^{}t^\beta ^{},\beta ,\beta ^{}𝒞(\alpha ),\alpha \mathrm{\Gamma }_0$$
(4.32)
Enfin la condition 4) du Lemme 25 montre que :
$$\epsilon _\alpha =\epsilon _{\theta (\alpha )},\alpha \mathrm{\Sigma }_1$$
(4.33)
Pour définir $`t^\alpha `$, on distingue plusieurs cas. On note $`\mathrm{\Sigma }_1`$, un sous-ensemble de $`\mathrm{\Sigma }_1`$ tel que tout $`\beta \mathrm{\Sigma }_1`$ soit élément d’un $`𝒞(\alpha )`$ pour un unique $`\beta `$ appartenant à $`\mathrm{\Sigma }_1\theta (\mathrm{\Sigma }_1)`$, et tel que les éléments de l’intersection de $`\mathrm{\Sigma }_1`$ et $`\theta (\mathrm{\Sigma }_1)`$ soient fixés par $`\theta `$. Soit $`\alpha \mathrm{\Sigma }_1`$ : 1) Si $`\alpha \mathrm{\Gamma }_0`$, et $`\theta (\alpha )=\alpha `$, (4.31) implique que $`t^\alpha `$ est réel et on pose :
$$t^\alpha =|t^\alpha |^{1/2}$$
2) Si $`\alpha \mathrm{\Gamma }_0`$, et $`\theta (\alpha )\alpha `$, on note $`z`$ une racine carrée de $`t^\alpha `$, et on pose :
$$j^\alpha =z,j^{\theta (\alpha )}=\overline{z}$$
3) Si $`\alpha \mathrm{\Gamma }_0`$ et $`\theta (𝒞(\alpha ))=𝒞(\alpha )`$, on a, d’ après (4.32) :
$$t^{\theta (\alpha )}=\epsilon _{\theta (\alpha )}\epsilon _\alpha t^\alpha $$
Tenant compte de (4.31), cela implique que $`t^\alpha `$ est réel. On note $`z`$ une racine carrée de $`|t^\alpha |`$, et on pose :
$$t^\beta =z,\beta 𝒞(\alpha )$$
4) Si $`\alpha \mathrm{\Gamma }_0`$ et $`\theta (C(\alpha ))C(\alpha )=\mathrm{}`$, on note $`z`$ une racine carrée de $`t^\alpha `$, et on pose :
$$t^\beta =z,t^{\theta (\beta )}=\overline{z},\beta C(\alpha )$$
On voit que les relations précédentes définissent $`t^{}`$, qui vérifie :
$$t^{A^{}\alpha }=t^\alpha ,\alpha \mathrm{\Gamma }_+^{}$$
Le Lemme 25 et les conditions imposées à $`t^{}`$ montrent que $`u:=t^{}t^{\sigma _1}t`$ vérifie les conditions voulues. (ii) est un cas particulier du Lemme 25 (ii). Traitons le cas où $`𝔣_1`$ est l’algèbre de Lie d’un tore maximal compact de $`𝔤_1`$. Alors $`𝔣_1`$ est l’ensemble des éléments de $`𝔧_1`$ sur lesquels toutes les racines sont imaginaires pures. Tout $`𝔥_1`$-triple est conjugué sous $`G_1`$ à un triple réel fortement standard $`(B_1,𝔦_1,𝔦_1^{})`$. Alors $`𝔦_1𝔧_1`$, qui est conjugué par $`G_1`$ à $`𝔣_1`$, est l’algèbre de Lie d’un tore maximal compact de $`G_1`$, donc est égal à $`𝔣_1`$. Par ailleurs $`𝔥_1^{}𝔧_1`$ est une sous-algèbre de Cartan fondamentale de $`𝔥_1^{}`$. Son intersection avec $`𝔣_1`$ est donc non réduite à zéro sauf si $`𝔥_1^{}`$ est réduite à zéro. On en déduit que la sous-algèbre Lagrangienne $`𝔦_1^{}`$ est sous une algèbre de Borel. Par complexification, il en résulte que dans (i), on doit avoir $`\mathrm{\Gamma }_+^{}=\mathrm{}`$ La définition de $`\theta `$ et le fait que les racines soient imaginaires pures sur $`𝔣_1`$ montrent que $`\theta `$ est l’identité. Ceci achève la preuve du Théorème.
Références
\[BD\], BELAVIN A., DRINFELD G., Triangle equations and simple Lie algebras, Mathematical Physics Reviews, vol. 4, 93-165
\[Bor\], BOREL A., Linear algebraic groups, Second Enlarged Edition, Graduate Text in Math.126, 1991, Springer Verlag, New York, Berlin, Heidelberg.
\[Bou\], BOURBAKI N., Groupes et Algèbres de Lie, Chapitre I, Chapitres IV, V, VI, Chapitres VII, VIII, Actualités Scientifiques et Industrielles 1285, 1337, 1364, Hermann, Paris, 1960, 1968, 1975.
\[De\], DELORME P., Sur les triples de Manin pour une algèbre réductive complexe, Preprint 1999.
\[G\], GANTMACHER F., Canonical representation of automorphism of a semisimple Lie group, Math Sb., 47, (1939), 101-144.
\[K1\], KAROLINSKY E., A classification of Poisson homogeneous spaces of a compact Poisson Lie group, Math. Phys., Anal. and Geom., 3 (1996), 545-563.
\[K2\], KAROLINSKY E., A classification of Poisson homogeneous spaces of a compact Poisson Lie group, Dokl. Ak. Nauk, 359 (1998), 13-15.
\[K3\], KAROLINSKY E.,A classification of Poisson homogeneous spaces of a reductive complex Poisson Lie group, Preprint, 1999
\[M1\], MATSUKI T., The orbits of affine symmetric spaces under the action of minimal parabolic subgroups, J. Math. Soc. Japan, 31 (1979), 331-357.
\[M2\], MATSUKI T., Orbits of affine symmetric spaces under the action of parabolic subgroups, Hiroshima J. Math., 12 (1982), 307-320.
\[P1\], PANOV A.., Manin triples of real simple Lie algebras, Part 1, Preprint.
\[P2\], PANOV A., Manin triples of real simple Lie algebras, Part 2, Preprint, QA 9905028.
\[W\], WARNER G., Harmonic Analysis on semi-simple Lie groups, Grundleh ren der math. Wis. in Einz., Vol 188, Springer Verlag, Berlin-Heidelberg-New York 1972
Institut de Mathématiques de Luminy, U.P.R. 9016 du C.N.R.S. Université de la Méditerrannée, 163 Avenue de Luminy, Case 907, 13288, Marseille Cedex 09, France e-mail : delorme@iml.univ-mrs.fr |
warning/0003/cond-mat0003512.html | ar5iv | text | # Columnar Defects and Scaling Behavior in Quasi 2D Type II Superconductors
## Abstract
Persistent scaling behavior of magnetization in layered high $`T_c`$ superconductors with short–range columnar defects is explained within the Ginzburg Landau theory. In the weak field region, the scaling function differs from that of a clean sample and both the critical and crossing temperatures are renormalized due to defects. In the strong field region, defects are effectively suppressed and scaling function, as well as critical and crossing temperatures are the same as in a clean superconductor. This picture is consistent with recent experimental results
Layered high-temperature superconducting (HTSC) materials, such as Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+δ</sub> and Bi<sub>2</sub>Sr<sub>2</sub>Ca<sub>2</sub>Cu<sub>3</sub>O<sub>10</sub>, are known to exhibit 2D scaling magnetic properties around the mean field transition line $`H_{c2}(T)`$. It is manifested by inspecting the magnetization $`M_0`$ as a function of temperature $`T`$ (measured in energy units) and the (external) magnetic field $`H`$:
$$\frac{s\mathrm{\Phi }_0}{A\sqrt{TH}}M_0(T,H)=(\sqrt{x^2+2}x),$$
(1)
where $`s`$ is an effective interlayer spacing, $`\mathrm{\Phi }_0`$ is the flux quantum, $`x=AH_{c2}^{}[TT_{c2}(H)]/\sqrt{TH}`$ is the scaling variable, $`H_{c2}^{}dH_{c2}(T)/dT|_{T=T_{c0}},`$ and $`T_{c0}`$ is the zero field critical temperature. For a superconductor with Ginzburg-Landau (GL) parameter $`\kappa `$ and Abrikosov geometric factor $`\beta _A`$ the constant $`A=\sqrt{s\mathrm{\Phi }_0/p}`$, where $`p=16\pi \kappa ^2\beta _A.`$ The form of the scaling function implies the existence of a crossing point: at some temperature $`T_0^{}=T_{c0}(1+1/(2A^2H_{c2}^{}))^1,`$ the sample magnetization is independent on $`H`$, $`M_0^{}M_0(T_0^{},H)=T_0^{}/(s\mathrm{\Phi }_0)`$.
Recently, the influence of linear defects (columnar defects, artificial holes etc.) on the magnetic properties of superconductors has been studied experimentally and theoretically. In particular, experiments by van der Beek et. al. showed that in HTSC with columnar defects the reversible magnetization of the sample is drastically affected, and that there are now two scaling regimes pertaining to relatively weak $`H<H_\mathrm{\Phi }`$ and strong $`H>H_\mathrm{\Phi }`$ magnetic fields (here the matching field $`H_\mathrm{\Phi }=n_d\mathrm{\Phi }_0`$ is proportional to the 2D density of defects $`n_d`$). These two scaling regimes correspond to two different critical temperatures (used in Ref. as fitting parameters) and crossing points.
aaaIn this Letter we propose an explanation of these results. Let us commence by presenting some intuitive arguments. Consider the quantity $`c=H_\mathrm{\Phi }/H`$ which, in a macroscopic sample, is the number of defects divided by the number of vortices. The magnetic field then serves as a control parameter for tuning the effective concentration $`c`$ of defects. In the weak field region, $`c`$ is large, each vortex is affected by a force emanating from many defects, and the fluctuations of this force play the main role. Short-range defects could be taken into account perturbatively. In first order they retain the same form of scaling function as that of a clean sample but renormalize the critical temperature $`T_c`$. Second order corrections indeed destroy the scaling behavior but in the vicinity of the crossing temperature scaling is approximately maintained. In the strong field region, $`c`$ is small, and the standard concentration expansion can be used. Here, even the first order correction (with respect to small concentration) destroys the scaling behavior. However, a strong field effectively suppresses the defects, thus restoring the scaling behavior of a clean superconductor with the initial critical temperature $`T_{c0}.`$ Identifying the two fitting temperatures of Ref. with the renormalized critical temperature $`T_c`$ and the initial one $`T_{c0}`$ respectively, one finds for the dimensionless defect strength $`\theta _1=0.49`$, well inside its allowed range $`0\theta _11`$. This indicates a full consistence between the description constructed below and the experimental results of Ref..
aaaOur quantitative discussion employs an approach proposed and successfully used for arbitrary fields in clean superconductors and for very low fields in disordered superconductors. Here we use it for disordered superconductors in much higher fields. Consider an irradiated thin superconducting film (or one layer in a layered superconductor) with area $`S`$ subject to perpendicular magnetic field (thus parallel to the defects). The effective interlayer separation $`s`$ is assumed to be much larger than the effective superconducting coherence length $`\xi (H,T)`$ in the magnetic field direction but much smaller than the magnetic penetration depth. Then the problem becomes effectively two dimensional. Columnar defects can be described as a local reduction of the critical temperature $`\delta T_c(𝐫)=T_{c0}t_j\mathrm{exp}((𝐫𝐫_j)^2/2L^2).`$ Here $`𝐫`$ is a two dimensional vector in the film plane, $`L`$ is the defect radius, and the positions $`𝐫_j`$ of defects are uniformly and independently distributed over the film plane with density $`n_d.`$ The value of $`n_d`$ is assumed to be moderate so that for the pertinent region of temperature the matching field $`H_\mathrm{\Phi }`$ is always much smaller than $`H_{c2}(T).`$ The dimensionless amplitudes of defects $`t_j1`$ are also independent random quantities distributed with some probability density. The thermodynamic properties of a type-II superconductor with $`\kappa 1`$ containing $`N_v`$ vortices are described by its partition function
$$𝒵𝒟\{\mathrm{\Psi }\}\mathrm{exp}(N_vg[\mathrm{\Psi }]),$$
(2)
where $`\mathrm{\Psi }`$ is the corresponding order parameter. The dimensionless GL free energy $`g[\mathrm{\Psi }]`$ of an irradiated superconductor is given by an expression
$$g=x\overline{|\mathrm{\Psi }|^2}+(4\beta _A)^1\overline{|\mathrm{\Psi }|^4}+\overline{\tau |\mathrm{\Psi }|^2},$$
(3)
where bar denotes averaging over the sample area. The scaling variable $`x`$ and the local temperature $`\tau (𝐫)`$ are defined below for each region of the magnetic field.
aaaFollowing, we replace in Eq.(3) $`\overline{\left|\mathrm{\Psi }(𝐫)\right|^4}`$ by $`\beta _A\left(\overline{\left|\mathrm{\Psi }(𝐫)\right|^2}\right)^2,`$ where $`\beta _A1.16`$ is the Abrikosov factor for a triangular lattice. This replacement is based on the assumption that the distribution of vortices is almost uniform in both regions of the magnetic field considered here. It is supported by noticing a remarkable difference between the number of vortices and the number of defects in both regions of fields. This substitution, together with the simplest version of the Hubbard-Stratonovich transformation (introduction of an additional integration over some auxiliary field $`\gamma `$) turns the problem to be an exactly solvable one. Then project the order parameter on the lowest Landau level (LLL) subspace,
$$\mathrm{\Psi }(𝐫)=\underset{m=0}{\overset{N_v}{}}C_mL_m(𝐫),$$
(4)
where $`L_m(𝐫)`$ are normalized LLL eigenfunctions with orbital momentum $`m.`$ As was recently demonstrated, the LLL approximation works quite well even down to $`HH_{c2}(T)/13`$. After integration over the expansion coefficients $`C_m`$ the partition function (2) reads,
$$𝒵_\mathrm{\Gamma }\mathrm{exp}\left\{N_v(\gamma ,x)\right\}𝑑\gamma ,$$
(5)
where
$$(\gamma ,x)=\gamma ^2+N_v^1\text{tr}\mathrm{ln}\left[(x+\gamma )\widehat{\text{I}}+\widehat{\tau }\right]$$
(6)
and $`\widehat{\tau }`$ is a random matrix with elements:
$$\tau _{mn}=_SL_m^{}(𝐫)\tau (𝐫)L_n(𝐫)d^2𝐫.$$
(7)
The contour $`\mathrm{\Gamma }`$ in Eq.(5) is parallel to the imaginary axis and stretches from $`\gamma ^{}i\mathrm{}`$ to $`\gamma ^{}+i\mathrm{}.`$ To assure convergence of the integrals over the coefficients $`\{C_m\}`$ the real constant $`\gamma ^{}`$ should satisfy the inequality $`\gamma ^{}+x+\mathrm{min}\tau _n>0,`$ where $`\tau _n`$ is the $`n`$-th eigenvalue of the matrix $`\tau _{mn}`$.
aaaIn the thermodynamic limit $`S\mathrm{}`$ with $`n_d=N_v/S`$ fixed, the partition function (5) could be calculated in a saddle point approximation. This results in the following form for the magnetization
$$\frac{s\mathrm{\Phi }_0}{A\sqrt{HT}}M(T,H)=\left(N_v^𝒵\right)^1𝒵/x=2\gamma (x),$$
(8)
where $`\gamma (x)`$ is the solution of the saddle point equation $`(\gamma ,x)/\gamma =0`$. For a clean superconductor ($`\widehat{\tau }=0`$) one gets two possible saddle points but only one of them
$$\gamma _0(x)=\frac{1}{2}(\sqrt{x^2+2}x)$$
(9)
can be reached by an allowed deformation of the contour $`\mathrm{\Gamma }`$. Substitution of Eq.(9) into (8) yields the magnetization $`M_0(T,H)`$ of a clean sample (1) obtained in Ref.. Note that $`2\gamma _0(x)`$ serves as the appropriate scaling function. To study the disordered case, we consider separately two regions of the magnetic field.
aaaIn the weak field region $`H<H_\mathrm{\Phi }`$ we, from the onset, take into account the renormalization of the critical temperature caused by defects. As a result, the scaling variable $`x`$ is defined in the same way as for a clean superconductor albeit with renormalized critical temperature $`T_c=T_{c0}\delta T_c,`$ where
$$\delta T_c=<\delta T_c(𝐫)>=2\pi \theta _1n_dL^2T_{c0}$$
(10)
and $`\theta _n<t^n>`$ where here and below $`<..>`$ implies ensemble average. The function $`\tau (𝐫)`$ in this field region, defined as $`\tau (𝐫)=(\delta T_c(𝐫)\delta T_c)AH_{c2}^{}/\sqrt{TH}`$, represents temperature fluctuations caused by short-range defects. They are small and can be accounted for perturbatively. Then, in the thermodynamic limit, the last term on the r.h.s. of Eq.(6) has an explicit self-averaged structure $`N_v^1\text{tr}(\mathrm{})`$ and can be replaced by its average. This procedure modifies the saddle point equation and therefore results in a modified magnetization
$$M(T,H)=M_0(T,H)(1+\epsilon (T)\frac{2\gamma _0(x)}{\sqrt{x^2+2}},),$$
(11)
where
$$\epsilon (T)=\frac{\mathrm{tr}\widehat{\tau }^2}{N_v}=\frac{\theta _2}{p}n_dL^2\frac{(2\pi H_{c2}^{}T_{c0})^2sL^2}{T}.$$
(12)
Note that the parameter $`\epsilon (T)`$ is proportional to the fourth power of the defect radius $`L`$ thus justifying the perturbation approach for short-range defects.
aaaIn the zeroth approximation with respect to $`\epsilon (T)`$ the magnetization (11) has exactly the same form, as for a clean sample (1) thus retaining both the scaling property and the existence of a crossing point. However, due to renormalization of the critical temperature, the crossing temperature $`T^{}=T_0^{}\delta T^{}`$ differs from its value $`T_0^{}`$ in a clean sample: $`\delta T^{}=\delta T_c(1+(2A^2H_{c2}^{})^1)^1.`$ In the next order, scaling is virtually destroyed, since the correction term (within the parenthesis in Eq.(11)) depends not only on the scaling variable $`x`$ but also on temperature. Yet, in a sufficiently narrow region around some temperature $`T`$, the deviation from scaling is negligibly small, but the scaling function itself is modified to be $`2\gamma (x,T)`$. At temperature $`T^{}`$ the magnetization reads
$$M(T^{},H)=M_0(T^{})\left(1+\epsilon (T^{})\frac{2H^{}}{H+H^{}}\right),$$
(13)
where $`H^{}=H_{c2}(T^{})=T^{}/(2A^2).`$ Therefore if the field is weak enough, $`HH^{},`$ then the crossing point is restored, $`T^{}`$ serves as a true crossing temperature and the magnetization at the crossing temperature differs from its unperturbed form $`2\gamma _0(x)`$ merely by a multiplicative constant $`1+2\epsilon (T^{}).`$
aaaWhen the magnetic field increases, the approach used above becomes inapplicable. Firstly, it fails in the vicinity of the matching field where the Abrikosov factor becomes very sensitive to the details of defect configuration. Secondly, higher order terms in the perturbation expansion for the saddle point equation (which are omitted), grow with magnetic field. Fortunately, we have here a new small parameter, that is, the dimensionless concentration $`c`$ of defects. It is then natural to use the concentration expansion. In such a case there is no sense in renormalizing the critical temperature, and the dimensionless temperature $`\tau (𝐫)`$ is now defined as $`\tau (𝐫)=\delta T_c(𝐫)AH_{c2}^{}/\sqrt{TH}.`$
aaaAs mentioned above, the second term in the r.h.s. of Eq.(6) is self-averaging and can be calculated using the limiting form of the density of states $`\rho (\tau )`$ of the matrix (7), which, for short-range defects in linear approximation with respect to $`c`$, reads
$$\rho (\tau )=(1c)\delta (\tau )+\frac{c}{\lambda }p(\frac{\tau }{\lambda }),$$
(14)
where $`\lambda =2\pi L^2T_{c0}AH_{c2}^{}\sqrt{H}(\mathrm{\Phi }_0\sqrt{T})^1`$ and $`p(t)`$ is probability distribution of the dimensionless temperature $`t_j`$. Indeed, the matrix $`\tau _{mn}`$ is nothing but the Hamiltonian of a particle with charge $`2e`$ in a $`2D`$ system subject to a perpendicular magnetic field and containing short-range defects (projected on the LLL). The first and second terms in Eq.(14) correspond, respectively, to those states whose energy is stuck to the LLL (despite the presence of zero-range defects (see e.g.)) and those states whose energies are lifted from the LLL by these defects. For sufficiently narrow distribution $`p(t)`$, the corresponding saddle-point equation leads to the magnetization
$$M=M_0\left(1\frac{c\lambda \theta _1}{(1+2\lambda \theta _1\gamma _0(x))\sqrt{x^2+2}}\right),$$
(15)
were $`M_0(T,H)`$ is given by Eq.(1) with an initial critical temperature $`T_{c0}.`$
aaaRigorously speaking, scaling is destroyed since both the concentration $`c`$ and the shifted eigenvalue $`\theta _1\lambda `$ depend explicitly on $`H`$ and $`T`$. However, at strong field the correction term in Eq.(15) becomes negligibly small. This implies a restoration of the crossing point. Indeed, at temperature $`T_0^{}`$ the magnetization $`M^{}=M(T^{},H)`$ assumes the form
$$M^{}=M_0^{}\left(1\frac{1}{1+\eta }\frac{H_\mathrm{\Phi }}{H+H^{}}\right),$$
(16)
with $`\eta ^1=2\pi L^2H_{c2}^{}T_{c0}\theta _1/\mathrm{\Phi }_0.`$ Therefore in the entire strong field region $`H_\mathrm{\Phi }HH^{}`$ the crossing temperature coincides with its initial value $`T_0^{}`$ and the magnetization in the crossing point practically coincides with its value $`M_0^{}`$ in a clean superconductor.
aaaLet us now discuss the limits of applicability of our results and their relation to the experiment of Ref.. Note that the first two moments $`\theta _{1,2}`$ of the random dimensionless temperature $`t`$ satisfy the inequality $`0\theta _1^2\theta _21.`$ In the pertinent region of fields $`0.2÷5\mathrm{T}`$, a typical defect radius $`L3.5\mathrm{nm}`$ is at least one order of magnitude smaller than the magnetic length, hence the defects can definitely be taken as short-range ones. In the weak field region, the important small parameters are then $`\epsilon (T)`$ (which enters the magnetization (13)) and $`\epsilon (T)/(x+\gamma _0(x))^2`$ (which enters the saddle point equation). Using parameters from the experimental setup $`s=1.5\mathrm{nm}`$, $`k_B\mu _0H_{c2}^{}=1.15\mathrm{TK}^1`$, $`\kappa =100`$, $`n_d=5\times 10^{10}\mathrm{cm}^2`$, $`\mu _0H_\mathrm{\Phi }=1\mathrm{T},`$ $`T=75÷85\mathrm{K},`$ $`T^{}=78.9\mathrm{K},`$ we find from Eq.(12) $`\epsilon (T^{})=0.5\theta _2`$ and $`\epsilon (T^{})/(x^{}+\gamma _0(x^{}))^20.25`$ (the latter figure is obtained for $`\mu _0H=0.2\mathrm{T}`$). For quite plausible value $`\theta _2=0.5`$ one then finds $`\epsilon (T^{})=0.5\theta _2=0.25.`$ The condition of convergence of the integral over the expansion coefficients $`\{C_m\}`$ can be written as $`H>0.25H_\mathrm{\Phi }\theta _1^2/\theta _2`$ and even in the worst case $`\theta _1^2=\theta _2`$ it reads $`\mu _0H0.25\mathrm{T}.`$ Finally, one has $`\mu _0H^{}6.4\mathrm{T}`$ and applicability of the LLL projection requires $`\mu _0H>0.5\mathrm{T}.`$ The weak field region of Ref. corresponds to $`\mu _0H=0.2÷0.02\mathrm{T}.`$ Thus, in the weak field region, the condition for applicability of the LLL projection is slightly violated, but the deviation is not dramatic. In the strong field region we find $`\eta 2.9`$ and therefore the correction term in parenthesis of equation (16) is less than three percents. Hence, in this region our assumptions are fully satisfied.
aaaUsing the same set of parameters we display in figure 1 the quantity $`M/\sqrt{TH}`$ as a function of the scaling variable for weak field (inset) and strong field (main part). We used here the maximal value $`\theta _1=1`$. In the strong field region, the deviation form clean sample scaling behavior is negligibly small for all three values of strong magnetic field.
In the weak field region, the scaling functions for three different fields can hardly be distinguished. This means that scaling is undoubtedly valid in a vicinity of the crossing temperature. At the same time the scaling function differs from its form in a clean sample (1) by a multiplicative constant (see the parenthesis in Eq.(11)). Note that scaling in the weak field region (which was experimentally established) is less pronounced than that in the strong field region. Apparently, the reason is that the experimental data are fitted to account for the clean sample scaling function. Nevertheless if we identify the fitted temperature $`82.6\mathrm{K}`$ (found in Ref. in the weak field region) with the renormalized critical temperature $`T_c=T_{c0}\delta T_c`$, and the fitted critical temperature $`84.2\mathrm{K}`$ in the strong field region with $`T_{c0}`$, then, even within such a rough approximation, we obtain $`\theta _10.5.`$ Recalling that $`\theta _1`$ should be positive and less than unity, the above result strongly supports the applicability of our theory to the pertinent experiment.
aaaIn summary, we calculated the magnetization of an irradiated superconductor below the mean–field transition line $`H_{c2}(T)`$, using the approach developed in Refs.. It was shown that, from a rigorous point of view, disordered short-range defects are expected to destroy the scaling behavior and prevent the existence of crossing point in both regions of weak and strong magnetic fields (with respect to matching field $`H_\mathrm{\Phi }`$). And yet, in the framework of the experimental setup the deviation from scaling behavior appears to be negligibly small and crossing points exist in both field regions, in complete agreement with the experimental findings. The two fitting critical temperatures introduced in Ref. for the strong and weak field regions correspond, in our formalism, to the initial and renormalized critical temperatures.
aaaThe authors would like to thank Z. Tesanović for helpful discussions and P.H. Kes who advised us some parameters of experimental setup of Ref..
aaaThis work was supported by MINERVA Foundation (G.B.), by grants from Israel Academy of Science “Mesoscopic effects in type II superconductors with short-range pinning inhomogeneities” (S.G.), and “Center of Excellence” (Y.A.), and by DIP grant for German-Israel collaboration (Y.A.). |
warning/0003/astro-ph0003473.html | ar5iv | text | # Pulsational 𝑀_𝑉 versus [𝐹𝑒/𝐻] relation(s) for globular cluster RR Lyrae variables
## 1 Introduction
The intrinsic luminosity of RR Lyrae variables has been for a long time a very popular way to give reasonable estimates of the distance to globular clusters (GCs) both in the Milky Way and in Local Group galaxies (Magellanic Clouds, M31) and, in turn, to constrain the age of these very old stellar systems. However, notwithstanding the large body of work, a general consensus on a precise evaluation of such a luminosity has been not yet achieved. One may notice that a firm knowledge of RR Lyrae luminosities would be of paramount relevance, since it would provide an independent test of the Cepheid distance scale as well as a reliable calibration of several secondary distance indicators (as, e.g., the GC luminosity function or the Red Giant Tip) for external galaxies, thus providing important clues on the value of the Hubble constant $`H_0`$. On these grounds, RR Lyrae variables could represent relevant milestones on the path to set both a lower and an upper limit to the age of the Universe, playing a fundamental role in several astrophysical problems ranging from stellar evolution to cosmological models.
From the observational side, studies dealing with the absolute magnitude $`M_V(RR)`$ of RR Lyraes and with the dependence of these magnitudes on the heavy element content $`[Fe/H]`$ has yielded to the well known debate between the so-named ”short” and ”long” distance scales. As recently reviewed by Cacciari (1999), empirical estimates of $`M_V(RR)`$ for RR Lyrae stars at $`[Fe/H]`$=-1.6 actually range from about 0.4 mag to 0.7 mag, thus leaving an uncertainty of $`\pm `$ 0.2 mag on the derived distance moduli (see also Popowski & Gould 1999).
Different estimates have been also given for the dependence of these magnitudes on the star metallicity. As a matter of the fact, for the often assumed linear relation
$$<M_V(RR)>=a+b[Fe/H]$$
one finds in the literature evaluations of the coefficient $`\mathrm{"}b\mathrm{"}`$ mainly in the range $`b0.18\pm 0.03`$ to $``$0.30, where the former value is based on the Baade-Wesselink method (see, e.g., Fernley et al. 1998b \[Fn98b\]) and the latter value has been early suggested by Sandage (1993 \[Sa93\]) when discussing the period-metallicity relation for field and GC RR Lyrae pulsators. However, an even milder slope has been suggested by Fusi Pecci et al. (1996), who investigated eight globular clusters in M31 to derive, over the range -1.8$`<[Fe/H]<`$-0.4,
$$<M_V(HB)>=(0.13\pm 0.07)[Fe/H]+(0.95\pm 0.09)$$
.
The recent release of HIPPARCOS statistical and trigonometric parallaxes for halo RR Lyraes ($`[Fe/H]`$-1.30) has not clarified the issue: one may indeed recall that Fernley et al. 1998a \[Fn98a\] and Groenewegen & Salaris (1999 \[GS99\]), both assuming the same slope $`b`$=0.18, give a zero-point of $`1.05\pm 0.15`$ mag and 0.77$`\pm 0.26`$ mag, respectively, as derived from an identical sample of variables but using different approaches (statistical parallaxes or reduced parallaxes, respectively). In the meantime, McNamara 1999 \[MN99\] claims that Baade-Wesselink results for variables with $`[Fe/H]>1.5`$ yield a quite different relation as given by
$$<M_V(RR)>=1.06+0.32[Fe/H]$$
On the theoretical side, the literature already contains several sets of horizontal branch (HB) evolutionary models computed for wide ranges of the overall metallicity ($`Z`$ in the range of 0.0001 to 0.02) which provide the ”theoretical route” to the calibration of the $`M_V(RR)`$ versus $`[Fe/H]`$ relation. One finds that almost all the recent theoretical predictions concerning the absolute magnitude $`M_V(ZAHB)`$ of the zero age horizontal branch (ZAHB) sequence at the RR Lyrae instability strip confirm the non-linear dependence of $`M_V(ZAHB)`$ on $`\mathrm{log}Z`$ formerly suggested by Castellani, Chieffi & Pulone (1991 \[CCP\]). However, the scaling of the overall metallicity $`Z`$ to the measured $`[Fe/H]`$ values could be a tricky matter since the classical assumption of solar-scaled chemical mixtures is likely inappropriate to GC stars. There is indeed a growing observational evidence for a significant enhancement of $`\alpha `$-elements with respect to iron ($`[\alpha /Fe]0.3`$) in GC and field metal-poor stars (see Carney et al. 1997, Gratton et al. 1997). Moreover, one has to bear in mind that observed RR Lyrae samples do contain stars evolved off their original ZAHB position. Thus, realistic predictions on the average magnitude $`<M_V(RR)>`$ require the evaluation of the evolutionary effects, possibly through synthetic HB simulations (SHB).
The wide grids of SHBs so far published (e.g., Lee, Demarque & Zinn 1990, Lee 1991, Bencivenni et al. 1991, Caputo et al. 1993) have already shown that the predicted mean magnitude of RR Lyrae stars $`<M_V(RR)>`$ significantly depends, with everything else being constant, on the HB morphology. Simulations based on slightly modified CCP models yielded Caputo (1997) to suggest
$$<M_V(RR)>=1.19+0.19\mathrm{log}Z$$
for RR Lyrae-rich metal-poor GCs with $`\mathrm{log}Z3.0`$, whereas for larger metallicities the theory gives
$$<M_V(RR)>=1.57+0.32\mathrm{log}Z$$
Similar results have been more recently found by Demarque et al. (1999), who definitively reject the existence of a unique linear relation covering the metallicities spanned by GCs, confirming that the slope of the predicted $`M_V(RR)\mathrm{log}Z`$ relation depends on the metallicity range and that the HB morphology of each cluster must be taken in the due account when using RR Lyrae stars as distance indicators.
On these grounds, one is tempted to conclude that theoretical and observational investigations do show a sort of consistency: the former give warnings against a ”universal” linear $`M_V(RR)[Fe/H]`$ relation, the latter fail to reach an agreement on both its slope and zero-point!
Within such a confusing scenario, one has to mark the seminal attempts made by Sandage (Sa93 and references therein) to use RR Lyrae periods to constrain the luminosity of these stars. This appears a quite relevant approach, since periods are firm and safe observational parameters, independent of distance and reddening. To discuss Sandage’s philosophy, one has to recall that since the pioneering work by Christy (1966) and Stellingwerf (1975, 1984) pulsating models have suggested the existence within the instability strip of a region where both fundamental ($`RR_{ab}`$) and first overtone ($`RR_c`$) modes are stable (see Bono & Stellingwerf 1994, Bono et al. 1997a, Bono et al. 1997c). The boundaries of this ”either-or” region, namely the fundamental blue edge (FBE) at the higher temperature side and the first overtone red edge (FORE) at the lower temperature side, encompass for each given luminosity the range of temperatures (or colours) where the mode-shift (i.e. the transition from $`RR_{ab}`$ to $`RR_c`$) may occur. Assuming that for both Oosterhoff type I (OoI) and Oosterhoff type II (OoII) globular clusters this transition occurs at the blue edge for fundamental pulsation and using periods and $`BV`$ colours of the shortest period $`RR_{ab}`$ in clusters and in the field, Sa93 derives the star luminosity from the well established period-mass-luminosity-temperature relation. In this way he obtains the relation
$$M_V(RR)=0.94+0.30[Fe/H]$$
which accounts for the Oosterhoff dichotomy in Galactic globular clusters as mainly due to a luminosity effect. However, a re-analysis by Fernley (1993 \[Fn93\]), using $`VK`$ colours and a limited sample of clusters with low and well-known reddening, yields
$$M_V(RR)=0.84+0.19[Fe/H]$$
More recently the assumption of a unique $`RR_{ab}/RR_c`$ transition line has been questioned by Bono, Caputo & Marconi (1995), who concluded that the Oosterhoff dichotomy is largely the result of different transition lines in OoI (near FBE) and OoII (near FORE) clusters, as early suggested by van Albada and Baker (1973). However, the pulsation theory predicts the limits of the whole instability strip, as given by the first overtone blue edge (FOBE) and the fundamental red edge (FRE), without any ambiguity about the actual pulsation mode. On this basis, Caputo (1997) used a preliminary set of pulsating models to show that theoretical predictions on the pulsator distribution in the period-absolute magnitude $`M_V\mathrm{log}P`$ plane can constrain the distance to RR Lyrae-rich globular clusters.
In recent times, the RR Lyrae pulsating models have been updated and extended to wide ranges of mass and chemical composition, shedding light on the dependence of the instability strip on the metal content. In this paper we will take advantage of these improvements to reconsider the Caputo (1997) analysis. The updated $`M_V\mathrm{log}P`$ relations at FOBE and FRE are discussed in the following Sect. 2, while Sect. 3 presents the comparison with observation and the derived ”pulsational” distance moduli for a selected sample of well-studied GCs. The resulting dependence of our $`<M_V(RR)>`$ values on the cluster metallicity is discussed in Sec. 4 in comparison with both empirical relations and recent theoretical HB models. Some concluding remarks will close the paper.
## 2 The predicted $`M_V\mathrm{log}P`$ diagram
The non-linear convective hydrodinamical code used for pulsating models has been already presented in a series of papers (Bono & Stellingwerf 1994, Bono, Caputo & Marconi 1995, Bono et al. 1997c) and it will be not further discussed. With respect to previous computations, the new models differ in the adopted opacity tables, using the most updated compilations by Iglesias & Rogers (1996) and extending in such a way the preliminary results presented in Caputo et al. (1999) for $`Z`$=0.001 and $`M=0.65M_{}`$.
Table 1 presents temperatures, absolute magnitudes and periods of stars at FOBE or FRE, for selected choices on $`Z`$ and suitable values of the star mass ($`M`$), luminosity ($`L`$) and helium abundance ($`Y`$). Absolute magnitudes are derived using the bolometric corrections provided by Castelli, Gratton & Kurucz (1997a,b). The adopted $`Y`$-values reasonably account for the extra helium brought to the stellar surface by the first dredge-up as well as for a galactic enrichment as given by $`\mathrm{\Delta }Y/\mathrm{\Delta }Z2.5`$. As discussed in Caputo, Marconi & Santolamazza (1998), mild variations of $`Y`$, as due also to uncertainties in the efficiency of element sedimentation (see Cassisi et al. 1998, 1999 \[Cs99\]), have negligible effects on the temperature of the instability edges and, in turn, on the related pulsational periods. Stellar masses and luminosities have been chosen in such a way to reasonably encompass available expectations about these evolutionary parameters for GC RR Lyrae pulsators.
It should be noted that, for each given set of entry parameters, the computations have been performed by steps of 100 K and that we adopt as limits of the instability region the average effective temperature between the last pulsating model and the first non pulsating one. It follows that the intrinsic uncertainty of the FOBE and FRE temperatures in Table 1 is $`\pm `$ 50 K, which in terms of period means $`\delta \mathrm{log}P\pm `$0.01 (see Caputo et al. 1998). From the data in Table 1 one derives analytical expressions connecting the absolute V-magnitude at the instability edges with the pulsator period, mass and metallicity, as given by
$`M_V(FOBE)=`$ $`0.6852.255\mathrm{log}P(FOBE)+`$ (1)
$`1.259\mathrm{log}M/M_{}+0.058\mathrm{log}Z`$
$`M_V(FRE)=`$ $`+0.5522.018\mathrm{log}P(FRE)+`$ (2)
$`1.348\mathrm{log}M/M_{}+0.108\mathrm{log}Z,`$
with a rms scatter $`\sigma _V`$=0.027 mag.
However, a further source of uncertainty on the predicted pulsation edges is given by the efficiency of convection in the external layers. The lack of a rigorous treatment of superadiabatic convection is indeed a well known fault in the whole stellar evolution theory and almost all the evolutionary sequences are calculated within the so-called ”mixing-length” scenario which involves an adjustable parameter $`l/H_p`$, the ratio of the mixing-length to the pressure scale height. Our pulsating models are consistent with this scenario and, in order to close the system of convective and dynamic equations, adopt $`l/H_p=1.5`$, in reasonable agreement with the values generally used for evolutionary computations. Since the effect of convection is to quench pulsation, variations of $`l/H_p`$ lead to variations in the effective temperature of the boundaries for instability, with the amount of this effect decreasing from the red to the blue edge.
To have light on such an uncertainty, we performed suitable numerical experiments, finding that decreasing $`l/H_p`$ down to 1.1 (i.e. decreasing the efficiency of convection and thus increasing the local temperature gradients) the FOBE periods decrease by $`\delta \mathrm{log}P0.029`$, while with $`l/H_p`$=2.0 these periods increase by $`\delta \mathrm{log}P0.017`$. On this ground one can estimate that mixing-length values in the range $`l/H_p=1.31.8`$, as widely adopted in the relevant literature, yield an additional uncertainty on FOBE periods by $`\delta \mathrm{log}P\pm 0.01`$. Since at the red side of the instability strip the mixing-length affects much more significantly the predicted periods, in the following we will rely on theoretical predictions concerning FOBE only, temptatively putting the red edge of the pulsation region at $`\mathrm{\Delta }\mathrm{log}P=0.45`$ with respect to FOBE.
As a relevant point, when constraining the luminosities of ZAHB pulsators in the range covered by current evolutionary predictions for various metallicities and for ages in the range from 8 to 18 Gyr, one finds that data in Table 1 predict FOBE effective temperatures very close to $`\mathrm{log}T_e=3.85`$, without significant variation with the metal content. This constant value is due to the balancing effects of metallicity on the FOBE temperature (which decreases if increasing only the luminosity) and ZAHB luminosity: when decreasing $`Z`$ only, the FOBE would become hotter, but the contemporary increase of the ZAHB luminosity eventually leaves unchanged the effective temperature. More in general, we can use such an evidence to safely take from evolutionary theories the predicted masses for HB pulsators at the blue side of the instability strip. Luckily enough, at variance with luminosities, the evolutionary masses have passed substantially unchanged the many improvements affecting HB models in the last years, thus representing a rather firm and trustworthy prediction. As a relevant point, one finds that such evolutionary prediction appears in close agreement with independent mass estimates from the period ratios of double-mode RR Lyrae pulsators (see Cox 1991, Bono et al. 1996).
Thus, by inserting into eq. (1) the predicted mass of the ZAHB model at $`\mathrm{log}T_e=3.85`$ as presented in Table 2 (from Bono et al. 1997b), one finally gets the period-luminosity-metallicity relation for evolutionary FOBE pulsators, as given by
$`M_V(FOBE)=`$ $`0.1782.255\mathrm{log}P(FOBE)+`$ (3)
$`+0.151\mathrm{log}Z,`$
with a total intrinsic dispersion (including the above uncertainty of $`\pm `$50 K, the mixing-length effects and mass variations by 5% the values in Table 2) of $`\sigma _V=0.065`$ mag.
Here one should note that the mass of HB models does depend on the abundance of $`\alpha `$-elements. However, one can benefit by the principle of correspondence for which the evolutionary behaviour of HB stars (and the predicted pulsator masses) depends on the overall metallicity $`Z`$, independently of the internal ratio between $`\alpha `$\- and heavy elements (see, e.g., Bencivenni et al. 1991, Salaris, Chieffi & Straniero 1993). In other words, the comparison between the predicted period-luminosity-metallicity relation given in eq. (3) and the GC RR Lyrae distribution in the observed $`V\mathrm{log}P`$ plane only requires that the scaling between the measured $`[Fe/H]`$ and the overall metallicity $`Z`$ is properly evaluated as
$$\mathrm{log}Z=[Fe/H]1.61+\mathrm{log}(0.638f+0.362),$$
(4)
where $`f`$ is the $`\alpha `$-enhancement factor with respect to iron (Salaris et al. 1993).
## 3 The observed $`V\mathrm{log}P`$ diagram
Since our analysis will be focused on the predicted FOBE, we selected well-studied clusters with statistically significant numbers of $`RR_c`$ stars. The sample of the used GCs is presented in Table 3, which gives for each cluster the adopted iron content $`[Fe/H]`$ (from Carretta & Gratton (1997) or from Zinn & West (1984) and Rutledge, Hesser & Stetson (1997) values transformed into the Carretta & Gratton metallicity scale), HB type (Harris 1996) and mean visual magnitude $`<V_{RR}>`$ of RR Lyrae stars.
For each assumption about the globular cluster distance modulus one obtains the distribution of the cluster RR Lyraes in the $`M_V\mathrm{log}P`$ plane. We derive a “pulsational” evaluation of the cluster distance modulus by constraining the observed $`RR_c`$ distribution to match the predicted blue limit of the pulsation region in order to have no variables in the hot stable region. Figure 1 shows the result of such a procedure, by assuming for the cluster sample a solar-scaled chemical composition, i.e. $`f`$=1 in eq. (4). As already stated, our analysis is focused on $`RR_c`$ stars and the right edge of the instability strip has been simply placed at $`\mathrm{\Delta }\mathrm{log}P=0.45`$ with respect to the left edge. However, it seems worthy of notice the fair agreement found also between the predicted FRE and the $`RR_{ab}`$ distribution.
The derived GC apparent distance moduli $`DM_V`$ are summarised in Table 4, together with the resulting mean absolute magnitude $`<M_V(RR)>`$ of RR Lyrae stars. The total errors on $`<M_V(RR)>`$ listed in Table 4 account for the observed dispersion $`\sigma _{<V(RR>}`$ (see Table 3) and the predicted total uncertainty ($`\pm `$0.07 mag) of the FOBE period-luminosity-metallicity relation. The last column gives the weight $`W`$ to our $`<M_V(RR)>`$ estimates, as simply derived from the number of $`RR_c`$ stars matching the predicted FOBE.
Note that the results in Table 2 refer to solar-scaled chemical compositions. If the chemical mixtures are $`\alpha `$-enhanced, then for each cluster the nominal metallicity Z increases \[see eq. (4)\] and the derived distance modulus decreases \[see eq. (3)\]. As a matter of example, with $`f`$=3 all the distance moduli in Table 4 have to be decreased by 0.05 mag, with a consequent increase of $`<M_V(RR)>`$.
## 4 The $`<M_V(RR)>`$ versus metallicity relation(s)
The final correlation between our $`<M_V(RR)>`$ and the cluster metallicity $`[Fe/H]`$ is presented in Fig. 2 for the two cases $`f`$=1 and $`f`$=3, assuming for each $`[Fe/H]`$ an error of $`\pm `$ 0.15 dex. The same figure shows the already quoted observational calibrations based on RR Lyrae periods (Sa93, Fn93), HIPPARCOS data for field RR Lyraes (Fn98a) and the Baade-Wesselink method (Fn98b). We add the result by Carretta et al. (1999 \[Cr99\]), as based on HIPPARCOS parallaxes for field subdwarfs and Main-Sequence fitting procedure, while the GS99 relation, which is only 0.03 mag fainter than Cr99, is not presented for the sake of clearness.
Inspection of Fig. 2 reveals that there is a general agreement between pulsational and other empirical calibrations, except the Fn98a relation which definitively suggests too faint magnitudes. If $`f`$=1, then also the Fn98b relation appears fainter with respect to our results. However, one derives that none of the empirical linear calibrations is able to fully match our pulsational results over the whole range of metal content. The lack of a full agreement is largely due to the fact that our data foresee a non linear relation between $`<M_V(RR)>`$ and metallicity, in agreement with stellar evolution theoretical predictions. We show in Table 5 that a bare linear best fit to the data in Table 4, starting from the four metal-poorest clusters with $`[Fe/H]`$-2.0 (NGC 4590, NGC 5053, NGC 5466, NGC 7078) and regularly increasing the metallicity range, yields a $`M_V(RR)[Fe/H]`$ relation which becomes steeper and steeper when moving towards metal-rich clusters, suggesting a change in slope at $`[Fe/H]`$-1.5. It seems worthy of notice that this result agrees with the recent analysis of RR Lyrae variables in the field (MN99) and the globular cluster $`\omega `$ Centauri (Rey et al. 2000).
On this ground, the least squares solutions performed through our weighted $`<M_V(RR)>`$ values with $`[Fe/H]<`$-1.5 and $`[Fe/H]>`$-1.5 yield
$`<M_V(RR)>=`$ $`0.71(\pm 0.10)+(0.17\pm 0.04)[Fe/H]+`$ (5)
$`+0.03f`$
and
$`<M_V(RR)>=`$ $`0.92(\pm 0.12)+(0.27\pm 0.06)[Fe/H]+`$ (6)
$`+0.03f,`$
respectively, in agreement with the empirical calibrations by Fn93, GS99 and Cr99 (metal-poor clusters) and MN99 (metal-rich clusters).
As shown in Fig. 3, where OoI and OoII clusters are depicted with different symbols, around $`[Fe/H]`$-1.5 the Oosterhoff dichotomy shows off, leading us to guess that the two above $`<M_V(RR)>[Fe/H]`$ relations hold for the two Oosterhoff groups. More interestingly, the same figure suggests that OoII clusters have brighter RR Lyraes than OoI clusters with similar metal content, an evidence which coupled with their blue HB morphology (see HB types in Table 4) confirms that the RR Lyrae evolutionary stage is more important than metallicity in triggering the Oosterhoff dichotomy, as early suggested by Lee, Demarque & Zinn (1990) and recently supported by independent investigations (Lee & Carney 1999, Clement & Shelton 1999).
Figure 4 finally presents the pulsational $`<M_V(RR)>`$ results as a function of log$`Z`$, as derived through eq. (4), together with selected theoretical predictions based on stellar evolution theory. The lines drawn in the figure refer to recent $`M_V(ZAHB)\mathrm{log}Z`$ calibrations as given by Cassisi & Salaris (1997 \[CS\]), Cs99, Caloi, D’Antona & Mazzitelli (1997 \[CDM\]) and Ferraro et al. (1999 \[Fr99\]), with the predicted ZAHB magnitude decreased by 0.1 mag to account for the luminosity excess of actual RR Lyrae stars over the ZAHB level.
Inspection of the figure reveals that none of the theoretical predictions fully agrees with our pulsational results with $`f`$=1, as the luminosities provided by Cs99 and CDM are systematically too bright, while those by CS and Fr99 match our results only with log$`Z<`$-3.0. If $`f`$=3, then the CS and Fr99 calibrations appear the most consistent with our data, but with a tendency to overestimate the luminosity of the most metal-rich variables. It seems worth noticing that the Cs99 relation is well reproducing all our data with $`f`$=3, but with an overluminosity of about 0.08 mag.
## 5 Concluding remarks
In this paper we have used results from the most recent and updated computations of non-linear convective pulsating models to constrain the distance modulus of Galactic globular clusters through the observed periods of $`RR_c`$ pulsators. The resulting $`<M_V(RR)>[Fe/H]`$ relation appears in the range of several empirical linear calibrations, but with evidence for a non linearity which suggests that the slope of the relation increases when moving towards the metal-richer variables. On observational grounds, a similar behavior seems present among RR Lyrae stars in $`\omega `$ Centauri (Rey et al. 2000). Moreover, we notice that over the range of metal-poor stars ($`[Fe/H]<`$-1.5) our pulsational calibration is in good agreement with the relations given by Fernley (1993), Groenewegen & Salaris (1999) and Carretta et al. (1999), while with $`[Fe/H]>`$-1.5 it agrees with MacNamara (1999) results.
Application of our results to RR Lyrae stars of the metal-poor globular clusters in the Large Magellanic Cloud (see data in GS) would give a distance modulus of 18.61$`\pm `$0.12 mag ($`f`$=1) and 18.56$`\pm `$0.12 mag ($`f`$=3), thus supporting the “long” distance scale (see also Romaniello et al. 1999).
By relying on the present pulsational RR Lyrae absolute magnitudes, one derives that the non linearity of our $`<M_V(RR)>[Fe/H]`$ relation is well reproduced by current predictions based on stellar evolution theory. However, in the case of solar-scaled chemical compositions, none of the evolutionary predictions published in the recent literature appears in satisfactory agreement, supporting observational evidence for $`\alpha `$-enhanced chemical mixtures in metal-poor stars. With the $`\alpha `$-elements enhanced by a factor of 3 with respect to iron, the predictions by CS and Fr99 agree with our pulsational magnitudes even though with a tendency of overestimating the luminosity of metal-rich pulsators. Interesting enough, one finds that the Cs99 relation is well reproducing the general dependence of $`M_V(RR)`$ on log$`Z`$, but with an overluminosity of about 0.08 mag. Holding CS99 results, a beautiful agreement with our data would be achieved by sistematically increasing the cluster metallicity by $``$ 0.2 dex, an occurrence hardly to be accepted.
To further discuss this point, one has to remind that differences in stellar models are mainly, if not only, the result of differences in the adopted input physics. Discussing RR Lyrae stars in the globular cluster M5 (Caputo et al. 1999) we have already reported pulsational evidence suggesting that models with the ”most updated” input physics (as in Cs99) give too luminous HB stars. Such an evidence has been further supported by independent estimates based on HIPPARCOS parallaxes for clumping field He burning stars (Castellani et al. 1999). Data in the previous Fig. 4 reinforce such an evidence, suggesting that the ”most updated physics” is probably far from being the most adequate one. As a whole, we remain with the tantalising evidence that Cs99 models give the rightest metal dependence but not right luminosities, whereas those by Fr99 and CS give much better luminosities but slightly worst slope.
The role played by the various physical ingredients in determining the predicted luminosity of HB structures has been recently discussed in several papers (see Cassisi et al. 1999, Castellani and Degl’Innocenti 1999, Castellani 1999) and cannot be repeated here. However, one may notice that the most recent theoretical predictions displayed in Fig. 4 all agree within a range of luminosity of about $`\pm `$ 0.05 mag. This in our feelings should be regarded as an evidence of the high standard reached by evolutionary theories, as well as a warning that better precision should require a corresponding level of accuracy in the input physics not yet reached by currently available evaluations.
Acknowledgment:
It is a pleasure thank the referee, B. Carney, for his valuable report. We deeply thank B. Carney and M. Corwin for providing us with data on NGC 5466 before publication. Thanks are also due to Santi Cassisi for several warm discussions with one of us (F.C.) during an icy week in Teramo. Financial support for this work was provided by the italian Ministero dell’Università e della Ricerca Scientifica e Tecnologica (MURST) under the scientific project “Stellar Evolution”. |
warning/0003/astro-ph0003446.html | ar5iv | text | # High resolution HST STIS spectra of C I and CO in the 𝛽 Pictoris (catalog HD39060) circumstellar disk
## 1 Introduction
$`\beta `$ Pictoris is the most extensively studied of the young planetary systems discovered in the last decade and a half. It is a bright Southern hemisphere star (type A5 V), located about 19.3 pc distant from the Sun, with a systemic radial velocity of 20 km s<sup>-1</sup> (for a review of the $`\beta `$ Pictoris system, see Vidal-Madjar, Lecavelier des Etangs, & Ferlet (1998)). It was observed in 1983 by the IRAS satellite to have a large excess of emission at infrared wavelengths. This was referred to as the Vega-like phenomenon and was identified as arising from an edge-on circumstellar (hereafter CS) dust disk, presumed to be associated with planetary formation (Smith & Terrile, 1984). It was soon determined through absorption spectroscopy that there was CS gas associated with the dust as well.
A large body of evidence has accumulated indicating that there are comet-like bodies present in the $`\beta `$ Pictoris CS disk. Collisions between dust particles are expected to produce submicron fragments which should be expelled from the system by radiation pressure on time scales much shorter than any plausible stellar age. Thus there must be a secondary source of particles; one model for the production of dust and gas in the CS disk focusses on evaporating comets and is called the Orbiting Evaporating Bodies model (OEB). In this picture, the comets orbit the star at several tens of AU, and thus, the $`\beta `$ Pictoris CS disk is a kind of “gigantic multi-cometary tail with its natural constituents: gas and dust” (Lecavelier des Etangs, Vidal-Madjar, & Ferlet, 1996). Spectra of $`\beta `$ Pictoris show variable redshifted absorption features arising from gas infalling toward the star at high velocities (and infrequently, blueshifted features as well); these features are best attributed to the evaporation of star-grazing comets, called the Falling Evaporating Bodies scenario (FEB) (Beust et al., 1990). Also, gas at close to 20 km s<sup>-1</sup> (zero radial velocity relative to the star) is identified in all observations and is called the stable gas component. This gas is difficult to understand, as modeling indicates that it should be expelled from the system by radiation pressure; a continuous source for this gas is required.
Neutral carbon and carbon monoxide have been observed in HST-GHRS UV absorption spectra of $`\beta `$ Pictoris; carbon monoxide is the only molecule detected in the CS disk to date (Vidal-Madjar et al., 1994; Jolly et al., 1998). Since CO and C I can be dissociated and ionized by interstellar UV photons on time scales of the order of 200 years, both must be continuously replenished. Carbon monoxide, in particular, is difficult to reform after dissociation in the $`\beta `$ Pictoris environment. Thus, the presence of these species indicates that a secondary source for this gas should exist, just as for the CS dust (Vidal-Madjar et al., 1994). Jolly et al. (1998) found the column densities of C I and CO to be comparable, around $`10^{15}`$ cm<sup>-2</sup>; since their rates of destruction are also comparable, this was taken as evidence that the C I is produced by photodissociation of CO, which evaporates from comets orbiting at various distances and velocities. However, the C I column density was determined from a heavily saturated multiplet and is therefore quite uncertain. In the hopes of further constraining the characteristics of C I and CO in the $`\beta `$ Pictoris system, we have reinvestigated the transitions observed in the GHRS data, as well as some that were not seen due to the relatively low spectral resolution of GHRS compared with that of the STIS high resolution echelle.
## 2 Observations
HST STIS high resolution echelle spectra of $`\beta `$ Pictoris were obtained on 1997 December 6 and 1997 December 19, covering the wavelength range from 1459 Å to 2888 Å in six exposures each day. Table 1 shows the log of observations. All the absorption features discussed in this paper appear in either the first or second data set listed for each day (o4g001010/o4g002010 or o4g001020/o4g002020).The data were initially reduced and calibrated using the STScI IRAF package *calstis v1.8*. Spectra with a signal-to-noise ratio of around 10 were achieved. Examination of the errors in the flux values showed that the error propagation calculation had not been performed correctly and that thestated errors were too small. Therefore the data were re-calibrated using *calstis v1.9a*, correcting the underestimate of the measurement errors.
The particular advantage of this data set over previous comparable ones is the very high spectral resolution achieved. The instrumental line spread function using the E140H grating is well described by a Gaussian with a FWHM of 1 pixel, corresponding to FWHM = $`\lambda `$/220,000 (Sahu et al., 1999). However, since the detector undersamples the line spread function, the effective resolution to separate two adjacent lines with this grating is only $`R`$ = 110,000. Using the E230H grating, the FWHM is about 2 pixels, corresponding to FWHM = $`\lambda `$/110,000 (Sahu et al., 1999).
## 3 Analysis
### 3.1 C I\] $`\lambda `$1613 line
This previously undetected spin-forbidden transition $`(^3`$P$`{}_{0}{}^{}_{}^{1}`$P$`{}_{1}{}^{})`$ at 1613.376 Å was clearly observed on both days. It is unsaturated, which allows for an improved determination of the column density of C I. As there was no change in the line between the two days of observation, the spectra were averaged together to improve the S/N; the result is shown in Figure 1. Unfortunately, the other two fine structure lines in the multiplet were not reliably detected. Atomic data were taken from Morton (1991) and the continuum around the line was fit with a sixth-degree polynomial. Voigt line profiles were used to generate transmission functions, which were then convolved with a Gaussian instrumental line spread function with FWHM = $`\lambda `$/220,000 to create model spectra. The $`\chi ^2`$ statistic between the model and the data was then minimized to determine the best velocity centroid, $`v`$, column density of C I in the ground level, $`N(^3\mathrm{P}_0)`$, and Doppler broadening parameter, $`b`$; error bars for these parameters were determined from the contours of $`\chi ^2`$. Since the line was unsaturated, the column density was also determined from the equivalent width of the line and $`b`$ and $`v`$ were determined from a simple Gaussian fit to the line, as a check on the results of the $`\chi ^2`$ minimization.
### 3.2 C I(<sup>3</sup>P) $`\lambda \lambda `$1561 and 1657 multiplets
The central portion of the 1561 Å and 1657 Å multiplets, arising from the <sup>3</sup>P ground term, were heavily saturated in the STIS data. However, this allowed us to examine the smaller variable red and blueshifted features in the spectra which are the signatures of infalling comets in the CS disk. In Figures 2 and 3, one can easily see absorption lying at higher redshift in the December 6 data that is not visible in the December 19 data. A model of the multiplet which contains only one velocity component is unable to reproduce the absorption dip at 1561.5 Å or the absorption at higher redshift in the December 6 data.
Due to the fact that the multiplets were extremely saturated and that the multiple velocity components were blended together, we were not able to perform $`\chi ^2`$ minimizations to find unique best values for the model parameters. Fe II absorption features in our data set were analyzed and the velocities of their multiple velocity components found. Models of the C I $`\lambda `$1561 multiplet were constructed, containing the same number of components as were found in the Fe II features, with roughly the same velocities. The parameters of the stable component at $`v`$ = 20 km s<sup>-1</sup> were set to those determined from the analysis of the unsaturated 1613.376 Å line. The remaining parameters were then adjusted to find the best model by eye. This model was compared to the 1657 Å multiplet to confirm that the values found in this way were reasonable. A small constant value ($`2\times 10^{13}`$ erg s<sup>-1</sup> cm<sup>-2</sup> Å<sup>-1</sup>) was subtracted from the 1657 Å data to bring the baselines of the totally saturated features to zero.
### 3.3 C I(<sup>1</sup>D) $`\lambda \lambda `$1931 and 1463 lines
These lines arise from the excited, metastable <sup>1</sup>D state of the ground configuration and have not previously been observed in spectra of $`\beta `$ Pictoris. The <sup>1</sup>D level has a lifetime of about 4000 s, lying about 10,000 cm<sup>-1</sup> above the ground level; C atoms in this state may be produced during the photodissociation of CO in solar system comet comae (Tozzi, Feldman, & Festou, 1998). The 1931 Å lines were analyzed using the same model generating and $`\chi ^2`$ minimization procedures described in § 3.1; the atomic data used were from Hibbert et al. (1993). Although there may be some reason to suspect that the strong central absorption features near 1931.05 Å contain multiple velocity components, since this behavior is seen in lines arising from excited levels of other atomic species, the $`\chi ^2`$ minimizations indicated that stable unique solutions containing more than one component in the central absorption features could not be found. Thus, best models were found using only one velocity component. However, the difficulty involved in modeling a saturated, blended line is such that there may well be other undetected components present, so models with two components in the central absorption feature were compared to the data by eye. The 1463 Å spectra, shown in Figure 5, were so noisy that they served only to roughly confirm the values found from the 1931 Å line. A small constant flux value ($`1\times 10^{13}`$ erg s<sup>-1</sup> cm<sup>-2</sup> Å<sup>-1</sup>) was subtracted from the 1463 Å data to bring the baselines of the saturated lines to zero.
### 3.4 CO Fourth Positive band system
The (0-0), (1-0), and (2-0) bands of the Fourth Positive system of CO ($`A^1\mathrm{\Pi }X^1\mathrm{\Sigma }^+`$) did not vary between the two days of observation. The spectra were therefore averaged together to improve the S/N; the result appears in Figure 6. Only a single velocity component was observed. Note the detection of bands arising from <sup>13</sup>CO, the perturbation band $`e^3\mathrm{\Sigma }^{}X^1\mathrm{\Sigma }^+`$ (1-0), not previously observed in spectra of $`\beta `$ Pictoris, and the strong perturbation band $`d^3\mathrm{\Delta }X^1\mathrm{\Sigma }^+`$ (5-0).
Models were generated using wavelengths and oscillator strengths from Morton & Noreau (1994). Energies of the ground state levels were calculated using the Dunham coefficients from Farrenq et al. (1991) and LTE assumed in order to determine the population of the rotational levels. The parameters of the model are the rotational excitation temperature, $`T`$, the Doppler broadening parameter, $`b`$, the column density of <sup>12</sup>CO in the ground vibrational state, $`N(^{12}\mathrm{CO})`$, the column density of <sup>13</sup>CO in the ground vibrational state, $`N(^{13}\mathrm{CO})`$, and the velocity centroid, $`v`$. $`\chi ^2`$ minimization was then performed on all three bands simultaneously.
## 4 Results
### 4.1 C I(<sup>3</sup>P) stable component
The C I\] $`\lambda `$1613.376 line showed only one velocity component at 20 km s<sup>-1</sup>, the systemic velocity of the star; the nominal uncertainty in the absolute wavelength calibration of STIS leads to an error in velocity determinations of about 1 km s<sup>-1</sup> (Sahu et al., 1999). This velocity and the fact that the line did not change between the two days of observation identifies the line as arising from stable gas. The results of the $`\chi ^2`$ minimizations for all the features analyzed appear in Table 2. The column density for stable carbon in the ground level, <sup>3</sup>P<sub>0</sub>, determined from $`\chi ^2`$ minimization was found to agree with that determined from the equivalent width of the line. Similarly the $`v`$ and $`b`$ determined from $`\chi ^2`$ minimization agreed with the values found from fitting a simple Gaussian to the unsaturated line.
Using a 3-$`\sigma `$ upper limit on the column density of C I in the <sup>3</sup>P<sub>2</sub> level (from the non-detection of the fine structure line at 1614.5068 Å) and assuming LTE, a firm upper limit of 100 K on the excitation temperature of <sup>3</sup>P carbon in the stable component was found. From analysis of the 1561 Å multiplet, discussed below, it was found that the excitation temperature of the stable component must be greater than about 50 K, or the multiplet could not be reasonably modeled. This range in temperature allows us to determine that the total column density of stable C I in the <sup>3</sup>P ground term is $`(24)\times 10^{16}`$ cm<sup>-2</sup>.
This column density is more than an order of magnitude larger than the total <sup>3</sup>P column density found by Jolly et al. (1998) from GHRS data taken in November 1994, $`N(^3\mathrm{P})`$$`=2\times 10^{15}`$ cm<sup>-2</sup>. Either the abundance of stable C I has varied over the three years between observations or, as is more likely, the difficult modeling of the heavily saturated C I $`\lambda `$1561 multiplet in the GHRS data led to an inaccurate column density. Models with column densities of $`10^{16}`$ cm<sup>-2</sup> in the stable component were compared to the 1994 GHRS spectrum of the 1561 Å multiplet and found to fit the data equally well as models with the lower column density. We found the Doppler broadening parameter of the stable C I(<sup>3</sup>P) to be $`1.3\pm 0.5\mathrm{km}\mathrm{s}^1`$. The previous work on the GHRS data found a Doppler broadening parameter of 4.2 km s<sup>-1</sup>, which was much greater than the parameter found for the other atomic species and CO (Jolly et al., 1998). Our smaller $`b`$ value is equal to that found for CO in our data; we thus do not see any excess kinetic energy in the motions of the C I atoms.
### 4.2 C I(<sup>3</sup>P) multiple velocity components
The difficulty in modeling a saturated, blended multiplet is formidable, and the values determined for the multiple velocity components of C I(<sup>3</sup>P) from the 1561 Å multiplet are not reliable. However, this analysis did confirm that models with column densities of $`10^{16}`$ cm<sup>-2</sup> in the stable component could reasonably fit the 1561 Å spectra. We can also conclude that the total column density of C I(<sup>3</sup>P) in the variable components is about $`10^{14}`$ cm<sup>-2</sup> on December 6 and about $`10^{15}`$ cm<sup>-2</sup> on December 19, to within an order of magnitude or so. It is also clear that the variable C I features are generally better fit with higher excitation temperatures and larger $`b`$ parameters, but these values are very poorly determined.
### 4.3 C I(<sup>1</sup>D)
Using models with only one velocity component in the strong central absorption features, two components arising from the excited <sup>1</sup>D level of C I were found in the December 6 data, only one in the December 19 data. The best models are overplotted on the spectra in Figures 4 and 5. Note that the central absorption features do not appear at the same velocity on the two days, that they also do not appear at the systemic velocity of the star, and that the column density of C I(<sup>1</sup>D) changes significantly between the two days of observation. This would seem to indicate that there is no stable component for C I(<sup>1</sup>D). Models with two velocity components in the central absorption feature, which were compared to the data by eye, indicate that there could be a “stable” component which has the same column density on both days. However, it must be at about 22-23 km s<sup>-1</sup>, which is significantly different from the systemic velocity of $`\beta `$ Pictoris. These data are not able to conclusively determine the velocity structure present in the C I(<sup>1</sup>D) gas; observation of unsaturated lines at similarly high resolution will be necessary.
### 4.4 CO
The CO bands showed no multiple velocity component structure and no change between the two days of observation; the absorbing gas is thus entirely associated with the stable component at the systemic velocity of the star. The best model is shown overplotted on the data in Figure 6. The rotational excitation temperature and the ratio of $`N(^{12}\mathrm{CO})`$ to $`N(^{13}\mathrm{CO})`$ show no significant change from the values found by Jolly et al. (1998). $`R(^{12}\mathrm{CO}/^{13}\mathrm{CO})`$ = $`15\pm 2`$ is quite small compared to typical values found in the ISM, e.g. $`R`$ = $`150\pm 27`$ for the diffuse clouds toward $`\zeta `$ Ophiuchi (Sheffer et al., 1992). However, chemical fractionation at low kinetic temperatures can explain this unusual ratio, as discussed below. The column densities of both <sup>12</sup>CO and <sup>13</sup>CO were found to be about 1/3 the values found from the GHRS data. Note that the column density of <sup>12</sup>CO is about 50 times smaller than the column density of stable C I(<sup>3</sup>P).
## 5 Discussion
Although CO in solar system comet comae is photodissociated by solar FUV photons, the CO in the $`\beta `$ Pictoris CS disk is primarily destroyed by interstellar photons. The dissociation energy for CO is 11.1 eV; thus CO may only be dissociated by photons with wavelengths shortward of $``$ 1100 Å. The type A5 star $`\beta `$ Pictoris lacks the strong FUV emission lines created in the Sun’s chromosphere; thus it emits very little flux in the FUV. The only source for CO-dissociating photons at $`\beta `$ Pictoris is therefore the interstellar UV radiation field. Since CO photodissociates primarily through discrete line absorptions, self-shielding can have a strong effect on the abundance of CO in interstellar clouds (Van Dishoeck & Black, 1988). However, this is unlikely to occur in the $`\beta `$ Pictoris CS disk because of the very small transverse dimension of the disk. Thus, the photodissociation rate for <sup>12</sup>CO in the $`\beta `$ Pictoris disk should be equal to the unshielded value and should also be equal to the value for <sup>13</sup>CO. Since <sup>13</sup>CO is not selectively dissociated in this situation, its high abundance was explained by the reaction
$${}_{}{}^{13}\mathrm{C}_{}^{+}+^{12}\mathrm{CO}^{13}\mathrm{CO}+^{12}\mathrm{C}^++35\mathrm{K}$$
which favors the production of <sup>13</sup>CO at gas kinetic temperatures below 35 K (Jolly et al., 1998). Assuming chemical equilibrium and that isotopic exchange is much more important than photodissociation for both <sup>12</sup>CO and <sup>13</sup>CO (Sheffer et al., 1992),
$$\frac{n(^{12}\mathrm{CO})}{n(^{13}\mathrm{CO})}=\mathrm{exp}\left(\frac{35}{T_{kin}}\right)\left(\frac{{}_{}{}^{12}\mathrm{C}}{{}_{}{}^{13}\mathrm{C}}\right)=15\pm 2.$$
The rotational excitation temperature of CO often does not accurately describe the gas kinetic temperature in diffuse environments (Wannier, Penprase, & Andersson, 1997). However, assuming a $`(^{12}\mathrm{C}/^{13}\mathrm{C})`$ ratio, we may use the above expression to estimate the true gas kinetic temperature of the CO.
The average carbon isotopic ratio in the local ISM is $`6070`$ (Langer & Penzias, 1993) and the typical solar system value found in comets is $`(^{12}\mathrm{C}/^{13}\mathrm{C})=89`$ (Jewitt et al., 1997). Using the range $`(^{12}\mathrm{C}/^{13}\mathrm{C})=8960`$, we find that the gas kinetic temperature of the carbon monoxide is 20 K $``$ 25 K, indicating that the CO gas is indeed colder than the stable C I gas. This suggests that the C I and CO are not located in the same regions of the disk. The assumption that isotopic exchange is more important than photodissociation should be reasonable. For this assumption to apply,
$$\mathrm{\Gamma }k^f\mathrm{exp}\left(\frac{35}{T_{kin}}\right)n(^{12}\mathrm{C}^+),$$
where $`\mathrm{\Gamma }`$ is the unshielded photodissociation rate, $`2\times 10^{10}\mathrm{s}^1`$ (Van Dishoeck & Black, 1988), k<sup>f</sup> is the forward reaction rate, $`6.8\times 10^{10}\mathrm{cm}^3\mathrm{s}^1`$ at 80 K (Smith & Adams, 1980), and n$`(^{12}\mathrm{C}^+)`$ is the volume density of $`{}_{}{}^{12}\mathrm{C}_{}^{+}`$. The average volume density of C, $`n(^{12}\mathrm{C})=N({}_{}{}^{3}\mathrm{P})/r20\mathrm{cm}^3`$, assuming that the carbon extends over a distance $`r`$ = 100 AU. Further assuming that $`n(^{12}\mathrm{C}^+)n({}_{}{}^{12}\mathrm{C})`$, as is likely, the right hand side of the above expression is greater than or equal to about $`2\times 10^9\mathrm{s}^1`$, an order of magnitude larger than $`\mathrm{\Gamma }`$. However, a more complete treatment of the relationship between $`(^{12}\mathrm{CO}/^{13}\mathrm{CO})`$ and $`(^{12}\mathrm{C}/^{13}\mathrm{C})`$ takes into account photodissociation but requires exact knowledge of the density of $`{}_{}{}^{12}\mathrm{C}_{}^{+}`$ in the CS disk (Sheffer et al., 1992).
The much larger column density of C I(<sup>3</sup>P) compared to that of CO leads us to believe that photodissociation of CO cannot be the only source of stable C I in the $`\beta `$ Pictoris CS disk. Since there is no evidence that CO is produced by infalling comets (no variability or red and blueshifted features), it has been postulated that the CO gas slowly evaporates from the OEBs at several tens of AU from the star and is photodissociated to produce the stable C I (Lecavelier des Etangs, 1998). Obviously some portion of the C I gas must be produced directly from the FEBs (the portion giving rise to the variable red and blueshifted absorption); perhaps this C I gas is decelerated somehow and accumulates in the stable component before being destroyed by photoionization. In this scenario, the equilibrium column density of C I is
$$N(^3\mathrm{P})=\frac{n\times N_{\mathrm{𝐹𝐸𝐵}}}{\mathrm{\Gamma }}$$
where $`n`$ is the mean number of infalling comets per year, $`10^2`$ per year (Vidal-Madjar, Lecavelier des Etangs, & Ferlet, 1998), $`N_{FEB}`$ is the total column density of C I gas produced by an infalling comet, and $`\mathrm{\Gamma }`$ is the photoionization rate for C I, 0.004 yr<sup>-1</sup>. This expression, which assumes that C I atoms are lost only through photoionization by interstellar UV photons, gives $`N_{FEB}10^{11}`$ cm<sup>-2</sup>, which should easily be produced by infalling comets like the ones giving rise to the variable C I(<sup>3</sup>P) components in our data. However, this is a rough treatment and the number of infalling comets per year varies significantly over time scales of a few years (Vidal-Madjar, Lecavelier des Etangs, & Ferlet, 1998). Also, this treatment does not take into account reformation of C I by radiative recombination; we cannot take this process into account properly without knowledge of the amount of C II in the CS disk and the electron density.
Considering the C I(<sup>1</sup>D), when such atoms are produced by photodissociation of CO, O I(<sup>1</sup>D) atoms must be produced also to conserve spin. The minimum total photon energy needed for this dissociation is 14.33 eV, corresponding to a threshold wavelength of 865 Å. Since this threshold is below the Lyman limit, there are virtually no interstellar UV photons capable of producing C I(<sup>1</sup>D) by photodissociation of CO. Thus, the C I(<sup>1</sup>D) atoms in the $`\beta `$ Pictoris CS disk cannot be produced by photodissociation of CO by stellar or interstellar photons and must be produced by a collisional process involving ground state carbon atoms. Since the energy of the <sup>1</sup>D state relative to the ground state is high, the collisional process must be a very energetic one and therefore is likely to be closely associated with the infalling bodies. Consequently, with a short C I(<sup>1</sup>D) lifetime, there would be no stable component in this gas; this behavior has not been previously seen in any constituent of the $`\beta `$ Pictoris CS disk. But this result is tentavive; examination of an unsaturated line at high resolution is needed to determine the velocity structure of the C I(<sup>1</sup>D) gas.
## 6 Concluding Remarks
The very high resolution, low scattered light contamination, and good order separation of this STIS echelle data set has provided some clear advantages over previous observations of $`\beta `$ Pictoris. The rotational lines of CO have been resolved, allowing for a much more precise determination of the physical parameters of the gas. The column density of CO is $`N(\mathrm{CO})=(6.3\pm 0.3)\times 10^{14}\mathrm{cm}^2`$ and the ratio $`R(^{12}\mathrm{CO}/^{13}\mathrm{CO})=15\pm 2`$ is found. The absence of transient red or blueshifted components in the high resolution CO spectra supports the suggestion that this gas evaporates from cometary bodies orbiting far (several tens of AU) from the star. But the fact that the column density of CO is only about 2% of the total column density of C I in the <sup>3</sup>P ground term implies that photodissociation of this CO is not the primary source for C I gas. It could perhaps be produced directly from infalling comets close to the star, but the mechanism by which it comes to zero velocity relative to the star and accumulates before being photoionized is unclear (although see Lagrange et al. (1998)). The C I(<sup>1</sup>D) gas may not have a stable component at 20 km s<sup>-1</sup>; this unique species could prove to be a valuable tracer of FEB activity in the $`\beta `$ Pictoris CS disk.
Despite the advantages of this data set, our lack of success in modeling the heavily saturated C I multiplets indicates that in order to really determine the characteristics of the variable components of the C I gas, we need to observe an unsaturated line or multiplet, with an oscillator strength between that of the 1561 Å multiplet and that of the spin-forbidden 1613.376 Å line. A number of suitable multiplets and lines lie in the FUV, shortward of $``$ 1300 Å. Also, an unsaturated line arising from the <sup>1</sup>D level would allow us to confirm the velocity structure in this gas and to determine if the velocities of the <sup>1</sup>D gas components correspond with any of the velocities of the variable components in the <sup>3</sup>P gas. Three likely lines lie between 1311 Å and 1359 Å. Measurement of the densities of C II and O I would greatly help to unravel the carbon chemistry of the $`\beta `$ Pictoris disk. Again, potentially useful multiplets of these species lie in the FUV below 1340 Å. Thus, although this data set has vastly increased our knowledge about the important species C I and CO in the $`\beta `$ Pictoris disk, our understanding would probably benefit greatly from investigation of $`\beta `$ Pictoris at shorter ultraviolet wavelengths.
We thank Jason McPhate for his work on our absorption line profile codes and John Debes for his work on the Fe II lines. We also thank B-G Andersson and our reviewer, X. Tielens, for their fruitful comments. This work is based on observations with the National Aeronautics and Space Administration – European Space Agency HST obtained at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Incorporated, under NASA contract NAS5-26555. Support for this work at JHU was provided by grant GO-07512.01-96A from the Space Telescope Science Institute. |
warning/0003/hep-ph0003250.html | ar5iv | text | # A STERILE NEUTRINO NEEDED FOR HEAVY-ELEMENT NUCLEOSYNTHESIS
## 1 Introduction
While there is strong evidence that the heaviest elements are produced in the neutrino-heated material ejected relatively long ($`10`$ s) after the explosion of a Type II or Type I b/c supernova, present calculations show conditions which would prevent this rapid-neutron-capture (or $`r`$) process from occurring. Though general relativistic effects and multi-dimensional hydrodynamic outflow have been invoked to solve these problems, these solutions are at best exceedingly finely tuned. In contrast, the solution presented here is extremely robust, and the neutrino mass-mixing scheme it requires is exactly that needed if one is to explain all present evidence for neutrino mass.
In the next section, the particle physics motivations for that neutrino scheme are discussed, with special emphasis on recent results of the LSND experiment, since those are particularly important for the needs of the $`r`$ process. That is followed by a section on the $`r`$ process and its difficulties, after which the solution is presented.
## 2 Particle Physics Evidence for a Four-Neutrino Scheme
The need for at least one light sterile (i.e., not having the usual weak interaction) neutrino in addition to the known three active neutrinos was proposed as a way to explain the solar $`\nu _e`$ deficit, the anomalous $`\nu _\mu /\nu _e`$ ratio from atmospheric neutrinos, and the apparent need for appreciable hot dark matter. The atmospheric anomaly, due to $`\nu _\mu \nu _\tau `$, requires a mass-squared difference between the $`\nu _\mu `$ and $`\nu _\tau `$ of $`\mathrm{\Delta }m_{\mu \tau }^2=0.003`$ eV<sup>2</sup> and maximal mixing ($`\mathrm{sin}^22\theta _{\mu \tau }=1.0`$), with the latter property being quite important to solving the $`r`$-process problem, as shown below. The solar $`\nu _e`$ deficit is explained by $`\nu _e\nu _s`$ with $`\mathrm{\Delta }m_{es}^210^5`$ eV<sup>2</sup> and a small mixing angle ($`\mathrm{sin}^22\theta _{es}0.01`$) or $`\mathrm{\Delta }m_{es}^210^{10}`$ eV<sup>2</sup> and $`\mathrm{sin}^22\theta _{es}1`$. The $`\nu _e`$$`\nu _s`$ pair is of lower mass than the $`\nu _\mu `$$`\nu _\tau `$ pair, which provide the hot dark matter. For the originally favored critical-mass-density universe the $`\nu _\mu `$ and $`\nu _\tau `$ needed masses of around 2 eV each.
This phenomenology was subsequently given theoretical bases in two 1993 papers, and has been utilized since in a large number of publications. Since that time the scheme has received support from the LSND experiment, the results of which provide some measure of the mass difference between the $`\nu _e`$$`\nu _s`$ and $`\nu _\mu `$$`\nu _\tau `$ pairs and hence of the neutrino contribution to dark matter.
In its 1996 publication, LSND claimed a signal in $`\overline{\nu }_\mu \overline{\nu }_e`$ on the basis that 22 events of the type $`\overline{\nu }_epe^+n`$ were seen, using a stringent criterion to reduce accidental coincidences between $`e^{}`$ or $`e^+`$ and $`\gamma `$ rays mimicking the 2.2-MeV $`\gamma `$ from $`npd\gamma `$, whereas only $`4.6\pm 0.06`$ events were expected. The probability of this being a fluctuation is $`4\times 10^8`$. Note especially that these data were restricted to the energy range 36 to 60 MeV to stay below the $`\overline{\nu }_\mu `$ endpoint and to stay above the region where backgrounds are high due to the $`\nu _e^{12}\mathrm{C}e^{}X`$ reaction. In plotting $`\mathrm{\Delta }m^2`$ vs. $`\mathrm{sin}^22\theta `$, however, events down to 20 MeV were used to increase the range of $`E/L`$, the ratio of the neutrino’s energy to its distance from the target to detection. This plot was intended to show the favored regions of $`\mathrm{\Delta }m^2`$, and all information about each event was used. The likelihood analysis applied did not have a Gaussian likelihood distribution, since its integral is infinite, but the likelihood contour labeled “90%” was obtained by going down a factor of 10 from the maximum, as in the Gaussian case. The contours in the LSND plot have been widely misinterpreted as confidence levels—which they certainly are not—because they were plotted along with confidence-level limits from other experiments.
Recently the difficult, computer-intensive analysis in terms of real confidence levels has been done. The likelihood for a grid in ($`\mathrm{sin}^22\theta `$, $`\mathrm{\Delta }m^2`$) space, including backgrounds, has been computed and compared with numerous Monte Carlo experiments to obtain a 90% confidence region. While the equivalency varies from point to point in the $`\mathrm{\Delta }m^2`$$`\mathrm{sin}^22\theta `$ plane, a typical value for the 90% confidence level is down a factor of 20 from the likelihood maximum. Thus the LSND allowed regions are considerably broader in $`\mathrm{sin}^22\theta `$ than in the plots published so far, and other experiments constrain allowed $`\mathrm{\Delta }m^2`$ regions less.
The confusion of comparing likelihood levels for LSND with confidence levels from other experiments may be exacerbated by using the 20–36 MeV region for the LSND data. While this higher background energy range makes some difference for the 1993–5 data, it could have had an appreciable effect for the parasitic 1996–8 runs, which were at a low event rate, increasing the effect of cosmic ray background. This could raise the low end of the supposed signal energy spectrum, especially as the one LSND distribution which was statistically worrisome was the ratio ($`R`$) of real to accidental events. Some accidental events in this 20–36 MeV region would favor low values of $`\mathrm{\Delta }m_{e\mu }^2`$ making the higher $`\mathrm{\Delta }m^2`$ values desirable for dark matter appear less likely.
Nevertheless, when a joint analysis is made of the LSND and KARMEN experiments even using the 20–36 MeV range for LSND, the region around 5.5 eV<sup>2</sup> is as probable as the banana-shaped region at lower $`\mathrm{\Delta }m^2`$, as shown in Fig. 1.
Frequently ignored by theorists, this higher mass region is favored by the $`\nu _\mu \nu _e`$ LSND data. Of course in the $`\nu _\mu \nu _e`$ case, using $`\nu _\mu `$ from $`\pi ^+`$ decay in flight and detecting $`\nu _e`$ by $`\nu _e^{12}\mathrm{C}e^{}X`$, the backgrounds are higher and hence yield much poorer statistics than for $`\overline{\nu }_\mu \overline{\nu }_e`$ with $`\overline{\nu }_\mu `$ from $`\mu ^+`$ at rest. In addition to the $`\mathrm{\Delta }m^2`$ issue, the important point of Fig. 1 is that although the KARMEN data are consistent with background, the joint analysis of the $`\overline{\nu }_\mu \overline{\nu }_e`$ data from the two experiments shows an appreciable region for a signal. KARMEN is continuing to take data, and LSND will have an improved analysis available soon. This new analysis has produced an excellent $`R`$ distribution for all the data and an energy distribution with reduced contributions at the low end, favoring higher $`\mathrm{\Delta }m_{e\mu }^2`$ values than does Fig. 1.
## 3 Problems with Synthesis of the Heaviest Elements
While in the next section we will find that the $`r`$ process of rapid neutron capture in supernovae provides strong support for the double doublet of neutrinos, initially the reverse appeared to be true, with the $`r`$ process apparently placing stringent limits on $`\nu _\mu `$$`\nu _e`$ mixing. The origin of these limits is that energetic $`\nu _\mu `$ ($`E25`$ MeV) coming from deep in the supernova core could convert via an MSW transition to $`\nu _e`$ inside the region of the $`r`$-process, producing $`\nu _e`$ of much higher energy than the thermal $`\nu _e(E11`$ MeV). The latter, because of their charged-current interactions, emerge from farther out in the supernova where it is cooler. Since the cross section for $`\nu _ene^{}p`$ rises as the square of the energy, these converted energetic $`\nu _e`$ would deplete neutrons, stopping the $`r`$-process. Calculations of this effect limit $`\mathrm{sin}^22\theta `$ for $`\nu _\mu \nu _e`$ to $`<10^4`$ for $`\mathrm{\Delta }m_{e\mu }^2>2`$ eV<sup>2</sup>, in conflict with at least the higher mass region of the LSND results, which will be of particular interest here.
More recently, serious problems have been found with the $`r`$ process itself. First, simulations have revealed the $`r`$-process region to be insufficiently neutron-rich, since about $`10^2`$ neutrons is required for each seed nucleus, such as iron. This was bad enough, but the recent realization of the full effect of $`\alpha `$-particle formation has created a disaster for the $`r`$ process. At a radial region inside where the $`r`$ process should occur, all available protons swallow up neutrons to form the very stable $`\alpha `$ particles, following which $`\nu _ene^{}p`$ reactions reduce the neutrons further and create more protons which make more $`\alpha `$ particles, and so on. The depletion of neutrons rapidly shuts off the $`r`$ process, and essentially no nuclei above $`A=95`$ are produced.
To solve this problem the $`\nu _e`$ flux has to be removed before the $`r`$ process site, while leaving a very large $`\nu _e`$ flux at a smaller radius for material heating and ejection. The obvious difficulty of accomplishing this has led to searches for other possible sites for the $`r`$ process, such as neutron star mergers.
## 4 Neutrino Solution for a Successful $`r`$ Process
The apparent miracle of having a huge $`\nu _e`$ flux disappear before it reaches the radius of the supernova where $`\alpha `$ particles form can be accomplished if there is (1) a sterile neutrino, (2) approximately maximal $`\nu _\mu \nu _\tau `$ mixing, (3) $`\nu _\mu \nu _e`$ mixing $`>10^4`$, and (4) an appreciable ($`>1`$ eV<sup>2</sup>) mass-squared difference between $`\nu _s`$ and the $`\nu _\mu `$$`\nu _\tau `$. This is precisely the neutrino mass pattern required to explain the solar and atmospheric anomalies and the LSND result, plus providing some hot dark matter!
Such a mass-mixing pattern creates two level crossings. The inner one, which is outside the neutrinosphere (beyond which neutrinos can readily escape) is near where the $`\nu _{\mu ,\tau }`$ potential $`(n_{\nu _e}n_n/2)`$ goes to zero. Here $`n_{\nu _e}`$ and $`n_n`$ are the numbers of $`\nu _e`$ and neutrons, respectively. The $`\nu _{\mu ,\tau }\nu _s`$ transition which occurs depletes the dangerous high-energy $`\nu _{\mu ,\tau }`$ population. Outside of this level crossing, another occurs where the density is appropriate for a matter-enhanced MSW transition corresponding to whatever $`\mathrm{\Delta }m_{e\mu }^2`$ LSND is observing. Because of the $`\nu _{\mu ,\tau }`$ reduction at the first level crossing, the dominant process in the MSW region reverses from the deleterious $`\nu _{\mu ,\tau }\nu _e`$, becoming $`\nu _e\nu _{\mu ,\tau }`$ and dropping the $`\nu _e`$ flux. For an appropriate value of $`\mathrm{\Delta }m_{e\mu }^2`$, the two level crossings are separate but sufficiently close so that the transitions are coherent. Then with adiabatic transitions (as calculations show) and maximal $`\nu _\mu `$$`\nu _\tau `$ mixing, the neutrino flux emerging from the second level crossing is 1/4 $`\nu _\mu `$, 1/4 $`\nu _\tau `$, and 1/2 $`\nu _s`$, and no $`\nu _e`$.
A more exact way to explain this in the four-neutrino formalism is to transform the four mass eigenstates into the flavor states via the mixing angles $`\varphi `$ for solar $`\nu _e\nu _s`$, $`\omega `$ for LSND $`\nu _\mu \nu _e`$, and $`\pi /4`$ (maximal mixing) for atmospheric $`\nu _\mu \nu _\tau `$. Symbolically,
$$\left(\begin{array}{c}|\nu _s\\ |\nu _e\\ |\nu _\mu \\ |\nu _\tau \end{array}\right)=\left(\begin{array}{c}\text{Atmospheric}\\ \nu _\mu \nu _\tau \end{array}\right)\left(\begin{array}{c}\text{Solar}\\ \nu _e\nu _s\end{array}\right)\left(\begin{array}{c}\text{LSND}\\ \nu _\mu \nu _e\end{array}\right)\left(\begin{array}{c}|\nu _1\\ |\nu _2\\ |\nu _3\\ |\nu _4\end{array}\right),$$
giving
$$\left(\begin{array}{c}|\nu _e\\ |\nu _s\\ |\nu _\mu ^{}\\ |\nu _\tau ^{}\end{array}\right)=\left(\begin{array}{cccc}\mathrm{cos}\varphi & \mathrm{sin}\varphi \mathrm{cos}\omega & \mathrm{sin}\varphi \mathrm{sin}\omega & 0\\ \mathrm{sin}\varphi & \mathrm{cos}\varphi \mathrm{cos}\omega & \mathrm{cos}\varphi \mathrm{sin}\omega & 0\\ 0& \mathrm{sin}\omega & \mathrm{cos}\omega & 0\\ 0& 0& 0& 1\end{array}\right)\left(\begin{array}{c}|\nu _1\\ |\nu _2\\ |\nu _3\\ |\nu _4\end{array}\right),$$
where
$`|\nu _\mu ^{}`$ $`=1/\sqrt{2}(|\nu _\mu |\nu _\tau )=\mathrm{sin}\omega |\nu _2+\mathrm{cos}\omega |\nu _3`$
$`|\nu _\tau ^{}`$ $`=1/\sqrt{2}(|\nu _\mu +|\nu _\tau )=|\nu _4,\text{a mass eigenstate}.`$
In this formalism what occurs at the first level crossing is $`\nu _\mu ^{}\nu _s`$, and at the second, $`\nu _e\nu _\mu ^{}`$, while $`\nu _\tau ^{}`$ being a mass eigenstate goes through both regions unaffected. Again this gives 1/4 $`\nu _\mu `$, 1/4 $`\nu _\tau `$, 1/2 $`\nu _s`$, and no $`\nu _e`$ at all.
Note that the $`\overline{\nu }_e`$ flux is also unaffected at the level crossings, so $`\overline{\nu }_epe^+n`$ enhances the neutron number in the $`r`$ process region, since the protons have not been depleted by $`\alpha `$ particle formation. It should be emphasized that this mechanism is quite robust, not depending on details of the supernova dynamics, especially as it occurs quite late in the explosive expansion.
It is essential that the two level crossings be in the correct order, and this provides a requirement on $`\mathrm{\Delta }m_{e\mu }^2`$, since the MSW transition depends on density and hence on radial distance from the protoneutron star. Detailed calculations have been made for $`\mathrm{\Delta }m_{e\mu }^26`$ eV<sup>2</sup>, which works very well. Possibly $`\mathrm{\Delta }m_{e\mu }^2`$ as low as 2 eV<sup>2</sup> or maybe even 1 eV<sup>2</sup> would work, but that is speculative. At any rate, the mass difference needed in this scheme, which is the only one surely consistent with all manifestations of neutrino mass and which rescues the $`r`$ process, implies appreciable hot dark matter.
## 5 Conclusions
It is quite remarkable that the profound problems of producing the heaviest elements by supernovae can be solved in a manner which requires no adjustment of parameters if the arrangement of masses and mixings of neutrinos is exactly that required to explain the solar $`\nu _e`$ deficit, the atmospheric neutrino anomaly, and the observations of the LSND experiment (or alternatively the need for hot dark matter). This is achieved via an active-sterile level crossing in the supernova, followed by an active-active transition. The total independence of this supernova information strongly enhances the case for this four-neutrino scheme.
## Acknowledgments
This paper is based largely on work done with G.M. Fuller and Y.-Z. Qian, to whom I am grateful, and was supported in part by the U.S. Department of Energy under contract DE-FG03-91ER40618. |
warning/0003/nucl-th0003017.html | ar5iv | text | # Rho-Nucleon Tensor Coupling and Charge-Exchange Resonances
(13 March 2000
Revised 15 July 2000 )
## Abstract
The Gamow-Teller resonance in <sup>208</sup>Pb is discussed in the context of a self-consistent RPA, based on the relativistic mean field theory. We inquire on the possibility of substituting the phenomenological Landau-Migdal force by a microscopic nucleon-nucleon interaction, generated from the rho-nucleon tensor coupling. The effect of this coupling turns out to be very small when the short range correlations are not taken into account, but too large when these correlations are simulated by the simple extraction of the contact terms from the resulting nucleon-nucleon interaction.
PACS: 21.60.-n; 21.60.Jz; 21.30.Fe; 24.30.Cz
Keywords: Relativistic mean field theory; Random phase approximation; One-boson-exchange models; Charge-exchange resonances
The quantum hadrodynamics (QHD) aims to describe the nuclear many-body system in terms of nucleons and mesons . Proposed initially as a full-fledged renormalizable quantum field theory, nowadays it is seen as an effective field theory, derivable, in principle, from the quantum chromodynamics .
The relativistic mean field theory (RMFT), which can be thought as a mean field (Hartree) approximation to the QHD, has been applied with great success during the last few decades. For instance, it accounts for both i) the nuclear matter saturation, and ii) the ground state properties of finite nuclei along the whole periodic table . More recently, the RMFT has also been exploited for the description of unstable nuclei all up to the nucleon drip lines .
Through a relativistic version of the random phase approximation (RRPA), various excited states and resonances have been studied in the context of the RMFT as well. Quite recently we have also reported the first calculation of this type for the Gamow-Teller (GT) and isobaric analogue (IA) resonances, excited from the ground states of <sup>48</sup>Ca, <sup>90</sup>Zr and <sup>208</sup>Pb nuclei.
Because of its pseudoscalar nature, the pion does not participate in the description of the the ground states in the RMFT. Thus, besides the nucleon and the Coulomb fields, only the $`\sigma `$, $`\omega `$ and $`\rho `$ mesons are usually involved in the calculations. Yet, in dealing with isovector excitations it is essential to include, together with the $`\rho `$ meson, the $`\pi `$ meson as well. This has already been done in our previous work , with the pseudovector pion-nucleon coupling $`f_\pi `$ fixed at its experimental value. For the remaining mesons, only the nonderivative couplings to the nucleon were included, as usually done in RMFT. With this prescription we were not able to reproduce the excitation energies of the just mentioned resonances. This has been possible only after introducing the repulsive Landau-Migdal (LM) delta force
$`V_{LM}(1,2)`$ $`=`$ $`g^{}\left({\displaystyle \frac{f_\pi }{m_\pi }}\right)^2𝝉_1𝝉_2𝝈_1𝝈_2\delta (𝒓_1𝒓_2),`$ (1)
of the same magnitude ($`g^{}=0.7`$) as the one used in the nonrelativistic calculations .
Here we wish to analyze whether the tensor (derivative) coupling of the $`\rho `$ meson to the nucleon could generate a sufficiently repulsive nucleon-nucleon force in order to locate the GT resonance at the correct experimental energy and in this way substitute the phenomenological LM force. The IA resonance is practically not affected by this part of the $`\rho `$-meson-exchange potential and therefore it will not be discussed so exhaustively as we do with the GT resonance.
As mentioned above, it is not usual to include the tensor coupling of the vector mesons to the nucleon in RMFT. This is because its effect on the ground state is (rightly) thought to be small. On a more general perspective, however, there are two good reasons why one should do so. For one, according to the rules of effective field theory such terms should appear in the effective QHD Lagrangian . For another, and perhaps more important reason for the phenomenological stand we are taking, it is well known that the tensor $`\rho `$-nucleon coupling gives a large contribution to the spin-isospin component of the nucleon-nucleon interaction , and as such it could have an important effect on the dynamics of the GT resonance.
Our Lagrangian density is now
$``$ $`=`$ $`\overline{\psi }(i\gamma _\mu ^\mu M)\psi `$ (2)
$`+{\displaystyle \frac{1}{2}}_\mu \sigma ^\mu \sigma {\displaystyle \frac{1}{2}}m_{\sigma }^{}{}_{}{}^{2}\sigma ^2{\displaystyle \frac{1}{3}}g_2\sigma ^3{\displaystyle \frac{1}{4}}g_3\sigma ^4g_\sigma \overline{\psi }\psi \sigma `$
$`{\displaystyle \frac{1}{4}}\mathrm{\Omega }_{\mu \nu }\mathrm{\Omega }^{\mu \nu }+{\displaystyle \frac{1}{2}}m_{\omega }^{}{}_{}{}^{2}\omega _\mu \omega ^\mu g_\omega \overline{\psi }\gamma _\mu \psi \omega ^\mu `$
$`+{\displaystyle \frac{1}{2}}_\mu 𝝅^\mu 𝝅{\displaystyle \frac{1}{2}}m_{\pi }^{}{}_{}{}^{2}𝝅𝝅{\displaystyle \frac{f_\pi }{m_\pi }}\overline{\psi }\gamma _5\gamma _\mu 𝝉\psi ^\mu 𝝅`$
$`{\displaystyle \frac{1}{4}}𝑹_{\mu \nu }𝑹^{\mu \nu }+{\displaystyle \frac{1}{2}}m_{\rho }^{}{}_{}{}^{2}𝝆_\mu 𝝆^\mu g_\rho \overline{\psi }\gamma _\mu 𝝉\psi 𝝆^\mu `$
$`{\displaystyle \frac{f_\rho }{2M}}\overline{\psi }\sigma _{\mu \nu }𝝉\psi ^\mu 𝝆^\nu `$
$`{\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }e\overline{\psi }\gamma _\mu {\displaystyle \frac{1+\tau _3}{2}}\psi A^\mu ,`$
where
$`\mathrm{\Omega }^{\mu \nu }`$ $`=`$ $`^\mu \omega ^\nu ^\nu \omega ^\mu ,`$
$`𝑹^{\mu \nu }`$ $`=`$ $`^\mu 𝝆^\nu ^\nu 𝝆^\mu 2g_\rho 𝝆^\mu \times 𝝆^\nu ,`$ (3)
$`F^{\mu \nu }`$ $`=`$ $`^\mu A^\nu ^\nu A^\mu .`$
This Lagrangian is identical to that of reference , except for the $`\rho `$-nucleon tensor coupling term (the one proportional to $`f_\rho `$).<sup>1</sup><sup>1</sup>1There is a minor correction to be made in . One must replace $`g_\rho `$ by $`2g_\rho `$ in eqs. (1) and (2) of that reference for consistency with the remaining equations. Therefore, following the same route one arrives at identical equations for the mean boson fields, except for that of the $`\rho `$ meson (only the component $`\rho _3^0`$ survives for spherical, definite-charge nuclei), which now takes the form
$$\left(^2+m_{\rho }^{}{}_{}{}^{2}\right)\rho _3^0=g_\rho \rho _3(r)+\frac{f_\rho }{2M}𝝆_{t3}(r),$$
(4)
where the (vector) isovector density $`\rho _3`$ is as defined in and we have introduced the tensor isovector density
$$𝝆_{t3}=\overline{\psi }i𝜶\tau _3\psi =\underset{\alpha =1}{\overset{A}{}}\overline{𝒰}_\alpha i𝜶\tau _3𝒰_\alpha .$$
(5)
The summation is over all the occupied single-particle, positive-energy states $`𝒰_\alpha `$, which obey the mean-field Dirac equation. This is also modified to
$`\left\{i𝜶+\beta \left[M+V_s(r)\right]+V_v(r)+(i\beta 𝜶𝒓/r)V_t(r)\right\}𝒰_\alpha `$ $`=`$ $`E_\alpha 𝒰_\alpha .`$
Again the scalar ($`V_s`$) and vector ($`V_v`$) potentials are as defined in , while the tensor potential,
$$V_t=\frac{f_\rho }{2M}\frac{d\rho _3^0}{dr}\tau _3,$$
(7)
is the contribution from the tensor-coupling term in (2).
The general structure and derivation of the RRPA for charge-exchange excitations, in the discretized spectral version we use, has been delineated in . An alternative, more detailed account can be found in . The main ingredient is the residual interaction $`V`$. For a self-consistent calculation, this must be obtained from the same Lagrangian (1) used for the mean field. Also, since Fock terms are ignored in RMFT, we must consider only the direct matrix elements of $`V`$. Hence only the isovector mesons contribute, and we get $`V=V_\pi +V_\rho `$, with, in the instantaneous approximation,
$`V_\pi (1,2)`$ $`=`$ $`\left({\displaystyle \frac{f_\pi }{m_\pi }}\right)^2𝝉_1𝝉_2(𝝈_1_1𝝈_2_2)Y(m_\pi ,r_{12}),`$ (8)
$`V_\rho (1,2)`$ $`=`$ $`𝝉_1𝝉_2[(g_\rho {\displaystyle \frac{f_\rho }{2M}}i\beta 𝜶)_1(g_\rho {\displaystyle \frac{f_\rho }{2M}}i\beta 𝜶)_2`$
$`(g_\rho 𝜶+{\displaystyle \frac{f_\rho }{2M}}\beta 𝝈\times )_1(g_\rho 𝜶+{\displaystyle \frac{f_\rho }{2M}}\beta 𝝈\times )_2]Y(m_\rho ,r_{12}),`$
where $`r_{12}=\left|𝒓_1𝒓_2\right|`$ and $`Y(m,r)=\mathrm{exp}(mr)/(4\pi r)`$.
For the numerical values of the parameters we follow mostly the philosophy of , adopting the parameter set NL1 . Yet, in view of the difficulties encountered by Ma et al. in accounting for the E1 and E0 giant resonances with the NL1 parameters, a few results for the TM1 model, worked out by Sugahara and Toki , will be presented as well.<sup>2</sup><sup>2</sup>2In the latter case, the $`\omega `$-meson self-interaction term, not appearing in (2), was also included in the numerical calculations. Taking experimental values for the pion, the only new parameter is the $`\rho `$-nucleon tensor coupling constant $`f_\rho `$. As mentioned in , the vector dominance model predicts for the ratio $`f_\rho /g_\rho K_\rho `$ a value equal to the isovector magnetic moment of the nucleon, i.e., $`\mu _p\mu _n1=3.7`$. On the other hand, most meson-exchange models for the nuclear force use $`K_\rho =6.6`$ . The former choice was preferred in the description of the ground state properties in closed shell nuclei within the relativistic Hartree-Fock approximation . Thus, the discussion that follows will mainly rely on the lower value for $`K_\rho `$, even though we are aware of the fact that the inclusion of Fock terms can considerably change the adjusted values of the QHD parameters .
Another point to consider is whether the inclusion of the tensor coupling term in the Lagrangian (2) does not sensitively affect the values of the remaining parameters. Fortunately, while the contribution of this term is not strictly zero in RMFT, its effects on the single particle energies as well as on the ground state properties are certainly very small. We therefore feel justified in keeping the remaining parameters fixed at their NL1 or TM1 values. With $`K_\rho =3.7`$, for instance, the spin-orbit splitting is modified in less than $`150`$ KeV. <sup>3</sup><sup>3</sup>3For identical particles, the NL1 paramerization yields significantly larger spin-orbit splittings than the TM1 model, while the opposite happens for nonidentical particles. Similarly tiny effects on the energy per particle and the root-mean-square radii are displayed in Table 1. An interesting side remark can be made concerning the latter observables. It is well known that, at variance with the nonrelativistic calculations, it is a common feature of the relativistic models to overestimate the neutron skin thickness . But, as seen from the results shown in Table 1, the tensor $`\rho `$-N coupling has the tendency to correct the RMFT for this handicap. This fact, in turn, could have very important consequences on the estimates of the atomic parity nonconservation .
The GT and IA resonances in <sup>208</sup>Pb were computed in RRPA, for both the NL1 and TM1 sets of parameters and within the same model space as that of , i.e., including only $`0\mathrm{}\mathrm{\Omega }`$ and $`2\mathrm{}\mathrm{\Omega }`$ excitations, and only those single-particle states that are bound at least for neutrons. For simplicity, we ignored the negative-energy states, although it has been shown that they are required in principle even if the no-sea approximation is made for RMFT, since one needs a complete single-particle basis to develop a perfectly consistent RRPA. In fact, the transitions from Fermi- to Dirac-sea states are essential to ensure certain desirable features, such as current conservation and the removal of the spurious $`J^\pi =1^{}`$ translational state. However, such issues are not crucial for our present purposes and, furthermore, Ma et al. have shown in a recent calculation that the contribution of the negative-energy states is of decisive importance only for the isoscalar modes. We therefore feel safe to leave their inclusion for a future, more sophisticated and detailed treatment of those isovector resonances.
In Fig. 1 are shown the NL1 results for the GT strength distribution, both in terms of the individual strengths,
$$s_\lambda =|\underset{p\overline{n}}{}X_{p\overline{n}}^\lambda p||𝝈||\overline{n}+\underset{n\overline{p}}{}Y_{n\overline{p}}^\lambda \overline{p}||𝝈||n|^2,$$
(10)
and of a “strength function” obtained by replacing the spikes by Lorentzians of conveniently chosen widths $`\mathrm{\Delta }`$ , i.e.,
$$S(E)=\frac{\mathrm{\Delta }}{\pi }\underset{\lambda }{}\frac{s_\lambda }{(EE_\lambda )^2+\mathrm{\Delta }^2},$$
(11)
where $`X_{p\overline{n}}^\lambda `$ and $`Y_{n\overline{p}}^\lambda `$ are, respectively, the forward and backward going RPA amplitudes for the state at excitation-energy $`E_\lambda `$. The upper, middle and lower panels correspond, respectively, to: (a) $`K_\rho =0,g^{}=0`$; (b) $`K_\rho =3.7,g^{}=0`$ and (c) $`K_\rho =3.7,g^{}=0.7`$. From these results one is induced to conclude that the tensor $`\rho `$-N coupling has a very small effect on the GT resonance. That is, it seems as though this coupling could merely redistribute the GT strength in the energy region between 5 and 15 MeV, but was incapable of promoting it to the correct experimental energy. The latter is only achieved after introducing an LM force of the same magnitude that has been used in the previous calculation, where the just mentioned coupling has not been considered at all . The issue of the NN-force generated by the $`\rho `$-N coupling is, however, not so simple and deserves further discussion, which is presented below. Before proceeding, let us just mention that we have not noticed large differences between the NL1 and TM1 results for the IA and GT resonances. For instance, in the case (c) we get that these excitations are localized at: $`E_{IA}`$(NL1) = 18.6 MeV and $`E_{GT}`$(NL1) = 19.5 MeV and $`E_{IA}`$(TM1) = 18.7 MeV and $`E_{GT}`$(TM1) = 20.3 MeV, while the experimental results are: $`E_{IA}`$(exp) = 18.8 MeV and $`E_{GT}`$(exp) = 19.2 MeV. Thus, henceforth only the parametrization NL1 will be used.
In the upper panel of Fig. 2 are confronted several diagonal $`J^\pi =1^+`$ proton-particle neutron-hole matrix elements for the $`V_{LM}`$, $`V_\pi `$, $`V_\rho ^{VV}`$, $`V_\rho ^{VT}`$ and $`V_\rho ^{TT}`$ potentials. (The meaning of the upper indices is self-explanatory.) One can see, in particular, that the matrix elements of $`V_\rho ^{TT}`$ are very small in comparison with those coming from $`V_{LM}`$. However, when we rewrite $`V_\rho ^{TT}`$ in the form
$`V_\rho ^{TT}(1,2)`$ $`=`$ $`\left({\displaystyle \frac{f_\rho }{2M}}\right)^2𝝉_1𝝉_2\beta _1\beta _2\{(𝜶)_1(𝜶)_2Y(m_\rho ,r_{12})`$ (12)
$``$ $`{\displaystyle \frac{1}{3}}m_\rho ^2\left({\displaystyle \frac{3}{m_\rho ^2r_{12}^2}}+{\displaystyle \frac{3}{m_\rho r_{12}}}+1\right)Y(m_\rho ,r_{12})S_{12}`$
$`+`$ $`{\displaystyle \frac{2}{3}}[m_\rho ^2Y(m_\rho ,r_{12})\delta (𝒓_1𝒓_2)]𝝈_1𝝈_2\}`$
and evaluate different parts separately, we find out that the Yukawa and contact pieces in the last term engender, each one, very large matrix elements. In fact, as shown in the lower panel, their individual values are larger than those of $`V_{LM}`$, but the overall contribution to $`V_\rho ^{TT}`$ is small, because they cancel each other very strongly. (A similar cancelation, though not so pronounced, also occurs in the case of $`V_\pi `$.)
It should be remembered that the contact terms in $`V_\pi `$ and $`V_\rho ^{TT}`$ would be smeared over a finite region if finite-nucleon-size effects (FNSE) were introduced, and they would be totally killed by realistic short range correlations (SRC). <sup>4</sup><sup>4</sup>4Note, however, that the contributions of the contact terms are nonzero when, both the FNSE, and the SRC are considered simultaneously . Yet, none of these two effects is considered in a mean field treatment, such as the present one. In return, it is common practice to extract the contact parts from (8) and (LABEL:9) by adding to the residual interaction the correction term $`\delta V=\delta V_\pi +\delta V_\rho `$, with
$`\delta V_\pi (1,2)`$ $`=`$ $`{\displaystyle \frac{1}{3}}\left({\displaystyle \frac{f_\pi }{m_\pi }}\right)^2𝝉_1𝝉_2𝝈_1𝝈_2\delta (𝒓_1𝒓_2),`$
$`\delta V_\rho (1,2)`$ $`=`$ $`{\displaystyle \frac{1}{3}}\left({\displaystyle \frac{f_\rho }{2M}}\right)^2𝝉_1𝝉_2\beta _1\beta _2\left(𝜶_1𝜶_2+2𝝈_1𝝈_2\right)\delta (𝒓_1𝒓_2).`$
For consistency, one must also perform such an extraction in the mean field part. Since, differently from the Hamiltonian formalism followed in , we are working within a Lagrangian formalism, we did this extraction in the baryon self-energy computed in the Hartree approximation (which is equivalent to RMFT). As a consequence the replacement $`V_tV_t+\delta V_t`$ has to be done in the Dirac equation (LABEL:6), with
$$\delta V_t=\frac{1}{3}\left(\frac{f_\rho }{2M}\right)^2\frac{𝝆_{t3}𝒓}{r}\tau _3$$
(14)
being a correction that arises upon the extraction from the baryon self-energy of the contact part due to this derivative coupling in eq. (2). But, when this recipe is implemented in the numerical calculation we get too much repulsion and the GT resonance is pushed up very high in energy. This comes from the fact that $`\delta V`$ is basically a $`\delta `$-force of the type (1), with
$$g_{\pi +\rho }^{}\frac{1}{3}+\frac{2}{3}\left(\frac{f_\rho }{f_\pi }\right)^2\left(\frac{m_\pi }{2M}\right)^2=1.6,$$
(15)
which is significantly larger than $`g^{}=0.7`$.
Note that in the nonrelativistic approximation the contact term also appears in $`V_\rho ^{VV}`$ and $`V_\rho ^{VT}`$, and instead of (15) one would have
$$g_{\pi +\rho }^{}\frac{1}{3}+\frac{2}{3}\left(\frac{g_\rho +f_\rho }{f_\pi }\right)^2\left(\frac{m_\pi }{2M}\right)^2=2.3.$$
(16)
There is no consensus on whether one should proceed in the same way in the relativistic case. Some authors exclude the contact terms only from $`V_\pi `$ and $`V_\rho ^{TT}`$ , while others do that for the full $`\pi +\rho `$ interaction . That the potentials $`V_\rho ^{VV}`$ and $`V_\rho ^{VT}`$ also contain a contact term follows from the substitution
$$\gamma _\mu \frac{1}{2M}(2P_\mu +\sigma _{\mu \nu }^\nu )$$
(17)
for the vector $`\rho `$-N coupling.
It is worth noting that Toki and Weise have interpreted microscopically the LM force as arising from the $`\pi +\rho `$ meson-exchange model combined with the SRC and FNSE. In the static limit, which is used here, the result is :
$`g_{LM}^{}(\omega =q=0)`$ $``$ $`{\displaystyle \frac{1}{3}}\left({\displaystyle \frac{\mathrm{\Lambda }^2m_\pi ^2}{\mathrm{\Lambda }^2+m_0^2}}\right)^2{\displaystyle \frac{m_0^2}{m_0^2+m_\pi ^2}}`$ (18)
$`+`$ $`{\displaystyle \frac{2}{3}}\left({\displaystyle \frac{g_\rho +f_\rho }{f_\pi }}\right)^2\left({\displaystyle \frac{m_\pi }{2M}}\right)^2{\displaystyle \frac{m_0^2}{m_0^2+m_\rho ^2}},`$
where $`\mathrm{\Lambda }`$ is the cut-off mass for the pion-nucleon vertex and $`m_0^1`$ is the correlation length. For $`\mathrm{\Lambda }=1`$ GeV, $`m_0=m_\rho `$ and $`g_\rho +f_\rho =17.2`$ this leads to $`g_{LM}^{}(\omega =q=0)=0.67`$ . (In the present work $`g_\rho +f_\rho =23.4`$.)
Our results can be summarized as follows:
1. When the short range correlations are not considered, the tensor $`\rho `$-nucleon coupling plays only a minor role in the description of the GT resonances.
2. If one tries to take these correlations into account by merely extracting the contact terms from the NN interaction, the GT resonance is pushed up too high in energy.
Thus, the simulation of the short range correlations by the simple-minded extraction of the contact terms alone is not a satisfactory procedure; at least not in the case of the heavier mesons. The explanation is that the contact terms in the $`\pi +\rho `$ NN-interaction are not the only ones to be strongly modified by the short range correlations. In particular, because of the large $`\rho `$-meson mass, also the Yukawa terms generated in (LABEL:9) should be strongly reduced. We conclude hence that the implementation of, both realistic short range correlations, and finite-nucleon-size effects, in the context of the relativistic RPA, is required. Presently, we are working on this issue.
Finally, let us mention that the tensor $`\rho `$-nucleon coupling plays an important role in the transverse spin response, and that some progress in assessing this through a relativistic many-body calculation has been made quite recently by Yoshida and Toki .
The authors wish to thank Peter Ring for the use of his spherical RMFT code, and one of us (F.K.) also thanks him for discussions and warm hospitality at the Technischen Universität München. The work of C.D.C. was supported by CAPES (Brazil) and FAPESP (S. Paulo, Brazil), and that of (F.K.) by CONICET and FONCyT (Argentina) under project N 03-04296. A.P.G. and F.K. acknowledge partial financial support by ICTP (Trieste). |
warning/0003/hep-th0003154.html | ar5iv | text | # 1 Introduction
## 1 Introduction
For building up self-consistent string models with $`N=(4,4)`$ worldsheet supersymmetry (SUSY) it is of primary importance to explore in full the structure of the relevant worldsheet conformal supergravity (SG), both on and off shell, as well as its couplings to $`N=(4,4)`$, $`2D`$ superconformal sigma models. In components and in the standard $`N=(4,4)`$, $`2D`$ superspace these issues were addressed in refs. -. Recently, there was a revival of interest to $`N=(4,4)`$ superconformal $`2D`$ theories caused by the fact that they describe the low-energy limits of some string theory compactifications (see. e.g., -). This makes it urgent to revert to the problem of finding out an adequate superspace description of $`N=(4,4)`$ SG and listing all possible versions of the latter.
Here we present the basics of conformal $`N=(4,4)`$, $`2D`$ SG in the analytic harmonic $`SU(2)\times SU(2)`$ superspace with two independent sets of harmonic variables (for the left and right light-cone sectors). This kind of harmonic superspace is indispensable for the off-shell description of $`N=(4,4)`$ supersymmetric torsionful sigma models, with all supersymmetries being manifest. Our construction in its starting points follows the analogous one for conformal SG in the analytic subspace of $`N=2`$, $`4D`$ harmonic superspace -, but eventually we find a few essential differences from the latter theory. These differences amount to a number of novel features of our construction compared to the existing approaches to $`N=(4,4)`$, $`2D`$ SG.
First, in the $`SU(2)\times SU(2)`$ harmonic superspace three different SG groups containing local $`SU(2)_L\times SU(2)_R`$ symmetry can be defined ($`L`$ and $`R`$ stand for the left- and right-handed $`2D`$ light-cone sectors). Two of them have as the rigid limits two different infinite-dimensional $`N=4`$, $`SU(2)`$ superconformal groups the realization of which in the flat harmonic superspace was given in . A closure of these two rigid superconformal groups is the large $`N=4`$, $`2D`$ superconformal group with the $`SO(4)\times U(1)`$ affine Kac-Moody group as internal symmetry (in each of two $`2D`$ light-cone sectors) -. The most general SG group which can be defined in the analytic $`SU(2)\times SU(2)`$ harmonic superspace yields in the flat limit just this large $`N=(4,4)`$ superconformal group. The corresponding SG can be treated as a “master theory” producing two $`N=(4,4)`$, $`SU(2)`$ SG theories as its proper truncations. Another, more elegant way of getting $`N=(4,4)`$, $`SU(2)`$ SG theories from the master $`N=(4,4)`$ SG is to couple the latter to appropriate harmonic superfield compensators. We explicitly demonstrate how one of $`N=(4,4)`$, $`SU(2)`$ SG groups can be recovered using this compensation procedure. The relevant compensator is one of the $`SU(2)\times SU(2)`$ harmonic superfields defined in (it contains $`(32+32)`$ off-shell components and generalizes the so-called nonlinear supermultiplet ). It should be stressed that the most characteristic feature of the master $`N=(4,4)`$ SG group is the presence of four local $`SU(2)`$ symmetries (corresponding to gauging left and right $`SO(4)`$) and two local $`U(1)`$ symmetries (corresponding to gauging left and right $`U(1)`$). The versions of off-shell conformal $`N=(4,4)`$ SG known until now contained at most two local $`SU(2)`$ symmetries and no local $`U(1)`$ symmetries at all.
One more difference from the $`N=2`$, $`4D`$ case stems from the presence of two independent sets of $`SU(2)`$ harmonic variables in the $`SU(2)\times SU(2)`$ harmonic superspace. This peculiarity gives rise, on the one hand, to the property that the relevant groups of analytic superdiffeomorphisms are more powerful than their $`N=2`$, $`4D`$ counterpart, in the sense that they allow to gauge away more fields from the basic geometric objects of the theory, analytic vielbeins which covariantize two analyticity-preserving harmonic derivatives. On the other hand, prior to any gauge fixing, we are led to impose the constraints on the analytic vielbeins reflecting the commutativity of two independent analyticity-preserving harmonic derivatives in the flat case. The constraints and the original SG gauge group together work in such a way that in the WZ gauge we are left with no auxiliary fields at all. Besides, the number of gauge fields coincides with that of the remaining independent gauge parameters in the left and right sectors. Thus, the analytic vielbeins in the considered case actually describe a sum of two pure gauge Weyl multiplets. This sum can be naturally called the $`N=(4,4)`$ Beltrami-Weyl (BW) multiplet (the $`SU(2)`$ or $`SO(4)\times U(1)`$ one, depending on from which superdiffeomorphism group one proceeds). For the $`N=(4,4)`$, $`SU(2)`$ case our results agree with those of Schoutens , who constructed the corresponding SG in the component approach by directly gauging the product of left and right $`N=4`$, $`SU(2)`$ superconformal groups. The standard conformal $`N=(4,4)`$ SG group corresponds to gauging the maximal finite-dimensional subgroup $`SU(1,1|2)\times SU(1,1|2)`$ of this product , - and also gives rise to Weyl multiplet containing no off-shell degrees of freedom. A novel point is that this phenomenon of the one-to-one correspondence between the gauge fields and residual gauge invariances is continued as well to the more general case of $`N=(4,4)`$, $`SO(4)\times U(1)`$ SG group. The supermultiplet of what is usually referred to as “the minimal off-shell Poincaré $`N=4`$ SG” arises after coupling $`N=4`$, $`SU(2)`$ BW multiplet to a compensating superfield which represents one of twisted chiral multiplets in the analytic harmonic $`SU(2)\times SU(2)`$ superspace. Thus the minimal $`N=(4,4)`$ SG representation corresponds to the two successive compensations: firstly, the $`N=(4,4)`$, $`SO(4)\times U(1)`$ SG group is compensated down to its $`N=(4,4)`$, $`SU(2)`$ subgroup by using some special harmonic compensator and, secondly, this subgroup is further compensated down to the group corresponding to the minimal representation via coupling to a twisted $`N=(4,4)`$ supermultiplet. The existence of a dual formulation of the twisted multiplet with an infinite number of auxiliary fields implies the existence of new off-shell version of $`N=(4,4)`$ Poincaré SG with an infinite number of auxiliary fields.
In the present paper we do not aim to present the whole formalism of $`N=(4,4)`$ SG in harmonic superspace. We concentrate on describing the analytic superspace geometry of the $`SU(2)`$ and $`SO(4)\times U(1)`$, $`N=(4,4)`$ BW supermultiplets: define the relevant groups, the analyticity-preserving harmonic derivatives and the covariant constraints on the latter, and show that after choosing appropriate WZ gauges and solving the constraints we are left with the needed irreducible field contents. We present the invariant couplings of $`SO(4)\times U(1)`$, $`N=(4,4)`$ BW multiplet to the compensating $`N=(4,4)`$ multiplets, such that the residual gauge freedom is just one of the $`N=(4,4)`$, $`SU(2)`$ SG groups. Then we extend this coupling to include an arbitrary number of self-interacting twisted multiplets. We also show, at the linearized level, how to extract another $`N=(4,4)`$, $`SU(2)`$ SG group from the $`N=(4,4)`$, $`SO(4)\times U(1)`$ one. We discuss various truncations and schemes of compensation of $`N=(4,4)`$, $`SO(4)\times U(1)`$ SG down to its $`N=(4,4)`$, $`SU(2)`$ superconformal descendants and, further, to different versions of Poincaré SG. A few novel possibilities are found. More detailed considerations with passing to component actions, etc, will be given elsewhere.
## 2 Flat $`SU(2)\times SU(2)`$ analytic harmonic superspace
To proceed, we need some facts about the flat analytic harmonic $`SU(2)\times SU(2)`$ superspace. In our notation we will basically follow ref. with minor deviations.
This superspace is spanned by the following set of coordinates
$$𝐀^{(1+2,1+2|2,2)}=(z^{++},z^{},\theta ^{(1,0)\underset{¯}{k}+},\theta ^{(0,1)\underset{¯}{b}},u_i^{(\pm 1,0)},v_a^{(0,\pm 1)})(\zeta ^\mu ,u_i^{(\pm 1,0)},v_a^{(0,\pm 1)}).$$
(2.1)
Here, the $`+,`$ indices of the $`z`$ and $`\theta `$ coordinates are the left and right light-cone $`SO(1,1)`$ ones, while $`i,\underset{¯}{k},a,\underset{¯}{b}`$ are doublet indices of four commuting $`SU(2)`$ groups which constitute the full automorphism group $`SO(4)_L\times SO(4)_R`$ of $`N=(4,4)`$, $`2D`$ Poincaré superalgebra. In what follows we will omit the light-cone indices of Grassmann coordinates, keeping in mind that the indices $`\underset{¯}{i}`$ and $`\underset{¯}{a}`$ are always accompanied by the indices $`+`$ and $``$. The harmonic part of $`𝐀^{(1+2,1+2|2,2)}`$ is parametrized by two independent sets of harmonic variables $`u_i^{(\pm 1,0)},v_a^{(0,\pm 1)}`$, each associated with one of the $`SU(2)`$ factors of $`SO(4)_L`$ and $`SO(4)_R`$, respectively (we denote these “harmonized” $`SU(2)`$ groups as $`SU(2)_L`$ and $`SU(2)_R`$):
$$u^{(1,0)i}u_i^{(1,0)}=1,v^{(0,1)a}v_a^{(0,1)}=1.$$
(2.2)
The harmonics $`u`$ and $`v`$, as well as the left and right odd coordinates, carry two independent $`U(1)`$ charges “$`(n,0)`$”, “$`(0,m)`$” which are assumed to be strictly conserved (like in the $`N=2`$, $`4D`$ harmonic superspace approach ). This requirement restricts $`u`$ and $`v`$ to parametrize 2-spheres $`SU(2)_L/U(1)_L`$ and $`SU(2)_R/U(1)_R`$. The superfields given on $`𝐀^{(1+2,1+2|2,2)}`$ (analytic $`N=(4,4)`$ superfields), $`\mathrm{\Phi }^{(p,q)}(\zeta ,u,v)`$, are also labelled by a pair of such $`U(1)`$ charges “$`(p,q)`$” and are assumed to admit expansions in the double harmonic series on the above 2-spheres. It should be stressed that the “harmonized” subgroups $`SU(2)_L,SU(2)_R`$ and the two remaining $`SU(2)`$ factors of $`SO(4)_L,SO(4)_R`$ are realized in essentially different ways. Namely, the “harmonized” $`SU(2)`$ symmetries are hidden, in the sense that they manifest themselves only in the existence of the double harmonic series; on the other hand, two extra $`SU(2)`$ symmetries are explicit, as they rotate the underlined doublet indices of the analytic Grassmann coordinates and the related indices of component fields in the $`\theta `$ expansion of $`\mathrm{\Phi }^{(p,q)}`$. Note that the latter in general can carry indices of any linear representation of these explicit $`SU(2)`$ symmetries.
The analytic superspace (2.1) is real with respect to the generalized involution “$``$” which is the product of ordinary complex conjugation and an antipodal map of the 2-spheres $`SU(2)_L/U(1)_L`$ and $`SU(2)_R/U(1)_R`$
$$\stackrel{~}{(\theta ^{(1,0)\underset{¯}{i}})}=\theta _{\underset{¯}{i}}^{(1,0)},\stackrel{~}{(u^{(\pm 1,0)i})}=u_i^{(\pm 1,0)},$$
(2.3)
(and similarly for $`\theta ^{(0,1)\underset{¯}{a}},v_a^{(0,\pm 1)}`$). The analytic superfields $`\mathrm{\Phi }^{(p,q)}`$ can be chosen real with respect to this involution, provided $`|p+q|=2n`$
$$\stackrel{~}{(\mathrm{\Psi }^{(p,q)})}=\mathrm{\Psi }^{(p,q)},|p+q|=2n.$$
(2.4)
In what follows we will need the fact of the existence of the mutually commuting sets of harmonic derivatives $`D^{(2,0)}`$, $`D_u^{(0,0)}D_u^0`$ and $`D^{(0,2)}`$, $`D_v^{(0,0)}D_v^0`$ which preserve $`N=(4,4)`$ Grassmann harmonic analyticity, i.e. yield an analytic superfield when acting on some analytic superfield. They are given by the expressions
$`D^{(2,0)}`$ $`=`$ $`^{(2,0)}+i(\theta ^{(1,0)})^2_{++},D^{(0,2)}=^{(0,2)}+i(\theta ^{(0,1)})^2_{}`$ (2.5)
$`D_u^0`$ $`=`$ $`_u^0+\theta ^{(1,0)\underset{¯}{i}}{\displaystyle \frac{}{\theta ^{(1,0)\underset{¯}{i}}}},D_v^0=_v^0+\theta ^{(0,1)\underset{¯}{a}}{\displaystyle \frac{}{\theta ^{(0,1)\underset{¯}{a}}}},`$ (2.6)
$`[D_u^0,D^{(2,0)}]`$ $`=`$ $`2D^{(2,0)},[D_v^0,D^{(0,2)}]=\mathrm{\hspace{0.33em}2}D^{(0,2)}.`$ (2.7)
Here $`_{\pm \pm }=/z^{\pm \pm }`$ and
$$^{(2,0)}=u^{(1,0)i}\frac{}{u^{(1,0)i}},_u^0=u^{(1,0)i}\frac{}{u^{(1,0)i}}u^{(1,0)i}\frac{}{u^{(1,0)i}},$$
(2.8)
(the same formulas are valid for $`^{(0,2)}`$ and $`_v^0`$ with the change $`uv`$). The operators $`D_u^0`$, $`D_v^0`$ count the $`U(1)`$ charges of the analytic superfields
$$D_u^0\mathrm{\Phi }^{(p,q)}(\zeta ,u,v)=p\mathrm{\Phi }^{(p,q)}(\zeta ,u,v),D_v^0\mathrm{\Phi }^{(p,q)}(\zeta ,u,v)=q\mathrm{\Phi }^{(p,q)}(\zeta ,u,v).$$
(2.9)
In the analytic superspace (2.1) one can realize two different infinite-dimensional groups of superconformal transformations. Each group consists of two commuting light-cone branches, the left and right ones, having as the algebra the classical $`N=4`$, $`SU(2)`$ superconformal algebra (SCA) . Without loss of generality we can specialize, e.g., to the left sector. It turns out that the form of the relevant superconformal transformations is basically specified by the transformation law of the analyticity-preserving covariant harmonic derivative $`D^{(2,0)}`$ (or $`D^{(0,2)}`$ in the right sector).
The basic distinguishing feature of the first group is that it does not touch the harmonics
$$\delta _Iu_i^{(\pm 1,0)}=0.$$
(2.10)
Its realization in $`𝐀^{(1+2,1+2|2,2)}`$ is completely fixed by the requirement that $`D^{(2,0)}`$ is invariant
$$\delta _ID^{(2,0)}=0.$$
(2.11)
The second superconformal group has the same Lie bracket structure as the first one, but it acts on all the left coordinates of $`𝐀^{(1+2,1+2|2,2)}`$, including the harmonic ones $`u^{(\pm 1,0)}`$. We give here only the generic form of transformations of harmonics and the derivative $`D^{(2,0)}`$
$`\delta _{II}u_i^{(1,0)}`$ $`=`$ $`\mathrm{\Lambda }_I^{(2,0)}(z^{++},\theta ^{(1,0)},u)u_i^{(1,0)},\delta _{II}u_i^{(1,0)}=\mathrm{\hspace{0.33em}0}`$
$`\delta _{II}D^{(2,0)}`$ $`=`$ $`\mathrm{\Lambda }^{(2,0)}D_u^0,\mathrm{\Lambda }^{(2,0)}=D^{(2,0)}\mathrm{\Lambda }_L,D^{(2,0)}\mathrm{\Lambda }^{(2,0)}=\mathrm{\hspace{0.33em}0}.`$ (2.12)
The main difference between these two $`N=4`$, $`SU(2)`$ superconformal groups lies in the realization of their affine $`SU(2)`$ subgroups: the second one acts on the indices $`i,j`$ and affects both the Grassmann and harmonic coordinates, while the first one acts only on the underlined indices and so affects only $`\theta `$’s. These groups do not commute; their closure is the “large” $`N=4`$, $`SO(4)\times U(1)`$ group . For our further purposes it will be important that the latter involves an extra $`U(1)`$ affine (Kac-Moody) symmetry with the dimensionless holomorphic parameter $`\lambda _L(z^{++})`$ (or $`\lambda _R(z^{})`$ in the right sector). It is realized, e.g., on $`u^{(1,0)i}`$ as
$$\delta _{U(1)}u^{(1,0)i}=(D^{(2,0)}\lambda _L(z))u^{(1,0)i}=i(\theta ^{(0,1)})^2_{++}\lambda _L(z)u^{(1,0)i}.$$
(2.13)
The “large” superconformal algebra corresponds to the most general solution of the constraints on $`\mathrm{\Lambda }^{(2,0)}`$ in eq. (2.12), while two of its $`SU(2)`$ subalgebras (SCA-I and SCA-II in what follows) are singled out by some additional conditions. Here we will not present the explicit form of the coordinate transformations of all these superconformal groups (see for details), since we will recover them as flat limits of the appropriate SG groups in the next Sections. Notice the following important property: both $`N=4`$, $`SU(2)`$ superconformal groups, and hence their closure, leave invariant the analytic superspace integration measure $`\mu ^{(2,2)}=d^2zd^2\theta ^{(1,0)}d^2\theta ^{(0,1)}[du][dv]`$:
$$\delta _I\mu ^{(2,2)}=\delta _{II}\mu ^{(2,2)}=0.$$
(2.14)
The last topic of this introductory Section is the harmonic superspace description of some important $`N=(4,4)`$ multiplets. We start with one of the possible $`N=(4,4)`$ twisted chiral multiplets , namely, the one having a simple description in $`SU(2)\times SU(2)`$ harmonic analytic superspace. It is represented by a real analytic $`(4,4)`$ superfield $`q^{(1,1)}(\zeta ,u,v)`$ subject to the constraints
$$D^{(2,0)}q^{(1,1)}=D^{(0,2)}q^{(1,1)}=0.$$
(2.15)
They leave in $`q^{(1,1)}`$ $`8+8`$ independent components , just the off-shell field content of $`N=(4,4)`$ twisted multiplet. The superfield $`q^{(1,1)}`$ is scalar with respect to the first $`N=4`$, $`SU(2)`$ superconformal group but it is transformed with the weight 1 under the second one (this is necessary for preserving the constraints (2.15))
$$\delta _Iq^{(1,1)}=0,\delta _{II}q^{(1,1)}=\mathrm{\Lambda }_Lq^{(1,1)}$$
(2.16)
(the transformations from the right-handed branches are similar). The physical dimension components of $`q^{(1,1)}`$ (four dimension 0 bosons and eight dimension 1/2 fermions) behave in different ways under these two kinds of $`N=(4,4)`$, $`SU(2)`$ transformations. In particular, the $`SU(2)`$ affine transformations from the first superconformal group act only on fermions, while those from the second group act both on bosons and fermions. The physical bosonic fields are naturally combined, with respect to the latter transformations and their right-handed counterparts, into a 2$`\times `$2 matrix $`q^{ia}(z^{++},z^{})`$ on which the left (right) conformal $`SU(2)`$ acts as a left (right) multiplication . So the purely $`SU(2)`$ part of $`q^{ia}`$ represents the coset $`SU(2)_L\times SU(2)_R/SU(2)_{diag}`$, and it is not too surprising that the $`q^{(1,1)}`$ action invariant under the second superconformal group is none other than $`N=(4,4)`$ extension of the $`SU(2)`$ WZW sigma model action. Indeed, it is just the $`N=4`$, $`SU(2)\times U(1)`$ WZW sigma model action of ref. . The $`SU(2)\times SU(2)`$ analytic superspace form of this action reads
$$S_{wzw}=\frac{1}{4\gamma ^2}\mu ^{(2,2)}\widehat{q}^{(1,1)}\widehat{q}^{(1,1)}\left(\frac{1}{(1+X)X}\frac{\text{ln}(1+X)}{X^2}\right),$$
(2.17)
where
$$\widehat{q}^{(1,1)}=q^{(1,1)}c^{(1,1)},X=c^{(1,1)}\widehat{q}^{(1,1)},c^{(\pm 1,\pm 1)}=c^{ia}u_i^{(\pm 1,0)}v_a^{(0,\pm 1)},c^{ia}c_{ia}=2,$$
(2.18)
and $`\gamma `$ is a dimensionless sigma model coupling constant. Despite the presence of an extra quartet constant $`c^{ia}`$ in the analytic superfield Lagrangian, the action actually does not depend on $`c^{ia}`$ .
We wish to stress that the action (2.17) is unique (up to adding full harmonic derivatives) in the sense that it is the only possible action of a single superfield $`q^{(1,1)}`$ invariant under the second $`N=(4,4),SU(2)`$ superconformal group. As we will see later, in the curved case the superfield $`q^{(1,1)}`$ serves as a compensator which breaks the appropriate $`N=(4,4)`$, $`SU(2)`$ SG group (having as the rigid limit the second $`N=(4,4)`$, $`SU(2)`$ superconformal group) down to the supergroup of minimal $`N=(4,4)`$, $`2D`$ SG .
As for the first superconformal group, an arbitrary action of the superfield $`q^{(1,1)}`$,
$$S_q=\mu ^{(2,2)}^{(2,2)}(q^{(1,1)M}(\zeta ,u,v),u,v),$$
(2.19)
is invariant with respect to it. As a consequence of this property, the particular $`q^{(1,1)}`$ action (2.17) is invariant under both $`N=(4,4)`$, $`SU(2)`$ superconformal groups and, hence, under their closure, i.e. the “large” $`N=(4,4)`$, $`SO(4)\times U(1)`$ superconformal group. Note that $`q^{(1,1)}`$ transforms under the left affine $`U(1)`$ transformations (2.13) as
$$\delta _{U(1)}q^{(1,1)}=\lambda _L(z^{++})q^{(1,1)}$$
(2.20)
(and analogously under their right-handed counterparts). The full transformation law of $`q^{(1,1)}`$ under the left “large” group looks like the second law in eq. (2.16), with $`\mathrm{\Lambda }_L=\lambda _L(z^{++})+\lambda ^{(ik)}(z^{++})u_i^{(1,0)}u_k^{(1,0)}+\mathrm{}`$. Further details will be given in Sect. 4. It is worth mentioning that the general action (2.19) always yelds the sigma model with torsion in the sector of physical bosons, just of the same kind as in the $`N=(4,4)`$ supersymmetric subclass of general $`N=(2,2)`$ chiral and twisted chiral superfield sigma models explored in . The actions of other matter multiplets in $`SU(2)\times SU(2)`$ harmonic superspace reveal the same characteristic feature. This is the radical difference of the considered case from the dimensionally-reduced off-shell sigma model actions of hypermultiplets in the standard harmonic superspace with one set of the $`SU(2)`$ harmonic variables : for physical bosons they yield the torsionless hyper-Kähler sigma model actions.
Note that there exist other types of twisted $`N=4`$ multiplets, with the same number of off-shell components, but with different realizations of various $`SU(2)`$ factors of the full $`SO(4)_L\times SO(4)_R`$ automorphism group of rigid $`N=(4,4),\mathrm{\hspace{0.33em}\hspace{0.33em}2}D`$ SUSY . Respectively, the above two $`N=(4,4)`$, $`SU(2)`$ superconformal groups are realized in different ways on these multiplets. In particular, there exists a sort of twisted multiplet on which the first and second superconformal groups act in the way just opposite to their action on $`q^{(1,1)}`$. <sup>1</sup><sup>1</sup>1In such a multiplet is called TM-I as opposed to $`q^{(1,1)}`$ which is TM-II in this classification. Such a classification makes sense with respect to a fixed $`N=(4,4),SU(2)`$ SCA: if the affine $`SU(2)_L\times SU(2)_R`$ subgroup acts both on the physical bosons and fermions, one deals with TM-II, whereas if it acts only on fermions, one faces TM-I. Conversely, $`q^{(1,1)}`$ is TM-I with respect to the first of the two $`N=(4,4),SU(2)`$ SCAs defined above, but it is TM-II with respect to the second one.
The $`SU(2)\times SU(2)`$ harmonic superspace description of these complementary twisted multiplets is somewhat more complicated. Nevertheless, all of them can be coupled to the $`N=(4,4)`$ Beltrami-Weyl SG multiplets to be defined below and so can serve as compensators. We are planning to present these couplings in a future work.
Finally, we mention one more analytic $`SU(2)\times SU(2)`$ harmonic supermultiplet which will be used in Sect. 5 as a compensator reducing the $`N=(4,4)`$, $`SO(4)\times U(1)`$ SG group to one of its $`N=(4,4)`$, $`SU(2)`$ subgroups. It is represented by a pair of analytic superfields $`N^{(2,0)}`$, $`N^{(0,2)}`$ satisfying the constraints
$`D^{(2,0)}N^{(2,0)}+N^{(2,0)}N^{(2,0)}`$ $`=`$ $`0,D^{(0,2)}N^{(0,2)}+N^{(0,2)}N^{(0,2)}=0,`$
$`D^{(2,0)}N^{(0,2)}D^{(0,2)}N^{(2,0)}`$ $`=`$ $`0.`$ (2.21)
These constraints are analogous to those defining the so-called nonlinear supermultiplet in the $`N=2`$, $`4D`$ harmonic superspace (the latter goes into $`N=(4,4)`$, $`SU(2)_{diag}`$ harmonic superspace upon reduction to $`2D`$). They are obviously covariant under the first $`N=(4,4)`$, $`SU(2)`$ superconformal group, if $`N^{(2,0)},N^{(0,2)}`$ are assumed to transform as scalars with respect to it. They are also covariant under the second group, provided $`N^{(2,0)},N^{(0,2)}`$ transform according to
$$\delta _{II}N^{(2,0)}=\mathrm{\Lambda }^{(2,0)},\delta _{II}N^{(0,2)}=\mathrm{\Lambda }^{(0,2)}.$$
(2.22)
The simplest invariant action (with the correct sign of the kinetic terms of the physical fields) is as follows:
$$S_N\mu ^{(2,2)}N^{(2,0)}N^{(0,2)}.$$
(2.23)
To see that it is invariant (up to surface terms) under (2.22), one should take into account the invariance of the analytic superspace integration measure and the properties
$$\mathrm{\Lambda }^{(2,0)}=D^{(2,0)}\mathrm{\Lambda }_L,\mathrm{\Lambda }^{(0,2)}=D^{(0,2)}\mathrm{\Lambda }_R,D^{(2,0)}\mathrm{\Lambda }_R=D^{(0,2)}\mathrm{\Lambda }_L=0.$$
(2.24)
The pair $`N^{(2,0)},N^{(0,2)}`$ describes $`32+32`$ off-shell degrees of freedom and is dual-equivalent to four $`q^{(1,1)}`$ superfields .
Having the multiplet $`N^{(2,0)},N^{(0,2)}`$, one can define further consistent non-linear multiplets $`G^{(2,0)},G^{(0,2)}`$ which are zero-weight scalars under both $`N=(4,4),SU(2)`$ superconformal groups
$$\delta _{I,II}G^{(2,0)}=\delta _{I,II}G^{(0,2)}=0.$$
(2.25)
The corresponding constraints (covariant with respect to both superconformal groups) are a slight modification of (2.21)
$`(D^{(2,0)}+2N^{(2,0)})G^{(2,0)}+\alpha G^{(2,0)}G^{(2,0)}`$ $`=`$ $`0,`$
$`(D^{(0,2)}+2N^{(0,2)})G^{(0,2)}+\alpha G^{(0,2)}G^{(0,2)}`$ $`=`$ $`0,`$
$`D^{(2,0)}G^{(0,2)}D^{(0,2)}G^{(2,0)}`$ $`=`$ $`0,`$ (2.26)
where $`\alpha `$ is an arbitrary dimensionless parameter (it can be equal to zero). All such representations comprise $`32+32`$ off-shell degrees of freedom. Their Lagrangians are bilinears like in (2.23).
## 3 Curved $`SU(2)\times SU(2)`$ analytic superspace and N=(4,4) Beltrami-Weyl multiplet
By analogy with the $`N=2`$, $`4D`$ case we assume that the fundamental group of $`N=(4,4)`$, $`2D`$ conformal supergravity is represented by the following diffeomorphisms of the analytic harmonic $`SU(2)\times SU(2)`$ superspace
$`\delta \zeta ^\mu =\mathrm{\Lambda }^\mu (\zeta ,u,v),\delta u_i^{(1,0)}=\mathrm{\Lambda }^{(2,0)}(\zeta ,u,v)u_i^{(1,0)},\delta v_a^{(0,1)}=\mathrm{\Lambda }^{(0,2)}(\zeta ,u,v)v_a^{(0,1)},`$
$`\delta u_i^{(1,0)}=\delta v_a^{(0,1)}=0.`$ (3.1)
Here $`\zeta ^\mu =(z^{++},z^{},\theta ^{(1,0)\underset{¯}{k}+},\theta ^{(0,1)\underset{¯}{b}})`$ as in (2.1) and the gauge parameters $`\mathrm{\Lambda }^\mu `$, $`\mathrm{\Lambda }^{(2,0)}`$, $`\mathrm{\Lambda }^{(0,2)}`$ are arbitrary functions over the whole harmonic analytic superspace $`𝐀^{(1+2,1+2|2,2)}`$. These transformation laws preserve the defining relations of harmonic variables (2.2) and the reality of $`𝐀^{(1+2,1+2|2,2)}`$ with respect to the “$``$” conjugation. The analyticity-preserving harmonic derivatives $`D^{(2,0)}`$ and $`D^{(0,2)}`$ are covariantized by introducing appropriate analytic vielbeins
$`D^{(2,0)}^{(2,0)}`$ $`=`$ $`D^{(2,0)}+H^{(2,0)\mu }_\mu +H^{(4,0)}^{(2,0)}+H^{(2,2)}^{(0,2)}`$
$``$ $`D^{(2,0)}+H^{(2,0)M}_M,`$
$`D^{(0,2)}^{(0,2)}`$ $`=`$ $`D^{(0,2)}+H^{(0,2)\mu }_\mu +\stackrel{~}{H}^{(2,2)}^{(2,0)}+H^{(0,4)}^{(0,2)}`$ (3.2)
$``$ $`D^{(0,2)}+H^{(0,2)M}_M,`$
where we used the notation
$`M=(\mu ,(2,0),(0,2)),_M=(_\mu ,^{(2,0)},^{(0,2)}),`$
$`^{(2,0)}=u^{(1,0)i}{\displaystyle \frac{}{u^{(1,0)i}}},^{(0,2)}=v^{(0,1)a}{\displaystyle \frac{}{v^{(0,1)a}}}`$ (3.3)
and separated the flat parts of the vielbein components in front of $`_{++}`$ in $`^{(2,0)}`$ and $`_{}`$ in $`^{(0,2)}`$. In eqs. (3.2) all the vielbeins are analytic $`N=(4,4)`$, $`2D`$ superfields,
$$H^{(2,0)M}=H^{(2,0)M}(\zeta ,u,v),H^{(0,2)M}=H^{(0,2)M}(\zeta ,u,v).$$
The flat limit is achieved by putting them equal to zero. The $`U(1)`$ charge-counting operators $`D_u^0`$ and $`D_v^0`$ retain their flat form (2.6).
Again in analogy with refs. , we postulate for $`^{(2,0)}`$, $`^{(0,2)}`$ the following transformation law under the $`N=(4,4)`$ SG group (3.1)
$$\delta ^{(2,0)}=\mathrm{\Lambda }^{(2,0)}D_u^0,\delta ^{(0,2)}=\mathrm{\Lambda }^{(0,2)}D_v^0,$$
(3.4)
whence
$`\delta H^{(2,0)++}`$ $`=`$ $`^{(2,0)}\mathrm{\Lambda }^{++}2i\mathrm{\Lambda }^{(1,0)}\theta ^{(1,0)},\delta H^{(2,0)}=^{(2,0)}\mathrm{\Lambda }^{},`$
$`\delta H^{(3,0)\underset{¯}{i}}`$ $`=`$ $`^{(2,0)}\mathrm{\Lambda }^{(1,0)\underset{¯}{i}}\mathrm{\Lambda }^{(2,0)}\theta ^{(1,0)\underset{¯}{i}},\delta H^{(2,1)\underset{¯}{a}}=^{(2,0)}\mathrm{\Lambda }^{(0,1)\underset{¯}{a}},`$
$`\delta H^{(4,0)}`$ $`=`$ $`^{(2,0)}\mathrm{\Lambda }^{(2,0)},\delta H^{(2,2)}=^{(2,0)}\mathrm{\Lambda }^{(0,2)},`$ (3.5)
$`\delta H^{(0,2)++}`$ $`=`$ $`^{(0,2)}\mathrm{\Lambda }^{++},\delta H^{(0,2)}=^{(0,2)}\mathrm{\Lambda }^{}2i\mathrm{\Lambda }^{(0,1)}\theta ^{(0,1)},`$
$`\delta H^{(1,2)\underset{¯}{i}}`$ $`=`$ $`^{(0,2)}\mathrm{\Lambda }^{(1,0)\underset{¯}{i}},\delta H^{(0,3)\underset{¯}{a}}=^{(0,2)}\mathrm{\Lambda }^{(0,1)\underset{¯}{a}}\mathrm{\Lambda }^{(0,2)}\theta ^{(0,1)\underset{¯}{a}},`$
$`\delta \stackrel{~}{H}^{(2,2)}`$ $`=`$ $`^{(0,2)}\mathrm{\Lambda }^{(2,0)},\delta H^{(0,4)}=^{(0,2)}\mathrm{\Lambda }^{(0,2)}.`$ (3.6)
From now on, the similarity with the $`N=2`$, $`4D`$ construction ceases to be literal and the specificity of the $`N=(4,4)`$ case comes into play.
First of all, we wish to generalize the notion of the twisted analytic superfield $`q^{(1,1)}`$ to the curved case and hence need to find a correct generalization of the defining constraints (2.15) and the superconformal transformation laws (2.16). As we have started with the most general diffeomorphism group of the analytic superspace, we expect it to yield, in the flat limit, the product of the left and right “large” $`SO(4)\times U(1)`$ superconformal groups, including their $`U(1)`$ affine subgroups with the parameters $`\lambda _L(z^{++})`$, $`\lambda _R(z^{})`$. However, a close inspection of the analytic superfield gauge parameters $`\mathrm{\Lambda }^\mu (\zeta ,u,v)`$, $`\mathrm{\Lambda }^{(2,0)}(\zeta ,u,v)`$ and $`\mathrm{\Lambda }^{(0,2)}(\zeta ,u,v)`$ shows that there is no place in them for such dimensionless parameters (these can appear only with their $`z`$ derivatives). To generalize the transformation laws of $`q^{(1,1)}`$ (2.16), (2.20) to the curved case, we are led to introduce two extra independent analytic gauge functions
$$\mathrm{\Lambda }_L(\zeta ,u,v)=\lambda _L(z^{++},z^{})+\mathrm{},\mathrm{\Lambda }_R(\zeta ,u,v)=\lambda _R(z^{++},z^{})+\mathrm{}$$
and to ascribe the following transformation laws to $`q^{(1,1)}`$
$$\delta q^{(1,1)}=(\mathrm{\Lambda }_L+\mathrm{\Lambda }_R)q^{(1,1)}.$$
(3.7)
We call these transformations the “$`U(1)`$ weight” ones, to distinguish them from the harmonic $`U(1)`$ phase transformations. We normalize the left and right $`U(1)`$ weights $`J_L`$ and $`J_R`$ as
$$J_Lq^{(1,1)}=J_Rq^{(1,1)}=q^{(1,1)}.$$
(3.8)
At this stage, the $`U(1)`$ weight analytic parameters $`\mathrm{\Lambda }_L`$, $`\mathrm{\Lambda }_R`$ are entirely unrelated to those of the coordinate transformations.
Such a relation naturally comes out, as a result of choosing the appropriate transformation law for the $`U(1)`$ weight-covariantized harmonic derivatives and fixing a proper gauge.
We covariantize $`^{(2,0)}`$, $`^{(0,2)}`$ by introducing four analytic superfield $`U(1)`$ connections $`H_L^{(2,0)}(\zeta ,u,v)`$, $`H_R^{(2,0)}(\zeta ,u,v)`$, $`H_L^{(0,2)}(\zeta ,u,v)`$, $`H_R^{(0,2)}(\zeta ,u,v)`$
$`^{(2,0)}`$ $``$ $`𝒟^{(2,0)}=^{(2,0)}+H_L^{(2,0)}J_L+H_R^{(2,0)}J_R`$
$`^{(0,2)}`$ $``$ $`𝒟^{(0,2)}=^{(0,2)}+H_L^{(0,2)}J_L+H_R^{(0,2)}J_R,`$ (3.9)
and postulate the following transformation laws for $`𝒟^{(2,0)}`$, $`𝒟^{(0,2)}`$
$`\delta 𝒟^{(2,0)}`$ $`=`$ $`\mathrm{\Lambda }^{(2,0)}(D_u^0J_L)^{(2,0)}\mathrm{\Lambda }_LJ_L^{(2,0)}\mathrm{\Lambda }_RJ_R,`$
$`\delta 𝒟^{(0,2)}`$ $`=`$ $`\mathrm{\Lambda }^{(0,2)}(D_v^0J_R)^{(0,2)}\mathrm{\Lambda }_LJ_L^{(0,2)}\mathrm{\Lambda }_RJ_R.`$ (3.10)
The transformation laws of the vielbeins in $`^{(2,0)}`$, $`^{(0,2)}`$ do not change, while the newly introduced $`U(1)`$ connections are transformed as
$`\delta H_L^{(2,0)}`$ $`=`$ $`\mathrm{\Lambda }^{(2,0)}^{(2,0)}\mathrm{\Lambda }_L,\delta H_R^{(2,0)}=^{(2,0)}\mathrm{\Lambda }_R,`$
$`\delta H_L^{(0,2)}`$ $`=`$ $`^{(0,2)}\mathrm{\Lambda }_L,\delta H_R^{(0,2)}=\mathrm{\Lambda }^{(0,2)}^{(0,2)}\mathrm{\Lambda }_R.`$ (3.11)
The $`𝒟^{(2,0)}`$ and $`𝒟^{(0,2)}`$ derivatives of the analytic superfield $`\mathrm{\Phi }^{(p,q)}`$, with the left and right $`U(1)`$ weights equal to $`l`$ and $`r`$, are transformed as follows:
$`\delta 𝒟^{(2,0)}\mathrm{\Phi }^{(p,q)}`$ $`=`$ $`\mathrm{\Lambda }^{(2,0)}(pl)\mathrm{\Phi }^{(p,q)}+(l\mathrm{\Lambda }_L+r\mathrm{\Lambda }_R)𝒟^{(2,0)}\mathrm{\Phi }^{(p,q)},`$
$`\delta 𝒟^{(0,2)}\mathrm{\Phi }^{(p,q)}`$ $`=`$ $`\mathrm{\Lambda }^{(0,2)}(qr)\mathrm{\Phi }^{(p,q)}+(l\mathrm{\Lambda }_L+r\mathrm{\Lambda }_R)𝒟^{(0,2)}\mathrm{\Phi }^{(p,q)}.`$ (3.12)
We see that only provided $`p=l`$, $`q=r`$, these derivatives are actually covariant, i.e. they transform as the superfield $`\mathrm{\Phi }^{(p,q)}`$ itself. But this is precisely what happens for $`q^{(1,1)}`$, which possesses $`J_L=J_R=1`$. Therefore, as the appropriate curved generalization of the constraints (2.15), we choose the following ones:
$`𝒟^{(2,0)}q^{(1,1)}`$ $`=`$ $`(^{(2,0)}+H_L^{(2,0)}+H_R^{(2,0)})q^{(1,1)}=\mathrm{\hspace{0.33em}0},`$
$`𝒟^{(0,2)}q^{(1,1)}`$ $`=`$ $`(^{(0,2)}+H_L^{(0,2)}+H_R^{(0,2)})q^{(1,1)}=\mathrm{\hspace{0.33em}0}.`$ (3.13)
Before going further, let us adduce some reasoning in favor of the choice of the transformation laws of $`𝒟^{(2,0)}`$, $`𝒟^{(0,2)}`$ in the form (3.10). The primary reason for this choice is the desire to relate the coordinate transformations with the $`U(1)`$ weight transformations, so as to eventually ensure a correct flat limit. Indeed, from eqs. (3.11) it follows that the connections $`H_L^{(2,0)}`$, $`H_R^{(0,2)}`$ can be entirely gauged away, thereby establishing the sought relation
$$H_L^{(2,0)}=H_R^{(0,2)}=0\mathrm{\Lambda }^{(2,0)}=^{(2,0)}\mathrm{\Lambda }_L,\mathrm{\Lambda }^{(0,2)}=^{(0,2)}\mathrm{\Lambda }_R.$$
(3.14)
In what follows we will frequently stick to this gauge. One more argument why we should assume (3.10) is based on an analogy with the harmonic space description of quaternionic manifolds in . There, the analyticity-preserving harmonic derivative in the analytic basis necessarily involves an analytic connection $`\varphi ^{++}`$ associated with the so called “$`Sp(1)`$ weight”. Its transformation law literally mimics that of $`H_L^{(2,0)}`$, $`H_R^{(0,2)}`$, so it is natural to assume that the $`U(1)`$ weights $`J_L`$, $`J_R`$ and the associated analytic superfield parameters $`\mathrm{\Lambda }_L`$ and $`\mathrm{\Lambda }_R`$ are direct analogs of the just mentioned $`Sp(1)`$ weight and the related analytic parameter inherent to the quaternionic manifolds <sup>2</sup><sup>2</sup>2A deep analogy between the description of quaternionic manifolds in the harmonic space and that of conformal $`N=2,\mathrm{\hspace{0.33em}4}D`$ SG in the harmonic superspace was pointed out in .. Of course, the most direct way to justify the transformation law (3.10) would be to deduce it proceeding from the appropriate constraints in the standard $`N=(4,4)`$ superspace. An alternative way is to show that it leads to a self-consistent SG theory, still in the framework of the analytic superspace. This is just what we are going to demonstrate.
An important consequence of the presence of two independent harmonic constraints in the definition of the twisted superfield $`q^{(1,1)}`$, eqs. (3.13), is the integrability condition
$$[𝒟^{(2,0)},𝒟^{(0,2)}]q^{(1,1)}=0.$$
(3.15)
It is easy to see that the direct generalization of the flat condition $`[D^{(2,0)},D^{(0,2)}]=0`$, namely,
$$[𝒟^{(2,0)},𝒟^{(0,2)}]=0,$$
is not covariant under (3.10). The covariant version of this constraint is as follows:
$$[𝒟^{(2,0)},𝒟^{(0,2)}]=H^{(2,2)}(D_v^0J_R)+\stackrel{~}{H}^{(2,2)}(D_u^0J_L).$$
(3.16)
It is evident that eq. (3.15) is automatically satisfied as a consequence of (3.16) and (3.8). This constraint implies
$`^{(2,0)}H_L^{(0,2)}^{(0,2)}H_L^{(2,0)}+\stackrel{~}{H}^{(2,2)}`$ $`=`$ $`0,`$
$`^{(2,0)}H_R^{(0,2)}^{(0,2)}H_R^{(2,0)}H^{(2,2)}`$ $`=`$ $`0`$ (3.17)
and
$$[^{(2,0)},^{(0,2)}]=H^{(2,2)}D_v^0+\stackrel{~}{H}^{(2,2)}D_u^0.$$
(3.18)
From the latter relation one deduces the constraints on the analytic vielbeins
$`^{(2,0)}H^{(0,2)++}^{(0,2)}H^{(2,0)++}2iH^{(1,2)}\theta ^{(1,0)}`$ $`=`$ $`0,`$
$`^{(2,0)}H^{(0,2)}^{(0,2)}H^{(2,0)}+2iH^{(2,1)}\theta ^{(0,1)}`$ $`=`$ $`0,`$
$`^{(2,0)}H^{(1,2)\underset{¯}{i}}^{(0,2)}H^{(3,0)\underset{¯}{i}}\stackrel{~}{H}^{(2,2)}\theta ^{(1,0)\underset{¯}{i}}`$ $`=`$ $`0,`$
$`^{(2,0)}H^{(0,3)\underset{¯}{a}}^{(0,2)}H^{(2,1)\underset{¯}{a}}+H^{(2,2)}\theta ^{(0,1)\underset{¯}{a}}`$ $`=`$ $`0,`$
$`^{(2,0)}H^{(0,4)}^{(0,2)}H^{(2,2)}`$ $`=`$ $`0,`$
$`^{(2,0)}\stackrel{~}{H}^{(2,2)}^{(0,2)}H^{(4,0)}`$ $`=`$ $`0.`$ (3.19)
Thus we see that in the $`N=(4,4),SU(2)\times SU(2)`$ case the analytic vielbeins and $`U(1)`$ connections covariantizing $`D^{(2,0)}`$, $`D^{(0,2)}`$ are necessarily constrained. This is the crucial difference from the formulation of $`N=2`$, $`4D`$ conformal SG in the standard harmonic superspace , where the analogous quantities are unconstrained analytic superfields, i.e. the prepotentials of the theory. Of course, this peculiarity is a direct consequence of the presence of two independent sets of harmonic variables in the considered case.
For the time being, we do not know how to solve (3.17), (3.19) via unconstrained superfield prepotentials. To single out the irreducible field representation carried by vielbeins and $`U(1)`$ connections, we keep to another strategy. Namely, we use the initial gauge freedom to gauge away from these objects as many components as possible, then substitute the resulting expressions into the constraints and solve the latter in this WZ-type gauge. Eventually, it turns out that the solution exists, is unique and is not reduced to a pure gauge. The superfield constraints prove to be purely kinematic: indeed, they do not imply any differential conditions, nor equations of motion, for the remaining fields. At present we are aware of the full nonlinear solution of these constraints. Here, we limit ourselves to the linearized level. This is quite sufficient for revealing the irreducible field contents of the SG theory under consideration.
In the present case, one can choose the WZ gauge in several different ways, the basic criterion for one or another choice being the desire to simplify the constraints (3.17), (3.19) as much as possible. As a first step, we choose the gauge (3.14) and the following additional ones
$`H^{(2,0)++}`$ $`=`$ $`H^{(0,2)}=H^{(3,0)\underset{¯}{i}}=H^{(0,3)\underset{¯}{a}}=\mathrm{\hspace{0.33em}0},`$ (3.20)
$`H^{(4,0)}`$ $`=`$ $`H^{(0,4)}=\mathrm{\hspace{0.33em}0}.`$ (3.21)
These gauges restrict in a certain way the original gauge parameters. At the considered linearized level, (3.20) and (3.21) give rise to the following relations:
$`D^{(2,0)}\mathrm{\Lambda }^{++}2i\mathrm{\Lambda }^{(1,0)}\theta ^{(1,0)}`$ $`=`$ $`0,D^{(2,0)}\mathrm{\Lambda }^{(1,0)\underset{¯}{i}}D^{(2,0)}\mathrm{\Lambda }_L\theta ^{(1,0)\underset{¯}{i}}=\mathrm{\hspace{0.33em}0},`$
$`D^{(0,2)}\mathrm{\Lambda }^{}2i\mathrm{\Lambda }^{(0,1)}\theta ^{(0,1)}`$ $`=`$ $`0,D^{(0,2)}\mathrm{\Lambda }^{(0,1)\underset{¯}{a}}D^{(0,2)}\mathrm{\Lambda }_R\theta ^{(0,1)\underset{¯}{a}}=\mathrm{\hspace{0.33em}0},`$
$`(D^{(2,0)})^2\mathrm{\Lambda }_L=(D^{(0,2)})^2\mathrm{\Lambda }_R`$ $`=`$ $`0,`$ (3.22)
which strictly fix the $`u`$ or $`v`$ dependence of the relevant parameters (depending on which derivative, i.e. either $`D^{(2,0)}`$ or $`D^{(0,2)}`$, enters the given relation). After this, there still remains a freedom associated with the surviving harmonic dependence. This freedom can be used to further gauge away some of the components in the double harmonic expansion of the remaining vielbeins $`H^{(2,0)}`$, $`H^{(0,2)++}`$, $`H^{(2,1)\underset{¯}{a}}`$, $`H^{(1,2)\underset{¯}{i}}`$ and the $`U(1)`$ connections $`H_R^{(2,0)}`$, $`H_L^{(0,2)}`$. At this stage, the $`u`$ and $`v`$ dependence of all analytic superfield gauge parameters is completely fixed and we are left with a finite set of the component parameters. However, in the vielbeins and connections one still finds a non-trivial harmonic dependence which is entirely fixed only after imposing the constraints. The final expressions for the vielbeins, connections and superfield gauge parameters at the linearized level are as follows:
$`H^{(2,0)}`$ $`=`$ $`i(\theta ^{(1,0)})^2\{h_{++}^{}2i\theta _{\underset{¯}{a}}^{(0,1)}h_{++}^{a\underset{¯}{a}}v_a^{(0,1)}i(\theta ^{(0,1)})^2h_{++}^{(ab)}v_a^{(0,1)}v_b^{(0,1)}\},`$
$`H^{(2,1)\underset{¯}{a}}`$ $`=`$ $`i(\theta ^{(1,0)})^2\{h_{++}^{a\underset{¯}{a}}v_a^{(0,1)}+\theta ^{(0,1)\underset{¯}{b}}[h_{++\underset{¯}{b})}^{(\underset{¯}{a}}+{\displaystyle \frac{1}{2}}\delta _{\underset{¯}{b}}^{\underset{¯}{a}}(_{}h_{++}^{}2h_{++}^{(ab)}v_a^{(0,1)}v_b^{(0,1)})]`$
$`+(\theta ^{(0,1)})^2({\displaystyle \frac{1}{2}}t_{++}^{b\underset{¯}{a}}i_{}h_{++}^{b\underset{¯}{a}})v_b^{(0,1)}\},`$
$`H_R^{(2,0)}`$ $`=`$ $`i(\theta ^{(1,0)})^2\{h_{++}+h_{++}^{(ab)}v_a^{(0,1)}v_b^{(0,1)}\theta _{\underset{¯}{a}}^{(0,1)}t_{++}^{b\underset{¯}{a}}`$ (3.23)
$`i(\theta ^{(0,1)})^2_{}h_{++}^{(ab)}v_a^{(0,1)}v_b^{(0,1)}\},`$
$`\mathrm{\Lambda }^{}`$ $`=`$ $`\lambda ^{}2i\theta _{\underset{¯}{a}}^{(0,1)}\lambda ^{a\underset{¯}{a}}v_a^{(0,1)}+i(\theta ^{(0,1)})^2\lambda ^{(ab)}v_a^{(0,1)}v_b^{(0,1)},`$
$`\mathrm{\Lambda }^{(0,1)\underset{¯}{a}}`$ $`=`$ $`\lambda ^{a\underset{¯}{a}}v_a^{(0,1)}+\theta ^{(0,1)\underset{¯}{b}}[\lambda _{\underset{¯}{b})}^{(\underset{¯}{a}}+{\displaystyle \frac{1}{2}}\delta _{\underset{¯}{b}}^{\underset{¯}{a}}(_{}\lambda ^{}+2\lambda ^{(ab)}v_a^{(0,1)}v_b^{(0,1)})]`$
$`(\theta ^{(0,1)})^2(i_{}\lambda ^{a\underset{¯}{a}}+{\displaystyle \frac{1}{2}}\beta _{}^{a\underset{¯}{a}})v_a^{(0,1)},`$
$`\mathrm{\Lambda }_R`$ $`=`$ $`\lambda _R+\lambda ^{(ab)}v_a^{(0,1)}v_b^{(0,1)}\theta _{\underset{¯}{a}}^{(0,1)}\beta _{}^{a\underset{¯}{a}}v_a^{(0,1)}`$ (3.24)
$`i(\theta ^{(0,1)})^2_{}\lambda ^{(ab)}v_a^{(0,1)}v_b^{(0,1)},`$
and $`H^{(0,2)++}`$, $`H^{(1,2)\underset{¯}{i}}`$, $`H_L^{(0,2)}`$, $`\mathrm{\Lambda }^{++}`$, $`\mathrm{\Lambda }^{(1,0)\underset{¯}{i}}`$, $`\mathrm{\Lambda }_L`$ can be obtained from these expressions via the substitutions $`+`$, $`\theta ^{(1,0)\underset{¯}{i}}\theta ^{(0,1)\underset{¯}{a}}`$, $`uv`$, $`i,\underset{¯}{i}a,\underset{¯}{a}`$. In (3.23), (3.24) all the component fields and gauge parameters are functions of $`z^{++},z^{}`$ and we have explicitly indicated their $`2D`$ space-time indices. Note that in the chosen gauge the diagonal components of the world-sheet zweibein $`h_{++}^{++}`$, $`h_{}^{}`$ equal unity and the parameters of two independent Weyl rescalings of $`\theta ^{(1,0)\underset{¯}{i}}`$, $`\theta ^{(0,1)\underset{¯}{a}}`$ are fixed to be $`_{++}\lambda ^{++}`$, $`_{}\lambda ^{}`$, so the difference between the world and tangent indices of the involved fields actually disappears. Actually, we have used all the gauge symmetries with pure shifts in their transformation laws for gauging away the corresponding field components (rescalings are just of this kind). We ended up only with the transformations starting with $`z`$-derivatives of gauge parameters.
Looking at the above expressions we observe that the irreducible content of the original set of analytic vielbeins and connections includes only gauge fields: the two components of the world-sheet zweibein $`h_{}^{++},h_{++}^{}`$, the left and right gravitino components $`h_{}^{+i\underset{¯}{i}},h_{++}^{a\underset{¯}{a}}`$, the left and right components of the $`SO(4)_L\times U(1)_L`$ and $`SO(4)_R\times U(1)_R`$ gauge connections $`h_{}^{(ij)},h_{}^{(\underset{¯}{i}\underset{¯}{j})},h_{}`$ and $`h_{++}^{(ab)},h_{++}^{(\underset{¯}{a}\underset{¯}{b})},h_{++}`$, as well as the left and right components of the “conformal gravitino” $`t_+^{i\underset{¯}{i}}`$, $`t_{++}^{a\underset{¯}{b}}`$, with a total of (16 + 16) independent components. The remaining gauge freedom involves just the same number of gauge parameters, so locally all these gauge fields can be gauged away, though such a gauge is inadmissible globally (e.g., after coupling this multiplet to the $`N=(4,4)`$ string fields, the zweibein components should produce two Virasoro constraints). Therefore it is natural to call the obtained gauge multiplet, with no off-shell degrees of freedom, the “ $`N=(4,4),SO(4)\times U(1)`$ Beltrami-Weyl (BW) multiplet”. We shall see later that it admits truncations to two different $`N=(4,4),SU(2)`$ ones. We will also show that the off-shell (8+8) “minimal $`N=4`$, $`2D`$ SG multiplet” naturally comes out as the result of coupling one of the $`N=(4,4)`$, $`SU(2)`$ BW multiplets to one kind of twisted $`N=(4,4)`$ multiplet treated as a compensator.
Actually, in order to be able to construct manifestly invariant superfield couplings of $`N=(4,4)`$ BW multiplets to $`N=(4,4)`$ matter, we need one more ingredient. This is an analytic density which should transform so as to cancel the transformation of the analytic superspace integration measure $`\mu ^{(2,2)}`$. Indeed, as distinct from the flat superspace superconformal groups, the full local group (3.1) does not leave $`\mu ^{(2,2)}`$ invariant:
$$\delta \mu ^{(2,2)}=((1)^{P(\mu )}_\mu \mathrm{\Lambda }^\mu +^{(2,0)}\mathrm{\Lambda }^{(2,0)}+^{(0,2)}\mathrm{\Lambda }^{(0,2)})\mu ^{(2,2)}\stackrel{~}{\mathrm{\Lambda }}\mu ^{(2,2)},$$
(3.25)
where $`P(\mu )`$ is 0 for bosonic and 1 for fermionic indices.
Defining the objects
$$\mathrm{\Gamma }^{(2,0)}=(1)^{P(M)}_MH^{(2,0)M},\mathrm{\Gamma }^{(0,2)}=(1)^{P(M)}_MH^{(0,2)M},$$
(3.26)
one finds them to transform as
$$\delta \mathrm{\Gamma }^{(2,0)}=^{(2,0)}\stackrel{~}{\mathrm{\Lambda }},\delta \mathrm{\Gamma }^{(0,2)}=^{(0,2)}\stackrel{~}{\mathrm{\Lambda }}$$
(3.27)
and to satisfy, as a consequence of the constraints (3.19), the condition
$$^{(2,0)}\mathrm{\Gamma }^{(0,2)}^{(0,2)}\mathrm{\Gamma }^{(2,0)}=0.$$
(3.28)
It is easy to show that (3.28) implies
$$\mathrm{\Gamma }^{(2,0)}=^{(2,0)}\mathrm{\Sigma }(\zeta ,u,v),\mathrm{\Gamma }^{(0,2)}=^{(0,2)}\mathrm{\Sigma }(\zeta ,u,v).$$
(3.29)
Again, with making use of the constraints (3.19), $`\mathrm{\Sigma }(\zeta ,u,v)`$ can be expressed in terms of the original BW multiplet (up to an unessential additive constant) <sup>3</sup><sup>3</sup>3To the zeroth order in the $`\theta `$’s and the first order in the fields, one has $`\mathrm{\Sigma }=const+(h_{++}^{++}+h_{}^{})+\mathrm{}`$. and shown to transform as
$$\delta \mathrm{\Sigma }=\stackrel{~}{\mathrm{\Lambda }}.$$
(3.30)
Hence the quantity
$$\mathrm{\Omega }e^\mathrm{\Sigma },\delta \mathrm{\Omega }=\stackrel{~}{\mathrm{\Lambda }}\mathrm{\Omega }$$
(3.31)
is the sought object, compensating for the non-invariance of the measure. In what follows we will need only the property
$$(^{(2,0)}+\mathrm{\Gamma }^{(2,0)})\mathrm{\Omega }=0,(^{(0,2)}+\mathrm{\Gamma }^{(0,2)})\mathrm{\Omega }=0.$$
(3.32)
In particular, due to this property, one can still integrate by parts with respect to the covariantized harmonic derivatives. Indeed, for any analytic function $`F(\zeta ,u,v)`$, the integral
$$\mu ^{(2,2)}\mathrm{\Omega }^{(2,0)}F(\zeta ,u,v),$$
up to full ordinary derivatives, reduces to
$$\mu ^{(2,2)}(^{(2,0)}+\mathrm{\Gamma }^{(2,0)})\mathrm{\Omega }F(\zeta ,u,v)=0$$
(the same is true for $`^{(0,2)}`$).
## 4 Various limits and truncations
Inspecting the residual symmetry parameters (3.24), one observes that after constraining their $`z`$ dependence, in such a way that the left (right) parameters are functions solely of $`z^{++}(z^{})`$,
$`_{}\mathrm{\Lambda }^{++}=_{++}\mathrm{\Lambda }^{(1,0)\underset{¯}{i}}=_{}\mathrm{\Lambda }_L=0,`$
$`_{++}\mathrm{\Lambda }^{}=_{++}\mathrm{\Lambda }^{(0,1)\underset{¯}{a}}=_{++}\mathrm{\Lambda }_R=0,`$ (4.1)
they constitute the direct sum of two “large” $`N=(4,4)`$, $`SO(4)\times U(1)`$ superconformal algebras . To see this, one should study the Lie brackets of the transformations (3.1) into which these restricted parameters expanded in series in $`z^{\pm \pm }`$ are substituted. Then, e.g., for the right branch, one finds that the expansion of $`\lambda ^{}(z^{})`$ produces a Virasoro subsector, that of $`\lambda ^{(\underset{¯}{a}\underset{¯}{b})}(z^{}),\lambda ^{(ab)}(z^{})`$ yields two affine $`SU(2)`$ subalgebras, and that of $`\lambda ^{a\underset{¯}{b}}(z^{}),\beta _{}^{a\underset{¯}{b}}(z^{})`$ corresponds to the two types of SUSY generators present in this SCA, i.e. the canonical generators (in particular, the $`N=4`$ Poincaré SUSY and the special conformal SUSY generators) and the non-canonical ones. <sup>4</sup><sup>4</sup>4Strictly speaking, such expansions define that part of $`N=4`$ SCA which is regular at the origin. Just such subalgebras of the left and right $`N=4,SU(2)`$ SCAs were gauged in the component approach of ref. . The affine $`U(1)`$ parameters contained in $`\lambda _R(z^{})`$ appear in the closure of the canonical and non-canonical SUSY transformations (actually, the rigid $`U(1)`$ parameter $`\lambda _R(z^{})|_{z=0}`$ never appears in the closure on the superspace coordinates, but it does appear when one considers the closure on the superfield $`q^{(1,1)}`$ with the transformation law (3.7)). It is also easy to check that these restricted superparameters coincide with those appearing in the realizations of these $`N=4`$ SCAs in the flat $`SU(2)\times SU(2)`$ harmonic superspace .
Thus, we found that the original $`N=(4,4)`$ SG group (3.1), (3.5), (3.6), (3.7), (3.11) contains the direct sum of two $`N=4,SO(4)\times U(1)`$ SCAs as the essential invariance subalgebra of the residual gauge freedom associated with the superparameters (3.24) (and their left counterparts). It should be stressed that it is an invariance of the full nonlinear theory, not only of the linearized approximation (3.23). Indeed, it could be recovered from the general harmonic vielbein transformation laws (3.5), (3.6), (3.11), as the maximal subgroup preserving the flat limit
$$H^{(2,0)M}=H^{(0,2)M}=H_{L,R}^{(2,0)}=H_{L,R}^{(0,2)}=0.$$
(4.2)
Thus, the analytic superdiffeomorphism group of Sect. 3 can be regarded as the local, gauged version of this maximal rigid $`N=(4,4)`$ superconformal group, with the BW multiplet defined by eq. (3.23) (and by its left counterpart) as the corresponding gauge multiplet. Presumably, the latter can be alternatively recovered via direct gauging of this SCA following the procedure of ref. . The $`SU(2)\times SU(2)`$ harmonic superspace approach allows one to relate it to the fundamental objects of the analytic superspace geometry, the analytic harmonic vielbeins $`H^{(2,0)M},H^{(0,2)M}`$ and the analytic $`U(1)`$ connections $`H_{L,R}^{(2,0)},H_{L,R}^{(0,2)}`$.
Since $`N=(4,4),SO(4)\times U(1)`$ SCA contains as its infinite-dimensional subalgebras two $`N=(4,4),SU(2)`$ SCAs (SCA-I and SCA-II), it is natural to expect that its local extension also contains two smaller $`N=(4,4)`$ SG groups having these superconformal symmetries as the maximal “rigid” subgroups. They can naturally be called the $`N=(4,4),SU(2)`$ SG-I and SG-II groups. They should come out as appropriate truncations of (3.1), (3.5), (3.6), (3.11) implemented through imposing certain constraints on the group parameters. The analytic harmonic vielbeins comprising the relevant shortened BW multiplets should then arise upon setting certain relations among the original analytic vielbeins, in a way covariant under the truncated SG group.
One obvious truncation of the original group and vielbeins is as follows:
$`\mathrm{\Lambda }^{(2,0)}=\mathrm{\Lambda }^{(0,2)}=\mathrm{\Lambda }_L=\mathrm{\Lambda }_R=0,`$ (4.3)
$`H^{(4,0)}=H^{(0,4)}=H^{(2,2)}=\stackrel{~}{H}^{(2,2)}=H_{L,R}^{(2,0)}=H_{L,R}^{(0,2)}=0.`$ (4.4)
The resulting group is the group of general analytic diffeomorphisms of the coordinates $`\zeta ^\mu `$, with the inert harmonics
$$\delta \zeta ^\mu =\mathrm{\Lambda }^\mu (\zeta ,u,v),\delta u=\delta v=0.$$
(4.5)
The corresponding covariant harmonic derivatives read
$$^{(2,0)}=D^{(2,0)}+H^{(2,0)\mu }_\mu ,^{(0,2)}=D^{(0,2)}+H^{(0,2)\mu }_\mu .$$
(4.6)
The transformation laws of these derivatives and vielbeins, as well as the constraints the latter should satisfy, directly follow from those given in the previous Section, after taking into account the constraints (4.3), (4.4). Note that the harmonic derivatives now are inert,
$$\delta ^{(2,0)}=\delta ^{(0,2)}=0,$$
and the integrability condition (3.18) becomes
$$[^{(2,0)},^{(0,2)}]=0$$
(4.7)
$`^{(2,0)}H^{(0,2)++}^{(0,2)}H^{(2,0)++}2iH^{(1,2)}\theta ^{(1,0)}=0,`$
$`^{(2,0)}H^{(0,2)}^{(0,2)}H^{(2,0)}+2iH^{(2,1)}\theta ^{(0,1)}=0,`$
$`^{(2,0)}H^{(1,2)\underset{¯}{i}}^{(0,2)}H^{(3,0)\underset{¯}{i}}=0,`$
$`^{(2,0)}H^{(0,3)\underset{¯}{a}}^{(0,2)}H^{(2,1)\underset{¯}{a}}=0.`$ (4.8)
Comparing the above truncated transformations with those of the first rigid superconformal $`N=4,SU(2)`$ group (eqs. (2.11), (2.10)), one can suspect that the truncated SG group corresponds to gauging just this SCA-I. This is indeed the case. One can again choose the gauges
$$H^{(2,0)++}=H^{(0,2)}=H^{(0,3)\underset{¯}{a}}=H^{(3,0)\underset{¯}{i}},$$
(4.9)
as in (3.20), and repeat all the steps which led us to the irreducible field representation (3.23) and the residual gauge freedom (3.24). For the truncated SG case we finally get, at the linearized level,
$`H^{(2,0)}`$ $`=`$ $`i(\theta ^{(1,0)})^2\{h_{++}^{}2i\theta _{\underset{¯}{a}}^{(0,1)}h_{++}^{a\underset{¯}{a}}v_a^{(0,1)}\},`$
$`H^{(2,1)\underset{¯}{a}}`$ $`=`$ $`i(\theta ^{(1,0)})^2\{h_{++}^{a\underset{¯}{a}}v_a^{(0,1)}+\theta ^{(0,1)\underset{¯}{b}}[h_{++\underset{¯}{b})}^{(\underset{¯}{a}}+{\displaystyle \frac{1}{2}}\delta _{\underset{¯}{b}}^{\underset{¯}{a}}_{}h_{++}^{}]`$ (4.10)
$`i(\theta ^{(0,1)})^2_{}h_{++}^{b\underset{¯}{a}}v_b^{(0,1)}\},`$
$`\mathrm{\Lambda }^{}`$ $`=`$ $`\lambda ^{}2i\theta _{\underset{¯}{a}}^{(0,1)}\lambda ^{a\underset{¯}{a}}v_a^{(0,1)},`$
$`\mathrm{\Lambda }^{(0,1)\underset{¯}{a}}`$ $`=`$ $`\lambda ^{a\underset{¯}{a}}v_a^{(0,1)}+\theta ^{(0,1)\underset{¯}{b}}[\lambda _{\underset{¯}{b})}^{(\underset{¯}{a}}+{\displaystyle \frac{1}{2}}\delta _{\underset{¯}{b}}^{\underset{¯}{a}}_{}\lambda ^{}]i(\theta ^{(0,1)})^2_{}\lambda ^{a\underset{¯}{a}}v_a^{(0,1)}`$ (4.11)
and analogous relations for the left vielbeins and parameters. We observe that the same can be obtained simply by setting
$$H_R^{(2,0)}=0,\mathrm{\Lambda }_R=0,H_L^{(0,2)}=0,\mathrm{\Lambda }_L=0$$
in the relations (3.23), (3.24) (and their left counterparts). Thus we end up with the BW multiplet $`h_{++}^{},h_{}^{++}`$, $`h_{++}^{a\underset{¯}{a}},h_{}^{+i\underset{¯}{i}}`$, $`h_{++}^{(\underset{¯}{a}\underset{¯}{b})},h_{}^{(\underset{¯}{i}\underset{¯}{k})}`$ the field content of which basically coincides with that of the $`N=(4,4),SU(2)`$ gauge multiplet found by Schoutens (a slight difference comes from the fact that, on the way to this field representation, we have already gauge-fixed some local symmetries with pure shifts in the relevant gauge parameters, in particular, the local $`2D`$ Lorentz and scale invariances by setting $`h_{++}^{++}=h_{}^{}=1`$). The residual gauge group has the parameters $`\lambda ^{},\lambda ^{++}`$ (local translations), $`\lambda ^{a\underset{¯}{a}},\lambda ^{+i\underset{¯}{i}}`$ (local supertranslations), $`\lambda ^{(\underset{¯}{a}\underset{¯}{b})},\lambda ^{(\underset{¯}{i}\underset{¯}{k})}`$ (right and left $`SU(2)`$ groups). The number of these gauge invariances coincides with that of the gauge fields, so that the $`N=(4,4),SU(2)`$ BW multiplet (BW-I in what follows) contains no off-shell components like its parental $`N=(4,4),SO(4)\times U(1)`$ BW multiplet. Once again, the maximal subgroup of (4.5) preserving the flat limit
$$H^{(2,0)\mu }=H^{(0,2)\mu }=0$$
is just the $`N=(4,4),SU(2)`$ SCA-I. It is singled out by imposing the light-cone chirality conditions on the parameters of the residual gauge group.
While specializing to the $`N=(4,4),SU(2)`$ SG-I group, we may retain the standard defining constraint for the twisted superfield $`q^{(1,1)}`$,
$$^{(2,0)}q^{(1,1)}=^{(0,2)}q^{(1,1)}=0$$
(4.12)
(because of the commutativity property (4.7)), and the zero-weight scalar transformation rule
$$q^{(1,1)}{}_{}{}^{}(\zeta {}_{}{}^{},u,v)=q^{(1,1)}(\zeta ,u,v).$$
(4.13)
So, with respect to this SG-I group, $`q^{(1,1)}`$ is what is called TM-I in because its physical bosonic fields $`q^{ia}(z)`$ are not affected by the local $`SU(2)`$ symmetries (on the contrary, the auxiliary fields $`F^{\underset{¯}{i}\underset{¯}{a}}`$ are transformed). Thus, the general rigidly supersymmetric $`q^{(1,1)}`$ action (2.19) can be straightforwardly extended to the locally supersymmetric one
$$S_q^I=\mu ^{(2,2)}\widehat{\mathrm{\Omega }}^{(2,2)}(q^{(1,1)M},u,v),$$
(4.14)
where the density $`\widehat{\mathrm{\Omega }}`$ is still defined by eqs. (3.29), (3.31), with the truncation conditions (4.4) taken into account. In components and with the auxiliary fields eliminated, it gives the general locally supersymmetric $`N=(4,4)`$ sigma-model of ref. which is a modification of the sigma-model action of ref. by torsion terms in the sector of the physical bosons. For the rigid $`q^{(1,1)}`$ action (2.19), the general torsionful off-shell component action was presented in . The action (4.14) yields a locally supersymmetric version of the latter. In Appendix we present, as an example, the component form of a very simple particular case of (4.14).
What about the second $`N=(4,4),SU(2)`$ SCA, with respect to which $`q^{(1,1)}`$ is TM-II? How to extract the relevant $`N=(4,4)`$ SG group from the original “master” SG group? It is easy to answer these questions at the linearized level. The answer is prompted by the known realization of the $`N=(4,4),SU(2)`$ SCA-II in the $`SU(2)\times SU(2)`$ harmonic analytic superspace . In order to have this SCA as the maximal symmetry after imposing the light-cone chiral constraints (4.1), one must seek for restrictions on the residual gauge superparameters (3.24) and their left counterparts such, that: i) the $`U(1)`$ parameters $`\lambda _{L.R}`$ are identified with $`_{\pm \pm }\lambda ^{\pm \pm }`$; ii) the affine $`SU(2)`$ parameters $`\lambda ^{(\underset{¯}{a}\underset{¯}{b})},\lambda ^{(\underset{¯}{i}\underset{¯}{k})}`$ are eliminated. The unique possibility to obey these requirements, still leaving the “true” $`SU(2)`$ parameters $`\lambda ^{(ab)},\lambda ^{(ik)}`$ unconstrained, is to impose the following relations:
$`{\displaystyle \frac{\mathrm{\Lambda }^{(1.0)\underset{¯}{i}}}{\theta ^{(1,0)\underset{¯}{k}}}}=\delta _{\underset{¯}{k}}^{\underset{¯}{i}}(\mathrm{\Lambda }_L+_{++}\mathrm{\Lambda }^{++}),{\displaystyle \frac{\mathrm{\Lambda }^{(0,1)\underset{¯}{a}}}{\theta ^{(0,1)\underset{¯}{b}}}}=\delta _{\underset{¯}{b}}^{\underset{¯}{a}}(\mathrm{\Lambda }_R+_{}\mathrm{\Lambda }^{}),`$ (4.15)
whence
$`\lambda ^{(\underset{¯}{a}\underset{¯}{b})}=\lambda ^{(\underset{¯}{i}\underset{¯}{k})}=0,\lambda _L={\displaystyle \frac{1}{2}}_{++}\lambda ^{++},\lambda _R={\displaystyle \frac{1}{2}}_{}\lambda ^{},`$
$`\beta _+^{i\underset{¯}{k}}=2i_{++}\lambda ^{+i\underset{¯}{k}},\beta _{}^{a\underset{¯}{b}}=2i_{}\lambda ^{a\underset{¯}{b}}.`$ (4.16)
It is easy to explicitly check that the superparameters (3.24) (and their left counterparts) restricted in this way indeed span the sought $`N=(4,4),SU(2)`$ SCA-II after imposing the chirality conditions (4.1). Then, at the linearized level, it is a consistent truncation to set equal to zero those combinations of the analytic vielbeins, which are not shifted under the subgroup singled out by eqs.(4.15):
$`{\displaystyle \frac{H^{(1.2)\underset{¯}{i}}}{\theta ^{(1,0)\underset{¯}{k}}}}=\delta _{\underset{¯}{k}}^{\underset{¯}{i}}(_{++}H^{(0,2)++}H_L^{(0,2)}),{\displaystyle \frac{H^{(2,1)\underset{¯}{a}}}{\theta ^{(0,1)\underset{¯}{b}}}}=\delta _{\underset{¯}{b}}^{\underset{¯}{a}}(_{}H^{(2,0)}H_R^{(2,0)}).`$ (4.17)
These relations amount to the following linearized constraints on the gauge fields:
$`h_{++}^{(\underset{¯}{a}\underset{¯}{b})}=h_{}^{(\underset{¯}{i}\underset{¯}{k})}=0,h_{}={\displaystyle \frac{1}{2}}_{++}h_{}^{++},h_{++}={\displaystyle \frac{1}{2}}_{}h_{++}^{},`$
$`t_+^{i\underset{¯}{i}}=2i_{++}h_{}^{+i\underset{¯}{i}},t_{++}^{a\underset{¯}{a}}=2i_{}h_{++}^{a\underset{¯}{a}}.`$ (4.18)
They leave us with the representation $`h_{++}^{},h_{}^{++}`$, $`h_{}^{+i\underset{¯}{k}}`$, $`h_{++}^{a\underset{¯}{b}}`$, $`h_{}^{(ik)},h_{++}^{(ab)}`$, which is again a $`N=(4,4),SU(2)`$ BW multiplet, but with another chiral pair of $`SU(2)`$ gauge fields, compared to (4.10). We call it the $`N=(4,4),SU(2)`$ BW-II multiplet. For completeness, we explicitly quote here the counterparts of (4.10), (4.11) for the considered case
$`H^{(2,0)}`$ $`=`$ $`i(\theta ^{(1,0)})^2\{h_{++}^{}2i\theta _{\underset{¯}{a}}^{(0,1)}h_{++}^{a\underset{¯}{a}}v_a^{(0,1)}i(\theta ^{(0,1)})^2h_{++}^{(ab)}v_a^{(0,1)}v_b^{(0,1)}\},`$
$`H^{(2,1)\underset{¯}{a}}`$ $`=`$ $`i(\theta ^{(1,0)})^2\{h_{++}^{a\underset{¯}{a}}v_a^{(0,1)}+{\displaystyle \frac{1}{2}}\theta ^{(0,1)\underset{¯}{a}}(_{}h_{++}^{}2h_{++}^{(ab)}v_a^{(0,1)}v_b^{(0,1)})\},`$
$`H_R^{(2,0)}`$ $`=`$ $`i(\theta ^{(1,0)})^2\{{\displaystyle \frac{1}{2}}_{}h_{++}^{}+h_{++}^{(ab)}v_a^{(0,1)}v_b^{(0,1)}2i\theta _{\underset{¯}{a}}^{(0,1)}_{}h_{++}^{b\underset{¯}{a}}`$ (4.19)
$`i(\theta ^{(0,1)})^2_{}h_{++}^{(ab)}v_a^{(0,1)}v_b^{(0,1)}\},`$
$`\mathrm{\Lambda }^{}`$ $`=`$ $`\lambda ^{}2i\theta _{\underset{¯}{a}}^{(0,1)}\lambda ^{a\underset{¯}{a}}v_a^{(0,1)}+i(\theta ^{(0,1)})^2\lambda ^{(ab)}v_a^{(0,1)}v_b^{(0,1)},`$
$`\mathrm{\Lambda }^{(0,1)\underset{¯}{a}}`$ $`=`$ $`\lambda ^{a\underset{¯}{a}}v_a^{(0,1)}+{\displaystyle \frac{1}{2}}\theta ^{(0,1)\underset{¯}{a}}(_{}\lambda ^{}+2\lambda ^{(ab)}v_a^{(0,1)}v_b^{(0,1)}),`$
$`\mathrm{\Lambda }_R`$ $`=`$ $`{\displaystyle \frac{1}{2}}_{}\lambda ^{}+\lambda ^{(ab)}v_a^{(0,1)}v_b^{(0,1)}+2i\theta _{\underset{¯}{a}}^{(0,1)}_{}\lambda ^{a\underset{¯}{a}}v_a^{(0,1)}`$ (4.20)
$`i(\theta ^{(0,1)})^2_{}\lambda ^{(ab)}v_a^{(0,1)}v_b^{(0,1)}.`$
The left objects are obtained via the same substitutions as in the previous cases.
For the time being, we do not know how to go beyond the linearized level in this important case. It seems that it is more fruitful to descend to the above shortened versions of the BW multiplets (and further to the Poincaré SG), using a more convenient approach based on the concept of superconformal compensation.
## 5 Superconformal matter couplings
The basic idea of the compensation approach (see, e.g., ) is to start from the pure superconformal SG and then to couple to it, in a superconformally covariant way, appropriate matter multiplets with inhomogeneous (Goldstone type) transformation laws with respect to certain (super)conformal symmetries. Then, by properly fixing gauges (normally, in such a way that all inhomogeneously transforming components are fully gauged away), one gets as a net result the theory with a smaller number of local symmetries and supersymmetries, i.e. a sort of Poincaré SG. The auxiliary fields of the compensating superfield become in this gauge auxiliary fields of the relevant Poincaré SG gauge multiplet. If, from the beginning, a few matter superfields coupled to a given conformal SG are included, being one of them a compensator, we end up with the theory of the remaining matter multiplets in a Poincaré SG background. In this way, one can derive various Poincaré-type supergravities (with all, or a part of, the original conformal symmetries compensated for), different off-shell SG multiplets (depending on the choice of compensator), etc.
We believe that the $`N=(4,4),SO(4)\times U(1)`$ SG group defined in Sect. 3 is the maximal, “master” $`N=(4,4),\mathrm{\hspace{0.33em}2}D`$ conformal SG group. Then, the relevant gauge multiplet, $`N=(4,4),SO(4)\times U(1)`$ BW multiplet, is the “master” multiplet from which all other known $`N=(4,4)`$ SG multiplets should follow by the appropriate compensating procedure. To list all possibilities, we need to know all possible superconformal rigid off-shell matter multiplets which can be defined in $`SU(2)\times SU(2)`$ harmonic superspace, their off-shell actions, and the locally superconformal extensions of the latter. As it was already noticed earlier, not all known types of twisted superfields (and their variant representations) admit a simple formulation in $`SU(2)\times SU(2)`$ analytic harmonic superspace . In what follows, we shall deal with the superconformal off-shell matter multiplets which admit a description in terms of analytic $`SU(2)\times SU(2)`$ harmonic superfields and which were reviewed in Sect. 2. These are the nonlinear multiplets $`N^{(2,0)},N^{(0,2)}`$, $`G^{(2,0)},G^{(0,2)}`$ and the twisted chiral multiplets $`q^{(1,1)}`$ which can be either TM-I or TM-II, depending on the superconformal $`N=(4,4),SU(2)`$ group with respect to which one studies their transformation properties. We shall show that some of these superfields can be used to compensate the “master” $`N=(4,4)`$ conformal SG group down to its $`N=(4,4),SU(2)`$ subgroups and, further, to the Poincaré SG groups, including the group of minimal off-shell SG of refs. .
We start with a local extension of the set $`N^{(2,0)},N^{(0,2)}`$. The rigid superconformal transformation laws of this multiplet (2.22) naturally generalize to the whole “master”$`N=(4,4)`$ SG group as
$$\delta N^{(2,0)}=\mathrm{\Lambda }^{(2,0)},\delta N^{(0,2)}=\mathrm{\Lambda }^{(0,2)},$$
(5.1)
where the transformation parameters are now the general analytic superfunctions introduced in (3.1). The defining constraints (2.21) are covariantized as follows:
$`(a)^{(2,0)}N^{(2,0)}+N^{(2,0)}N^{(2,0)}=H^{(4,0)},^{(0,2)}N^{(0,2)}+N^{(0,2)}N^{(0,2)}=H^{(0,4)},`$
$`(b)^{(2,0)}N^{(0,2)}^{(0,2)}N^{(2,0)}=H^{(2,2)}\stackrel{~}{H}^{(2,2)}`$ (5.2)
(for a similar covariantization of the standard nonlinear multiplet in the conventional harmonic superspace, see ). It is obvious that the $`N`$-multiplet can be used to fully compensate all gauge invariances contained in $`\mathrm{\Lambda }^{(2,0)},\mathrm{\Lambda }^{(0,2)}`$, including two chiral $`SU(2)`$ symmetries acting on the harmonic variables. One can achieve this purpose, choosing the gauge
$$N^{(2,0)}=N^{(0,2)}=0(a)H^{(0,4)}=H^{(4,0)}=0,(b)H^{(2,2)}\stackrel{~}{H}^{(2,2)}=0.$$
(5.3)
Prior to any gauge-fixing, it is instructive to fully elaborate on the corollaries of the constraints (5.2). For the quantities
$$Q^{(2,0)}N^{(2,0)}H_L^{(2,0)}H_R^{(2,0)},Q^{(0,2)}N^{(0,2)}H_L^{(0,2)}H_R^{(0,2)},$$
(5.4)
eq. (5.2$`b`$), combined with eqs. (3.17) implies the following constraint:
$$^{(2,0)}Q^{(0,2)}^{(0,2)}Q^{(2,0)}=0Q^{(2,0)}=^{(2,0)}\mathrm{\Phi },Q^{(0,2)}=^{(0,2)}\mathrm{\Phi },$$
(5.5)
where $`\mathrm{\Phi }=\mathrm{\Phi }(\zeta ,u,v)`$ is a new analytic compensating superfield. Recalling the transformation properties (3.11), (5.1), we see that
$$\delta Q^{(2,0)}=^{(2,0)}(\mathrm{\Lambda }_L+\mathrm{\Lambda }_R),\delta Q^{(0,2)}=^{(0,2)}(\mathrm{\Lambda }_L+\mathrm{\Lambda }_R)\delta \mathrm{\Phi }=\mathrm{\Lambda }_L+\mathrm{\Lambda }_R.$$
(5.6)
Hence, the newly introduced analytic object $`\mathrm{\Phi }`$ can be fully gauged away using the analytic gauge parameter $`\mathrm{\Lambda }_L+\mathrm{\Lambda }_R`$
$$\mathrm{\Phi }=0\mathrm{\Lambda }_L=\mathrm{\Lambda }_R\mathrm{\Lambda }.$$
(5.7)
As a corollary of this choice, the following relations occur:
$$Q^{(2,0)}=Q^{(0,2)}=0N^{(2,0)}=H_L^{(2,0)}+H_R^{(2,0)},N^{(0,2)}=H_L^{(0,2)}+H_R^{(0,2)}.$$
(5.8)
At this stage, it is time to fix the gauge freedom associated with the superparameters $`\mathrm{\Lambda }^{(2,0)},\mathrm{\Lambda }^{(0,2)}`$, by imposing the gauge (5.3). As a result of this gauge choice, the original “master” $`N=(4,4)`$ SG group (3.1) proves to be compensated just down to its $`N=(4,4),SU(2)`$ SG-I subgroup (4.5). Eqs. (5.8), in this gauge, imply
$`H_L^{(2,0)}=H_R^{(2,0)}H^{(2,0)},H_L^{(0,2)}=H_R^{(0,2)}H^{(0,2)},`$ (5.9)
$`\delta H^{(2,0)}=^{(2,0)}\mathrm{\Lambda },\delta H^{(0,2)}=^{(0,2)}\mathrm{\Lambda }.`$ (5.10)
As a consequence of these relations and the gauge choice (5.7), the transformation law (3.7) of the twisted multiplet $`q^{(1,1)}`$ in the “master” SG group, as well as its defining constraints (3.13), are reduced to those covariant under the $`N=(4,4),SU(2)`$ SG group, i.e. (4.13) and (4.12). Nevertheless, the resulting theory is not yet identical to what we have got after truncation in Sect. 4. Indeed, the gauge-fixed covariant derivatives $`^{(2,0)},^{(0,2)}`$ differ from those defined by eq. (4.6)
$`^{(2,0)}=D^{(2,0)}+H^{(2,0)\mu }_\mu +H^{(2,2)}^{(0,2)},`$
$`^{(0,2)}=D^{(0,2)}+H^{(0,2)\mu }_\mu +H^{(2,2)}^{(2,0)},`$ (5.11)
$`H^{(2,2)}=^{(2,0)}H^{(0,2)}^{(0,2)}H^{(2,0)}.`$ (5.12)
Though $`H^{(2,2)}`$ as well as $`H^{(2,0)},H^{(0,2)}`$ transform as scalars under the remaining $`N=(4,4),SU(2)`$ SG-I group, the harmonic partial derivatives $`^{(0,2)},^{(2,0)}`$ are not covariant, due to the presence of a non-trivial $`u,v`$ dependence in the group parameters in (4.5). As a result, $`H^{(2,2)}`$ appears in the transformation laws of the vielbeins $`H^{(2,0)\mu },H^{(0,2)\mu }`$. The constraints (3.19) also do not go into the set (4.8), due to the presence of $`H^{(2,2)}`$. For this object, the original constraints (3.19) imply the following ones (recall that $`H^{(4,0)}=H^{(0,4)}=0`$ in the gauge (5.3)):
$$^{(2,0)}H^{(2,2)}=^{(0,2)}H^{(2,2)}=0.$$
(5.13)
This peculiarity comes out only at the nonlinear level. The linearized analysis goes as before and shows that $`H^{(2,0)\mu },H^{(0,2)\mu }`$, in the present case, carry the same set of fields forming the $`N=(4,4),SU(2)`$ BW-I multiplet. In other words, after fixing appropriate conformal gauges in the locally superconformal system of the original “master” BW multiplet and the compensator multiplet $`N^{(2,0)},N^{(0,2)}`$, we are left with a smaller $`N=(4,4),SU(2)`$ BW-I multiplet and an extra off-shell multiplet. The latter is carried by the superfields $`H^{(2,0)},H^{(0,2)}`$ which exhibit the gauge freedom (5.10) with an extra analytic gauge parameter $`\mathrm{\Lambda }(\zeta ,u,v)`$ and satisfy the constraints (5.13). This extended representation is not fully reducible, in the sense that the additional gauge superfields $`H^{(2,0)},H^{(0,2)}`$ are scalars with respect to the conformal $`N=(4,4),SU(2)`$ SG-I group (4.5) while the SG-I transformation laws of the analytic vielbeins $`H^{(2,0)\mu },H^{(0,2)\mu }`$ include these extra superfields.
Thus, we have found the previously unknown off-shell $`N=(4,4)`$ SG gauge multiplet. In the WZ gauge and at the linearized level, its part coming from the analytic vielbeins is the same BW-I gauge fields representation which was described in Sect. 4 and which has no off-shell degrees of freedom (the linearized structure (4.10) in this case is slightly modified by the fields from $`H^{(2,0)},H^{(0,2)}`$, because of the presence of $`H^{(2,2)}`$ in the constraints on $`H^{(0,2)\mu },H^{(2,0)\mu }`$). To examine the off-shell content of $`H^{(2,0)},H^{(0,2)}`$, we have chosen an appropriate WZ gauge with respect to the parameter $`\mathrm{\Lambda }(\zeta ,u,v)`$, so as to kill as much component fields in the $`\theta `$, $`u,v`$ expansions of these superfields as possible, and inserted the result into the linearized form of the constraints (5.13). Solving the latter (it does not put any field on shell), we have eventually found $`(32+32)`$ independent off-shell components listed below (the numerals in the parentheses on the right to the fields denote the “engineering” dimension and the number of independent real components, respectively):
$`\underset{¯}{\text{bosons}}:(h_{++},h_{})(1,1),l_{++}^{(ab)}(1,3),l_{}^{(ik)}(1,3),l_{\underset{¯}{i}\underset{¯}{a}}^{ia}(1,16),l^{(ik)(ab)}(0,9),`$
$`\underset{¯}{\text{fermions}}:l_{\underset{¯}{a}+}^b(3/2,4),l_{\underset{¯}{k}}^i(3/2,4),l_{\underset{¯}{b}}^{(ik)a}(1/2,12),l_{\underset{¯}{k}}^{(ab)i}(1/2,12).`$ (5.14)
The fields $`h_{\pm \pm }`$ are gauge fields for a $`U(1)`$ with the gauge parameter $`\lambda (z)`$ which is the first component in $`\mathrm{\Lambda }(\zeta ,u,v)`$. This $`U(1)`$ is the only residual gauge symmetry of the given WZ gauge. The fields $`l_{++}^{(ab)},l_{}^{(ik)}`$ and $`l_{\underset{¯}{a}+}^b,l_{\underset{¯}{k}+}^i`$ are “former” gauge fields for the symmetries with the parameters $`\lambda ^{(ab)},\lambda ^{(ik)}`$ and $`\beta _{\underset{¯}{a}}^b,\beta _{\underset{¯}{k}+}^i`$ in the “master” BW multiplet (eqs. (3.23), (3.24) and their left counterparts). Now these local symmetries have been entirely compensated by the appropriate compensating fields from $`N^{(2,0)},N^{(0,2)}`$. Note that the residual gauge group $`U(1)`$ is the diagonal in the product of two chiral gauge $`U(1)`$ groups realized on the “master” BW multiplet; the rest of these $`U(1)`$ symmetries has been compensated by a dimension-0, $`SO(4)`$ singlet field present in $`N^{(2,0)},N^{(0,2)}`$ (this is just the first component of the compensator $`\mathrm{\Phi }`$ introduced in (5.5)). The biggest flat limit symmetry of the extended gauge multiplet ($`N=(4,4),SU(2)`$ BW-I together with (5.14)) is $`N=(4,4),SU(2)`$ SCA-I augmented with an extra rigid $`U(1)`$ symmetry. Note that in the matter couplings we shall discuss in this Section, the superfields $`H^{(2,0)},H^{(0,2)}`$ always appear only through their analytic superfield strength $`H^{(2,2)}`$ containing, in particular, the field strength of the $`U(1)`$ gauge field $`h_{\pm \pm }`$. In other words, the residual local $`U(1)`$ group is hidden, and for the time being we do not see in which situations it could become active.
A comment is to the point here. In principle, we could completely eliminate the extra multiplet by treating it as pure gauge. This possibility corresponds to adding additional constraints to the set (5.2)
$$H^{(2,2)}=^{(2,0)}N^{(0,2)},\stackrel{~}{H}^{(2,2)}=^{(0,2)}N^{(2,0)}.$$
(5.15)
These constraints are manifestly covariant and compatible both with (5.2) and (3.19).The same reasoning which led us to eqs. (5.5), (5.9) implies that in the gauge (5.3) the pairs $`H_L^{(2,0)},H_L^{(0,2)}`$ and $`H_R^{(2,0)},H_R^{(0,2)}`$ become pure gauge, with respect to the $`U(1)`$ gauge groups with parameters $`\mathrm{\Lambda }_L`$ and $`\mathrm{\Lambda }_R`$. Hence, they can be gauged away, fully compensating this gauge freedom. As the result, the $`N=(4,4),SU(2)`$ SG-I group and the BW-I multiplet are finally reproduced. A deviation from the standard compensation point of view is that, after imposing (5.15), the compensators $`N^{(2,0)},N^{(0,2)}`$ cease to have a flat off-shell limit (when all vielbeins are put equal to zero): the resulting modified set of constraints proves to be too restrictive, it puts these superfield on shell . On the other hand, one can view the relations (5.2) and (5.15) merely as the covariant definition of particular harmonic vielbeins $`H^{(0,4)},H^{(4,0)},H^{(2,2)},\stackrel{~}{H}^{(2,2)}`$, such that it provides a covariant way to make some gauge fields in the “master” BW multiplet purely longitudinal and, so, globally removable by fixing appropriate gauges. Indeed, from the standpoint of the linearized WZ representation (3.23) for the “master” BW multiplet, these relations mean that all gauge fields except those comprising the $`N=(4,4),SU(2)`$ BW-I multiplet are postulated to be pure gauge.
Let us now turn to the issue of constructing matter actions invariant under the “master” conformal SG group.
We start by seeking for the appropriate generalization of the $`N`$-action (2.23). Somewhat surprisingly, it cannot be straightforwardly promoted to an invariant of the local superconformal group. The best we have reached, in our attempts to covariantize (2.23), is the action
$$S_N^{loc}=\mu ^{(2,2)}\mathrm{\Omega }\left(Q^{(2,0)}Q^{(0,2)}+2Q^{(2,0)}H_L^{(0.2)}+2Q^{(0,2)}H_R^{(2,0)}+2H_L^{(0,2)}H_R^{(2,0)}\right),$$
(5.16)
where $`Q^{(2,0)},Q^{(0,2)}`$ are defined in (5.4). It is shifted, up to surface terms, by the expression
$$\delta S_N^{loc}=2\mu ^{(2,2)}\mathrm{\Omega }\left(H_L^{(0,2)}^{(2,0)}\mathrm{\Lambda }_L+H_R^{(2,0)}^{(0,2)}\mathrm{\Lambda }_R\right),$$
(5.17)
which cannot be further cancelled in any way.
On the other hand, it is possible to construct invariant actions for the second type of superconformally invariant (32+32) nonlinear multiplet defined by (2.25), (2.26). The constraints (2.26) admit a direct covariantization
$`(^{(2,0)}+2N^{(2,0)})G^{(2,0)}+\alpha G^{(2,0)}G^{(2,0)}=0,`$
$`(^{(0,2)}+2N^{(0,2)})G^{(0,2)}+\alpha G^{(0,2)}G^{(0,2)}=0,`$
$`^{(2,0)}G^{(0,2)}^{(0,2)}G^{(2,0)}=0.`$ (5.18)
Indeed, it is easy to check their covariance under (3.1) provided that the superfields $`G`$ transform as scalars: $`\delta G^{(2,0)}=\delta G^{(0,2)}=0`$. Then the simplest manifestly invariant action of $`G^{(2,0)},G^{(0,2)}`$ in the background of the $`N=(4,4)`$ “master” conformal SG fields and compensators $`N^{(2,0)},N^{(0,2)}`$ is given by
$$S_G^{loc}=\mu ^{(2,2)}\mathrm{\Omega }G^{(2,0)}G^{(0,2)}.$$
(5.19)
Another possibility to construct an invariant off-shell action for the pair of compensators $`N^{(2,0)},N^{(0,2)}`$ is to take as the relevant Lagrangian density the constraints (5.2) with the appropriate analytic Lagrange multipliers $`\omega ^{(2,2)}`$, $`\omega ^{(2,2)}`$, $`\omega `$. Just an action of this kind describes the standard nonlinear multiplet coupled to a conformal $`N=2,\mathrm{\hspace{0.33em}4}D`$ SG in the conventional harmonic superspace . Its $`SU(2)\times SU(2)`$ analogue would also have no propagating degrees of freedom and, before varying with respect to Lagrange multipliers, contain an infinite number of auxiliary fields. This possibility requires a thorough analysis and we postpone discussing it to the future.
It is worth noting that there are no problems with extending the flat superspace actions of $`N^{(2,0)},N^{(0,2)}`$ and $`G^{(2,0)},G^{(0,2)}`$ to invariants of the $`N=(4,4)`$ SG-I group. The relevant constraints are obtained from the flat ones (2.21), (2.26) by the replacements $`D^{(2,0)},D^{(0,2)}^{(2,0)},^{(0,2)}`$, where $`^{(2,0)},^{(0,2)}`$ are given by eqs. (4.6), and the locally supersymmetric actions are obtained via the replacement $`\mu ^{(2,2)}\mu ^{(2,2)}\widehat{\mathrm{\Omega }}`$ in the flat superspace ones.
Let us now switch over to the twisted multiplets. We already constructed in Sect. 4 a locally supersymmetric $`q^{(1,1)}`$ action (4.14) invariant under the $`N=(4,4),SU(2)`$ SG-I group. An important question is how to construct the $`q^{(1,1)}`$ actions invariant under the full “master” $`N=(4,4)`$ group (3.1). The main difficulty here is related to the non-trivial transformation law (3.7) of $`q^{(1,1)}`$ in this group.
The simplest way to construct such a coupling is to consider $`q^{(1,1)}`$ together with the compensators $`N^{(2,0)},N^{(0,2)}`$. In this case, due to the existence of the analytic scalar compensator $`\mathrm{\Phi }`$ which is shifted by the sum $`\mathrm{\Lambda }_L+\mathrm{\Lambda }_R`$ (eq. (5.5)), one can redefine any $`q^{(1,1)}`$ with the transformation law (3.7) in such a way that it will transform as a scalar under the “master” SG group
$$q^{(1,1)}(\zeta ,u,v)\stackrel{~}{q}^{(1,1)}=e^\mathrm{\Phi }q^{(1,1)},\stackrel{~}{q}^{(1,1)}{}_{}{}^{}(\zeta ^{},u^{},v^{})=\stackrel{~}{q}^{(1,1)}(z,u,v).$$
(5.20)
The constraints (3.13) become
$$(^{(2,0)}+N^{(2,0)})\stackrel{~}{q}^{(1,1)}=0,(^{(0,2)}+N^{(0,2)})\stackrel{~}{q}^{(1,1)}=0.$$
(5.21)
Their covariance is evident. The general invariant action is similar to (4.14)
$$\stackrel{~}{S}_q^I=\mu ^{(2,2)}\mathrm{\Omega }^{(2,2)}(\stackrel{~}{q}^{(1,1)},\stackrel{~}{u},\stackrel{~}{v}),$$
(5.22)
where <sup>5</sup><sup>5</sup>5Within the conventional harmonic superspace, the necessity of analogous redefinitions of the harmonics explicitly appearing in the action of hypermultiplets coupled to conformal $`N=2`$, $`4D`$ SG was firstly shown in .
$`\stackrel{~}{u}^{(1,0)}=u^{(1,0)}N^{(2,0)}u^{(1,0)},\stackrel{~}{u}^{(1,0)}=u^{(1,0)},`$
$`\stackrel{~}{v}^{(0,1)}=v^{(0,1)}N^{(0,2)}v^{(0,1)},\stackrel{~}{v}^{(1,0)}=v^{(0,1)}.`$ (5.23)
In the gauge (5.3) the action (5.22) coincides with (4.14) modulo a modification of both the covariant harmonic derivatives and the constraints on the analytic vielbeins due to the presence of the $`U(1)`$ gauge multiplet $`H^{(2,0)},H^{(0,2)}`$. The effect of this modification is two-fold: first, the constraints defining $`q^{(1,1)}`$ are obscured by this extra multiplet and, second, the density $`\mathrm{\Omega }`$ differs from $`\stackrel{~}{\mathrm{\Omega }}`$ in (4.14) owing to the presence of the extra multiplet in the constraints for the analytic vielbeins. It would be interesting to see what is the precise impact of this modification on the component sigma-model action as compared to the action (4.14) which includes the $`N=(4,4),SU(2)`$ BW-I multiplet without any additional SG fields.
Since there exists the unique $`N=(4,4)`$ WZW $`q^{(1,1)}`$ action (2.17) invariant under the full rigid $`N=(4,4),SO(4)\times U(1)`$ superconformal symmetry, it is natural to seek for its direct coupling to the “master” $`N=(4,4)`$ BW multiplet without adding any extra compensators. If such a coupling can be set up, $`q^{(1,1)}`$ can be regarded, like $`N^{(2,0)},N^{(0,2)}`$, as a compensator extending the master BW multiplet to some SG multiplet with a smaller number of gauge symmetries and gauge fields. The corresponding SG group should be some subgroup of the master $`N=(4,4)`$ SG group. Indeed, the shifted superfield $`\widehat{q}^{(1,1)}`$ defined in (2.18) transforms inhomogeneously under (3.1), (3.7)
$$\delta \widehat{q}^{(1,1)}=(\mathrm{\Lambda }_L+\mathrm{\Lambda }_R)(\widehat{q}^{(1,1)}+c^{(1,1)})\mathrm{\Lambda }^{(2,0)}c^{(1,1)}\mathrm{\Lambda }^{(0,2)}c^{(1,1)},$$
(5.24)
and hence it can be employed as a compensator.
Unfortunately, we do not have yet any general recipe how to construct such a locally supersymmetric extension of (2.17). The main difficulty stems from the fact that the analytic superfield density in (2.17) is not a tensor: it is shifted by full harmonic derivatives under the rigid superconformal $`SO(4)\times U(1)`$ transformation. The most straightforward approach is to restore the full action order by order in the SG superfields, and this is what we shall undertake.
First, we make the replacement
$$\mu ^{(2,2)}\mu ^{(2,2)}\mathrm{\Omega }$$
in (2.17) in order to be able to integrate by parts with respect to $`^{(2,0)}`$, $`^{(0,2)}`$ (recall the discussion at the end of Sect. 3). We do not fix beforehand any gauges including (3.14). Thus we represent the sought $`S_{wzw}^{loc}`$ as a series in powers of the SG superfields
$$S_{wzw}^{loc}=S_{(0)}+S_{(1)}+S_{(2)}+\mathrm{}=\frac{1}{4\gamma ^2}\mu ^{(2,2)}\mathrm{\Omega }[_{(0)}^{(2,2)}+_{(1)}^{(2,2)}+_{(2)}^{(2,2)}+\mathrm{}],$$
(5.25)
where $`_{(0)}^{(2,2)}`$ is just the density in (2.17). Then, using the formula
$$\delta S_{(0)}=\frac{1}{4\gamma ^2}\mu ^{(2,2)}\mathrm{\Omega }\left(\widehat{q}^{(1,1)}\delta \widehat{q}^{(1,1)}\frac{1}{(1+X)^2}\right),$$
it is rather straightforward to restore the first correction term in (5.25):
$`S_{(1)}`$ $`=`$ $`{\displaystyle \frac{1}{4\gamma ^2}}{\displaystyle }\mu ^{(2,2)}\mathrm{\Omega }(\widehat{q}^{(1,1)}{\displaystyle \frac{1}{(1+X)^2}}[c^{(1,1)}H_L^{(2,0)}+c^{(1,1)}H_R^{(0,2)}`$ (5.26)
$`(c^{(1,1)}H_L^{(0,2)}+c^{(1,1)}H_R^{(2,0)})(2+X)]).`$
A problem is met at the next step, when trying to calculate the second term. Including from the beginning all possible appropriate structures, we finally found that almost all structures appearing in the first-order variation of $`S_{(0)}+S_{(1)}`$ can be cancelled by the zeroth-order variation of $`S_{(2)}`$. Only one term cannot be cancelled. It looks just the same as the term (5.17) appearing in the variation of the non-invariant $`N^{(2,0)},N^{(0,2)}`$ action (5.16)
$$\delta (S_{(0)}+S_{(1)}+S_{(2)})=\frac{1}{4\gamma ^2}\mu ^{(2,2)}\mathrm{\Omega }\left[2H_L^{(0,2)}^{(2,0)}\mathrm{\Lambda }_L+2H_R^{(2,0)}^{(0,2)}\mathrm{\Lambda }_R\right].$$
(5.27)
The origin of this anomaly can be inferred from the results of ref. where the problem of gauging isometries of bosonic sigma models with torsion was studied. As was shown there, in the case of group manifold WZW model associated with a group $`G`$ it is impossible to construct an action in which the full $`G\times G`$ symmetry of the rigid WZW action would be gauged (without adding extra copies of WZW fields). One can only gauge either the left, or right, or diagonal subgroups of $`G\times G`$. The bosonic sector of the above rigid $`q^{(1,1)}`$ action is just the $`SU(2)_L\times SU(2)_R/SU(2)_{diag}`$ WZW action, while the “master” $`N=(4,4)`$ SG group implies gauging both $`SU(2)_L`$ and $`SU(2)_R`$ symmetries. Thus, in view of the argument just adduced, a direct coupling of the WZW $`q^{(1,1)}`$ action (2.17) to the “master” BW multiplet does not exist and the “classical anomaly” (5.27) is just a manifestation of this fact. The unremovable piece in the gauge variation (5.17) is of the same origin, because the $`N^{(2,0)},N^{(0,2)}`$ action (2.23) also contains the $`SU(2)_L\times SU(2)_R`$ WZW model in its bosonic sector. The same reasoning implies the non-existence of similar straightforward $`N=(4,4)`$ SG-II group-invariant extensions of (2.17), (2.23), since this SG group still includes gauge $`SU(2)_L,SU(2)_R`$ symmetries which act on the physical bosons of $`q^{(1,1)}`$ ($`SU(2)`$ WZW fields). Note that no problems of this sort arise while promoting (2.17) to an invariant of the $`N=(4,4)`$ SG-I gauge group, or to that of the “master” SG group with making use of the $`N^{(2,0)},N^{(0,2)}`$ compensators at the intermediate step: such locally supersymmetric $`q^{(1,1)}`$ actions are particular cases of (4.14), (5.22).
Thus the construction of direct couplings of $`N=(4,4)`$ WZW action (2.17) or the $`N^{(2,0)},N^{(0,2)}`$ action (2.23) to the “master” conformal $`N=(4,4)`$ SG or $`N=(4,4)`$ SG-II is a non-trivial problem. It seems that the unique possibility to arrange such couplings is to consider a few copies of the superfields $`q^{(1,1)}`$, $`N^{(2,0)},N^{(0,2)}`$. Then one can construct invariant actions as sums of the individual actions of the type (5.16), (5.25), taking some of them with the wrong sign so as to cancel out the non-vanishing variations like (5.17), (5.27) coming from different actions. This is possible just because these anomalous variations involve only SG gauge fields.
The simplest possibility is to consider a pair of nonlinear multiplets, $`N_1^{(2,0)},N_1^{(0,2)}`$ and $`N_2^{(2,0)},N_2^{(0,2)}`$, each set being subjected to the constraints (5.2). Then the difference of two actions (5.16)
$`S_{N_1N_2}^{loc}=S_{N_1}^{loc}S_{N_2}^{loc}={\displaystyle }\mu ^{(2,2)}\mathrm{\Omega }[N_1^{(2,0)}N_1^{(0,2)}N_2^{(2,0)}N_2^{(0,2)}`$
$`+(N_1^{(2,0)}N_2^{(2,0)})(H_L^{(0,2)}H_R^{(0,2)})(N_1^{(0,2)}N_2^{(0,2)})(H_L^{(2,0)}H_R^{(2,0)})]`$ (5.28)
can be easily checked to be invariant under the “master” SG group. Each of these multiplets, or their sum $`N_1^{(2,0)}+N_2^{(2,0)},N_1^{(0,2)}+N_2^{(0,2)}`$ can be chosen as compensators reducing the “master” SG group to SG-I group via gauge-fixings like (5.3) (and, further, (5.7)). Note that the gauge-invariant combinations $`\stackrel{~}{N}^{(2,0)}=N_1^{(2,0)}N_2^{(2,0)},\stackrel{~}{N}^{(0,2)}=N_1^{(0,2)}N_2^{(0,2)}`$ obey the constraints
$`\left[^{(2,0)}+(N_1^{(2,0)}+N_2^{(2,0)})\right]\stackrel{~}{N}^{(2,0)}=0,\left[^{(0,2)}+(N_1^{(0,2)}+N_2^{(0,2)})\right]\stackrel{~}{N}^{(0,2)}=0,`$
$`^{(2,0)}\stackrel{~}{N}^{(0,2)}^{(0,2)}\stackrel{~}{N}^{(2,0)}=0,`$ (5.29)
which are recognized as the $`\alpha =0`$ version of (5.18). So one can from the beginning add the invariant piece
$$\stackrel{~}{N}^{(2,0)}\stackrel{~}{N}^{(0,2)}$$
(5.30)
to the Lagrangian density in (5.28). Finally, e.g., in the gauges
$$N_1^{(0,2)}+N_2^{(0,2)}=N_1^{(2,0)}+N_2^{(2,0)}=0$$
and (5.7), we are left with the action of the (32 + 32) matter multiplet $`\stackrel{~}{N}^{(2,0)},\stackrel{~}{N}^{(0,2)}`$ in the background of the BW-I multiplet augmented with the $`U(1)`$ multiplet (5.14). The invariant couplings of an arbitrary number of $`q^{(1,1)}`$ multiplets can be arranged with the help of the compensator $`N_1^{(2,0)}+N_2^{(2,0)},N_1^{(0,2)}+N_2^{(0,2)}`$ as explained above, and added to the $`N_1,N_2`$ action.
Another possibility to set up a direct coupling to the “master” $`N=(4,4)`$ conformal SG is to take the difference of the “almost-covariant” $`N`$ -superfields action (5.16) and the $`q^{(1,1)}`$ action (5.25)
$$S_{qN}^{loc}=S_{wzw}^{loc}\frac{1}{4\gamma ^2}S_N^{loc}.$$
(5.31)
In analogy with the $`N^{(2,0)},N^{(0,2)}`$ action (5.16), it is natural to assume that the only unremovable term in the “master” SG group variation of (5.25) is given by (5.27), while all higher-order variations can be cancelled by inserting into the Lagrangian density the appropriate higher-order structures composed out of the analytic vielbeins and the superfield $`q^{(1,1)}`$. This still has to be proved (it would be desirable to find out the geometric principle behind such a recursion procedure). <sup>6</sup><sup>6</sup>6The $`N=(4,4)`$, $`2D`$ WZW - SG couplings were earlier constructed using $`N=1`$ superfields and the conventional $`N=4`$ superfields in . If such an “almost-covariant” $`q^{(1,1)}`$ action exists, the action (5.31) is invariant like (5.28), and in the gauges (5.3), (5.7) we arrive at the SG-I group-invariant action of the (8+8) multiplet $`q^{(1,1)}`$ in the background consisting of the BW-I gauge multiplet and the extra multiplet (5.14). Adding other $`q^{(1,1)}`$ superfields in a way covariant under “master” SG group can be accomplished, like in the previous case, with making use of the compensators $`N^{(2,0)},N^{(0,2)}`$.
An essentially new situation comes out if one directly couples $`q^{(1,1)}`$ superfields to BW multiplet, without using $`N`$-superfields. For this purpose one should take at least two different WZW $`q^{(1,1)}`$ superfields with the same transformation law (5.24) (although with different sets of constants $`c^{ia}`$, generally speaking). Under the assumptions that the “almost-covariant” $`q^{(1,1)}`$ action $`S_{wzw}^{loc}`$ exists to all orders in SG fields and that its non-invariance is given only by the variation (5.27), the fully invariant action could be constructed as
$$S_{q_1q_2}^{loc}=S_{wzw_1}^{loc}S_{wzw_2}^{loc}.$$
(5.32)
More generally, one can take a sum of $`n`$ such actions and choose the coefficients in such a way that the anomaly variations (5.27) coming from different items in the sum are cancelled out. <sup>7</sup><sup>7</sup>7A similar trick was used in for construction of a gauge-invariant WZW action with the full $`G\times G`$ symmetry group gauged. Clearly, at least one of such actions should enter with a “wrong” sign, presumably indicating that the relevant $`q^{(1,1)}`$ is a sort of “Liouville coordinate” . In view of the inhomogeneous nature of the transformation law (5.24), one of the $`q^{(1,1)}`$ superfields will play the role of a compensator.
As it was already mentioned, for the time being we are not aware of the full nonlinear structure of the “almost-covariant” $`q^{(1,1)}`$ actions and even of the complete proof of their existence. Nevertheless, taking for granted that such actions can be constructed, let us inspect which kind of compensation of the “master” SG group can be achieved with the help of $`q^{(1,1)}`$. It will be enough to perform this analysis at the linearized level.
We shall start from the linearized WZ gauge content of BW gauge multiplet (3.23) and the corresponding form (3.24) of the residual symmetry. At the linearized level, the $`q^{(1,1)}`$ constraints (3.13) read (for the shifted superfield $`\widehat{q}^{(1,1)}=q^{(1,1)}c^{(1,1)}`$)
$`D^{(2,0)}\widehat{q}^{(1,1)}=c^{(1,1)}D^{(0,2)}H_R^{(2,0)}c^{(1,1)}H_R^{(2,0)},`$
$`D^{(0,2)}\widehat{q}^{(1,1)}=c^{(1,1)}D^{(2,0)}H_L^{(0,2)}c^{(1,1)}H_L^{(0,2)}.`$ (5.33)
They imply
$`\widehat{q}^{(1,1)}`$ $`=`$ $`\widehat{q}^{ia}(z)u_i^{(1,0)}v_a^{(0,1)}+\theta ^{(1,0)\underset{¯}{i}}\psi _{+\underset{¯}{i}}^a(z)v_a^{(0,1)}+\theta ^{(0,1)\underset{¯}{a}}\chi _{\underset{¯}{a}}^i(z)u_i^{(1,0)}`$ (5.34)
$`+i\theta ^{(1,0)\underset{¯}{i}}\theta ^{(0,1)\underset{¯}{a}}F_{\underset{¯}{i}\underset{¯}{a}}(z)+\mathrm{},`$
where dots stand for the terms involving the BW multiplet gauge fields and derivatives of the explicitly written physical dimension fields of $`q^{(1,1)}`$. The purely shift part of the transformation (5.24) (we need only the latter for our linearized analysis)
$$\delta \widehat{q}^{(1,1)}=c^{(1,1)}\left(\mathrm{\Lambda }_L+\mathrm{\Lambda }_R\right)c^{(1,1)}D^{(2,0)}\mathrm{\Lambda }_Lc^{(1,1)}D^{(0,2)}\mathrm{\Lambda }_R$$
(5.35)
amounts to the following transformations of the fields:
$`\delta \widehat{q}^{ia}=c^{ia}(\lambda _L+\lambda _R)c_j^a\lambda _L^{(ji)}c_b^i\lambda _R^{(ba)},`$
$`\delta \psi _{+\underset{¯}{i}}^a=c_i^a\beta _{+\underset{¯}{i}}^i,\delta \chi _{\underset{¯}{a}}^i=c_a^i\beta _{\underset{¯}{a}}^a,`$
$`\delta F_{\underset{¯}{i}\underset{¯}{a}}=0.`$ (5.36)
One observes that all the physical dimension fields can be gauged away by appropriate gauge parameters
$`\widehat{q}^{ia}=0(a)\lambda _L=\lambda _R\lambda ,(b)\lambda _L^{(ij)}={\displaystyle \frac{1}{c^2}}\lambda _R^{(ab)}c_a^ic_b^j\lambda ^{(ij)},`$ (5.37)
$`\psi _{\underset{¯}{i}}^a=\chi _{\underset{¯}{a}}^i=0\beta _{+\underset{¯}{i}}^i=\beta _{\underset{¯}{a}}^a=0,`$ (5.38)
where $`c^2=c^{ia}c_{ia}0`$. As it follows from (5.37), the product of two local $`U(1)`$ symmetries is compensated down to the diagonal $`U(1)`$, and the same occurs for the product $`SU(2)_L\times SU(2)_R`$ (this results in identifying the $`SU(2)`$ indices of the left and right harmonics, though still does not reduce two harmonic sets to each other). Eq. (5.38) implies the full compensation of the local non-canonical supersymmetries.
As the result, in the gauge (5.37), (5.38) the irreducible off-shell gauge representation comprises the (0+0) BW-I multiplet (4.10) as a submultiplet of the “master” BW multiplet we started with, as well as a new off-shell (8+8) gauge multiplet. The latter inherits a part of its fields from the original BW multiplet, and a part from the compensating $`\widehat{q}^{(1,1)}`$ multiplet
$`\underset{¯}{\text{bosons}}:(h_{++},h_{})(1,1),(h_{++}^{(ik)},h_{}^{(ik)})(1,3),F_{\underset{¯}{i}\underset{¯}{a}}(1,4),`$
$`\underset{¯}{\text{fermions}}:t_{++}^{i\underset{¯}{a}}(3/2,4),t_+^{k\underset{¯}{j}}(3/2,4).`$ (5.39)
Comparing it with the (8+8) “$`Sp(1)`$ vector multiplet” of ref. , we find almost full identity between the two representations, except for a minor distinction related to the fact that one bosonic degree of freedom in (5.39) is represented by the dimension 1 $`U(1)`$ gauge field $`h_{\pm \pm }`$, while in it is carried over by the dimension 2 auxiliary field. It is natural to identify the latter with the curl $`_{++}h_{}_{}h_{++}`$, in view of the well-known equivalence of the auxiliary scalar field and the curl of gauge vector field in two dimensions. Note that the $`Sp(1)`$ vector multiplet was introduced in “by hand”, in addition to the purely gauge SG multiplet which we call here BW-I, in order to be able to construct locally $`N=(4,4)`$ supersymmetric sigma models on quaternionic manifolds. In our scheme it naturally appears, along with the BW-I gauge multiplet, as a result of compensating the “master” $`N=(4,4)`$ SG group by the TM-II multiplet $`q^{(1,1)}`$. The gauge $`Sp(1)`$ symmetry of is recognized as the diagonal in the product of $`SU(2)_L`$ and $`SU(2)_R`$ symmetries realized as isometries of the WZW bosonic fields in $`q^{(1,1)}`$. It would be of interest to study this correspondence at the full nonlinear level and, in particular, to inquire how to construct superconformally-invariant couplings of some other matter $`q^{(1,1)}`$ superfields to this field representation (different from a simple sum of the “almost-covariant” actions). Because of the presence of the $`SU(2)_{diag}`$ gauge fields in (5.39) which couple to the physical bosonic fields of $`q^{(1,1)}`$, such couplings should be very restrictive.
Let us summarize the above ways of descending from the “master” conformal BW multiplet to the BW-I multiplet.
A. The (32 + 32) field representation. This option corresponds to the use of the pure gauge nonlinear multiplet $`N^{(2,0)}`$, $`N^{(0,2)}`$ as the conformal compensator. One imposes the covariant constraints (5.2) which imply some specific form for the analytic vielbeins $`H^{(4,0)}`$, $`H^{(0,4)},H^{(2,2)},\stackrel{~}{H}^{(2,2)}`$. After properly fixing the gauges, one ends up with the BW-I multiplet and an additional (32+32) off-shell $`U(1)`$ gauge multiplet (5.14) represented by the analytic superfield strength $`H^{(2,2)}`$ (5.12). The general action of the TM-II superfields $`q^{(1,1)}`$ in the background of this representation is given by eq. (5.22). No action for the compensator $`N^{(2,0)}`$, $`N^{(0,2)}`$ itself is assumed. A version with the additional constraints (5.15) yields the pure BW-I multiplet, with no extra multiplets.
B. The (64+64) field representation. This case corresponds to assuming an invariant action for the $`N^{(2,0)}`$, $`N^{(0,2)}`$ compensator. It is constructed using two copies of such superfields (eqs. (5.28), (5.30)). Only one set from this pair is the genuine compensator. As in the previous case, after gauging this compensator away, one ends up with the BW-I multiplet and the (32+32) multiplet (5.14). One more (32+32) off-shell multiplet is $`\stackrel{~}{N}^{(2,0)}`$, $`\stackrel{~}{N}^{(0,2)}`$ (5.29),which is the remnant of the two original copies of $`N`$-multiplets.
C. The (40+40) field representation. In this scheme, in order to construct the invariant action for the compensator $`N^{(2,0)}`$, $`N^{(0,2)}`$, one uses the hypothetical “almost-covariant” action for one TM-II multiplet $`q^{(1,1)}`$ which is a gauged extension of the $`N=(4,4)`$, $`SU(2)`$ WZW action (2.17). The invariant action of two multiplets is given by (5.31). After gauging away the $`N`$-compensator, one is left with the (0+0) BW-I multiplet, the (32+32) $`U(1)`$ multiplet (5.14) and the (8+8) TM-II multiplet $`q^{(1,1)}`$.
D. The (16+16) field representation. This option is different from the preceding ones, as it uses $`q^{(1,1)}`$ as a compensator for the SG-I group. The invariant action (5.32) is given by the difference of two “almost-covariant” $`q^{(1,1)}`$ actions. In the gauge with all possible symmetries of the “master” SG group being compensated for, the surviving field representation consists of the BW-I multiplet, the (8+8) $`SU(2)`$ gauge multiplet (5.39) and the extra (8+8) TM-II multiplet $`q^{(1,1)}`$ which was added to set up the action (5.32).
There still remain the questions as to, how to descend to another, smaller conformal $`N=(4,4)`$, $`SU(2)`$ SG group, i.e. the SG-II group, and how to reproduce the known - and, perhaps, the new $`N=(4,4)`$ Poincaré SG multiplets, by continuing the above process of compensation.
The answer to the first question is as follows. As was already mentioned, there should be a “democracy” between different $`SU(2)`$ factors in the automorphism group $`SO(4)_L\times SO(4)_R`$ of the $`N=(4,4)`$, $`2D`$ Poincaré superalgebra. This implies the existence of “mirror” counterparts of the superconformal matter multiplets discussed so far, such that the roles of the $`SU(2)`$ groups acting on the doublet indices $`i,a`$ and $`\underset{¯}{i},\underset{¯}{a}`$ are switched. An example of such a correspondence is the TM-I multiplet , the off-shell field content of which is given by $`q^{\underset{¯}{i}\underset{¯}{a}},\psi _k^{\underset{¯}{i}},\chi _b^{\underset{¯}{a}},F_{ia}`$ that should be compared with the field content of TM-II (5.34). A similar mirror counterpart should exist for the nonlinear multiplet $`N^{(2,0)},N^{(0,2)}`$. It is natural to call the latter NM-II, with respect to the $`N=(4,4),SU(2)`$ SG-II group acting on the harmonic variables. Then, with respect to the same SG group, the mirror counterpart can be called NM-I. It seems plausible to conjecture that these mirror TM-I and NM-I multiplets (being, in fact, TM-II and NM-II with respect to the $`N=(4,4),SU(2)`$ SG-I group), can be employed to compensate the “master” conformal SG group just down to the SG-II group, quite analogously to how TM-II and NM-II can be used for compensating the “master” group down to the SG-I group. Some of these mirror matter multiplets, in the rigid case, admit a description in the $`SU(2)\times SU(2)`$ harmonic superspace , so we can hope to find their locally supersymmetric versions, cousins of the actions considered above.
To clarify the second question, let us come back to the action (5.32) and assume that the “master” conformal SG group is reduced in it “by hand” to the SG-II one (taking for granted that a nonlinear version of the truncation conditions (4.17) exists). One of the $`q^{(1,1)}`$ superfields can still be used as a compensator. The linearized, purely shift part of the transformation laws of its components, under the action of the residual group (4.20), can be obtained by the substitution of (4.16) into (5.36)
$`\delta \widehat{q}^{ia}={\displaystyle \frac{1}{2}}c^{ia}(_{++}\lambda ^{++}+_{}\lambda ^{})c_j^a\lambda _L^{(ji)}c_b^i\lambda _R^{(ba)},`$
$`\delta \psi _{+\underset{¯}{i}}^a=2ic_i^a_{++}\lambda _{\underset{¯}{i}}^{+i},\delta \chi _{\underset{¯}{a}}^i=2ic_b^i_{}\lambda _{\underset{¯}{a}}^b,`$
$`\delta F_{\underset{¯}{i}\underset{¯}{a}}=0.`$ (5.40)
One sees that the two chiral $`SU(2)`$s are reduced to the diagonal $`SU(2)`$, like in the case (5.36), (5.37), (5.38), by gauging away the triplet part of $`\widehat{q}^{ia}`$. However, the singlet part cannot be gauged away; it becomes just the third component of zweibein. Analogously, $`\psi _{+\underset{¯}{i}}^a`$, $`\chi _{\underset{¯}{a}}^i`$, together with $`h_{}^{+i\underset{¯}{i}},h_{}^{a\underset{¯}{a}}`$ from the BW-II multiplet (4.19), are combined into the 16-component $`N=(4,4)`$ Poincaré SG gravitino (the indices $`i`$ and $`a`$ now refer to the same diagonal $`SU(2)`$). Eventually, bearing in mind the auxiliary field $`F_{\underset{¯}{i}\underset{¯}{a}}`$, we end up just with the $`(8+8)`$ off-shell content of the minimal $`N=(4,4)`$, $`2D`$ Poincaré SG representation . However, recalling that the invariant action (5.32) includes one more $`q^{(1,1)}`$, the total off-shell representation for this case is (16+16). This off-shell content coincides with that of the “TM $`N=4`$ superstring” considered in .
Analogously, one can use the nonlinear multiplet NM-II as a compensator from $`N=(4,4)`$, $`SU(2)`$ SG-II down to some Poincaré SG. The resulting version involves (32+32) off-shell components; its interesting feature is that both conformal $`SU(2)`$ symmetries turn out to be fully compensated for, and $`h_{++}^{(ik)},h_{}^{(ab)}`$ in (4.19) (and in its left counterpart) cease to be gauge fields. The full off-shell content, taking into account an additional $`q^{(1,1)}`$ multiplet needed to construct the invariant action as in (5.31), is (40+40). This coincides with the off-shell content of the “relaxed hypermultiplet $`N=4`$ superstring” of ref. . It is interesting to inquire whether the latter representation is indeed identical to ours.
At last, one can start from the action (5.28) and recover a version of Poincaré SG with (64+64) off-shell fields. Once again, in this version both $`SU(2)`$ symmetries are fully compensated for.
In accord with the previous discussion, various mirror versions of the Poincaré SG can be obtained, starting from the $`N=(4,4)`$, $`SU(2)`$ SG-I and making use of the multiplets TM-I and NM-I as compensators. The various patterns of descent from the “master” $`N=(4,4)`$ SG to the SG-I described above, as well as their hypothetical mirror cousins, seem also to admit further compensations down to the $`N=(4,4)`$ Poincaré SG representations, along similar lines. New possibilities can arise while simultaneously using both types of matter multiplets, i.e. the types I and II, as compensators.
Finally, let us note that there exists a dual version of the rigidly supersymmetric $`q^{(1,1)}`$ actions, including the $`N=(4,4)`$ WZW one (2.17), in terms of unconstrained $`SU(2)\times SU(2)`$ harmonic analytic superfields with infinite numbers of auxiliary fields . This should obviously generalize to the case of local SUSY, which in turn suggests the existence of new versions of Poincaré $`N=(4,4)`$, $`2D`$ SG with infinite sets of auxiliary fields.
A thorough analysis of all these possibilities can be a good program for a future study.
## 6 Conclusions
In this paper we constructed a new sort of $`N=(4,4)`$, $`2D`$ conformal SG gauge multiplet, i.e. the Beltrami-Weyl multiplet, starting from the group of diffeomorphisms in the $`SU(2)\times SU(2)`$ analytic harmonic superspace. This multiplet can be regarded as the result of gauging the most extensive rigid $`N=(4,4)`$ superconformal $`2D`$ group, i.e. the product of two light-cone copies of the infinite-dimensional “large” $`SO(4)\times U(1)`$, $`N=4`$ superconformal group. The previously known $`N=(4,4)`$ conformal SG groups and the corresponding Weyl multiplets were argued to follow from the new “master” SG group and BW multiplet upon their various truncations and compensations, with making use of the appropriate superconformal matter multiplets. Also, various versions of $`N=(4,4)`$ Poincaré SG can be recovered.
There still remain a few important conceptual and technical points to be fully elaborated on. This concerns, before all, constructing the full nonlinear version of the “almost-covariant” $`q^{(1,1)}`$ action (5.25) and the nonlinear completion of the constraints (4.17), (4.15), as well as revealing the component fields structure of the locally supersymmetric superfield actions presented. An important problem is to incorporate into the present scheme mirror counterparts of the superconformal multiplets employed in this paper and to study the relevant compensation patterns. Different $`N=(4,4)`$ SG-matter couplings correspond to various versions of $`N=(4,4)`$ superstrings . It would be interesting to inquire the quantum properties of the systems described here, e.g., along the lines of refs. , . Note that the rigid $`N=(4,4)`$ WZW action (2.17) admits an extension to the $`N=(4,4)`$ WZW-Liouville one , with breaking the $`N=(4,4),SO(4)\times U(1)`$ superconformal invariance down to the type-II $`N=(4,4),SU(2)`$ one. Such a Liouville extension plays an important role in the quantum case . It is of interest to inquire whether a locally supersymmetric extension of the Liouville term can be constructed in $`SU(2)\times SU(2)`$ harmonic superspace.
## Acknowledgments
S.B. wishes to thank JINR-Dubna for hospitality at the early stages of this research. E.I. thanks Jim Gates for enlightening correspondence and INFN-LNF for the hospitality multiply extended to him during the course of this long-term work. This research was supported in part by the Fondo Affari Internazionali Convenzione Particellare INFN-JINR, Project PAST-RI 99/01, RFBR Grant 99-02-18417, RFBR-CNRS Grant 98-02-22034, NATO Grant PST.CLG 974874 and INTAS Grants INTAS-96-0538, INTAS-96-0308.
## Appendix: A simple example of the component action
Here, just to give a feeling how the locally supersymmetric actions presented in this paper look in terms of component fields, we quote the free part of the general conformal SG-I group-invariant action (4.14)
$$S_q^{free}=\mu ^{(2,2)}\widehat{\mathrm{\Omega }}q^{(1,1)}q^{(1,1)}.$$
(A.1)
It will be convenient to choose a gauge for the analytic vielbeins which is slightly different from (4.9)
$`H^{(2,0)++}=H^{(3,0)\underset{¯}{i}}=0,H^{(2,0)}=i(\theta ^{(1,0)})^2\widehat{h}_{++}^{}(z,v,\theta ^{(0,1)}),`$
$`H^{(2,1)\underset{¯}{a}}=i(\theta ^{(1,0)})^2\widehat{h}_{++}^{(0,1)\underset{¯}{a}}(z,v,\theta ^{(0,1)}).`$ (A.2)
This gauge is also globally well-defined. To simplify the situation as soon as possible, we recall that all components of the BW-I multiplet are locally pure gauge, and we choose the additional gauge, which is admissible only locally,
$$h_{++}^{}=h_{++}^{(\underset{¯}{a}\underset{¯}{b})}=h_{++}^{a\underset{¯}{a}}=0$$
(A.3)
(we could alternatively choose the left counterparts of (A.3) to vanish). It is easy to show that the full solution of the constraints (4.8) in this gauge is given by
$`^{(2,0)}`$ $`=`$ $`^{(2,0)}+i(\theta ^{(1,0)})^2_{++},^{(0,2)}=^{(0,2)}+i(\theta ^{(0,1)})^2_{},`$
$`_{}`$ $`=`$ $`_{}+\{h_{}^{++}2i\theta _{\underset{¯}{i}}^{(1,0)}h_{}^{+k\underset{¯}{i}}u_k^{(1,0)}\}_{++}`$ (A.4)
$`+`$ $`\{h_{}^{+k\underset{¯}{i}}u_k^{(1,0)}+\theta ^{(1,0)\underset{¯}{k}}[h_{\underset{¯}{k})}^{(\underset{¯}{i}}+{\displaystyle \frac{1}{2}}\delta _{\underset{¯}{k}}^{\underset{¯}{i}}_{++}h_{}^{++}]`$
$``$ $`(\theta ^{(1,0)})^2_{++}h_{}^{+k\underset{¯}{i}}v_k^{(1,0)}\}{\displaystyle \frac{}{\theta ^{(1,0)\underset{¯}{i}}}}.`$
It is easy to explicitly check the integrability condition
$$[^{(2,0)},^{(0,2)}]=0.$$
The residual gauge symmetry of (A.4) is given by (4.11), with all parameters being functions of only $`z^{}`$ (this is just the right $`N=4,SU(2)`$ SCA-I), and by the left counterpart of (4.11), with the parameters still being general functions of both coordinates $`z^{\pm \pm }`$. It is easy to check that under this group
$$\delta \mu ^{(2,2)}=0,$$
so one can expect $`\widehat{\mathrm{\Omega }}\text{const}`$ in this gauge. This is indeed so, because it is easy to check that
$$\mathrm{\Gamma }^{(2,0)}=\mathrm{\Gamma }^{(0,2)}=0$$
(A.5)
for the vielbeins in (A.4). Then, the action (A.1) is
$$S_q^{free}=\mu ^{(2,2)}q^{(1,1)}q^{(1,1)},$$
(A.6)
with
$$^{(2,0)}q^{(1,1)}=^{(0,2)}q^{(1,1)}=0.$$
(A.7)
Being aware of the explicit expressions for $`^{(2,0)},^{(0,2)}`$, it is easy to directly solve these constraints in terms of the physical fields of $`q^{(1,1)}`$ defined in (5.34) and the SG fields. This is rather straightforward, so we quote only the final form of the action. It is obtained by substituting this solution into (A.6) and integrating there over the $`\theta `$’s and the harmonics:
$`S_q^{free}`$ $`=`$ $`{\displaystyle }d^2z\{_{++}q_{ia}(\widehat{}_{}q^{ia}+h_{}^{+i\underset{¯}{i}}\psi _{+\underset{¯}{i}}^a)+{\displaystyle \frac{i}{2}}\chi _{}^{i\underset{¯}{a}}_{++}\chi _{i\underset{¯}{a}}+{\displaystyle \frac{1}{4}}F^{\underset{¯}{i}\underset{¯}{a}}F_{\underset{¯}{i}\underset{¯}{a}}`$ (A.8)
$`+`$ $`{\displaystyle \frac{i}{2}}\psi _+^{a\underset{¯}{i}}[\widehat{}_{}\psi _{+a\underset{¯}{i}}+(h_{\underset{¯}{i})}^{(\underset{¯}{k}}+{\displaystyle \frac{1}{2}}\delta _{\underset{¯}{i}}^{\underset{¯}{k}}_{++}h_{}^{++})\psi _{+a\underset{¯}{k}}2ih_{\underset{¯}{i}}^{+k}_{++}q_{ka}]\},`$
where
$$\widehat{}_{}=_{}+h_{}^{++}_{++}.$$
Note that (A.8) is just the action of the $`N=4`$ chiral bosons constructed in ref. (up to switching the $`+`$ and $``$ light-cone indices), with the residual local $`N=4`$ SUSY as the relevant Siegel symmetry. |
warning/0003/astro-ph0003087.html | ar5iv | text | # Kinematics of Metal-Poor Stars in the Galaxy. III. Formation of the Stellar Halo and Thick Disk as Revealed from a Large Sample of Non-Kinematically Selected Stars
## 1 Introduction
Over the past few decades, studies of the luminous halo population of metal-deficient field stars and globular clusters have provided an increasingly detailed picture of the formation and evolution of the Galaxy. Because the time required for mixing of the initial phase-space distribution of these objects, via exchange of energies and angular momenta, is thought to exceed the age of the Galaxy, kinematic information obtained at the present enables one to elucidate the initial dynamical conditions under which these objects were born. To the extent one is able to estimate ages of these halo population objects, either directly (which is difficult at present), or indirectly (by postulating that the ensemble metallicities of these objects increases with time), their formation history is obtainable as well. Thus, the dynamical and chemical state of these halo-population objects provides important information on how the Galaxy has developed its characteristic structures during the course of its evolution.
Almost forty years ago, the “canonical” scenario of the early dynamical evolution of the Galaxy was put forward by Eggen, Lynden-Bell, & Sandage (1962, hereafter ELS) to explain what they believed to be an observed correlation between the orbital characteristics and metal abundances of stars in the solar neighborhood. Focusing on the lack of metal-poor stars with low eccentricity orbits in their (proper-motion selected) sample, ELS argued that the Galaxy must have undergone a rapid collapse, then formed a rotationally supported disk. Although criticism of the ELS model has been levied because of the potential influence of their kinematic selection bias, especially the extent to which this might alter the derived collapse timescale of the early Galaxy (Yoshii & Saio 1979; Norris, Bessell, & Pickles 1985; Norris 1986; Norris & Ryan 1991; Beers & Sommer-Larsen 1995, hereafter BSL; Chiba & Yoshii 1998, hereafter CY), the ELS collapse picture has long been influential for studies of disk galaxies like our own, and for elliptical galaxies as well (e.g., Larson 1974; van Albada 1982).
An alternative picture for the formation of the Galactic halo was proposed by Searle & Zinn (1978, hereafter SZ), who noted a number of difficulties in reconciling several observed properties of the halo globular cluster system with predictions of the ELS model. Among these, the existence of a large (several Gyr) spread in the inferred ages of the Galactic globulars, and the lack of an abundance gradient with distance from the Galactic center were thought to be the most crucial. SZ suggested that, in its earliest epochs, the halo component of the Galaxy may have experienced prolonged, chaotic accretion of subgalactic fragments. More recent studies of halo field stars also provide evidence which may support the SZ picture, including a possible age spread among halo subdwarfs (e.g., Schuster & Nissen 1989; Carney et al. 1996), a gradient in the inferred ages of field horizontal-branch (FHB) stars and RR Lyrae variables, in the sense that the outer halo FHB stars and RR Lyrae variables appear several Gyr younger than those of the inner halo (Preston, Shectman, & Beers 1991; Lee & Carney 1999), a report of the apparent clustering of FHB stars in the halo (Doinidis & Beers 1989), possible kinematic substructure in the halo (e.g., Doinidis & Beers 1989; Majewski, Munn, & Hawley 1994; 1996), and distinct changes in the kinematics of the field populations as one moves from the inner to outer halo (e.g., Majewski 1992; Carney et al. 1996; Sommer-Larsen et al. 1997).
In order to assess which picture, ELS or SZ (or both, e.g., Norris 1994; Freeman 1996; Carney et al. 1996), more correctly describes the early history of the Galaxy, we require a large and reliable set of data for halo-population objects chosen with criteria that do not unduly influence the subsequent analysis. As we show, the analysis of stars chosen with a kinematic selection bias can be particularly troublesome. Interestingly, in order to obtain adequately large samples of stars exhibiting a range of metallicities in the solar neighborhood, an abundance bias is actually required, otherwise the exceedingly rare very low-metallicity stars of the halo population will not be represented in sufficient numbers. One must exercise caution, however, that abundance estimates for stars in the sample under consideration are as accurate as possible, due to the presence of an overlap of the local halo with the relatively high density thick-disk population (e.g., Anthony-Twarog & Twarog 1994; BSL; Ryan & Lambert 1995; Twarog & Anthony-Twarog 1996; CY). It is similarly important to assemble a large and homogeneously analyzed sample, both to minimize statistical fluctuations in the derived kinematic quantities, and to reduce the effects of other systematic errors (such as might arise in estimates of stellar distances).
In this paper we present an analysis of the kinematics of metal-deficient field stars in the solar neighborhood, based on a large catalog of stars selected without kinematic bias (Beers et al. 2000, hereafter Paper II). This catalog, consisting of 2041 stars from the published literature with abundances \[Fe/H\]$`0.6`$, includes updated stellar positions, newly derived homogeneous distance estimates, revised radial velocities, and refined metal abundance estimates. Moreover, a subset of some 1200 stars in the catalog now have available proper motions, taken from a variety of recently completed proper motion catalogs. We note that this catalog is (by far) the largest sample of metal-deficient field stars with available proper motions among any previously assembled non-kinematically selected samples. Thus, it is now possible to draw a much more definitive picture of the early kinematic evolution of the Galaxy.
Our paper is organized as follows. In §2 we present a discussion of the detailed velocity distributions of our sample stars, concentrating on those presently located in the solar neighborhood. In §3 we analyze the orbital motions of these stars. In §4 we consider the global character of the halo of the Galaxy, as deduced from the kinematics of a local sample. In §5, we further examine evidence for kinematic substructure in the phase-space distribution of the halo. Finally, in §6, the results of the present work are summarized, and their implications for the formation and evolution of the Galaxy are discussed.
## 2 Velocity Distributions of the Metal-Poor Stars
### 2.1 Individual and Mean Space Velocities
Paper II of this series presented proper motions for 1214 stars with \[Fe/H\] $`0.6`$, as well as for a number of slightly more metal-rich stars. Within this sample, 1203 stars with \[Fe/H\] $`0.6`$ have distance estimates and radial velocities as well as proper motions, so that the full three-dimensional velocities are directly calculable. Figure 1 is a reproduction of Figure 8 from Paper II, and shows the local velocity components vs. \[Fe/H\] for these 1203 stars. The velocity components $`U`$, $`V`$, and $`W`$ are directed to the Galactic anticenter, rotation direction, and north Galactic pole, respectively, and have been corrected for the local solar motion $`(U_{},V_{},W_{})=(9,12,7)`$ km s<sup>-1</sup> with respect to the local standard of rest (LSR) (Mihalas & Binney 1981). Note the excellent coverage of this sample over the entire range of Galactic abundances, especially below \[Fe/H\] $`=2`$ and above \[Fe/H\] $`=1`$, where previous studies that made use of, for example, the sample of metal-poor stars studied with Hipparcos, lack sufficient numbers of stars (compare, e.g., with Figure 4 of CY).
To examine the characteristic local velocity distributions of our sample, we confine ourselves to a discussion of the stars for which the Galactocentric distance along the plane, $`R`$, is between 7 and 10 kpc, and those for which the distance from the Sun, $`D`$, is within 4 kpc. Six stars in this subsample have large rest-frame velocities, $`V_{RF}>550`$ km s<sup>-1</sup>, that are in excess of the canonical escape velocity in the solar neighborhood ($`V_{esc}500550`$ km s<sup>-1</sup>; Carney, Latham, & Laird 1988). Although some of the space velocities may indeed be this high, the majority of these stars probably have large $`V_{RF}`$ due to an over-estimation of their distances, which has artificially inflated their estimated tangential velocities. We choose to remove these extreme-velocity stars by placing an additional limit of $`V_{RF}550`$ km s<sup>-1</sup> on the sample. The stars satisfying the above selection criteria are referred to as the “Selected Sample” in the following discussion.
We first calculate the mean velocities $`(<U>,<V>,<W>)`$ and velocity dispersions $`(\sigma _U,\sigma _V,\sigma _W)`$ for the Selected Sample – values for five characteristic ranges in metal abundance are listed in Table 1a. Velocity dispersions are estimated from the standard deviations, after correction for the typical measurement errors in the velocities ($`10`$ km s<sup>-1</sup>). Figure 2 shows $`(\sigma _U,\sigma _V,\sigma _W)`$ as a function of \[Fe/H\]. In this Figure we have adopted a finer binning in metal abundance; the dispersion measurements in each bin are listed in Table 1b. The filled and open circles in Figure 2 denote the stars at $`|Z|<1`$ kpc and $`|Z|<4`$ kpc, respectively, where $`Z`$ is the height above the Galactic plane.
The most metal-deficient stars in the Selected Sample, those more metal-poor than \[Fe/H\] $`=2.2`$, are dominated by members of the halo population. For $`|Z|<1`$ kpc, these stars exhibit a radially elongated velocity ellipsoid $`(\sigma _U,\sigma _V,\sigma _W)=(141\pm 11,106\pm 9,94\pm 8)`$ km s<sup>-1</sup>, in good agreement with previous results (e.g., BSL; CY). With a slightly more metal-rich cut on the abundances, i.e., selecting stars with \[Fe/H\] $`<1.7`$, we obtain similar values, $`(148\pm 7,110\pm 5,92\pm 4)`$ km s<sup>-1</sup>, so it appears that the shape of the velocity ellipsoid remains essentially unchanged with varying \[Fe/H\] below \[Fe/H\] $`=1.7`$. In this regard Norris (1994) claimed, from his analysis of a sample of high proper-motion stars, that $`\sigma _W`$ continues to increase with decreasing \[Fe/H\], even at its lowest levels. This result was taken to indicate the possible existence of a dynamically “hot” proto-disk population at low abundances. Carney et al. (1996) disputed this result, as their analysis of a different set of high proper-motion stars indicated that the disk component is not dynamically hot, at least when membership is confined to the stars orbiting exclusively within the inner part of the Galaxy, $`R14`$ kpc. Figure 2c shows no evidence for an increase of $`\sigma _W`$ at low abundances.
The velocity dispersion components of the Selected Sample in the more metal-rich abundance ranges decrease as the contribution of the thick disk component progressively increases. In particular, for $`0.7`$ \[Fe/H\] $`<0.6`$ and $`|Z|<1`$ kpc, where the contribution of the halo component is expected to be negligible, the mean $`V`$ velocity, $`<V>`$, is $`20\pm 5`$ km s<sup>-1</sup>; the velocity dispersions are $`(\sigma _U,\sigma _V,\sigma _W)=(46\pm 4,50\pm 4,35\pm 3)`$ km s<sup>-1</sup> . This result is in agreement with previously derived kinematic parameters for the thick disk, which appears to be in rapid rotation ($``$200 km s<sup>-1</sup>), provided the rotational speed of the LSR is $`V_{LSR}=220`$ km s<sup>-1</sup> (BSL). With these values for the thick disk kinematics, it is possible to estimate the radial scale length of this component using the following formula (Binney & Tremaine 1987):
$$2V_{LSR}V_{lag}V_{lag}^2=\sigma _U^2\left(1+\frac{\sigma _V^2}{\sigma _U^2}+2\frac{R}{h_R}\right),$$
(1)
where $`V_{lag}`$ is the asymmetric drift given by $`V_{lag}=<V>`$, and $`h_R`$ is the scale length of the disk, provided its density varies as $`\mathrm{exp}(R/h_R)`$. By inserting $`V_{LSR}=220\pm 10`$ km s<sup>-1</sup>, the assumed solar radius $`R=8.5`$ kpc (Kerr & Lynden-Bell 1986) and derived parameters for $`0.7`$ \[Fe/H\] $`<0.6`$ in Equation (1), we obtain $`h_R=4.5\pm 0.6`$ kpc. This is in good agreement with $`h_R=4.7\pm 0.5`$ kpc obtained by BSL, and also with the lower limit of $`h_R4.5`$ kpc derived by Ratnatunga & Freeman (1989). We note that the second term on the left-hand side of Equation (1) has been omitted in some previous works, as it is small compared to other terms. If we were to exclude this term, we would obtain $`h_R=4.3\pm 0.6`$ kpc, thereby slightly underestimating $`h_R`$.
### 2.2 Rotational Character of the Selected Sample
We now examine the rotational character of the Selected Sample. Figure 3 shows the mean rotational velocities $`<V_\varphi >`$ (the rotation velocity in a cylindrical coordinate frame) as a function of \[Fe/H\], based on stars in the Selected Sample. The left-hand panel in Figure 3a displays the results for the abundance ranges listed in Table 2, for three subsets of the sample as a function of distance above the plane. In the right-hand panel of Figure 3a, the bins are obtained by passing a box of width $`N=100`$ stars, ordered by metallicity, with an overlap of 20 stars each. The latter approach is adopted to avoid any effects of the arbitrary placement of bins on the results.
Note that the panels in Figure 3a are obtained with full knowledge of the space motions of the stars in the Selected Sample. Figure 3b is based on the radial velocities alone, applying the methodology of Frenk & White (1980) (FW) (see also Norris 1986; Morrison, Flynn, & Freeman 1990, hereafter MFF; BSL). The solid lines in the left and right-hand panels of Figure 3b denote $`<V_\varphi >`$ as derived for the stars with available proper motions (denoted as $`<V_\varphi >_{pm}^{FW}`$ in Table 2), i.e., the same sample as in Figure 3a, whereas the dashed lines are for all of the stars with available radial velocities ($`<V_\varphi >_{all}^{FW}`$ in Table 2). Comparison between Figures 3a and 3b allows us to examine whether the Selected Sample is subject to any significant kinematic bias, since, if so, the $`<V_\varphi >`$ derived from the space motions is expected to be systematically smaller from that determined on the basis of radial velocities alone (Ryan & Norris 1991; Norris & Ryan 1991).
Figure 3a clearly indicates that the rotational properties of the Selected Sample change discontinuously at \[Fe/H\] $`1.7`$. Stars with \[Fe/H\] $`<1.7`$ exhibit no systematic variation of $`<V_\varphi >`$ with decreasing \[Fe/H\]. It is interesting to note that the subsample of low-abundance stars with $`|Z|<1`$ kpc show a rather large prograde rotation of $`<V_\varphi >=3050`$ km s<sup>-1</sup>, and that $`<V_\varphi >`$ decreases if stars at larger heights are considered, at least for the two abundance bins centered on \[Fe/H\] $`=1.9`$ and $`2.2`$, respectively. This behavior is not exhibited, however, for stars in the lowest abundance bin. To check the significance of this feature, we have combined the stars in our Selected Sample in the metallicity interval $`2.4[\mathrm{Fe}/\mathrm{H}]1.9`$ and obtained $`<V_\varphi >`$ by sweeping a box of 50 stars ordered by $`|Z|`$, with an overlap of 30 stars. The results are summarized in Table 3 and depicted in Figure 4. The lower solid line in Figure 4 is a least-squares fit to the data, which yields $`\mathrm{\Delta }<V_\varphi >/\mathrm{\Delta }|Z|=52\pm 6`$ km s<sup>-1</sup> kpc<sup>-1</sup>, indicating the presence of a significant vertical gradient in $`<V_\varphi >`$ at low abundances. Figure 4 also suggests that $`<V_\varphi >`$ beyond $`|Z|1.2`$ kpc has a nearly constant zero value. If we exclude the last point at $`|Z|=1.76`$ kpc from the fit, we obtain $`\mathrm{\Delta }<V_\varphi >/\mathrm{\Delta }|Z|=62\pm 5`$ km s<sup>-1</sup> kpc<sup>-1</sup> (shown as a dotted line). Note that Majewski (1992) reported evidence for a halo component which is in retrograde motion ($`V_\varphi =275\pm 16`$ km s<sup>-1</sup>), but which exhibited no gradient of rotation with distance from the plane, a result that is clearly at odds with our present result.
Figure 3a shows that for stars in the Selected Sample with \[Fe/H\] $`>1.7`$ there is a clear linear dependence of $`<V_\varphi >`$ on \[Fe/H\], a dependence that remains essentially unchanged even if the range of $`|Z|`$ is varied. This discontinuity has been seen before, of course, based on analysis of smaller samples (e.g., Norris 1986; Carney 1988; Zinn 1988; Norris & Ryan 1989; BSL; CY). The inescapable conclusion is that the transition from halo to disk must not have occurred in a continuous manner, as predicted in the ELS model. However, the vertical gradient in $`<V_\varphi >`$ for metal-poor stars noted above suggests that the halo was not formed in a totally chaotic, dissipationless manner as implied in the SZ hypothesis. Rather, dissipational processes may have played a role in the initial contraction of the halo, likely involving energy exchange with the gas phase (Carney et al. 1996).
It is also worth noting from Figure 3 that the stars of the thick disk, which dominate the Selected Sample for \[Fe/H\] $`>1`$, exhibit only a small change in $`<V_\varphi >`$ with increasing $`|Z|`$. Table 3 and Figure 4 summarize the change of $`<V_\varphi >`$ with $`|Z|`$ for stars likely to be dominated by the thick disk, i.e., in the abundance range $`0.8`$ \[Fe/H\] $`0.6`$. The least-squares fit to the data (as shown by upper solid line) yields $`\mathrm{\Delta }<V_\varphi >/\mathrm{\Delta }|Z|=30\pm 3`$ km s<sup>-1</sup> kpc<sup>-1</sup>, much smaller than the gradient obtained for halo stars with $`2.4`$ \[Fe/H\] $`1.9`$, but similar to previous estimates of the thick disk rotational velocity gradients reported by Majewski (1992).
The relation between $`<V_\varphi >`$ and \[Fe/H\] as shown in the panels of Figure 3a is also seen in the panels of Figure 3b. In particular, $`<V_\varphi >`$ obtained from the subsample having available proper motions (dashed lines) is essentially the same as that from the entire sample (solid lines) within standard errors in $`<V_\varphi >`$ (except for \[Fe/H\] $`1.6`$: see below). This is consistent with our argument given in Paper II that the subsample based on stars with available proper motions is not subject to any significant kinematic bias.
We note that near \[Fe/H\] $`=1.6`$ in Figure 3b, $`<V_\varphi >`$ obtained from consideration of the sample with full space motions is larger than that obtained from the sample using radial velocities alone. The apparent retrograde rotation of globular clusters in the similar metallicity range was also reported by Rodgers & Paltoglou (1984) based on radial velocities alone, whereas Dinescu, Girard, & van Altena (1999) found no sign of significant retrograde rotation using full space motions of globular clusters. It is worth noting that this difference in the rotational velocity using either of radial velocities or full space motions is not a signature of kinematic bias, since the result is in the opposite direction to that expected. This difference probably arises due to limitations of the FW methodology as applied to our sample. The FW method implicitly assumes that the angle, $`\psi `$, between the line-of-sight and the rotational direction is randomly distributed in the sample, which may not apply in this case. In addition, the apparent excursion to large retrograde rotation, $`<V_\varphi >40`$ km s<sup>-1</sup> near \[Fe/H\] $`=1.6`$ in the right-hand panel of Figure 3b (also noted by BSL), may be caused by a few outliers having large velocities as seen by an observer at rest with respect to the Galactic center, $`V_{gal}`$. If we exclude five stars having large $`V_{gal}`$ from the bin centered at \[Fe/H\] $`=1.6`$, we obtain $`<V_\varphi >=12`$ km s<sup>-1</sup>. We have verified that the effect of outliers on $`<V_\varphi >`$ is small in other abundance ranges.
The influence of a disk-like population for \[Fe/H\] $`>1.7`$ is certainly suggested by the appearance of Figure 3. The question of the limiting abundance of a so-called metal-weak thick disk (MWTD)<sup>1</sup><sup>1</sup>1These metal-poor stars with disk-like kinematics may also include a considerable portion of the thin disk if its metallicity distribution overlaps that of the thick disk (Wyse & Gilmore 1995). has been considered several times in the past (MFF; Rodgers & Roberts 1993; Layden 1995; BSL; Ryan & Lambert 1995; Twarog & Anthony-Twarog 1996; CY). Figure 5 shows the frequency distribution of $`V_\varphi `$ for the stars in the Selected Sample with available space motions, for subsets chosen to have (a) $`|Z|<1`$ kpc, and (b) $`|Z|1`$ kpc, respectively. At $`|Z|<1`$ kpc, where the disk-like kinematics are expected to be more evident than at larger heights above the plane, the metal-rich stars with \[Fe/H\] $`>1`$ are peaked at $`V_\varphi =200`$ km s<sup>-1</sup>, a rather high rotational velocity. One also sees the presence of a small contribution of the stars with halo-like kinematics. At lower abundances the halo-like kinematics become much more dominant; the contribution of the MWTD is apparently decreasing at lower abundances and higher $`|Z|`$.
To quantify the fraction of the MWTD in our local sample within the specified abundance ranges, we have fit the subset of stars with $`|Z|<1`$ kpc using a mixture of two Gaussian distributions for $`V_\varphi `$, representing the halo and disk populations. The halo kinematic parameters ($`<V_\varphi >_{halo},\sigma _{\varphi ,halo}`$) = (+33,106) km s<sup>-1</sup> are derived from stars with \[Fe/H\] $`2.2`$. The rotation velocity of the disk component, $`<V_\varphi >_{disk}=+200`$ km s<sup>-1</sup>, is obtained considering the stars in the metallicity range $`0.7<`$ \[Fe/H\] $`0.6`$. With these parameters fixed, we evaluate the most likely values of the velocity dispersion of the disk, $`\sigma _{\varphi ,disk}`$ and fraction $`F`$ of the MWTD, using a maximum likelihood analysis (see also MFF; CY). The likelihood function for the stars with $`V_\varphi ^i`$ is given by
$$\mathrm{log}f(F,\sigma _{\varphi ,disk})=\underset{i=1}{\overset{N}{}}\mathrm{log}[Ff_{disk}^i+(1F)f_{halo}^i],$$
(2)
where $`f_{disk}^i`$ ($`f_{halo}^i`$) denote Gaussian functions with mean velocities $`<V_\varphi >_{disk}`$ ($`<V_\varphi >_{halo}`$) and dispersions $`\sigma _{\varphi ,disk}`$ ($`\sigma _{\varphi ,halo}`$). The results of the likelihood analysis are tabulated in Table 4, and shown by the solid curves in Figure 5. The MWTD contributes about 30% of the metal-poor stars in the abundance range $`1.7<`$ \[Fe/H\] $`1`$, which is smaller than the fraction derived by MFF ($``$72%) and BSL ($``$60%), but larger than the result of CY ($``$10%). The fraction of the MWTD is quite modest in the more metal-poor ranges, in contrast to the suggestion of BSL, who argued for $``$30% even at \[Fe/H\] $`<2`$. One reason that our result may differ so strikingly from that of BSL is that the estimated abundances for many of the HK survey stars at \[Fe/H\] $`<2`$, listed in the original BSL catalog, are likely to have been underestimated by $`0.3`$ dex (see Figure 1 of Paper II).
The RR Lyrae stars in our Selected Sample (shown as shaded histograms in Figure 5) exhibit no clear disk-like kinematics, even in the intermediate abundance range of $`1.7<`$ \[Fe/H\] $`1`$. This confirms earlier results by MFF, Layden (1995), and CY, and may imply a somewhat younger age of the thick disk as compared to the halo. However, the numbers of RR Lyraes with available space motions is rather small, so this question, especially in conjunction with their period distributions to investigate different populations of RR Lyraes (Lee & Carney 1999), should be revisited when the sample size has been increased.
## 3 Orbital Properties of the Metal-Poor Stars
In this section we investigate the orbital properties of our sample of stars in a given Galactic potential. We adopt the analytic Stäckel-type potential developed by Sommer-Larsen & Zhen (1990, hereafter SLZ), which consists of a flattened, oblate disk and a nearly spherical massive halo. This model potential is consistent with the mass model of Bahcall, Schmidt, & Soneira (1982), exhibiting a flat rotation curve beyond $`R=4`$ kpc, and having a commensurate local mass density at $`R=R_{}`$. In contrast to a non-analytic potential, for which numerical integrations of orbits are required, the analytic nature of the adopted potential has the great advantage of maintaining clarity in the analysis, as demonstrated below. In the Appendix, we summarize the properties of the Stäckel mass model, and provide expressions for three integrals of motion in such a model (see also de Zeeuw 1985; Dejonghe & de Zeeuw 1988)
### 3.1 The Relationship between Orbital Eccentricity and Metal Abundance
We first compute orbital eccentricities, defined as $`e=(r_{ap}r_{pr})/(r_{ap}+r_{pr})`$, where $`r_{ap}`$ and $`r_{pr}`$ denote the apogalactic and perigalactic distances of the orbits, respectively. These orbital parameters are tabulated in Table 3 of Paper II for the stars under consideration. In Figure 6a, we show the relation between $`e`$ and \[Fe/H\]. As is evident, there is no strong correlation between these quantities, and the metal-poor stars below \[Fe/H\] $`=2`$ exhibit a diverse range in orbital eccentricities. This is in sharp contrast to the ELS result, and confirms previous suggestions from a number of workers, but in a much more definitive manner (Yoshii & Saio 1979; NBP; Carney & Latham 1986; Carney, Latham, & Laird 1990; Norris & Ryan 1991; CY). In addition to the diverse distribution of $`e`$ at all abundances, we note a small concentration of the stars at $`e0.9`$ and \[Fe/H\] $`1.7`$, which is somewhat reminiscent of the original ELS result. It is perhaps not coincidental that the excess number of high-$`e`$ stars occurs at an abundance which matches the sharp discontinuity of $`<V_\varphi >`$ found at \[Fe/H\] $`=1.7`$, where $`<V_\varphi >`$ is almost zero (Figure 3). This may suggest that a significant fraction of the metal-poor stars with abundances near \[Fe/H\] $`=1.7`$ formed from infalling gas of this metallicity during an early stage of Galaxy formation, in a manner similar to an ELS collapse.
In Figure 6b, we show the mean eccentricity, $`<e>`$, vs. \[Fe/H\], where the bins are obtained by passing a box of width $`N=100`$ stars, ordered by metallicity, with an overlap of 20 stars each. For comparison, the dashed line denotes the result of Carney et al. (1996) for their high proper motion sample. There is a clear difference from our results – the use of a kinematically selected sample overestimates the average orbital eccentricities at a given \[Fe/H\], by an amount up to 0.2.
Figure 7 shows the cumulative distributions of $`e`$, $`N(<e)`$, in two specific abundance ranges, (a) for \[Fe/H\] $`2.2`$, and (b) for $`1.4<`$ \[Fe/H\] $`1`$. Figure 7a clearly demonstrates that even at quite low abundance, roughly 20% of our stars have $`e<0.4`$. The different lines correspond to the cases when the range of $`|Z|`$ is changed. It is apparent that, for stars with \[Fe/H\] $`2.2`$, the cumulative distribution function of $`e`$ is unchanged when considering subsets of the data with a range of $`|Z|`$, suggesting the absence of any substantial disk-like component below this metallicity. By way of contrast, Figure 7b shows that stars with intermediate abundances exhibit (a) a higher fraction of orbits with $`e<0.4`$ than for the lower abundance stars, (b) a decrease in the relative fraction of low eccentricity stars as larger heights above the Galactic plane are considered, and (c) convergence at larger heights to a fraction which is close to the 20% obtained for the lower abundance stars. These results imply that the orbital motions of the stars in the intermediate abundance range are, in part, affected by the presence of thick-disk component with a finite scale height. We recall that CY and Chiba, Yoshii & Beers (1999), using a sample of metal-poor stars with Hipparcos measurements, found a further decrease of the fraction of the stars with $`e<0.4`$ at larger $`|Z|`$, without achieving the convergence noted here (see Figure 15 of CY). This was presumably due to the lack of a sufficient number of intermediate abundance stars at large $`|Z|`$ in the sample considered by these authors.
### 3.2 Structural Parameters of the MWTD Component
We now seek to quantitatively describe the abundance range, scale height, and fraction of the MWTD component in our sample. We apply a Kolmogorov-Smirnoff (KS) test of the null hypothesis that the differential distributions of $`e`$, $`n(e)`$, for stars in a specified abundance range, are drawn from the same parent population of eccentricities as stars belonging to a “pure” halo component. Based on our analysis above, we take the subsample of 78 stars with \[Fe/H\] $`2.2`$ and $`|Z|<1`$ kpc to represent the pure halo component. We then calculate the KS probabilities, $`P_{KS}`$, for the stars in various intermediate abundance ranges and with $`|Z|>Z_{lim}`$, where $`Z_{lim}`$ is the lower limit on the heights of the stars above the Galactic plane. We expect that, even if the specified abundance range is contaminated by stars with disk-like kinematics, it will be dominated by halo-like kinematics above some $`|Z|=Z_{lim}`$, with $`P_{KS}`$ exceeding 0.2 (i.e., the subsamples being consistent with draws from the same parent population of orbital eccentricities).
Figure 8 shows the results of the KS tests. In order to obtain an estimate for the value of $`Z_{lim}`$ above which the populations cannot be distinguished, Figure 8a depicts the results for the stars below \[Fe/H\] $`=1`$, but above the specified lower limit for the abundance. In the left-hand panel of Figure 8a, it is seen that $`P_{KS}`$ rapidly increases at $`Z_{lim}=0.5`$ to 1.3 kpc and then remains roughly constant at larger $`Z_{lim}`$. In the right-hand panel of Figure 8a, the distribution of $`P_{KS}`$ on $`Z_{lim}`$ changes dramatically as one passes from the inclusion of stars with \[Fe/H\] $`>1.9`$ to those with \[Fe/H\] $`>2.0`$. When stars with abundances as low as \[Fe/H\] $`=2.0`$ are included, there is no value of $`Z_{lim}`$ for which the distributions can be distinguished. To identify the lower limit on the abundance of stars which are members of the MWTD, Figure 8b shows the distribution of $`P_{KS}`$ for metal-poor stars above \[Fe/H\] $`=2.2`$, but below the specified upper limit on abundance. As is seen in the left-hand panel of this figure, $`P_{KS}`$ remains small (indicating that the populations can be distinguished) at all $`Z_{lim}`$. In the right-hand panel of Figure 8b, one sees that although there exists a region at small $`Z_{lim}`$ where the populations can be marginally distinguished when the upper limit on abundance is taken to be \[Fe/H\] $`=1.9`$, there is no such region when stars with an upper limit of \[Fe/H\] $`=2.0`$ is considered. These results suggest that the MWTD component has a characteristic scale height of roughly 1 kpc, above which halo-like orbital motions dominate, and a lower abundance limit near \[Fe/H\] $`=2.0`$. We note here that the number of stars employed in the KS tests is sufficiently large in the ranges of $`Z_{lim}`$ considered (e.g. $`N=240`$ for $`2<`$ \[Fe/H\] $`1`$ and $`N=54`$ for $`2.2<`$ \[Fe/H\] $`2`$ at $`Z_{lim}=1`$ kpc). However, it would be useful, especially at larger $`Z_{lim}`$, to boost the sample sizes so that more detailed investigations can be carried out.
We now estimate the contribution of the MWTD component in the solar neighborhood, $`F`$, using the distribution of $`e`$ in various abundance ranges. Following the methodology developed by CY, we perform a Monte Carlo simulation to predict the $`e`$-distribution from a mixture of stars contributed by the thick-disk and halo populations, adopting the kinematic parameters for these components derived in §2, and compare with the observed cumulative distribution functions of eccentricity in our sample with $`|Z|<1`$ kpc. Figure 9 shows the results of this exercise. As is clear, more metal-rich ranges are described by a larger $`F`$, but in the abundance range $`2.2<`$ \[Fe/H\] $`2`$, $`F0`$, in good agreement with the results obtained from comparison of the differential eccentricity distributions. Figure 10 shows the dependence of $`F`$ on \[Fe/H\] for stars with $`|Z|<1`$ kpc, where the fits are made with bins of $`0.2`$ dex for \[Fe/H\] and $`0.1`$ for $`F`$. The $`F`$(\[Fe/H\]) derived here, based on full knowledge of the stellar orbital motions, is rather similar to that found by BSL based on radial velocities alone, except for the abundance range below \[Fe/H\] $`<1.5`$, where the MWTD appears more modestly populated. A characteristic value of $`F=0.3`$ over the abundance range $`1.7<`$ \[Fe/H\] $`1`$ is obtained from our present analysis.
Having obtained the structural parameters of the MWTD component, we extract a set of likely members of the MWTD component in our sample. High-resolution spectroscopic observations of these stars, to obtain estimates of their individual elemental abundances, should provide valuable information concerning the nature of MWTD stars, and reveal differences, if any, in their compositions relative to similar metallicity stars of the halo population (e.g., Bonifacio, Centurion, & Molaro 1999). In Table 5 we list the stars satisfying (1) $`2.2`$ \[Fe/H\] $`1`$, (2) $`|Z|1`$ kpc, (3) $`V_\varphi <V_\varphi >_{disk}\sigma _{\varphi ,disk}`$, and (4) $`|V_R|\sigma _{R,disk}`$ and $`|V_Z|\sigma _{Z,disk}`$, where $`<V_\varphi >_{disk}=200`$ km s<sup>-1</sup> and $`(\sigma _{R,disk},\sigma _{\varphi ,disk},\sigma _{Z,disk})=(46,50,35)`$ km s<sup>-1</sup>, as derived above). Condition (3) corresponds to the high rotation velocity of the candidate members and condition (4) is placed so that their velocities in $`R`$ and $`Z`$ directions are confined within a 1 $`\sigma `$ range relative to the zero mean, i.e. within the velocity dispersions of the MWTD component. In Table 5, the fourth column denotes the classification of each stellar type. We follow the coding of paper II – D: main-sequence dwarf star; TO: main-sequence turnoff star; SG: subgiant star; G: giant star; AGB: asymptotic giant branch star; FHB: field horizontal-branch star; RRV: RR Lyrae variable star; V: variable star. Note that Table 5 supersedes Table 8 of BSL, as we now have much more complete kinematic information.
## 4 Global Dynamics and Structure of the Halo
Although our present sample is dominated by stars located in the vicinity of the Sun, the orbits of many of these stars explore regions well into the more distant halo of the Galaxy. Thus, their local kinematics provide information on the global dynamics and structure of the halo (May and Binney 1986). In this section we first investigate the rotational properties of the halo at large heights from the Galactic plane, then use this same sample to obtain a picture of the global density distribution of the halo.
### 4.1 Rotational Properties of the Halo
Majewski (1992) claimed, on the basis of measured proper motions for an in situ sample of halo subdwarfs located at $`Z>5`$ kpc (in a small field in the direction of the North Galactic Pole), that stars at such large heights above the Galactic plane exhibit a net retrograde rotation $`<V_\varphi >55\pm 16`$ km s<sup>-1</sup>, in contrast to the stars nearer the plane, which show a near-zero or slightly prograde rotation. Although there have been criticisms of this result (e.g., Ryan 1992), a number of workers have also reported observations of separate samples which seem to support this view, so the true situation has remained unclear. For example, Carney et al. (1996) reported evidence for a retrograde rotation in the subset of their local sample of high proper-motion stars whose orbits extend far above the plane. In their analysis, they divided the sample into those stars with $`Z_{max}2`$ kpc, and $`Z_{max}5`$ kpc, respectively, where $`Z_{max}`$ is the maximum distance of the derived orbit from the plane. Carney et al. showed that their “high halo” sample, with $`<`$ \[Fe/H\] $`>=2.04`$ and $`Z_{max}5`$ kpc, exhibited a net retrograde rotation ($`<V_\varphi >=45\pm 22`$ km s<sup>-1</sup>), whereas their “low halo” sample at $`Z_{max}2`$ kpc exhibited a net prograde rotation ranging from $`<V_\varphi >=12`$ km s<sup>-1</sup> to 44 km s<sup>-1</sup>, depending on the specific criteria chosen to avoid stars of the disk component.
Figure 11a reproduces the original data of Carney et al. (1996), for stars with \[Fe/H\] $`1.5`$. Since, in their estimates of $`Z_{max}`$, Carney et al. adopted a different Galactic potential than ours, we have also re-determined $`Z_{max}`$ for their sample with the same potential described in §3, and show the results in this same figure. Regardless of the adopted potential, it is apparent that the stars at $`Z_{max}5`$ kpc exhibit a net retrograde rotation, as compared to the stars at $`Z_{max}2`$ kpc, which are in prograde rotation. Carney et al. argued that this result might be explained by the presence of two distinct halo populations, a high halo formed via accretion of fragments (such as in the SZ model), and a low halo formed from an organized contraction (similar to the ELS model).
The Carney et al. sample is based on high proper-motion stars selected from the Lowell Proper Motion Catalog, where proper motions are measured to exceed 0.26” yr<sup>-1</sup>, and the New Luyten Two-Tenths Catalog with proper motions exceeding roughly 0.18” yr<sup>-1</sup>. Thus, their sample is a-priori biased against inclusion of stars with prograde rotation close to the velocity of the LSR. What remains to be evaluated is the effect that this kinematic bias may have on the observed kinematics of the high halo.
We consider the question of bias in the Carney et al. sample via a Monte Carlo simulation, with the following assumptions. “Stars” with \[Fe/H\] $`<1.5`$ are randomly distributed in Galactic coordinates $`(l,b)`$, and are assumed to have a Gaussian velocity distribution with $`(\sigma _U,\sigma _V,\sigma _W)=(141,108,94)`$ km s<sup>-1</sup>, but with no systematic rotation. Distances to the simulated stars are taken by adopting a Gaussian form with mean 0.18 kpc and dispersion 0.09 kpc, which reproduces well the distance distribution of their sample stars. We further assume that only the stars whose inferred proper motions exceed 0.26” yr<sup>-1</sup> are observed at the Sun with $`V_{LSR}=220`$ km s<sup>-1</sup>. The result is shown in panel Figure 11b. The simulated data sample exhibits a net retrograde rotation for stars in the high halo. We note that our simulation does not exhibit a net prograde rotation in the low halo, as seen in the Carney et al. sample. This may arise because their sample is not distributed randomly in $`(l,b)`$, as our simulation assumes, and/or because their disk sample may be mainly drawn from stars for which the proper motion limit is 0.18” yr<sup>-1</sup>. Nevertheless, the results of this simple simulation provide reason to be skeptical of their claim of the existence of a retrograde high halo, which clearly can be influenced by selection bias of the input sample<sup>2</sup><sup>2</sup>2Carney (1999) reported that, after making a statistical correction for the kinematic bias in the Carney et al. 1996 sample, he also obtained a net prograde rotation even in the high halo..
Our large non-kinematically selected sample provides the means to elucidate the rotational character of the high halo without the effect of an input selection bias. Figure 11c presents our results based on sample stars with \[Fe/H\] $`1.5`$. It is found that the stars at small $`Z_{max}`$ show a net prograde rotation, in agreement with the results presented in §2 and also with the Carney et al. (1996) result: we obtain $`<V_\varphi >=59\pm 7`$ km s<sup>-1</sup> for 230 stars at $`Z_{max}<2`$ kpc. However, at large $`Z_{max}`$, the stars exhibit no systematic rotation, in sharp contrast to the Carney et al. result: we obtain $`<V_\varphi >=0\pm 8`$ km s<sup>-1</sup> for 212 stars at $`Z_{max}4`$ kpc. We note that although there is a difference in the rotational velocities for the stars close to and farther from the plane, which may suggest two populations (accreted and contracted populations), the boundary between them is not obvious.
### 4.2 The Global Density Distribution of the Halo
May & Binney (1986) discussed the interesting possibility that, on the basis of Jeans’ theorem, the global structure of the stellar halo can be recovered from local kinematic information for a sufficiently large sample of stars observed in the solar neighborhood. The theorem states that, for a well-mixed stellar system, the six-dimensional phase-space distribution function of stars, $`f(𝐱,𝐯)`$, can be taken to be a function of the three isolating integrals of motion $`I_i`$, $`i=1,2,3`$, i.e. $`f(I_1,I_2,I_3)`$. Within this isolating integral space, the stars constitute a set of fixed points with no time evolution. May & Binney argued that the stars observed in the solar neighborhood actually occupy a large fraction of this phase space, and it is hence possible to reconstruct the global structure of the stellar system from the kinematic data of nearby stars. Following these strategies, SLZ developed a maximum likelihood method for recovering a global model of the halo based on a discrete sum of orbits, and applied it to a sample of 118 local stars with \[Fe/H\] $`1.5`$ selected without kinematic bias. SLZ found that the stellar halo at $`8<R<20`$ kpc may consist of two components – a main, nearly spherical component, and an overlapping, highly flattened component. We note here that the actual halo system is unlikely to be in a well mixed equilibrium state, as we will discuss in the next section. However, the relaxation process is very slow compared to the orbital periods of typical stars, so the Jeans theorem and the above approach based on it are at least approximately valid.
We now apply the SLZ methodology to the present sample of stars, which is both larger, and has more accurately determined kinematic information than was available to SLZ. To exclude the MWTD stars as much as possible, we select as representatives of the halo population the stars in our sample with \[Fe/H\] $`1.8`$, a more restrictive abundance cut than the \[Fe/H\] $`1.5`$ used by SLZ. We also select a sample of stars with $`1.6<`$ \[Fe/H\] $`1`$ in order to examine the characteristics derived for a halo population contaminated by the MWTD. To minimize the effects of distance errors we limit our samples to those stars satisfying $`D4`$ kpc. We also remove stars with inferred (and possibly incorrect) extreme space motions ($`V_{RF}550`$ km s<sup>-1</sup>). After applying these cuts, the samples we investigate include $`N=359`$ stars for \[Fe/H\] $`1.8`$, and $`N=302`$ stars for $`1.6<`$ \[Fe/H\] $`1`$ .
The method is summarized as follows (see SLZ for the complete description): (1) the $`N`$ sets of isolating integrals $`(I_{1,i},I_{2,i},I_{3,i})`$, $`i=1,\mathrm{},N`$ are calculated from the observed positions $`𝐱_i`$ and velocities $`𝐯_i`$ of the stars, within an assumed Galactic potential. Here, $`I_1`$ is the total orbital energy $`E`$, $`I_2`$ is proportional to the square of the angular momentum vector pointing in the $`Z`$-direction $`I_2=L_z^2/2`$ (which measures azimuthal angular momentum), and $`I_3`$ is the so-called third integral of motion as defined in the Appendix (see also equation 15 of SLZ). For the Galactic potential, we adopt the same Stäckel model as in §3. (2) At all locations of the stars $`𝐱_j`$, $`j=1,\mathrm{},N`$, the probability density $`\rho _i(𝐱_j)`$ of an orbit characterized by $`(E_i,I_{2,i},I_{3,i}`$) is calculated for $`i=1,\mathrm{},N`$. In other words, we calculate the $`N^2`$ matrix $`\rho _{ij}`$, $`i,j=1,\mathrm{},N`$ from knowledge of integrals and locations of the orbits. (3) By maximizing the probability that the star found at $`𝐱_{j=1}`$ is on orbit $`i=1`$, the star found at $`𝐱_{j=2}`$ is on orbit $`i=2`$, and so forth, the orbit weighting factors, $`c_i`$, are used to estimate the total density at $`𝐱`$, viz
$$\rho (𝐱)=\underset{i=1}{\overset{N}{}}c_i\rho _i(𝐱).$$
(3)
Thus, equation (3) provides an estimate of the density of the halo stars at any point $`𝐱`$. The method also permits one to derive the mean azimuthal velocities as:
$$<V_\varphi >(𝐱)=\frac{1}{\rho (𝐱)}\underset{i=1}{\overset{N}{}}c_i\rho _i(𝐱)V_{\varphi ,i}(𝐱).$$
(4)
We proceed to average the results from equations (3) and (4) over grids of finite area. Following SLZ, we define the grids in the meridional plane of the spheroidal coordinates $`(\lambda ,\nu )`$, which is a suitable choice for Stäckel mass models (see the Appendix for more details). The grids are defined as $`\lambda _k=k^2\alpha `$, $`k=1,\mathrm{},30`$ and $`\nu _l=(\gamma \alpha )\mathrm{cos}^2(\theta _l)\gamma `$, $`\theta _l=(\pi /2)(l/20)`$, $`l=0,\mathrm{}20`$, as shown in Figure 4 of SLZ, where $`\alpha `$ and $`\gamma `$ are constants. The spatial resolution of the grids is about 1 kpc.
Figure 12a shows a plot of the reconstructed density distribution at the Galactic plane (the averaged density over the area at $`l=20`$), for \[Fe/H\] $`1.8`$ (filled circles) and $`1.6<`$ \[Fe/H\] $`1`$ (open circles). As a comparison, the results using the SLZ sample with \[Fe/H\] $`1.5`$ are also shown (crosses). As in the analysis of SLZ, the density distribution for $`R>8`$ kpc is well described by a power-law model $`\rho R^\beta `$. For \[Fe/H\] $`1.8`$, we find that the power-law model with exponent $`\beta =3.55\pm 0.13`$ fits well at all radii beyond $`R=8`$ kpc, up to the grid point for the largest radius, $`R=35`$ kpc. Note that for the SLZ sample the density at the largest three radii appears to fall short of the power-law model (probably as a result of their smaller sample size). If we omit these outer points, we obtain a fit to the power-law index $`\beta =3.57\pm 0.16`$ at $`8<R<28`$ kpc, which is yet slightly steeper than the $`\beta =3.29`$ result of SLZ. It is of interest to note that a power-law model with exponent $`\beta 3.5`$ derived for our sample of field halo stars is similar to the radial density distribution of the halo globular cluster population derived by Harris (1976) and Zinn (1985) ($`\beta =3.5`$) and by Carney, Latham, & Laird (1990) ($`\beta =3.0`$). A similar density distribution, but with a slightly shallower slope, has been found for field RR Lyrae stars by Saha (1985) ($`\beta 3`$) and Hawkins (1984) ($`\beta 3.1`$). Preston et al. (1991) combined counts of RR Lyrae stars and FHB stars in several fields to obtain the exponent $`\beta =3.5\pm 0.3`$. We see, in our reconstructed density distribution, no clear evidence for a break in the density distribution at $`R=2025`$ kpc as was detected in the number counts of globular clusters and RR Lyraes. For the subsample of stars with $`1.6<`$ \[Fe/H\] $`1`$, we obtain $`\beta =3.47\pm 0.18`$ for $`8<R<25`$ kpc, thus contamination from the MWTD has little effect. Below $`R=8`$ kpc, the density distributions of all three samples clearly deviate from a single power-law model, a result which is likely caused by incomplete representation, in the solar neighborhood, of stars for which apocentric radii are below $`R=R_{}`$.
Figure 12b shows the mean azimuthal velocities, $`<V_\varphi >`$, for the same three samples of stars, projected onto the Galactic plane. The value of $`<V_\varphi >`$ is nearly zero for $`R>10`$ kpc, but there is a signature of increasing $`<V_\varphi >`$ with decreasing $`R`$. In particular, for the sample of stars with $`1.6<`$ \[Fe/H\] $`1`$, $`<V_\varphi >`$ rises rather discontinuously at $`R10`$ kpc, which may correspond to the radial limit of the rapidly rotating thick-disk component.
Figure 13 is a plot of the inferred global density distribution in the $`(R,Z)`$ plane, in the form of equidensity contours. The lack of stars at small $`R`$ and large $`Z`$ (which gives rise to the ill-formed contour levels in this portion of the diagram) is a consequence of the small probability that stars in the Galaxy that explore such a region are represented in the solar neighborhood (as argued in SLZ). Other than in this region, the inferred density distribution based on the 359 stars with \[Fe/H\] $`1.8`$ (panel b) appears to be very similar to that which SLZ obtained from their 118 stars (panel a): The outer part of the halo, at $`R>15`$ kpc, is round, in good agreement with inferences based on star counts (see Freeman 1987 and references therein; Preston et al. 1991).
While we reproduce the general sense of the SLZ results with our much larger sample, there are notable differences in the details. SLZ argued that there is a clear indication that the halo, at any given radius, consists of both a main, nearly spherical component and an overlapping, highly flattened component. There no clear evidence for this result in the density reconstruction based on our new data. To further examine this point, we have fit elliptical contours to the reconstructed density maps after, following SLZ, omitting the data points near the Galactic plane. Specifically, we obtain fits to ellipses of major axis $`a`$, and axial ratio $`q`$, over the polar angle $`40^{}<\theta <80^{}`$. Residuals to the fits obtained to these ellipses as a function of polar angle are shown in Figure 14a. Thick solid and dotted lines correspond to the fits with $`a=10.5`$ kpc and $`q=0.70`$ and with $`a=13.5`$ kpc and $`q=0.51`$, respectively, for our sample with \[Fe/H\] $`1.8`$. For comparison, we show the corresponding results using the SLZ sample fit over $`30^{}<\theta <80^{}`$ (thin solid and dotted lines). It follows that, while we reproduce the SLZ result that there is a large density excess near the Galactic plane when using their sample, it is not evident in our sample – an additional flat component is not required for the fitting.
Based on the above result, we proceed to make a fit including the density data near the plane, with a single value for the axial ratio $`q`$ at each major axis $`a`$, i.e., without taking into account an additional flat component having small $`q`$. The change of our estimate of $`q`$ as a function of radius is shown in Figure 14b. For the sample of stars with \[Fe/H\] $`1.8`$, the density distribution in the outer part of the halo, $`R20`$ kpc, is quite round. However, the axial ratio $`q`$ appears to decrease with decreasing $`R`$ over $`15<R<20`$ kpc, and the inner part, at $`R<15`$ kpc, exhibits $`q0.65`$. Thus, the halo can be described as nearly spherical in the outer part and highly flattened in the inner part, instead of the overlap of both components at all locations in the halo. This result is in good agreement with previous studies of the distribution of RR Lyrae (Hartwick 1987; Layden 1995) and FHB stars (Preston et al. 1991; Kinman, Suntzeff, & Kraft 1994), as well as with the flattening of the inner halo reported by Larsen & Humphreys (1994) based on counts of F- and G-type stars. It remains an open question as to whether or not there exists a distinct boundary between the outer spherical halo and the inner flattened halo.
As seen in Figure 13c, the density distribution in the inner part of the halo for stars of intermediate abundance ($`1.6<`$ \[Fe/H\] $`1`$) appears to be more flattened than is the case for stars with \[Fe/H\] $`1.8`$. This result is reflected in the ellipse fits for these stars shown in Figure 14b, which indicate that $`0.5<q<0.6`$ at $`R<12`$ kpc. This may be due to the contribution of the (by definition) flattened, MWTD population with a finite radial scale length. We note that our results imply the decrease of axial ratios at $`R>17`$ kpc. This may be an artifact due to the small numbers of stars employed at such large radii, so further analysis using much larger samples is necessary.
## 5 Kinematic Substructure of the Halo in the Solar Neighborhood
If the halo of the Galaxy was assembled from the merging and/or accretion of small subgalactic clumps, as argued by SZ, then one might hope to find signatures of those events, even now, in the form of kinematic substructures, because the mixing of phase space for such stars is expected to be incomplete (Helmi & White 1999). The Sagittarius dwarf galaxy is an ongoing merging event at the current epoch (Ibata, Gilmore, & Irwin 1994), and the Magellanic Clouds may ultimately follow the similar fate, as they lose energy via dynamical friction (Tremaine 1976). Signatures of past merging events in the halo were reported by Majewski et al. (1994) in their in situ sample of the stars at about 4.5 kpc above the Galactic plane. Helmi & White (1999) argue that the reported clumpiness in the velocity distribution of the halo stars from Majewski et al. is actually a superposition of two individual streams of stars, possibly arising from a common progenitor.
Clear “fossil evidence” in the solar neighborhood for a previous merger of the Galaxy with what may have been a “Seale & Zinn fragment” has recently been discovered by Helmi et al. (1999, hereafter HWdZZ). These authors examined a subset of the stars in samples from our previous work (BSL; CY), and identified a statistically significant clumping of stars in the angular momentum diagram $`L_z`$ vs $`L_{}=(L_x^2+L_y^2)^{1/2}`$. The substructure identified by HWdZZ consists of 7 stars (in a sample of 97 stars with \[Fe/H\] $`1.6`$ and $`D<1`$ kpc), or 12 stars (in a sample of 275 stars with \[Fe/H\] $`1`$ and $`D<2.5`$ kpc). HWdZZ suggest, based on the observed numbers of stars in this clump, that roughly 10% of the halo stars outside the solar radius may have arisen from a single coherent object with a total mass of about $`10^8`$ M, disrupted during the process of halo formation.
We now consider the HWdZZ result based on our revised catalog. Figure 15 shows the angular momentum diagram of HWdZZ as populated by the stars of our present sample within 2.5 kpc of the Sun. Panels (a) and (b) are for stars in the abundance ranges \[Fe/H\] $`1.6`$ and $`1.6<`$ \[Fe/H\] $`1`$, respectively. In both of these abundance ranges there exists a clump of stars in the region that HWdZZ pointed out, at $`L_{}2200`$ kpc km s<sup>-1</sup> and $`L_z1200`$ kpc km s<sup>-1</sup>. In addition to this clump, we identify a possible “trail” (in angular momentum space) which appears to connect the clump and the high $`L_z`$ region, most clearly evident among the higher abundance stars shown in panel (b). For the purposes of this discussion, we define the “clump” region to be comprised of stars with $`2100<L_{}<2600`$ kpc km s<sup>-1</sup> and $`800<L_z<1500`$ kpc km s<sup>-1</sup> (solid box), and the “trail” region to be comprised of stars having $`1250<L_{}<2000`$ kpc km s<sup>-1</sup> and $`1200<L_z<2000`$ kpc km s<sup>-1</sup> (dotted box), respectively. We include BPS CS 22876-0040 as a “trail” member (triangle in Fig.15a), which is somewhat outside the region defined above, $`(L_z,L_{})=(1037,1672)`$ kpc km s<sup>-1</sup>, because this star exhibits similar orbital motions to the “trail” stars examined below.
We note that the angular momentum, $`L_{}`$, is not an exact integral of motion in the currently adopted Galactic potential (which consists of a disk and halo component), thus its use may not be generally appropriate for the study of kinematic substructures in the halo. Rather, the so-called third integral, $`I_3`$, which is related to $`L_{}`$ in a non-spherical potential, should be used. Unfortunately, no general analytic expression exists for $`I_3`$. However, the Stäckel form of the currently adopted potential allows one to estimate $`I_3`$ in an explicit manner (de Zeeuw 1985; de Zeeuw, Peletier, & Franx 1986; Dejonghe & de Zeeuw 1988), as well as other integrals, such as the orbital energy $`E`$ and the angular momentum $`L_z`$. Expressions for these three integrals are described in the Appendix.
For a spherically symmetric potential, the quantity $`(2I_3)^{1/2}`$ is equivalent to $`L_{}`$. Figure 16 is a plot of $`(2I_3)^{1/2}`$ vs. $`L_z`$ (panel a), and $`|E|`$ vs $`L_z`$ diagram (panel b), respectively, for the 723 stars from our sample with \[Fe/H\] $`1`$ and $`D2.5`$ kpc. Filled and open circles denote the stars in the “clump” and “trail” regions of Figure 15, respectively. We note that one star in the “clump” region of Figure 15 (HD 214161) exhibits quite different orbital parameters than the others (see below), so this star is drawn as a cross in Figure 16. The fact that we see such similar structures in Figures 15 and 16 implies that, along the orbits of the stars constituting the “clump” and “trail”, the gravitational potential can be regarded as nearly spherical – the member stars spend the majority of their orbits far from the Galactic plane, where the effect of the disk potential is modest. Also, panel (b) suggests that the orbital energies $`|E|`$ of the stars in the “clump” are confined to the narrow range near $`|E|10^5`$ km<sup>2</sup> s<sup>-2</sup>, whereas the stars in the “trail” have a rather diverse range of $`|E|`$.
Another choice of integrals in the current mass model are the so-called action integrals, $`𝐉`$. Action integrals are adiabatic invariants, and thus remain unchanged even if the orbital energy $`E`$ changes via gravitational interaction, provided the time scale for the interaction is sufficiently long compared to the orbital period. One of the actions is $`J_\varphi `$, which is equivalent to $`L_z`$ in an axisymmetric potential. Other suitable actions in the current Stäckel potential, defined in terms of the spheroidal coordinates $`(\lambda ,\nu )`$, are $`(J_\lambda ,J_\nu )`$ (see the Appendix for complete definitions). Panels (c) and (d) of Figure 16 show the distribution of our sample stars in the $`J_\nu `$ vs $`J_\lambda `$ and $`J_\nu `$ vs $`L_z`$ diagrams, respectively. It can been seen that all of the “clump” stars, except for the star HD 214161 $`(J_\lambda =5058,J_\nu =1223)`$, exhibit a clear clumpiness in action space, whereas the “trail” stars show a broad distribution in $`J_\nu `$, but with a rather narrow distribution in $`J_\lambda `$ \[with the exception of the star CS Ser $`(J_\lambda =6957,J_\nu =587)`$\].
From these integrals of motion, we conclude that the “clump” consists of only 9 stars, instead of 12 stars in the HWdZZ result, and the “trail” consists of 9 stars. The kinematic quantities for these stars are listed in Table 6. The orbit of the “clump” is characterized by $`Z_{max}16`$ kpc, $`r_{ap}20`$ kpc, and $`r_{pr}7`$ kpc, which are in good agreement with the HWdZZ result. Note that even though we have tripled the numbers of stars considered in our present sample (728 stars), relative to that of the sample examined by HWdZZ (275 stars), the number of detected clump members has not increased. Thus, HWdZZ’s conclusion that as much as one-tenth of the halo stars presently located in the solar neighborhood originates from a single object may not apply. We have not considered, as of yet, the impact of the “trail” stars on this argument.
Figure 17 shows the metallicity distribution of the stars in the “clump” (solid histogram) and in the “trail” (dotted histogram) features. The two features seem to share a similar metallicity distribution, though the number of stars involved is still too small for any definite conclusion to be reached. It might be tempting to suggest that the “trail” was formed via the tidal interaction of the precursor object of the “clump” with the Galactic potential, with the “trail” stars gaining some of the orbital angular momentum lost by the precursor object. Such an interaction may have proceeded in a rapid manner, being comparable to the orbital period of the object, so that the actions of the “trail” differ from those of the “clump”. Progress on evaluation of this picture will require detailed numerical simulations, which we are presently investigating.
## 6 Discussion and Conclusion
We have analyzed both the local and global kinematics of 1203 metal-poor stars in the Galaxy with \[Fe/H\] $`0.6`$, based on a large, revised, catalog of stars selected without kinematic bias (Paper II). All of these stars have available distance estimates, radial velocities, proper motions, and abundance estimates over the full applicable range in the Milky Way. This is the largest non-kinematically selected sample yet assembled, so the derived kinematic properties are the least affected by systematics as well as statistical fluctuations. We summarize our results below, and discuss them in the context of the formation of the Galaxy.
### 6.1 Summary of the Results
The local kinematics of the halo population, based on the stars with \[Fe/H\] $`2.2`$ and $`|Z|<1`$ kpc, are characterized by a radially elongated velocity ellipsoid $`(\sigma _U,\sigma _V,\sigma _W)=(141\pm 11,106\pm 9,94\pm 8)`$ km s<sup>-1</sup> and a small prograde rotation $`<V_\varphi >=3050`$ km s<sup>-1</sup> (assuming $`V_{LSR}=220`$ km s<sup>-1</sup>). When additional halo stars at larger $`|Z|`$ are taken into account, the velocity ellipsoid remains essentially unchanged, but $`<V_\varphi >`$ exhibits a marked decrease (Figs. 2 and 3). We find no evidence for an increase of $`\sigma _W`$ at the lowest abundances, as had been previously suggested. At higher metallicities, the stars in our sample exhibit disk-like kinematics, and a higher mean rotation. Specifically, for stars in the abundance interval $`0.7`$ \[Fe/H\] $`<0.6`$ and with $`|Z|<1`$ kpc, we have obtained $`(\sigma _U,\sigma _V,\sigma _W)=(46\pm 4,50\pm 4,35\pm 3)`$ km s<sup>-1</sup> and $`<V_\varphi >=200`$ km s<sup>-1</sup>, which characterize the kinematic parameters of the thick disk. We have also confirmed previous results that there exists a remarkable discontinuity of the rotational properties of the Galaxy at \[Fe/H\] $`1.7`$ (Fig. 3).
Analysis of a large sample of non-kinematically selected stars provides clear evidence, supporting earlier suspicions based on much smaller samples, that there exists no correlation between metal abundances and orbital eccentricities for metal-poor stars of the Milky Way (Fig. 5). Even at the lowest abundances explored in our sample, \[Fe/H\] $`2.2`$, about 20% of the stars have $`e<0.4`$. In addition, there is a small concentration of high-$`e`$ stars at \[Fe/H\] $`1.7`$, which is possibly responsible for the near zero $`<V_\varphi >`$ at the same metallicity. We found that the fraction of the low-eccentricity stars with \[Fe/H\] $`2.2`$ remains the same, even as one changes the range of $`|Z|`$ (Fig. 6a), so such stars belong to the halo, not the MWTD component. On the other hand, stars in intermediate abundance ranges above $`2.2`$ dex exhibit a decrease of low-$`e`$ stars with increasing $`|Z|`$, and the fraction of such stars appears to converge at high $`|Z|`$ to that found for \[Fe/H\] $`2.2`$ (Fig. 6b). Both a KS test and a Monte Carlo simulation enable a determination of the structural parameters of the disk component in these abundance ranges. Specifically, the fraction of the disk component is about 30% for $`1.7<`$ \[Fe/H\] $`1`$, but is less than 10% for more metal-poor ranges (Fig. 8).
The global kinematics of the halo stars are summarized as follows. In contrast to the claims of Majewski (1992), and Carney et al. (1996), stars in our sample do not show a net retrograde rotation at large $`Z_{max}`$, but rather exhibit a near zero systematic rotation (Fig. 10). The difference between our result and that of Carney et al. (1996) probably arises from the (unavoidable) kinematic bias inherent in their sample selection criteria. The observed decrease of $`<V_\varphi >`$ with increasing $`Z_{max}`$ is continuous, so that it is not possible to conclude that the inner “contracted” population (with a positive $`<V_\varphi >`$) is distinct from an outer “accreted” population, based on the rotational properties of the metal-poor stars alone. Our analysis of the global density distribution of halo stars, based on the reconstruction method developed by SLZ, confirms SLZ’s conclusion that the outer halo is quite round (Fig. 12). However, we see no evidence of an overlapping flattened component in addition to the main, nearly spherical one, as was claimed by SLZ. Rather, the density distribution of the halo is better described as nearly spherical in the outer region (beyond $`R=1520`$ kpc) and highly flattened in the inner region.
We have confirmed a recent detection of kinematic substructure in the solar neighborhood by HWdZZ, based on a small number of stars which cluster together in the halo angular momentum diagram. We have also found an additional elongated “trail” which appears to connect between HWdZZ’s “clump” and the high $`L_z`$ region (Fig. 14). Further analysis, using several integrals of motion for the “clump,” does not result in a dramatic increase in the numbers of stars associated with it, even though the total number of our sample stars is three times as large as that available to HWdZZ.
### 6.2 Implications for the Formation of the Galaxy
The local and global kinematics of metal-poor stars provide valuable clues for understanding the formation process of the halo and thick disk components in the Galaxy, as well as in disk-type galaxies in general.
If the primordial collapse from the halo to the disk occurred in a monolithic manner, starting from an overdense homogeneous spheroid, one might expect (as predicted by the ELS model) to observe a continuous increase of $`<V_\varphi >`$ for the stars born from the infalling gas, as well as a continuous decrease of their orbital eccentricities with increasing \[Fe/H\] as the spheroid spins up in order to conserve angular momentum. The fact that we observe no correlation between \[Fe/H\] and $`e`$, and basically no change of $`<V_\varphi >`$ with abundance for stars with \[Fe/H\] $`1.7`$ conflicts with this scenario. The lack of an abundance gradient in the halo stars (Carney et al. 1990; CY) is also difficult to interpret in this context. The outer halo, if formed from a monolithic collapse, might be expected to be dominated by radially elongated motions of the stars, but this is actually opposite to the inferred tangentially anisotropic velocity ellipsoid at large distance from the Sun (see Sommer-Larsen et al. 1997). We also note that a small portion of the metal-poor stars having \[Fe/H\] $`1.7`$ may have been formed from the infalling part of gas, so as to explain both the nearly zero $`<V_\varphi >`$ and the excess number of high-$`e`$ stars found at \[Fe/H\] $`1.7`$.
If the halo is assembled from merging and/or accretion of numerous fragments falling into the Galaxy (SZ), one might expect little or no correlation between kinematic and chemical properties, as each fragment has its own chemical history, and the merging process may proceed in a chaotic manner. Our results for the $`<V_\varphi >`$ vs. \[Fe/H\] and \[Fe/H\] vs. $`e`$ relations are basically in agreement with this scenario. The SZ scenario is also consistent with a several Gyr age spread in globular clusters in the outer halo (see, e.g., Rosenberg et al. 1999), and even in field stars (Schuster & Nissen 1989), because the initiation and duration of star formation may not be coherent from fragment to fragment<sup>3</sup><sup>3</sup>3Harris et al. (1997) showed that the most metal-poor globular clusters, such as M92, have essentially the same age everywhere in the halo. As they argued, this result could also be explained within the precepts of the SZ scenario if all of the “SZ fragments” began building the first generation of clusters in the same time period.. However, the original SZ scenario seems unlikely to explain our observed vertical gradient of $`<V_\varphi >`$ for halo stars, as well as the highly flattened density distribution of the inner halo, in contrast to the nearly spherical outer halo. Totally incoherent, chaotic merging of SZ fragments would not be expected to produce these “internal” kinematic structures in the halo. It is also unclear as to how the rapidly rotating disk component subsequently formed out of the aftermath of merging (see also Freeman 1996).
SZ first suggested that at least the inner part of the halo may have undergone a coherent contraction in a manner similar to the ELS hypothesis, an idea which has been invoked by subsequent workers to explain the duality of the density, kinematics, and ages of the halo field stars (e.g., SLZ; Norris 1994; Carney et al. 1996; Sommer-Larsen et al. 1997), as well as the age difference between outer and inner globular clusters (Zinn 1996; Rosenberg et al. 1999). This sort of hybrid picture, combining aspects of both the ELS and SZ scenarios, proposes that the outer halo is made up from merging and/or accretion of subgalactic objects, such as dwarf-type satellite galaxies, whereas the inner part of the halo has undergone a dissipative contraction on relatively short timescales. This hybrid model might explain our identification of the inner, flattened, slowly rotating component of the halo with a finite spatial gradient in $`<V_\varphi >`$.
An alternative hypothesis to explain an observed “duality” of the Galactic halo relies on the existence of a thick-disk population of stars even at rather low abundances (MFF; Norris 1994; BSL). If stars with disk-like kinematics have a metallicity distribution which extends below \[Fe/H\] $`=2`$, then a finite fraction of their orbits would be characterized by low $`e`$, as found in our present investigation. One possible origin of this MWTD component may be the heating of the pre-existing thin disk triggered by the dissipationless merging of a satellite falling into the disk (Quinn, Hernquist, & Fullagar 1993). Under this hypothesis, the currently observed thin disk, with a vertical scale height of $`350`$ pc, could only have formed after the merging event was completed. However, our finding that few thick-disk stars exist with \[Fe/H\] $`1.7`$, and no observed increase of $`\sigma _W`$ with decreasing \[Fe/H\], belies the existence of a dynamically hot, proto-disk population at the lowest abundances (see also Norris 1994; Ryan & Lambert 1995; Twarog & Anthony-Twarog 1996). Furthermore, following the results presented in §4.2, we see no evidence for an overlapping flattened component of the halo in addition to a nearly spherical component. An indication that there might exist a significant vertical gradient in $`<V_\varphi >`$ for \[Fe/H\] $`1.7`$, compared with a much smaller gradient observed for the thick disk itself, also conflicts with this hypothesis. Thus, we conclude that the formation of the inner flattened halo possibly involves a dissipative contraction, not a dissipationless heating of the proto-disk.
If a hybrid halo formation picture, based on dissipationless merging in the outer halo and dissipative contraction in the inner halo, applies, the question arises as to whether there is a clear boundary distinguishing the two regions. The results presented in §4 suggest no clear distinction between the outer and inner regions of the halo, at least as seen from inspection of the $`<V_\varphi >`$ vs. $`Z_{max}`$ relation, and the inferred globafl density distribution of the halo. Furthermore, there presently exists no reasonable theoretical explanation for the existence of two distinct populations of stars in the halo. Thus, our current analysis implies that both dissipationless and dissipative processes in the outer and inner halo, respectively, may have occurred more or less in a simultaneous manner.
We now ask whether the above hybrid scenario is a natural consequence of the currently favored theory of galaxy formation based on hierarchical assembly of cold dark matter (CDM) halos (e.g., Peacock 1999). The CDM model postulates that initial density fluctuations in the early Universe have larger amplitudes on smaller scales. Thus, the initially overdense regions that end up forming large galaxies such as our own contain large density fluctuations on subgalactic scales. As a protogalaxy collapses from the general cosmological expansion, these small-scale fluctuations develop into numerous clumps of CDM particles, into which the interstellar gas falls from gravitational attraction. The protogalaxy is thus made up of numerous clumps comprised of a mixture of primordial gas and dark matter, interacting with one another via their mutual gravitational attraction. According to numerical simulations by Steinmetz & Müller (1994; 1995) and Bekki & Chiba (1999), most of the metal-poor stars which presently occupy the outer halo of our Galaxy form in these local, small-scale density fluctuations. Once star formation initiates, the gas inside of these small fragments quickly escapes due to energy feedback from supernovae. Later, these clumps begin to merge with one another, and the aftermath of these essentially dissipationless merging processes exhibits a nearly spherical density distribution with no abundance gradient.
The subsequent evolution of the system may be described in the following way (Bekki & Chiba 1999). As a consequence of the merging of low-mass fragments, a smaller number of more massive clumps develops – within each of these merged clumps one expects to find previously formed metal-poor stars as well as newly born stars. These large clumps continue to accrete gas from their immediate surroundings. These clumps gradually move toward the central region of the system due to both dynamical friction, and dissipative merging with smaller clumps. Then, when the last merging event between the largest clumps occurs, the stars which have been confined inside the clumps are disrupted and spread over the inner part of the halo, whereas a large fraction of the disrupted gas appears to end up in the center of the Galaxy and may form a bulge. As a consequence, the inner part of the halo should have a flattened density distribution with a finite prograde rotation, as reported here, and its angular momentum distribution may be similar to that of the bulge (Wyse & Gilmore 1992). Also, the stars born from this infalling stage of gas may explain the existence of high-$`e`$ stars at \[Fe/H\] $`1.7`$. The simulations conducted to date imply that the thick disk component is partially composed of debris stars, but it is mainly made from diluted gas which has been accreted from the outer part of the halo (Sommer-Larsen et al. 1997). Therefore, although more detailed simulation work is required, CDM models appear to reproduce, at least qualitatively, the overall kinematic properties of the metal-poor stars via both dissipationless merging in the outer halo, and dissipative merging in the inner halo.
It is unknown whether or not evidence for merging events of CDM clumps during the early evolution of the Galaxy might be still preserved as kinematic substructures at the current epoch. Within the currently available precision of space velocities for stars in our sample, typically of the order of 10 to 20 km s<sup>-1</sup>, the main body of the halo appears to be well mixed in phase space; higher precision measurements of proper motions by the planned astrometric satellites missions (e.g., FAME, SIM, GAIA) may be able to disentangle this complex mixture of halo stars (Helmi, Zhao & de Zeeuw 1999). Alternatively, the confirmed kinematic clumping of halo stars presented in §5 may originate from the recent accretion of a satellite galaxy, which has fallen into the Galaxy after the major part of the halo was formed.
Firmer conclusions on the formation of the Galaxy require the assembly and analysis of still larger numbers of stars with accurate distances and proper motions, especially at larger distances from the Sun. Exploration along this line is now in progress. More elaborate numerical modeling of the formation of large spiral galaxies such as the Milky Way is also needed in order to clarify the physical processes that lead to the currently observed dynamics and structure of the halo and disk components. It is of particular importance to model and understand the chemo-dynamical evolution of the system of subgalactic fragments in the course of the Galaxy’s collapse. Once a fundamental understanding of the formation and evolution of our Galaxy is established, it will then be possible to obtain additional insights into formation of disk-type galaxies in general, by combining our refined picture with the rapidly growing observational database of young galaxies becoming available in the deep realm of the Universe.
We are grateful to the referee, Bruce Carney, for his careful review of the paper and a number of thoughtful suggestions. MC acknowledges partial support from Grants-in-Aid for Scientific Research (09640328) from the Ministry of Education, Science, Sports and Culture of Japan. TCB acknowledges partial support for this work from grant AST 95-29454 from the National Science Foundation.
## Appendix A The Stäckel Potential and Integrals of Motion
We briefly describe the properties of the Stäckel potential adopted in this work, and present expressions for the associated integrals of motion. For more details, see, e.g., de Zeeuw (1985), Dejonghe & de Zeeuw (1988), and SLZ.
We define spheroidal coordinates $`(\lambda ,\varphi ,\nu )`$, where $`\varphi `$ corresponds to the azimuthal angle in cylindrical coordinates $`(R,\varphi ,Z)`$, and $`\lambda `$ and $`\nu `$ are the roots for $`\tau `$ of
$$\frac{R^2}{\tau +\alpha }+\frac{Z^2}{\tau +\gamma }=1,$$
(A1)
where $`\alpha `$ and $`\gamma `$ are constants, giving $`\gamma \nu \alpha \lambda `$. The coordinate surfaces are spheroids $`(\lambda =const.)`$ and hyperboloids of revolution $`(\nu =const.)`$ with the $`Z`$-axis as the rotation axis, where the focal distance $`\mathrm{\Lambda }=(\gamma \alpha )^{1/2}`$ fixes the coordinate system.
The gravitational potential of the Stäckel type is then written as
$$\phi (\lambda ,\nu )=\frac{(\lambda +\gamma )G(\lambda )(\nu +\gamma )G(\nu )}{\lambda \nu },$$
(A2)
where $`G(\tau )`$ is an arbitrary function. In this work, $`G(\tau )`$ is the sum of $`G_{disk}(\tau )`$ from the disk and $`G_{halo}(\tau )`$ from the massive dark halo. Following SLZ, we adopt the oblate perfect spheroid for $`G_{disk}(\tau )`$ and the $`s=2`$ model of de Zeeuw, Peletier, & Franx (1986) for $`G_{halo}(\tau )`$.
The Hamiltonian, $`H`$, per unit mass, for motion in this potential $`\phi (\lambda ,\nu )`$ is
$$H=\frac{p_\lambda ^2}{2P^2}+\frac{p_\varphi ^2}{2R^2}+\frac{p_\nu ^2}{2Q^2}+\phi (\lambda ,\nu ),$$
(A3)
where $`P`$ and $`Q`$ are the metric coefficients of the spheroidal coordinates, given by
$$P^2=\frac{\lambda \nu }{4(\lambda +\alpha )(\lambda +\gamma )},Q^2=\frac{\lambda \nu }{4(\nu +\alpha )(\nu +\gamma )},$$
(A4)
and $`p_\lambda `$, $`p_\varphi `$, and $`p_\nu `$ are the conjugate momenta to $`\lambda `$, $`\varphi `$, and $`\nu `$, respectively,
$$p_\lambda =P^2\dot{\lambda }=Pv_\lambda ,p_\varphi =R^2\dot{\varphi }=Rv_\varphi ,p_\nu =Q^2\dot{\nu }=Qv_\nu .$$
(A5)
The velocities $`v_\lambda `$, $`v_\varphi `$, and $`v_\nu `$ at a point $`(\lambda ,\varphi ,\nu )`$ are the components of the velocity $`𝐯`$ along the orthogonal axis defined locally by the spheroidal coordinate system.
The three integrals of motion, $`|E|`$, $`I_2`$, and $`I_3`$, are defined as
$`|E|`$ $`=`$ $`H`$ (A6)
$`I_2`$ $`=`$ $`{\displaystyle \frac{L_z^2}{2}}`$ (A7)
$`I_3`$ $`=`$ $`{\displaystyle \frac{1}{2}}(L_x^2+L_y^2)+\mathrm{\Delta }^2\left[{\displaystyle \frac{1}{2}}v_z^2Z^2{\displaystyle \frac{G(\lambda )G(\nu )}{\lambda \nu }}\right],`$ (A8)
and the action integrals $`J_\lambda `$, $`J_\varphi `$, and $`J_\nu `$ are defined as
$`J_\lambda `$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle p_\lambda 𝑑\lambda }={\displaystyle \frac{2}{\pi }}{\displaystyle _{\lambda _1}^{\lambda _2}}p_\lambda 𝑑\lambda `$ (A9)
$`J_\varphi `$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle p_\varphi 𝑑\varphi }=L_z`$ (A10)
$`J_\nu `$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle p_\nu 𝑑\nu }={\displaystyle \frac{2}{\pi }}{\displaystyle _\gamma ^{\nu _0}}p_\nu 𝑑\nu ,`$ (A11)
where ($`\lambda _1,\lambda _2)`$ and $`\nu _0`$ are the turning points of the orbit, defined as the values for which $`v_\lambda =0`$ and $`v_\nu =0`$, respectively, and $`\nu =\gamma `$ defines the equatorial plane. For the evaluation of $`J_\lambda `$, we have taken four times the integrals from $`\lambda _1`$ to $`\lambda _2`$, to maintain symmetry between $`J_\lambda `$ and $`J_\nu `$ and ensure continuity of the actions across transitions from one orbital family to another (de Zeeuw 1985). |
warning/0003/math0003154.html | ar5iv | text | # Lemma 1 (Li)
L-embedded Banach spaces and measure topology
H. Pfitzner
## Abstract
An L-embedded Banach spaace is a Banach space which is complemented in its bidual such that the norm is additive between the two complementary parts. On such spaces we define a topology, called an abstract measure topology, which by known results coincides with the usual measure topology on preduals of finite von Neumann algebras (like $`\mathrm{L}^1([0,1])`$). Though not numerous, the known properties of this topology suffice to generalize several results on subspaces of $`\mathrm{L}^1([0,1])`$ to subspaces of arbitrary L-embedded spaces.
§1 Introduction
This article continues the investigations made in on asymptotically isometric copies of $`l^1`$ in preduals of von Neumann algebras and in L-embedded Banach spaces. (For defintions see below.) In it has been proved that, roughly speaking, in the predual of a finite von Neumann algebra the only non-trivial bounded sequences that converge to $`0`$ with respect to the measure topology are essentially those that span $`l^1`$ asymptotically; for $`\mathrm{L}^1(\mu )`$, $`\mu `$ a finite measure, this characterization has been known for quite a time \[15, Th. 2\], \[25, Th. 3, Rem. 6bis\].
From the point of view of Banach space theory, L-embedded Banach spaces provide a natural frame for preduals of von Neumann algebras. So the starting point of this paper is on the one hand the definition of an abstract measure topology, Definition 3, patterned after the just mentionend characterization and on the other hand the easy but important observation, Theorem 4, that every L-embedded space admits such a topology. Although this topology does not come out easily with its properties \- at the time of this writing it is not clear whether it is Hausdorff let alone metrizable or whether addition is continuous - it allows to generalize several results on subspaces of $`\mathrm{L}^1(\mu )`$ to subspaces of arbitrary L-embedded spaces. Thus section §4 of the present paper is titled ”Section IV.3 of (partly) revisited”. For example, Theorem 10 generalizes a theorem of Buhvalov-Lozanovskii which describes the link between L-embeddedness and measure topology for subspaces $`Y`$ of $`\mathrm{L}^1(\mu )`$, $`\mu `$ finite: $`Y`$ is L-embedded if and only if its unit ball is closed in measure. (Note in passing that this criterion involves only the space $`Y`$ itself, not its bidual.) Moreover, as a consequence of this, the closedness in measure of the unit ball of $`Y`$ is a weak substitute for compactness which could be called ”convex sequential compactness”, see Corollary 9. We also reprove a result of Godefroy and Li concerning a criterion for L-embedded subspaces which are duals of M-embedded spaces, see Theorem 13. In this vein, that is by substituting arbitrary L-embedded spaces for $`\mathrm{L}^1(\mu )`$, we recover also some results of Godefroy, Kalton, Li in §5. Finally, in §6 it is proved that addition is $`\tau _\mu `$-continuous in preduals of von Neumann algebras.
§2 Notation, Background:
The results are stated for complex scalars. The dual of a Banach space $`X`$ is denoted by $`X^{}`$. $`\mathrm{B}_X`$ denotes the unit ball of $`X`$. Subspace of a Banach space means norm-closed subspace, bounded always means norm-bounded. As usual, we consider a Banach space as a subspace of its bidual omitting the canonical embedding. $`[x_n]`$ denotes the closed linear span of a (finite or infinite) sequence $`(x_n)`$.
Basic properties and definitions which are not explained here can be found in or in - for Banach spaces and in for C-algebras. The standard reference for M- and L-embedded spaces is the monograph .
Let $`(x_n)`$ be a sequence of nonzero elements in a Banach space $`X`$.
We say that $`(x_n)`$ spans $`l^1`$ isomorphically (or $`r`$-isomorphically to be more precise) \- $`(x_n)_{n\mathrm{I}N}\stackrel{r}{}l^1`$ or just $`x_n\stackrel{r}{}l^1`$ in symbols - if there exists $`r>0`$ (trivially $`r1`$) such that $`r(_{n=1}^{\mathrm{}}|\alpha _n|)_{n=1}^{\mathrm{}}\alpha _n\frac{x_n}{x_n}_{n=1}^{\mathrm{}}|\alpha _n|`$ for all scalars $`\alpha _n`$ (the second inequality being trivial).
We say that $`(x_n)`$ spans $`l^1`$ almost isometrically \- $`x_n\stackrel{\mathrm{alm}}{}l^1`$ in symbols - if there is a sequence $`(\delta _m)`$ in $`[0,1[`$ tending to $`0`$ such that $`(x_n)_{nm}\stackrel{1\delta _m}{}l^1`$ for all $`m\mathrm{I}N`$. Recall that the Banach-Mazur distance of two Banach spaces $`X`$ and $`Y`$ is defined by $`\text{dist}(X,Y)=infTT^1`$ where the infimum extends over all surjective isomorphisms $`T:XY`$. To avoid confusion, notice that $`\text{dist}(l^1,[x_n])=1`$ is not the same as $`x_n\stackrel{\mathrm{alm}}{}l^1`$.
Finally, a sequence $`(x_n)`$ is said to span $`l^1`$ asymptotically isometrically or just to span $`l^1`$ asymptotically \- $`x_n\stackrel{\mathrm{asy}}{}l^1`$ in symbols - if there is a sequence $`(\delta _n)`$ in $`[0,1[`$ tending to $`0`$ such that $`_{n=1}^{\mathrm{}}(1\delta _n)|\alpha _n|_{n=1}^{\mathrm{}}\alpha _n\frac{x_n}{x_n}_{n=1}^{\mathrm{}}|\alpha _n|`$ for all scalars $`\alpha _n`$.
Note that the present definitions of almost and asymptotically isometric differ slightly from those in , by the term $`x_n/x_n`$ but that, of course, for normalized sequences the definitions are the same. Note also the technical detail that because of this term one might have $`x_n0`$ for a sequence spanning $`l^1`$ isomorphically whereas sequences that are equivalent to the canonical $`l^1`$-basis (\[4, p. 43\]) are uniformly bounded away from $`0`$. We say that a Banach space is isomorphic (respectively almost isometric respectively asymptotically isometric) to $`l^1`$ if it has a basis with the corresponding property.
Let $`Y`$ be a subspace of a Banach space $`X`$ and $`P`$ be a projection on $`X`$. $`P`$ is called an L-projection provided $`x=Px+(\mathrm{id}_XP)x`$ for all $`xX`$. A subspace $`YX`$ is called an M-ideal in $`X`$ if its annihilator $`Y^{}`$ in $`X^{}`$ is the range of an L-projection on $`X^{}`$. $`Y`$ is called an L-summand in $`X`$ if it is the range of an L-projection on $`X`$. In the special case in which $`X=Y^{\prime \prime }`$ and in which $`Y`$ is an M-summand (respectively an L-summand) in $`Y^{\prime \prime }`$ we say that $`Y`$ is M-embedded (respectively L-embedded). As examples we only mention that preduals of von Neumann algebras, in particular $`l^1`$ and $`L^1`$-spaces, furthermore the Hardy space $`H_0^1`$ and the dual of the disc algebra are L-embedded. The sequence space $`c_0`$, the space of compact operators on a Hilbert space, and the quotient $`C/A`$ of the continuous functions on the unit circle by the disc algebra $`A`$ are examples among M-embedded spaces. It is not difficult but important to see that if there is an L-projection $`P`$ on a Banach space $`X`$ then each contractive projection on $`X`$ which has the same kernel as $`P`$ coincides with $`P`$, see \[13, Prop. I.1.2\]. It follows that if $`X`$ is M-embedded then the canonical inclusion of $`X^{}`$ in $`X^{\prime \prime \prime }`$ is an L-summand in $`X^{\prime \prime \prime }`$ that is $`X^{}`$ is L-embedded; the converse is false \[13, III.1.3\]; in fact, for an L-embedded space being the dual of an M-embedded space can be quite a restrictive condition: For instance, while the dual of any C-algebra is L-embedded only those C-algebras are M-embedded which are isometrically -isomorphic to the algebra of compact operators or to a $`c_0`$-sum of such algebras \[13, III.2.9\]. Throughout this note, if $`X`$ denotes an L-embedded Banach space (which is not always the case) we will write $`X_s`$ for the complement of (the canonical embedding of) $`X`$ in $`X^{\prime \prime }`$ that is $`X^{\prime \prime }=X_1X_s`$. In this case $`P`$ will denote the L-projection from $`X^{\prime \prime }`$ onto $`X`$.
We recall Godefroy’s fundamental result , \[13, IV.2.2\] that L-embedded Banach spaces are $`w`$-sequentially complete. This will be used mostly without reference; together with Rosenthal’s $`l^1`$-theorem a typical application is that each bounded sequence in an L-embedded space contains a subsequence which either spans $`l^1`$ or converges weakly. There is a useful criterion for L-embeddedness of subsapces of L-embedded spaces due to Li ( or \[13, Th. IV.1.2\]) which we state for easy reference:
###### Lemma 1 (Li)
For an L-embedded Banach space $`X`$ (with L-decomposition $`X^{\prime \prime }=X_1X_s`$ and L-projection $`P`$ on $`X^{\prime \prime }`$ with range $`X`$) and a closed subspace $`Y`$ of $`X`$ the following assertions are equivalent.
(i) $`Y`$ is L-embedded.
(ii) $`Y^{}=Y_1(Y^{}X_s)`$.
(iii) $`P\overline{\mathrm{B}_Y}^w^{}=\mathrm{B}_Y`$.
(iv) $`PY^{}=Y`$.
In particular if $`Y`$ is L-embedded and if one identifies $`Y^{\prime \prime }=Y_1Y_s`$ and $`Y^{}X^{\prime \prime }`$ then $`Y_s=Y^{}X_s`$.
Let us finally cite some technical results from which will be used in the sequel.
It is routine to show that sequences spanning $`l^1`$ asymptotically are stable by adding norm-null sequences \[24, Lem. 4\], to be more precise, let $`(x_n)`$, $`(y_n)`$ be two sequences in a Banach space $`X`$ such that $`(x_n)`$ spans $`l^1`$ asymptotically, $`infx_n>0`$, $`y_n0`$ and $`x_n+y_n0`$. Then $`(x_n+y_n)`$ spans $`l^1`$ asymptotically, too.
Although it has been proved in that there are almost isometric $`l^1`$-copies which do not contain asymptotic ones, both notions ”coincide up to subsequences” in L-embedded Banach spaces, more precisely, in L-embedded Banach spaces each sequence spanning $`l^1`$ almost isometrically contains a subsequence spanning $`l^1`$ asymptotically \[23, Cor. 3\].
The following lemma is fundamental for the rest of the paper. It is an immediate consequence of and says that within L-embedded spaces, sequences spanning $`l^1`$ almost or asymptotically isometrically behave like the standard basis of $`l^1`$ as to their $`w^{}`$-accumulation points. Mostly it will be used with $`M`$ being a countable set of normalized elements that span $`l^1`$ asymptotically.
###### Lemma 2
Let $`X`$ be L-embedded (with L-decomposition $`X^{\prime \prime }=X_1X_s`$), let $`MX`$ be a subset of the unit sphere of $`X`$. Then the following assertions are equivalent.
(i) $`M`$ contains a sequence spanning $`l^1`$ asymptotically.
(ii) $`M`$ admits a $`w^{}`$-accumulation point $`x_sX_s`$ of norm one.
In fact, each $`w^{}`$-accumulation point of a sequence spanning $`l^1`$ asymptotically lies in $`X_s`$.
Proof: (i)$``$(ii) and the second statement: see last paragraph of the proof of \[23, Lem. 1\]; (ii)$``$(i): \[23, Th. 2\]
§3 Abstract measure topology
To prepare the definiton of an abstract measure topology we recall some facts about sequential spaces (see for example \[8, 1.6-1.7\] or , cf. also ) because the topology will be defined by determining the class of its convergent sequences.
A topological space is called a sequential space if closedness and sequential closedness coincide. A topological space is called a Fréchet space if closure and sequential closure coincide. Clearly first countable spaces are Fréchet spaces and Fréchet spaces are sequential spaces. Recall that a topological space is a $`\text{T}_1`$-space if every one-point set is closed; this happens if each convergent sequence has a unique limit.
A $`Li^{}`$-space<sup>1</sup><sup>1</sup>1To avoid confusion with the letter L like in L-embedded, L-projection, L-structure we prefer the notation $`Li^{}`$instead of $`^{}`$ as in . is a triple $`(X,𝒞,li^{})`$ where $`X`$ is a set, $`𝒞X^{\mathrm{I}N}`$ a class of sequences of $`X`$ (called the convergence class) and $`li^{}:𝒞X`$ a map (called limit operator) satisfying the following conditions (L1)-(L3). We write $`li^{}x_n`$ instead of $`li^{}((x_n)_{n\mathrm{I}N})`$; the elements of $`𝒞`$ are called $`𝒞`$-convergent sequences.
Let $`(x_n)`$ be any sequence in $`X`$, let $`xX`$.
(L1) If $`x_n=x`$ for all $`n\mathrm{I}N`$ then $`(x_n)𝒞`$ and $`li^{}x_n=x`$.
(L2) If $`(x_n)𝒞`$ with $`li^{}x_n=x`$ then $`(x_{n_k})𝒞`$ and $`li^{}x_{n_k}=x`$ for each subsequence $`(x_{n_k})`$ of $`(x_n)`$.
(L3) If a sequence $`(x_n)`$ is such that there is $`xX`$ and any subsequence $`(x_{n_k})`$ contains a further subsequence $`(x_{n_{k_m}})`$ such that $`(x_{n_{k_m}})𝒞`$ and $`li^{}x_{n_{k_m}}=x`$ then $`(x_n)𝒞`$ and $`li^{}x_n=x`$.
On a $`Li^{}`$-space $`(X,𝒞,li^{})`$ one defines a topology $`\tau _{li^{}}`$ \- called the sequential topology induced by $`li^{}`$ \- by taking as the family of closed sets all $`li^{}`$-sequentially closed sets; here we call a set $`A`$ $`li^{}`$-sequentially closed if $`li^{}x_nA`$ for all $`𝒞`$-convergent sequences $`(x_n)`$ contained in $`A`$. It is elementary to verify that in this way one indeed obtains a topology and that the $`\tau _{li^{}}`$-convergent sequences are exactly the $`𝒞`$-convergent sequences, that is for every sequence $`(x_n)`$ in $`X`$ one has $`li^{}x_n=x`$ if and only if $`x_n\stackrel{\tau _{li^{}}}{}x`$.
An $`𝒮^{}`$-space is a $`Li^{}`$-space $`(X,𝒞,li^{})`$ satisfying additionally
(L4) If $`(x_m^{(n)})_m𝒞`$ for all $`n\mathrm{I}N`$ and $`(x_n)𝒞`$ such that $`li^{}x_m^{(n)}=x_n`$ for all $`n\mathrm{I}N`$ and $`li^{}x_n=x`$ then there exist two sequences $`(n_k)`$, $`(m_k)`$ such that $`li^{}x_{m_k}^{m_k}=x`$.
Endowed with $`\tau _{li^{}}`$, an $`𝒮^{}`$-space becomes a Fréchet space \[8, 1.7.18,19\].
As already mentioned in the introduction the following definition is patterned after a characterization of bounded measure-null sequences in the preduals of finite von Neumann algebras \[24, Th. 1\].
###### Definition 3
Let $`X`$ be a Banach space. A system $`\tau _\mu `$ of subsets of $`X`$ is called an abstract measure topology if it satisfies the following four conditions.
1. $`(X,\tau _\mu )`$ is a sequential space in which every convergent sequence has a unique limit.
2. $`\tau _\mu `$ is weaker than the norm topology.
3. $`\tau _\mu `$ is translation invariant for sequences more precisely, $`x_n\stackrel{\tau _\mu }{}x`$ if and only if $`x_nx\stackrel{\tau _\mu }{}0`$ for any sequence $`(x_n)`$ in $`X`$.
4. Each bounded sequence in $`X`$ that spans $`l^1`$ asymptotically $`\tau _\mu `$-converges to $`0`$,
and each sequence in $`X`$ that $`\tau _\mu `$-converges to $`0`$ is bounded and contains a subsequence which spans $`l^1`$ asymptotically or tends to $`0`$ in norm.
The following result is quite easy to prove. Nevertheless, because of its importance, we call it a theorem.
###### Theorem 4
Every L-embedded Banach space admits an abstract measure topology.
Proof: Let $`X`$ be an L-embedded Banach space with L-decomposition $`X^{\prime \prime }=X_1X_s`$. Set
$`𝒞_0=\{(x_n)|`$ $`(x_n)\text{ is bounded and every subsequence }(x_{n_k})\text{ of }(x_n)\text{ contains}`$
$`\text{a subsequence }(x_{n_{k_l}})\text{ such that }x_{n_{k_l}}\stackrel{\mathrm{asy}}{}l^1\text{ or }x_{n_{k_l}}0\},`$
$`𝒞=\{(x_n)|`$ $`\text{ there exists }xX\text{ such that }(x_nx)𝒞_0\}.`$
We define a limit operator $`li^{}:𝒞X`$ by $`li^{}x_n=x`$ where $`xX`$ is such that $`(x_nx)𝒞_0`$. To show that $`(X,𝒞,li^{})`$ is a $`Li^{}`$-space the only thing to verify is that $`li^{}`$ is well defined as a map. because then conditions (L1) - (L3) are immediate from the definiton of $`𝒞`$.
Suppose that there are $`x,yX`$, $`(x_n)𝒞`$ such that both $`(x_nx)𝒞_0`$ and $`(x_ny)𝒞_0`$. If $`(x_nx)`$ or $`(x_ny)`$ admits a subsequence tending to $`0`$ in norm then $`x=y`$ because the norm topology is Hausdorff. Otherwise, after passing to an appropriate subsequence, we suppose that both sequences are uniformely bounded away from $`0`$ in norm and that both $`x_nx\stackrel{\mathrm{asy}}{}l^1`$ and $`x_ny\stackrel{\mathrm{asy}}{}l^1`$. Since both sequences are bounded, by Lemma 2 they admit two $`w^{}`$-accumulation points $`x_s,y_sX_s`$ and there is a net $`(x_{n_\gamma })`$ such that $`x_{n_\gamma }x\stackrel{w^{}}{}x_s`$ and $`x_{n_\gamma }y\stackrel{w^{}}{}y_s`$. But this means that $`x_{n_\gamma }\stackrel{w^{}}{}x+x_s=y+y_s`$ whence $`x=y`$ (and $`x_s=y_s`$).
We define the abstract measure topology $`\tau _\mu `$ as the sequential topology induced by $`li^{}`$. It is immediate from the definition of $`𝒞`$ that $`\tau _\mu `$ satisfies the conditions of Definition 3.
Lemma 5 shows some elementary properties of $`\tau _\mu `$ which in the sequel will be used mostly without reference.
###### Lemma 5
If a Banach space $`X`$ admits an abstract measure topology $`\tau _\mu `$ then $`\tau _\mu `$ has the following properties.
(a) $`(X,\tau _\mu )`$ is a $`\text{T}_1`$-space.
(b) The relative topology of $`\tau _\mu `$ on a subspace of $`X`$ is again an abstract measure topology.
(c) Closedness and sequential closedness coincide for $`\tau _\mu `$. Sequentially continuous maps on $`X`$ are continuous.
(d) If $`X`$ does not contain a copy of $`l^1`$ the norm topology is an abstract measure topology and is the only one.
(e) $`\tau _\mu `$ is unique.
(f) If $`X`$ is the predual of a finite von Neumann algebra then $`\tau _\mu `$ coincides on bounded sets with the usual measure topology; on unbounded sets it does not in general.
(g) Multiplication by scalars is $`\tau _\mu `$-continuous.
(h) If $`X`$ is L-embedded (i.e. $`X^{\prime \prime }=X_1X_s`$) and if a net $`(x_\gamma )`$ in $`X`$ $`w^{}`$-converges to $`x_sX_s`$ such that $`x_\gamma x_s`$ then $`x_\gamma \stackrel{\tau _\mu }{}0`$.
(i) If $`X`$ is L-embedded then for any $`x^{\prime \prime }\overline{\mathrm{B}_{X_s}}^w^{}`$ there exists a net $`(x_\gamma )`$ in $`\mathrm{B}_X`$ such that both $`x_\gamma \stackrel{\tau _\mu }{}0\text{ and }x_\gamma \stackrel{w^{}}{}x^{\prime \prime }`$.
Proof: (a) - (d) are clear from the definition.
(e) The topology of sequential spaces is determined by its convergent sequences. Thus the abstract measure topology is unique because the conditions of Definition 3 determine all convergent sequences.
(f) The assertion concerning bounded sets follows from \[24, Th. 1\]. For the assertion concerning unbounded sets we consider the usual measure (=pointwise) topology and the $`w^{}`$-topology on $`l^1`$: These two and $`\tau _\mu `$ coincide on bounded sets. But the unbounded sequence $`(ne_n)`$ converges in the usual measure topology while it does not with respect to $`\tau _\mu `$. (Here $`(e_n)`$ denotes the standard basis of $`l^1`$.)
(g) Let $`\lambda _n\lambda `$ in $`\mathrm{C}\text{ }`$, and $`x_n\stackrel{\tau _\mu }{}x`$ in $`X`$. If $`\lambda _n0`$ or $`x_nx0`$ then $`\lambda _nx_n\lambda x`$ with respect to the norm-topology and thus also with respect to $`\tau _\mu `$. If $`inf|\lambda _n|>0`$, $`x_nx\stackrel{\mathrm{asy}}{}l^1`$, and $`infx_nx>0`$ then $`\lambda _nx_n\lambda _nx\stackrel{\mathrm{asy}}{}l^1`$ whence $`\lambda _nx_n\lambda x\stackrel{\mathrm{asy}}{}l^1`$ because $`\lambda _nx\lambda x0`$ and because sequences spanning $`l^1`$ asymtotically are stable by adding norm-null sequences. Hence $`\lambda _nx_n\stackrel{\tau _\mu }{}\lambda x`$. Up to the usual reasoning on subsequences of subsequences this proves that $`\lambda _n\lambda `$ and $`x_n\stackrel{\tau _\mu }{}x`$ imply $`\lambda _nx_n\stackrel{\tau _\mu }{}\lambda x`$.
(h) We assume to the contrary that $`(x_\gamma )`$ does not $`\tau _\mu `$-converge to $`0`$. Then there is a $`\tau _\mu `$-neighborhood $`𝒪`$ of $`0`$ and a subnet $`(x_\gamma ^{})`$ which does not meet $`𝒪`$ and which still $`w^{}`$-converges to $`x_s`$. But by Lemma 2 there is a sequence $`(x_{\gamma _n^{}})`$ that spans $`l^1`$ almost isometrically i.e. $`x_{\gamma _n^{}}\stackrel{\tau _\mu }{}0`$ hence $`(x_{\gamma _n^{}})`$ does meet $`𝒪`$.
(i) Let $`𝒰`$ be a $`w^{}`$-neighbourhood of $`x^{\prime \prime }`$. Then there exists a $`x_s\mathrm{B}_{X_s}𝒰`$ and $`𝒰`$ is also a $`w^{}`$-neighbourhood of $`x_s`$. Let $`(x_\gamma )\mathrm{B}_X`$ be a net with $`w^{}`$-limit $`x_s`$ and such that $`x_\gamma =x_s`$. Let $`𝒱`$ be a $`\tau _\mu `$-neighbourhood of $`0`$. Then $`x_\gamma \stackrel{\tau _\mu }{}0`$ by (h) and there is $`\gamma _0`$ such that $`x_\gamma 𝒰`$ and $`x_\gamma 𝒱`$ for all $`\gamma \gamma _0`$. That is $`𝒱𝒰\mathrm{B}_X\mathrm{}`$ whence the assertion.
Remarks:
1. Part (h) of the Lemma above corresponds to \[13, IV.3.7\], part (i) to \[10, Lem. 2.2\]. Since in general $`\mathrm{B}_{X_s}`$ is not $`w^{}`$-closed there are $`\tau _\mu `$-null sequences which differ from the ones described in part (h). It is tempting to suppose that the $`\tau _\mu `$-null nets are exactly those that admit $`w^{}`$-limits in $`\overline{\mathrm{B}_{X_s}}^w^{}`$ but at the time of this writing this is not at all clear.
2. There exists a non-Hausdorff Fréchet space in which every sequence has at most one limit \[8, 1.6E\]. Therefore the question whether $`\tau _\mu `$ is Hausdorff is not necessarily trivial.
3. If $`X,Y`$ are two L-embedded Banach spaces, $`YX`$ a subspace of $`X`$, then by (b) and (e) one can identify the intrinsic abstract measure topology of $`Y`$ with the relative topology of the abstract measure topology of $`X`$. Therefore, in theorems like Theorem 13 or Proposition 16 it is no longer necessary to consider a sourrounding L-embedded Banach space like $`\mathrm{L}^1(\mu )`$ in the corresponding theorems \[13, IV.3.10\] or \[10, Prop. 2.1\]. This observation sheds also some light on the remarks after \[13, IV.3.5\] and after \[13, Def. IV.4.2\] on the use of ”nicely placed” and ”L-embedded”.
4. It might be usefull for other purposes to modify Definition 3. For example one could replace ”asymptotically” by ”almost isometrically” in Definition 3; by \[23, Cor. 3\] this would give the same topology for L-embedded spaces. As a less trivial modification one could first restrict Definition 3 to bounded subsets of a Banach space $`X`$ and then define the abstract measure topology on the whole of $`X`$ as the inductive limit of the family $`\tau _\mu |_{n\mathrm{B}_x}`$. In this case, the results of this paper would remain valid up to some minor modifications because in all proofs except for Lemma 5 only the restriciton of $`\tau _\mu `$ to the unit ball is considered. In passing we note that we restrict our attention to bounded sets mainly because the characterization of measure-null sequences in does not work for unbounded sequences, see the second remark after the proof of Theorem 1 in .
We end this section with a modest attempt to get closer to the sequential structure of $`\tau _\mu `$. In case the addition is $`\tau _\mu `$-continuous, at least on bounded sets $`\tau _\mu `$ gives a Fréchet space. This applies to von Neumann preduals which in §6 below will be shown to have $`\tau _\mu `$-continuous addition.
###### Lemma 6
Let $`X`$ be an L-embedded Banach space. If the addition is $`\tau _\mu `$-continuous then the restriction of $`\tau _\mu `$ to a bounded subset of $`X`$ makes this set a Fréchet space.
Proof: To avoid trivialities we consider a bounded sequence $`x_n\stackrel{\mathrm{asy}}{}l^1`$ in $`X`$, and a uniformely bounded sequence of sequences $`(x_m^{(n)})_{m\mathrm{I}N}`$ such that $`x_m^{(n)}x_n\stackrel{\mathrm{asy}}{}l^1`$ for all $`n\mathrm{I}N`$. It is enough to show the existence of a sequence $`(m_n)`$ such that $`x_{m_n}^{(n)}\stackrel{\tau _\mu }{}0`$ because this will prove that $`(X,𝒞,li^{})`$ of the proof of Theorem 4 satisfies (L4).
We set $`y_m^{(n)}=x_m^{(n)}x_n`$. Since all $`x_m^{(n)}`$ are uniformely bounded there is no loss of generality if we suppose that $`y_m^{(n)}=1`$ for all $`m`$ and all $`n`$. By Lemma 2 each universal net $`(y_{m_\gamma }^{(n)})_{\gamma \mathrm{\Gamma }}`$ $`w^{}`$-converges to a limit $`y_s^{(n)}X_s`$ of norm one. What follows is a straightforward modification of the proof of \[23, Th. 2\]. Let $`(\delta _n)`$ be a sequence of strictly positive numbers in $`]0,1]`$ converging to $`0`$. Set $`\eta _1=\frac{1}{3}\delta _1`$ and $`\eta _{n+1}=\frac{1}{3}\mathrm{min}(\eta _n,\delta _{n+1})`$ for $`n\mathrm{I}N`$. By induction over $`n\mathrm{I}N`$ one constructs $`m_n\mathrm{I}N`$ such that
$`\left({\displaystyle \underset{i=1}{\overset{n}{}}}(1\delta _i)|\alpha _i|\right)+\eta _n{\displaystyle \underset{i=1}{\overset{n}{}}}|\alpha _i|{\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _iy_{m_i}^{(i)}\text{ for all }n\mathrm{I}N,\alpha _i\mathrm{C}.`$ (1)
The first induction step is settled by setting $`m_1=1`$. For the induction step $`nn+1`$ fix an element $`\alpha =(\alpha _i)_{i=1}^{n+1}`$ in the unit sphere of $`l_{n+1}^1`$ such that $`\alpha _{n+1}0`$. The $`w^{}`$-convergence (along $`\gamma `$) of $`(_{i=1}^n\alpha _iy_{m_i}^{(i)})+\alpha _{n+1}y_{m_\gamma }^{(n+1)}`$ to $`(_{i=1}^n\alpha _iy_{m_i}^{(i)})+\alpha _{n+1}y_s^{(n+1)}`$ yields
$`\underset{\gamma }{lim\; inf}\left({\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _iy_{m_i}^{(i)}\right)+\alpha _{n+1}y_{m_\gamma }^{(n+1)}`$ $``$ $`\left({\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _iy_{m_i}^{(i)}\right)+\alpha _{n+1}y_s^{(n+1)}`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _iy_{m_i}^{(i)}+|\alpha _{n+1}|`$
$`\stackrel{(\text{1})}{}`$ $`\left({\displaystyle \underset{i=1}{\overset{n+1}{}}}(1\delta _i)|\alpha _i|\right)+\mathrm{min}(\eta _n,\delta _{n+1})`$
whence
$`\left({\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _iy_{m_i}^{(i)}\right)+\alpha _{n+1}y_m^{(n+1)}\left({\displaystyle \underset{i=1}{\overset{n+1}{}}}(1\delta _i)|\alpha _i|\right)+3\eta _{n+1}`$
for infinitely many $`m`$. This extends to a finite $`\eta _{n+1}`$-net $`(\alpha ^l)_{l=1}^{L_{n+1}}`$ in the unit sphere of $`l_{n+1}^1`$ with $`2\eta _{n+1}`$ instead of $`3\eta _{n+1}`$, to all $`\alpha `$ in the unit sphere of $`l_{n+1}^1`$ with $`\eta _{n+1}`$ instead of $`2\eta _{n+1}`$, and finally to all $`\alpha l_{n+1}^1`$. The details are the same as in the proof of \[23, Th. 2\].
This ends the induction. By (1) we have $`y_{m_n}^{(n)}\stackrel{\tau _\mu }{}0`$. Hence $`x_{m_n}^{(n)}=y_{m_n}^{(n)}+x_n\stackrel{\tau _\mu }{}0`$ because addition is supposed to be $`\tau _\mu `$-continuous.
§4 Section IV.3 of (partly) revisited
Proposition 7 generalizes some well known facts, see e.g. \[7, Th. IV.8.12\] for (a) and \[13, p. 202\] for (b).
###### Proposition 7
Let $`X`$ be an L-embedded Banach space with its abstract measure topology $`\tau _\mu `$. Then the following statements hold.
(a) A sequence converges in norm if (and only if) it converges both weakly and with respect to $`\tau _\mu `$, and all limits coincide.
(b) A norm closed subspace $`YX`$ is reflexive if and only if $`\tau _\mu `$ and the norm topology coincide on the unit ball of $`Y`$.
Proof: (a) The statement is almost immediate from the definiton of $`\tau _\mu `$: First remark that a sequence which $`\tau _\mu `$-converges to $`0`$ contains a subsequence which either converges to $`0`$ in norm or is uniformely bounded away from 0 in norm and spans $`l^1`$ (asymptotically); but the latter case is excluded if the sequence also converges weakly (to whatever limit) because $`l^1`$-bases do not converge weakly. Now let $`(x_n)`$ be a sequence in $`X`$, let $`x,yX`$ be such that both $`x_n\stackrel{\tau _\mu }{}x`$ and $`x_n\stackrel{w}{}y`$. Then by what has just been remarked, for each subsequence $`(x_{n_k})`$ there is a subsequence $`(x_{n_{k_l}})`$ such that $`x_{n_{k_l}}x0`$ in norm for any subsequence $`(x_{n_k})`$ whence $`x=y`$ and the assertion follows.
(b) By $`\tau _{}`$ we denote the norm topology on $`X`$. Let $`Y`$ be reflexive. To show that $`\tau _\mu `$ and $`\tau _{}`$ coincide on the unit ball $`\mathrm{B}_Y`$ of $`Y`$ it is enough to show that each subsequence of a $`\tau _\mu `$-convergent sequence in $`\mathrm{B}_Y`$ admits a subsequence which converges in norm to the same limit. But if $`Y`$ is reflexive then each bounded sequence contains a weakly convergent subsequence which if the sequence is also $`\tau _\mu `$-convergent converges in norm by (a). Thus $`\tau _\mu `$ and $`\tau _{}`$ coincide on $`\mathrm{B}_Y`$.
Conversely suppose that $`\tau _\mu `$ and $`\tau _{}`$ coincide on $`\mathrm{B}_Y`$. In order to prove that $`Y`$ is reflexive it is enough to prove that $`Y`$ does not contain isomorphic copies of $`l^1`$ because by Rosenthal’s theorem \[4, Ch. XI\] in the absence of $`l^1`$ each bounded sequence contains a weak Cauchy subsequence which then converges weakly by the weak sequential completeness of $`X`$. But if $`Y`$ contained an isomorphic copy of $`l^1`$ then by James’ distortion theorem it would contain also an almost isometric copy of $`l^1`$. By \[23, Cor. 3\] $`Y`$ would contain an asymptotic copy of $`l^1`$. It’s canonical normalized basis would $`\tau _\mu `$-converge to $`0`$; finally, since $`\tau _\mu `$ and $`\tau _{}`$ coincide on $`\mathrm{B}_Y`$ this basis would even converge in norm to $`0`$, a contradiction which proves that $`Y`$ is reflexive.
Lemma 8 is the technical key for the other results of this section. It corresponds to \[13, IV.3.1\].
###### Lemma 8
Let $`X`$ be an L-embedded Banach space with L-projection $`P`$ from $`X^{\prime \prime }`$ onto $`X`$. Then for every net $`(x_\gamma )_{\gamma \mathrm{\Gamma }}`$ in $`X`$ $`w^{}`$-converging to $`x^{\prime \prime }X^{\prime \prime }\backslash X`$ there is a bounded sequence $`(y_n)`$ in $`\text{co}\{x_\gamma |\gamma \mathrm{\Gamma }\}`$ such that the sequence $`(y_nPx^{\prime \prime })`$ spans $`l^1`$ asymtotically isometrically.
Proof: Set $`x=Px^{\prime \prime }`$, $`x_s=x^{\prime \prime }Px^{\prime \prime }`$. Choose a net $`(z_\gamma )_{\gamma \mathrm{\Gamma }}`$ in $`X`$ such that $`z_\gamma \stackrel{w^{}}{}x_s`$ and $`z_\gamma =x_s`$. We assume without loss of generality that both nets $`(x_\gamma )`$, $`(z_\gamma )`$ are indexed by the same directed set $`\mathrm{\Gamma }`$. Then $`x_\gamma z_\gamma \stackrel{w}{}x`$.
The idea of the proof is that on one hand by the theorem of Hahn-Banach the net $`(x_\gamma z_\gamma )`$ admits convex combinations which converge to $`x`$ in norm and that on the other hand by a slight modification of Godefroy’s construction the corresponding convex combinations of the $`x_\gamma `$ can be chosen so to span $`l^1`$ asymtotically. Here are the details.
Since $`x^{\prime \prime }X`$ we have $`x_s0`$ and thus without loss of generality we suppose $`x_s=z_\gamma =1`$. Let $`(\delta _n)`$ be a sequence of numbers in $`]0,1[`$ convergent to $`0`$. Set $`\eta _1=\frac{1}{4}\delta _1`$ and $`\eta _{n+1}=\frac{1}{4}\mathrm{min}(\eta _n,\delta _{n+1})`$ for $`n\mathrm{I}N`$. By induction over $`n\mathrm{I}N`$ we will construct finite sets $`A_n\mathrm{I}N`$, finite sequences $`(\lambda _k)_{kA_n}`$ in $`[0,1]`$ and $`(\gamma _k)_{kA_n}`$ in $`\mathrm{\Gamma }`$ such that
$`{\displaystyle \underset{kA_n}{}}\lambda _k`$ $`=`$ $`1,G_iG_n=\mathrm{},G_n\mathrm{\Gamma }_ni<n`$ (3)
$`(y_nx){\displaystyle \underset{kA_n}{}}\lambda _kz_{\gamma _k}<\eta _n.`$
$`{\displaystyle \underset{i=1}{\overset{n}{}}}|\alpha _i|(1+\eta _i)`$ $``$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _i(y_ix)\left({\displaystyle \underset{i=1}{\overset{n}{}}}(1\delta _i)|\alpha _i|\right)+\eta _n{\displaystyle \underset{i=1}{\overset{n}{}}}|\alpha _i|\alpha _i\mathrm{C}`$ (4)
where
$`y_n={\displaystyle \underset{kA_n}{}}\lambda _kx_{\gamma _k},G_n=\{\gamma _k|kA_n\},\mathrm{\Gamma }_n=\mathrm{\Gamma }\backslash {\displaystyle \underset{i=1}{\overset{n1}{}}}G_i(=\mathrm{\Gamma }\text{ if }n=1).`$
For the first induction step $`n=1`$ we choose $`x^{}\mathrm{B}_X^{}`$ such that $`1=x_s\mathrm{Re}x_s(x^{})>1\delta _1+2\eta _1`$. The $`w^{}`$-convergence of $`(z_\gamma )`$ to $`x_s`$ yields $`\beta _1\mathrm{\Gamma }`$ such that
$`\mathrm{Re}x^{}(z_\gamma )>1\delta _1+2\eta _1\gamma \beta _1.`$ (5)
The net $`((x_\gamma x)z_\gamma )_{\gamma \beta _1}`$ $`w`$-converges to $`0`$ thus by the theorem of Hahn-Banach we find a convex combination $`y_1=_{kA_1}\lambda _kx_{\gamma _k}`$ such that
$$(y_1x)\underset{kA_1}{}\lambda _kz_{\gamma _k}<\eta _1.$$
Thus (4, $`n=1`$) follows from $`y_1x_{kA_1}\lambda _kz_{\gamma _k}+\eta _11+\eta _1`$ and from
$`y_1x`$ $`>`$ $`{\displaystyle \underset{kA_1}{}}\lambda _kz_{\gamma _k}\eta _1`$
$``$ $`\eta _1+\mathrm{Re}{\displaystyle \underset{kA_1}{}}\lambda _kx^{}(z_{\gamma _k})`$
$`>`$ $`\eta _1+{\displaystyle \underset{kA_1}{}}\lambda _k(1\delta _1+2\eta _1)=1\delta _1+\eta _1.`$
For the induction step $`nn+1`$ we suppose $`A_i\mathrm{\Gamma }`$, $`(\lambda _k)_{kA_i}[0,1]`$, $`G_i\mathrm{\Gamma }`$ to be constructed according to (3) - (4) for $`i=1,\mathrm{},n`$.
First we consider the $`w`$-convergent net $`(x_\gamma z_\gamma )_{\gamma \mathrm{\Gamma }_{n+1}}`$. By the theorem of Hahn-Banach we can choose a finite set $`A_{n+1}\mathrm{I}N`$, numbers $`(\lambda _k)_{kA_{n+1}}[0,1]`$ and indices $`(\gamma _k)_{kA_{n+1}}\mathrm{\Gamma }_{n+1}`$ such that (3) and (3) hold for $`n+1`$. Together with (4, $`n`$) this gives the first inequality of (4, $`n+1`$).
We fix an element $`\alpha =(\alpha _i)`$ in the unit sphere of $`l_{n+1}^1`$ such that $`\alpha _{n+1}0`$ and use the L-decomposition of $`X^{\prime \prime }=X_1X_s`$ in order to get
$`\left({\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _i(y_ix)\right)+\alpha _{n+1}x_s`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _i(y_ix)+\alpha _{n+1}x_s`$
$`\stackrel{(\text{4})}{}`$ $`\left({\displaystyle \underset{i=1}{\overset{n}{}}}(1\delta _i)|\alpha _i|\right)+\eta _n\left({\displaystyle \underset{i=1}{\overset{n}{}}}|\alpha _i|\right)+|\alpha _{n+1}|`$
$`=`$ $`\left({\displaystyle \underset{i=1}{\overset{n+1}{}}}(1\delta _i)|\alpha _i|\right)+\eta _n(\eta _n\delta _{n+1})|\alpha _{n+1}|`$
$``$ $`\left({\displaystyle \underset{i=1}{\overset{n+1}{}}}(1\delta _i)|\alpha _i|\right)+\mathrm{min}(\eta _n,\delta _{n+1})`$
$`=`$ $`\left({\displaystyle \underset{i=1}{\overset{n+1}{}}}(1\delta _i)|\alpha _i|\right)+4\eta _{n+1}`$
because $`\alpha =1`$ and $`|\alpha _{n+1}|1`$. Thus there is $`x^{}\mathrm{B}_X^{}`$ (depending on $`\alpha `$) such that
$`\left({\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _i(y_ix)\right)+\alpha _{n+1}x_s`$ $``$ $`\text{Re}\left(\left({\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _i(y_ix)\right)+\alpha _{n+1}x_s\right)(x^{})`$
$`>`$ $`\left({\displaystyle \underset{i=1}{\overset{n+1}{}}}(1\delta _i)|\alpha _i|\right)+3\eta _{n+1}.`$
Then the $`w^{}`$-convergence (along $`\gamma \mathrm{\Gamma }_{n+1}`$) of $`\left((_{i=1}^n\alpha _i(y_ix))+\alpha _{n+1}z_\gamma \right)`$ to $`(_{i=1}^n\alpha _i(y_ix))+\alpha _{n+1}x_s`$ yields $`\beta \mathrm{\Gamma }_{n+1}`$ (depending on $`\alpha `$ and $`x^{}`$) such that
$`\text{Re}\left(x^{}\left(({\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _i(y_ix))+\alpha _{n+1}z_\gamma \right)\right)>\left({\displaystyle \underset{i=1}{\overset{n+1}{}}}(1\delta _i)|\alpha _i|\right)+3\eta _{n+1}\gamma \beta .`$
Choose a finite $`\eta _{n+1}`$-net $`(\alpha ^l)_{l=1}^{L_{n+1}}`$ in the unit sphere of $`l_{n+1}^1`$ in the sense that for each $`\alpha `$ in the unit sphere of $`l_{n+1}^1`$ there is $`lL_{n+1}`$ such that $`\alpha \alpha ^l=_{i=1}^{n+1}|\alpha _i\alpha _i^l|<\eta _{n+1}`$. Then we may repeat the reasoning above finitely many times for $`l=1,\mathrm{},L_{n+1}`$ in order to get $`\beta _{n+1}\mathrm{\Gamma }_{n+1}`$ and $`x_l^{}\mathrm{B}_X^{}`$ for $`lL_{n+1}`$ such that
$`\text{Re}\left(x_l^{}\left(({\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _i^l(y_ix))+\alpha _{n+1}^lz_\gamma \right)\right)>\left({\displaystyle \underset{i=1}{\overset{n+1}{}}}(1\delta _i)|\alpha _i^l|\right)+3\eta _{n+1}lL_{n+1},\gamma \beta _{n+1}.`$ (6)
For each $`lL_{n+1}`$ we get that
$`{\displaystyle \underset{i=1}{\overset{n+1}{}}}\alpha _i^l(y_ix)`$ $`\stackrel{(\text{3})}{}`$ $`\eta _{n+1}+\left({\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _i^l(y_ix)\right)+\alpha _{n+1}^l{\displaystyle \underset{kA_{n+1}}{}}\lambda _kz_{\gamma _k}`$ (7)
$`=`$ $`\eta _{n+1}+{\displaystyle \underset{kA_{n+1}}{}}\lambda _k\left(({\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _i^l(y_ix))+\alpha _{n+1}^lz_{\gamma _k}\right)`$ (8)
$``$ $`\eta _{n+1}+{\displaystyle \underset{kA_{n+1}}{}}\lambda _k\mathrm{Re}x_l^{}\left(({\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _i^l(y_ix))+\alpha _{n+1}^lz_{\gamma _k}\right)`$ (9)
$`\stackrel{(\text{6})}{}`$ $`\left({\displaystyle \underset{i=1}{\overset{n+1}{}}}(1\delta _i)|\alpha _i^l|\right)+2\eta _{n+1}`$ (10)
For an arbitrary $`\alpha `$ in the unit sphere of $`l_{n+1}^1`$ choose $`lL_{n+1}`$ such that $`\alpha \alpha ^l<\eta _{n+1}`$. Then
$`{\displaystyle \underset{i=1}{\overset{n+1}{}}}\alpha _i(y_ix)`$ $``$ $`{\displaystyle \underset{i=1}{\overset{n+1}{}}}\alpha _i^l(y_ix){\displaystyle \underset{i=1}{\overset{n+1}{}}}(\alpha _i\alpha _i^l)(y_ix)`$
$``$ $`\left({\displaystyle \underset{i=1}{\overset{n+1}{}}}(1\delta _i)|\alpha _i|\right)+2\eta _{n+1}\alpha \alpha ^l`$
$``$ $`\left({\displaystyle \underset{i=1}{\overset{n+1}{}}}(1\delta _i)|\alpha _i|\right)+\eta _{n+1}`$
$`=`$ $`\left({\displaystyle \underset{i=1}{\overset{n+1}{}}}(1\delta _i)|\alpha _i|\right)+\eta _{n+1}{\displaystyle \underset{i=1}{\overset{n+1}{}}}|\alpha _i|.`$
This extends to all scalars $`\alpha _i\mathrm{C}\text{ }`$ and thus ends the induction. The sequence $`(y_n)`$ is bounded because of (3). By (4) the sequence $`(y_nx)`$ is easily seen to span $`l^1`$ asymptotically. This ends the proof.
Remark: The proof yields not only $`y\text{co}\{x_\gamma |\gamma \mathrm{\Gamma }\}`$ but separated blocks $`y_n=_{kA_n}\lambda _kx_{\gamma _k}`$ where the sets $`\{x_{\gamma _k}|kA_n\}`$ are pairwise disjoint. Moreover one can obtain, given a sequence $`(\gamma _n^{})`$ in $`\mathrm{\Gamma }`$, that $`x_{\gamma _k}\gamma _n^{}`$ for $`kA_n`$.
In general the unit ball of $`X`$ is not $`\tau _\mu `$-compact; the Rademacher functions $`r_n`$ in $`\mathrm{L}^1([0,1])`$ which are bounded without having a measure convergent subsequence provide a counterexample. \[If $`(r_n)`$ contained a measure convergent subsequence this subsequence would admit a norm convergent subsequence by Proposition 7 (a) because $`(r_n)`$ spans $`l^2`$ isomorphically.\] But we have:
###### Corollary 9
Every bounded sequence in an L-embedded Banach space admits a sequence of convex combinations which converges with respect to the measure topology.
Proof: Let $`(x_n)`$ be a bounded sequence in an L-embedded space $`X`$ and let $`(x_{n_\gamma })`$ be a universal net that $`w^{}`$-converges to $`x^{\prime \prime }X^{\prime \prime }`$ by the $`w^{}`$-compactness of $`\mathrm{B}_{X^{\prime \prime }}`$. If $`x^{\prime \prime }X`$ then this net $`w`$-converges, admits norm-convergent convex combinations and we are done in this case. Otherwise $`x^{\prime \prime }`$ lies in $`X^{\prime \prime }\backslash X`$ and one applies Lemma 8 to get a sequence $`(y_n)`$ in $`\text{co}\{x_n|n\mathrm{I}N\}`$ which $`\tau _\mu `$-converges to $`Px^{\prime \prime }`$.
Corollary 9 corresponds to \[13, p. 202\]. In this context there is a natural
Question: Does Komlos’ theorem hold accordingly? More precisely, given a bounded sequence in an L-embedded space, does it admit a aubsequence whose Cesaro (=arithmetic) means converge with respect to the measure topology?
Note that by (a) of Proposition 7 Komlos’ theorem implies the weak Banach-Saks property (which by definition claims that a $`w`$-convergent sequence admits a subsequence whose Cesaro means converge in norm, see for example , \[4, p. 112, 121\], ). By Rosenthal’s $`l^1`$-theorem the weak Banach-Saks property is also half a converse to Komlos’ theorem, that is by Rosenthal’s $`l^1`$-theorem a bounded sequence in an L-embedded space admits a subsequence which is either equivalent to the standard basis of $`l^1`$ or converges weakly; but if one supposes the weak Banach-Saks property to hold then in the second case of a $`w`$-convergent sequence there are Cesaro means that converge in norm whence with respect to the measure topology.
There is another related
Question: Does the Kadec-Pełczyński subsequence decomposition (sometimes also called the Kadec-Pełczyński splitting lemma) hold accordingly? This lemma says that a bounded sequence $`(f_n)`$ in $`\mathrm{L}^1([0,1])`$ admits a subsequence $`(f_{n_k})`$ which can be decomposed in the following sense: there are two bounded sequences $`(g_k)`$, $`(h_k)`$ in $`\mathrm{L}^1([0,1])`$ such that $`f_{n_k}=g_k+h_k`$, the $`g_k`$ are pairwise disjoint, and $`(h_k)`$ converges weakly. In a recent preprint Randrianantoanina showed the Kadec-Pełczyński subsequence decomposition for preduals of von Neumann algebras. Since the latter are known to have the weak Banach-Saks property , Komlos’ theorem follows almost immediately from Randrianantoanina’s result for von Neumann preduals, see Proposition 24 below.
The following theorem generalizes a theorem of Buhvalov-Lozanovskii (, \[13, IV.3.4\]). As in \[13, IV.3.4\] the implication (ii)$``$(i) holds also for unbounded $`C`$.
###### Theorem 10
Let $`X`$ be an L-embedded Banach space with L-projection $`P`$ (on $`X^{\prime \prime }`$ with range $`X`$) and endowed with its abstract measure topology $`\tau _\mu `$. For a norm closed bounded convex set $`CX`$ the following two assertions are equivalent.
(i) $`P\overline{C}^w^{}=C`$ where $`w^{}`$ refers to the $`w^{}`$-topology of $`X^{\prime \prime }`$.
(ii) $`C`$ is $`\tau _\mu `$-closed.
Proof: (i)$``$(ii) Take $`c_nC`$, $`xX`$ with $`c_n\stackrel{\tau _\mu }{}x`$. It is enough to show that $`xC`$ because closedness and sequential closedness coincide for $`\tau _\mu `$. If the $`\tau _\mu `$\- null sequence $`(c_nx)`$ contains a norm convergent subsequence then we are done because in the norm topology $`C`$ is closed. Otherwise an appropriate subsequence $`(c_{n_k}x)`$ spans $`l^1`$ asymptotically and $`infc_{n_k}x>0`$. By $`w^{}`$-compactness of $`\overline{C}^w^{}`$ there is a net $`(c_{n_{k_\gamma }}x)_{\gamma \mathrm{\Gamma }}`$ on $`𝒮=\{c_{n_k}x|k\mathrm{I}N\}`$ that $`w^{}`$-converges to $`\overline{c}x\overline{C}^w^{}x`$. By Lemma 2 we have that $`\overline{c}xX_s`$ because $`C`$ is bounded. Thus $`P(\overline{c}x)=0`$ and $`x=Px=P\overline{c}C`$ by hypothesis.
(ii)$``$(i) This implication is essentially Lemma 8: If $`(c_\gamma )_{\gamma \mathrm{\Gamma }}`$ is a $`w^{}`$-convergent net in $`C`$ with limit $`\overline{c}`$ it is enough to prove that $`P\overline{c}C`$ because the inclusion $`CP\overline{C}^w^{}`$ is trivial. If $`\overline{c}X`$ then there is a sequence of convex combinations of the $`c_\gamma `$ convergng to $`\overline{c}`$ in norm whence $`\overline{c}C`$. Otherwise, if $`\overline{c}X^{\prime \prime }\backslash X`$, by Lemma 8 there is a sequence $`(d_n)\text{co}\{c_\gamma |\gamma \mathrm{\Gamma }\}C`$ which $`\tau _\mu `$-converges to $`P\overline{c}`$ hence $`P\overline{c}C`$ because $`C`$ is $`\tau _\mu `$-closed.
Corollary 11 corresponds to \[13, IV.3.5\]) for the case $`X=\mathrm{L}^1(\mathrm{\Omega },\mathrm{\Sigma },\mu )`$.
###### Corollary 11
Let $`X`$ be an L-embedded Banach space endowed with its abstract measure topology $`\tau _\mu `$. Then a norm closed subspace $`YX`$ is L-embedded if and only if its unit ball $`\mathrm{B}_Y`$ is $`\tau _\mu `$-closed.
The proof is immediate from Theorem 10 with $`C=\mathrm{B}_Y`$ and from Li’s criterion Lemma 1.
Let $`X`$ be a Banach space admitting an abstract measure topology $`\tau _\mu `$. Then we define
$$X^\mathrm{\#}=\{x^{}X^{}|x^{}|_{\mathrm{B}_X}\text{is}\tau _\mu \text{continuous}\}.$$
Remarks:
1. $`X^\mathrm{\#}`$ is a closed subspace of $`X^{}`$. (The proof is left to the reader.)
2. If $`X`$ is a subspace of an L-embedded Banach space then one has $`X^\mathrm{\#}=X^{}`$ if and only if $`X`$ does not contain copies of $`l^1`$. For, in the absence of $`l^1`$, $`\tau _\mu `$ coincides with the norm topology hence $`X^\mathrm{\#}=X^{}`$. Conversely, if $`X`$ contains a copy of $`l^1`$ then by James’ destortion theorem it contains also an almost isometric copy $`U`$ of $`l^1`$ spanned by a normalized basis $`(u_n)`$. Since $`X`$ is contained in an L-embedded space, by \[23, Cor. 3\] we (may pass to an appropriate subsequence and) suppose that $`u_n\stackrel{\tau _\mu }{}0`$. Let $`x^{}X^{}`$ be a Hahn-Banach extension of the functional on $`U`$ which corresponds to $`1l^{\mathrm{}}`$. Then $`x^{}`$ is not $`\tau _\mu `$-continuous on $`\mathrm{B}_X`$ since $`u_n\stackrel{\tau _\mu }{}0`$ but $`x^{}(u_n)1`$. (Compare also with \[13, Rem. (b) p. 186\].)
###### Proposition 12
Let $`X`$ be an L-embedded Banach space (with L-decomposition $`X^{\prime \prime }=X_1X_s`$) endowed with its abstract measure topology $`\tau _\mu `$. Then
$`X^\mathrm{\#}=(X_s)_{}\text{ (}X^{}\text{)}.`$
Proof: ”$``$” Take $`x^{}X^\mathrm{\#}`$ and $`x_sX_s`$. To prove the inclusion we show that $`x_s(x^{})=0`$.
Let $`(x_\gamma )`$ be a net that $`w^{}`$-converges to $`x_s`$ with $`x_\gamma =x_s`$. But then, by Lemma 2 there is a sequence $`(x_{\gamma _n})`$ that spans $`l^1`$ asymptotically. Hence $`x_{\gamma _n}\stackrel{\tau _\mu }{}0`$ and $`x^{}(x_{\gamma _n})0`$ since $`x^{}`$ is $`\tau _\mu `$-continuous on bounded sets. This proves that $`x^{\prime \prime }(x^{})=0`$. (In passing we note that $`x_\gamma \stackrel{\tau _\mu }{}0`$ by (h) of Lemma 5 but that one can not infer from this that $`x^{\prime \prime }(x^{})=0`$ because it is not clear whether a $`\tau _\mu `$-convergent net has a unique limit.)
$``$” Assume that there is $`x^{}(X_s))_{}`$ that is not $`\tau _\mu `$-continuous on $`\mathrm{B}_X`$. Then by the definiton of $`\tau _\mu `$ there are $`\epsilon >0`$ and a sequence $`(x_n)`$ in $`\mathrm{B}_X`$ such that $`x_n\stackrel{\tau _\mu }{}0`$ but $`|x^{}(x_n)|>\epsilon `$ for all $`n\mathrm{I}N`$. Still by definiton of $`\tau _\mu `$ and because $`x^{}`$ is norm-continuous we suppose that $`(x_n)`$ spans $`l^1`$ almost isometrically. By $`w^{}`$-compactness of $`\mathrm{B}_{X^{\prime \prime }}`$ there exists a $`w^{}`$-accumulation point $`x^{\prime \prime }\mathrm{B}_{X^{\prime \prime }}`$ of $`\{x_n|n\mathrm{I}N\}`$. Let $`(x_{n_\gamma })`$ be a net $`w^{}`$-converging to $`x^{\prime \prime }`$. By Lemma 2, $`x^{\prime \prime }X_s`$. Thus $`x^{}(x_{n_\gamma })x^{\prime \prime }(x^{})=0`$ by hypothesis. This contradicts $`|x^{}(x_n)|>\epsilon `$ and proves that $`x^{}`$ is $`\tau _\mu `$-continuous on $`\mathrm{B}_X`$.
Theorem 13 (see also \[13, IV.3.10\]) was proved in for nicely placed (=L-embedded) subspaces $`X`$ of $`\mathrm{L}^1(\mathrm{\Omega },\mathrm{\Sigma },\mu )`$, $`\mu `$ finite. For its proof Proposition 12 plays the same rôle as \[13, IV.3.9\] in the proof of \[13, IV.3.10\]. Recall that if an L-embedded space admits a predual then this predual need not be M-embedded (\[13, p. 102\]).
###### Theorem 13
Let $`X`$ be an L-embedded Banach space endowed with its abstract measure topology $`\tau _\mu `$. The following assertions are equivalent.
(i) $`X`$ is (isometrically isomorphic to) the dual of an M-embedded Banach space $`Z`$.
(ii) $`X^\mathrm{\#}`$ separates $`X`$.
If (i) and (ii) are satisfied $`Z`$ is (isometrically isomorphic to) $`X^\mathrm{\#}`$.
Proof: (i)$``$(ii): By \[13, III.1.3\] the L-projections on $`X^{\prime \prime }`$ and on $`Z^{\prime \prime \prime }`$ with kernel $`X_s=Z^{}X^{\prime \prime }`$ can be identified. Thus $`Z=(Z^{})_{}=(X_s)_{}=X^\mathrm{\#}`$ by Proposition 12; in particular, $`X^\mathrm{\#}`$ separates $`X`$.
(ii)$``$(i): By \[13, IV.1.9\] it is enough to show that $`X_s`$ is $`w^{}`$-closed in $`X^{\prime \prime }`$; then an M-embedded predual of $`X`$ exists and is isometrically isomorphic to $`(X_s)_{}`$ whence to $`X^\mathrm{\#}`$ by Proposition 12. To see that $`X_s`$ is $`w^{}`$-closed we take an element
$`\overline{x}=x+x_s\overline{X_s}^w^{}=((X_s)_{})^{}=(X^\mathrm{\#})^{}`$
with $`xX`$, $`x_sX_s`$. Then $`x=\overline{x}x_s(X^\mathrm{\#})^{}X=(X^\mathrm{\#})_{}=\{0\}`$ where the latter equality comes from the fact that $`X^\mathrm{\#}`$ separates $`X`$. Thus $`\overline{x}=x_sX_s`$ which proves that $`X_s`$ is $`w^{}`$-closed in $`X^{\prime \prime }`$.
Analogously to the $`l^1`$-case we say that a sequence $`(x_n)`$ of nonzero elements in a Banach space $`X`$ spans $`c_0`$ almost (respectivley asymptotically) isometrically if there exists a sequence $`(\delta _m)`$ in $`[0,1[`$ tending to $`0`$ such that $`(1\delta _m)\mathrm{max}_{mnm^{}}|\alpha _n|_{n=m}^m^{}\alpha _n\frac{x_n}{x_n}(1+\delta _m)\mathrm{max}_{mnm^{}}|\alpha _n|`$ for all $`mm^{}`$ (respectively such that $`\mathrm{max}_{nm}(1\delta _n)|\alpha _n|_{n=1}^m\alpha _n\frac{x_n}{x_n}\mathrm{max}_{nm}(1+\delta _n)|\alpha _n|`$ for all $`m\mathrm{I}N`$).
It follows from the proof of \[6, Th. 2\] that the dual of an asymptotic $`c_0`$-copy is an asymptotic $`l^1`$-copy. A similar argument shows that this remains true if ”asymptotic” is replaced by ”almost isometric”. Analogously we get
###### Lemma 14
Let $`X=[x_n]`$ be an almost isometric copy of $`l^1`$ spanned by the normalized basis $`(x_n)`$. Let $`x_n^{}X^{}`$ be the biorthogonal functionals (that is $`x_n^{}(x_k)=\delta _{n,k}`$). Then $`(x_n^{})`$ spans $`c_0`$ almost isometrically.
The same holds with ”almost” replaced by ”asymptotically”.
Proof: First we deal with the case where $`x_n\stackrel{\mathrm{alm}}{}l^1`$.
By hypothesis there is a null sequence $`(\delta _m)[0,1[`$ sucht that
$`(1\delta _m){\displaystyle \underset{m}{\overset{\mathrm{}}{}}}|\beta _n|{\displaystyle \underset{m}{\overset{\mathrm{}}{}}}\beta _nx_n{\displaystyle \underset{m}{\overset{\mathrm{}}{}}}|\beta _n|`$
for all scalars $`\beta _n`$. Let $`mm^{}`$ be arbitrary in $`\mathrm{I}N`$. Since $`x_n=1`$ we have the first inequality of
$`\underset{mnm^{}}{\mathrm{max}}|\alpha _n|{\displaystyle \underset{n=m}{\overset{m^{}}{}}}\alpha _nx_n^{}(1\delta _m)^1\underset{mnm^{}}{\mathrm{max}}|\alpha _n|`$ (11)
for all scalars $`\alpha _n`$. For the second inequality we take any $`x\mathrm{B}_X`$ of the form $`x=\beta _nx_n`$; then
$`\left|{\displaystyle \underset{n=m}{\overset{m^{}}{}}}\alpha _nx_n^{}(x)\right|`$ $`=`$ $`\left|{\displaystyle \underset{n=m}{\overset{m^{}}{}}}\alpha _n\beta _n\right|\underset{mnm^{}}{\mathrm{max}}|\alpha _n|{\displaystyle \underset{n=m}{\overset{m^{}}{}}}|\beta _n|`$
$``$ $`(1\delta _m)^1\underset{mnm^{}}{\mathrm{max}}|\alpha _n|`$
whence the second inequality of (11). It follows from (11) that $`(x_n^{})`$ spans $`c_0`$ almost isometrically.
The case in which $`x_n\stackrel{\mathrm{asy}}{}l^1`$ is proved similarly.
###### Proposition 15
An L-embedded almost isometric copy of $`l^1`$ is the dual of an M-embedded space which is almost isometric to $`c_0`$.
The statement remains true when ”almost” is replaced by ”asymptotically”.
Proof: Let $`X`$ be an L-embedded almost isometric $`l^1`$-copy with a normalized canonical basis $`(x_n)`$. Let $`(x_n^{})`$ be the biorthogonal functionals of $`(x_n)`$ that is $`x_n^{}(x_m)=\delta _{n,m}`$ for $`n,m\mathrm{I}N`$.
Sublemma $`x_n^{}X^\mathrm{\#}`$ for all $`n\mathrm{I}N`$.
Proof of the Sublemma: Suppose there is $`n_0\mathrm{I}N`$ such that $`x^{}=x_{n_0}X^\mathrm{\#}`$. Then there is a sequence $`(y_n)\mathrm{B}_X`$ and there is $`\epsilon >0`$ such that
$`y_n=1,y_n\stackrel{\mathrm{asy}}{}l^1,|x^{}(y_n)|>\epsilon \text{ for all }n\mathrm{I}N.`$
By \[23, Lem. 1\] $`Y=[y_n]`$ is L-embedded and $`Y^{}=Y_1Y_s`$ with $`Y_s=Y^{}X_s`$. Let $`\delta >0`$ be arbitrary for the moment. Set $`X_0=[x_n]_{nm_0}`$ where $`m_0`$ is such that $`(x_n)_{nm_0}\stackrel{1\delta }{}l^1`$. $`X_0`$ is L-embedded by \[23, Lem. 1\] and we write $`X_0^{}=X_0_1(X_0)_s`$ with $`(X_0)_s=X_0^{}X_s`$. Let $`y_sX^{\prime \prime }`$ be a $`w^{}`$-accumulation point of $`\{y_n|n\mathrm{I}N\}`$. By Lemma 2 we have $`y_s=1`$ and $`y_sY_s`$. We have $`y_s(X_0)_s`$ because $`X_0`$ is co-finite-dimensional in $`X`$; furthermore, $`|y_s(x^{})|\epsilon `$. Let $`(z_\gamma )\mathrm{B}_{X_0}`$ be a normalized net $`w^{}`$-converging to $`y_s`$. After passing to an appropriate subnet we suppose that $`|(y_sz_\gamma )(x^{})|<\epsilon /2`$ for all $`\gamma `$. By Lemma 2 one can extract a sequence $`(z_{\gamma _n})`$ that spans $`l^1`$ asymptotically. Then $`|x^{}(z_{\gamma _n})|\epsilon /2`$ for all $`n`$. We define an isomorphism $`T:X_0[e_n]_{nm_0}`$ by $`x_ne_n`$ where $`(e_n)`$ is the standard basis of $`l^1`$. Then $`T(1\delta )^1`$ and $`T^11`$. There is $`m_1\mathrm{I}N`$ such that $`(z_{\gamma _n})_{nm_1}\stackrel{1\delta }{}l^1`$ since $`z_{\gamma _n}\stackrel{\mathrm{alm}}{}l^1`$. Hence, with the notation $`f_n=Tz_{\gamma _n}`$, we get $`(f_n)_{nm_1}\stackrel{12\delta }{}l^1`$ because
$`(12\delta ){\displaystyle \underset{n=m_1}{\overset{\mathrm{}}{}}}|\alpha _n|`$ $`<`$ $`(1\delta )^2{\displaystyle \underset{n=m_1}{\overset{\mathrm{}}{}}}|\alpha _n|(1\delta ){\displaystyle \underset{n=m_1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{|\alpha _n|}{f_n}}`$
$``$ $`{\displaystyle \underset{n=m_1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\alpha _n}{f_n}}z_{\gamma _n}{\displaystyle \underset{n=m_1}{\overset{\mathrm{}}{}}}\alpha _n{\displaystyle \frac{f_n}{f_n}}.`$
Then by \[5, Th. B\] or \[24, L. 10, 6\] there is a sequence $`(\stackrel{~}{f}_k)l^1`$ of pairwise disjointly supported elements of $`l^1`$ and a subsequence $`(f_{n_k})`$ such that
$`f_{n_k}\stackrel{~}{f}_k<\delta ^{}`$
where $`\delta ^{}0`$ as $`\delta 0`$. Set $`\stackrel{~}{z}_k=T^1\stackrel{~}{f}_k`$. Then
$`\stackrel{~}{z}_kz_{\gamma _{n_k}}=T^1\stackrel{~}{f}_kT^1f_{n_k}\stackrel{~}{f}_kf_{n_k}<\delta ^{}.`$
Now we choose $`\delta >0`$ small enough in order to have $`\delta ^{}x^{}<\epsilon /4`$. Hence
$`|x^{}(\stackrel{~}{z}_k)|{\displaystyle \frac{\epsilon }{2}}|x^{}(\stackrel{~}{z}_kz_{\gamma _{n_k})}|{\displaystyle \frac{\epsilon }{2}}\delta ^{}x^{}{\displaystyle \frac{\epsilon }{4}}`$
for all $`k\mathrm{I}N`$. On the other hand, for $`e^{}=e_{n_0}^{}=(T^1)^{}(x^{})`$ we have $`e^{}(e_n)=\delta _{n_0,n}`$ that is
$`x^{}(\stackrel{~}{z}_k)=e^{}(\stackrel{~}{f}_k)=0`$
for all but possibly one $`k\mathrm{I}N`$ because the $`\stackrel{~}{f}_k`$ are pairwise disjoint. This contradiction proves the Sublemma.
Since the biorthogonal functionals separate $`X`$ the Sublemma and Theorem 13 imply that $`X^\mathrm{\#}`$ is M-embedded and $`X=(X^\mathrm{\#})^{}`$. The Sublemma states that $`[x_n^{}]X^\mathrm{\#}`$. In fact one has $`[x_n^{}]=X^\mathrm{\#}`$. To see this let $`z^{}X^\mathrm{\#}\backslash [x_n^{}]`$. Since $`X^{}`$ is isomorphic to $`l^{\mathrm{}}`$ there is an infinite set $`N^{}\mathrm{I}N`$ and there is $`\epsilon >0`$ such that $`|z^{}(x_n)|>\epsilon `$ for all $`nN^{}`$. But $`(x_n)_{nN^{}}\stackrel{\mathrm{alm}}{}l^1`$ whence $`(x_n)_{nN}\stackrel{\mathrm{asy}}{}l^1`$ for an appropriate infinite set $`NN^{}`$ by \[23, Cor. 3\]. This means that $`(x_n)_{nN}`$ $`\tau _\mu `$-converges to $`0`$ hence $`z^{}X^\mathrm{\#}`$.
We have proved that if $`X`$ is L-embedded and almost isometric to $`l^1`$ then $`X`$ is the dual of the M-embedded space $`X^\mathrm{\#}`$ and $`X^\mathrm{\#}=[x_n^{}]`$. To end the proof it just remains to apply Proposition 14.
The arguments are similar for the case in which $`x_n\stackrel{\mathrm{asy}}{}l^1`$.
§5 On some results of Godefroy, Kalton, Li
First we deal with \[10, Prop. 2.1\]. The first part of that proposition says that $`X^\mathrm{\#}`$ where $`X`$ is a nicely placed (=L-embedded) subspace of $`\mathrm{L}^1(\mu )`$ is always M-embedded, not only in the situation of Theorem 13.
###### Proposition 16
Let $`X`$ be an L-embedded Banach space. Then $`X^\mathrm{\#}`$ is M-embedded.
Proof: With the usual notation $`X^{\prime \prime }=X_1X_s`$ and with Proposition 12 we have
$$X^\mathrm{\#}=(X_s)_{}X^{}.$$
We set
$$Z=\overline{X_s}^w^{}=(X^\mathrm{\#})^{}X^{\prime \prime }$$
and
$$Y=X\overline{X_s}^w^{}=(X^\mathrm{\#})_{}X.$$
By Corollary 11, $`Y`$ is L-embedded because its unit ball
$`\mathrm{B}_Y=\mathrm{B}_X{\displaystyle \underset{x^{}X^\mathrm{\#}}{}}\text{ker}x^{}`$
is $`\tau _\mu `$-closed. Hence $`Y^{}=Y_1(Y^{}X_s)`$ by Lemma 1. Now the fact $`Z=Y_1X_s`$ and the fact that $`X^\mathrm{\#}`$ is an M-ideal in its bidual $`Y^{}`$ can be deduced exactly as in the proof of \[10, Prop. 2.1\].
For the proof of Proposition 17 we recall property $`(m_1^{})`$. In a separable Banach space $`Z`$ is defined to have property $`(m_1^{})`$ if for all $`z^{},z_n^{}Z^{}`$
$`lim\; supz^{}+z_n^{}=z^{}+lim\; supz_n^{}`$ (12)
whenever $`z_n^{}\stackrel{w^{}}{}0`$. Analogously a separable Banach space $`X`$ is defined to have property $`(m_1)`$ if for all $`x,x_nX`$
$`lim\; supx+x_n=x+lim\; supx_n`$ (13)
whenever $`x_n\stackrel{w}{}0`$.
The second part of \[10, Prop. 2.1\] reads as follows in our context.
###### Proposition 17
Let $`X`$ be a separable L-embedded Banach space. If $`\mathrm{B}_X`$ is $`\tau _\mu `$-sequentially compact then for any $`\epsilon >0`$ there is a subspace $`X_\epsilon `$ of $`c_0`$ such that $`\text{dist}(X^\mathrm{\#},X_\epsilon )<1+\epsilon `$.
Proof: Exactly as in one distinguishes three steps: Firstly one proves that $`X^\mathrm{\#}`$ has property $`(m_1^{})`$, secondly one deduces from this property $`(m_{\mathrm{}}^{})`$ and thirdly it remains to apply \[16, Th. 3.5\]. Only the first step must be modified a bit.
From the proof of Proposition 16 we know that $`(X^\mathrm{\#})^{}=X/Y`$ where $`Y=X\overline{X_s}^w^{}=(X^\mathrm{\#})_{}X`$. Let $`(u_n)X/Y=(X^\mathrm{\#})^{}`$ be a $`w^{}`$-null sequence. We denote by $`Q:XX/Y`$ the quotient map. Let $`(x_n)X`$ be a bounded sequence such that $`Qx_n=u_n`$. By hypothesis there is a $`\tau _\mu `$-convergent subsequence - still denoted by $`(x_n)`$ \- such that $`x_nx_0\stackrel{\mathrm{asy}}{}l^1`$ where $`x_0=\tau _\mu limx_n`$ and such that $`limx_0x_n`$ exists. We have $`x_0Y`$ because for any $`x^{}X^\mathrm{\#}`$ one has
$`x^{}(x_0)=limx^{}(x_m)=limx^{}(u_m)=0.`$
We have furthermore that
$`lim\; supy+x+(x_nx_0)y+x+limx_nx_0`$ (14)
for all $`xX`$, $`yY`$. To see this, recall that $`X`$ is L-embedded, and that by Lemma 2 each universal net $`(x_{n_\gamma }x_0)`$ $`w^{}`$-converges to a limit $`x_sX_s`$ such that $`x_s=lim_\gamma x_{n_\gamma }x_0`$ and
$`\underset{\gamma }{lim}y+x+(x_{n_\gamma }x_0)y+x+x_s=y+x+x_s`$
by $`w^{}`$-continuity of the norm whence (14).
Since $`Qx=inf_{yY}y+x`$ we deduce from (14) that
$`lim\; supy+x+(x_nx_0)Qx+limu_n`$
and
$`lim\; supQx+u_nQx+limu_n`$
which proves that $`X^\mathrm{\#}`$ has property $`(m_1^{})`$. The deduction of property $`(m_{\mathrm{}}^{})`$ and the conclusion via \[16, Th. 3.5\] do not depend on the measure topology and coincide therefore with the arguments in .
The following remark gives a characterization of property $`(m_1^{})`$ in L-embedded Banach spaces. We will not need it in the sequel and state it only because the way we prove it by constructing asymptotic $`l^1`$-sequences fits naturally in the main theme of this paper. Note that the implication (i)$``$(ii) holds for arbitrary Banach spaces $`Z`$, that the implications (i)$``$(ii)$``$(iii) hold whenever $`Z`$ is such that its dual admits an abstract measure topology, and that the implication (iii)$``$(iv) does not need the M-embeddedness of $`Z`$. Note furthermore that in Remark 18 the separation assumption on $`Z`$ could be omitted because the definition of properties $`(m_1^{})`$ and $`(m_1)`$ makes sense also for non-separable spaces.
###### Remark 18
(a) Let $`Z`$ be a Banach space such that its dual is L-embedded. Then the following assertions are equivalent:
(i) $`Z`$ has property $`(m_1^{})`$.
(ii) Each $`w^{}`$-null sequence in $`Z^{}`$ admits of a subsequence that converges to $`0`$ in norm or spans $`l^1`$ asymptotically.
(iii) Each $`w^{}`$-null sequence in $`Z^{}`$ is $`\tau _\mu `$-null.
If $`Z`$ is even separable and M-embedded then the assertions above are equivalent to
(iv) $`\mathrm{B}_Z^{}`$ is $`\tau _\mu `$-sequentially compact.
(b) An arbitrary Banach space $`X`$ has property $`(m_1)`$ if and only if it has the Schur property.
Proof: (a) We set $`X=Z^{}`$. Sketch of
(i)$``$(ii): The proof ressembles the one of \[23, Th. 2\] the only difference being that (12) replaces the $`w^{}`$-lower semicontinuity of the norm and the L-embeddedness of $`Z^{}`$. Let $`(x_m)X`$ be a $`w^{}`$-null sequence, suppose without loss of generality that $`x_m=1`$. Let $`(\delta _n)`$ be a sequence of strictly positive numbers converging to $`0`$. Set $`\eta _1=\frac{1}{6}\delta _1`$ and $`\eta _{n+1}=\frac{1}{6}\mathrm{min}(\eta _n,\delta _{n+1})`$ for $`n\mathrm{I}N`$. By induction over $`n\mathrm{I}N`$ one constructs $`m_n\mathrm{\Gamma }`$ such that
$`\left({\displaystyle \underset{i=1}{\overset{n}{}}}(1\delta _i)|\alpha _i|\right)+\eta _n{\displaystyle \underset{i=1}{\overset{n}{}}}|\alpha _i|{\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _ix_{m_i}{\displaystyle \underset{i=1}{\overset{n}{}}}|\alpha _i|n\mathrm{I}N,\alpha _i\mathrm{C}.`$ (15)
For the induction step $`nn+1`$ fix an element $`\alpha =(\alpha _i)_{i=1}^{n+1}`$ in the unit sphere of $`l_{n+1}^1`$ such that $`\alpha _{n+1}0`$. Then (12) yields
$`\left({\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _ix_{m_i}\right)+\alpha _{n+1}x_m+{\displaystyle \frac{1}{2}}\mathrm{min}(\eta _n,\delta _{n+1})`$
$``$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _ix_{m_i}+|\alpha _{n+1}|`$
$`\stackrel{(\text{15})}{}`$ $`\left({\displaystyle \underset{i=1}{\overset{n}{}}}(1\delta _i)|\alpha _i|\right)+\eta _n\left({\displaystyle \underset{i=1}{\overset{n}{}}}|\alpha _i|\right)+|\alpha _{n+1}|`$
$``$ $`\left({\displaystyle \underset{i=1}{\overset{n+1}{}}}(1\delta _i)|\alpha _i|\right)+\mathrm{min}(\eta _n,\delta _{n+1})`$
whence
$`\left({\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _iy_{m_i}^{(i)}\right)+\alpha _{n+1}y_m^{(n+1)}\left({\displaystyle \underset{i=1}{\overset{n+1}{}}}(1\delta _i)|\alpha _i|\right)+3\eta _{n+1}`$
for infinitely many $`m`$. This gives (15); the details are the same as in the proof of \[23, Th. 2\] or of Lemma 6.
Note for the proof of part (b) below that we used the $`w^{}`$-convergence of $`(x_n)`$ only in order to apply property $`(m_1^{})`$ but not for the construction of the $`l^1`$-basis itself.
(ii)$``$(i): Suppose that $`Z`$ satisfies (ii) without having $`(m_1^{})`$ and suppose that $`Z^{}`$ is L-embedded. Then there are $`xX`$, $`\epsilon >0`$ and a $`w^{}`$-null sequence $`(x_n)X`$ such that $`limx+x_n`$ and $`limx_n`$ exist and
$`\epsilon +limx+x_n<x+limx_n.`$ (16)
Since (16) excludes the case $`x_n0`$ we have - after passing to an appropriate subsequence - that $`x_n\stackrel{\mathrm{asy}}{}l^1`$. Let $`(x_{n_\gamma })`$ be a universal net. By Lemma 2 it $`w^{}`$-converges to a point $`x_sX_s`$ and
$`\underset{\gamma }{lim}x+x_{n_\gamma }x+x_s=x+\underset{\gamma }{lim}x_{n_\gamma }`$
by $`w^{}`$-continuity of the norm which contradicts (16) and proves property $`(m_1^{})`$.
(ii)$``$(iii) is immediate from the definition of $`\tau _\mu `$.
(iii)$``$(iv): If $`Z`$ is separable then $`\mathrm{B}_X`$ is $`w^{}`$-sequentially compact. Hence $`\mathrm{B}_X`$ is $`\tau _\mu `$-sequentially compact if (iii) holds.
(iv)$``$(i): If $`Z`$ is M-embedded then $`Z=X^\mathrm{\#}`$. In this case, if $`\mathrm{B}_X`$ is $`\tau _\mu `$-sequentially compact then the proof of Proposition 17 shows that $`Z`$ has $`(m_1^{})`$.
(b) The Schur property clearly implies $`(m_1)`$. Conversely, suppose a Banach space $`X`$ has $`(m_1)`$ but fails to have the Schur property. Then there exists a normalized $`w`$-null sequence $`(x_n)`$ in $`X`$. Exactly as in (i)$``$(ii) of part (a), using (13) instead of (12), we can extract a subsequence $`(x_{n_k})`$ which spans $`l^1`$ asymptotically. But a normalized sequence spanning $`l^1`$ cannot converge weakly. This contradiction proves that $`(m_1)`$ implies the Schur property.
###### Sublemma 19
For any L-embedded space $`X`$ one has the following inclusions
$`X\overline{\mathrm{B}_{X_s}}^w^{}`$ $``$ $`{\displaystyle \{\overline{V}^w|V\tau _\mu \text{open in }\mathrm{B}_X,\mathrm{\hspace{0.17em}0}V\}}`$
$``$ $`{\displaystyle \{\overline{\text{co}}^{}(V)|V\tau _\mu \text{open in }\mathrm{B}_X,\mathrm{\hspace{0.17em}0}V\}}.`$
Proof: The first inclusion is immediate from (i) of Lemma 5, the second inclusion follows from $`\overline{\text{co}}^w(V)=\overline{\text{co}}^{}(V)`$.
That at the time of this writing we are not able to give the whole analogue of \[10, L. 2.6\] is essentially due to the fact that we know only the behaviour of sequences but not of nets in $`\tau _\mu `$, more specifically if we knew that one can extract a $`\tau _\mu `$-convergent sequence from a $`\tau _\mu `$-convergent net then we could also prove that $`\overline{\text{co}}^{}(V)X\overline{\mathrm{B}_{X_s}}^w^{}`$. In any case, our reduced version suffices for Lemma 20 which corresponds to \[10, L. 2.7\]. Note that (ii) of Lemma 20 is equivalent to the two conditions of Theorem 13.
###### Lemma 20
Let $`X`$ be an L-embedded Banach space. The following assertions are equivalent.
(i) There exists a locally convex Hausdorff topology on $`X`$ which is coarser than $`\tau _\mu `$ on $`\mathrm{B}_X`$.
(ii) $`\mathrm{B}_X`$ is compact Hausdorff with respect to $`\sigma (X,X^\mathrm{\#})`$.
(iii) $`\{0\}`$ is the intersection of the convex $`\tau _\mu `$-neighbourhoods of $`0`$ in $`\mathrm{B}_X`$.
If moreover $`\mathrm{B}_X`$ is $`\tau _\mu `$-compact Hausdorff, then the above conditions are also equivalent to
(iv) The weak topology of $`X`$ is finer than $`\tau _\mu `$ on $`\mathrm{B}_X`$.
(v) There exists a locally convex Hausdorff topology - namely $`\sigma (X,X^\mathrm{\#})`$ \- which coincides with $`\tau _\mu `$ on $`\mathrm{B}_X`$.
(vi) $`0`$ has a basis of $`\tau _\mu `$-neighbourhoods in $`\mathrm{B}_X`$ consisting of convex sets.
Proof: The proof is word-for-word the same as the one of \[10, L. 2.7\] (with Sublemma 19 replacing \[10, L. 2.6\]).
Suppose that with the notation of (i) of Lemma 5 one has $`x=x^{\prime \prime }X`$; then $`x_\gamma \stackrel{\tau _\mu }{}0`$ and $`x_\gamma \stackrel{w}{}x`$. If (iv) of Lemma 20 holds then this means that both $`x_\gamma \stackrel{\tau _\mu }{}0`$ and $`x_\gamma \stackrel{\tau _\mu }{}x`$. Since it is not clear whether $`\tau _\mu `$ is Hausdorff we cannot deduce that $`x=0`$. Therefore we claimed $`\mathrm{B}_X`$ to be $`\tau _\mu `$-Hausdorff for (iv) - (vi) in Lemma 20 (in particular for the proof of (iv)$``$(ii)).
We end this section with some remarks on the continuity of the canonical L-projection $`P`$ on the bidual of an L-embedded space $`X`$.
In \[10, Prop. 3.8\] it was shown that $`P`$ is ($`w^{}`$-$`\tau _\mu `$)-sequentially continuous if and only if $`X`$ has the Schur property. In the proof the authors of use the Grothendieck property of $`\mathrm{L}^{\mathrm{}}`$ in order to obtain the ($`w^{}`$-$`w`$)-sequential continuity of the L-projection on $`(\mathrm{L}^{\mathrm{}})^{}`$. But in the general case it is not known whether $`P`$ is always ($`w^{}`$-$`w`$)-sequentially continuous although by an example of Johnson it is known that in general the dual of an L-embedded Banach space does not have the Grothendieck property. Therefore it is not clear whether the first part of \[10, Prop. 3.8\] can be generalized in our context; its second part can:
###### Proposition 21
Let $`X`$ be L-embedded and $`P`$ be the canonical L-projection on $`X^{\prime \prime }`$. Suppose that the abstract measure topology $`\tau _\mu `$ on $`X`$ is Hausdorff. Then $`P`$ is ($`w^{}`$-$`\tau _\mu `$)-continuous if and only if the restriction of $`\tau _\mu `$ to $`\mathrm{B}_X`$ is compact Hausdorff locally convex.
Proof: See \[10, Prop. 3.8(b)\].
§6 The special case of the predual of a von Neumann algebra
Recall that two elements $`\varphi ,\psi 𝒩_{}`$ in the predual of a von Neumann algebra $`𝒩`$ are called orthogonal - $`\varphi \psi `$ in symbols - if they have orthogonal right and orthogonal left support projections. It is well know that $`\varphi \psi `$ if and only if the linear span of $`\varphi `$ and $`\psi `$ is isometrically isomorphic to the two-dimensional $`l^1(2)`$.
###### Sublemma 22
Let $`𝒩`$ be a von Neumann algebra with predual $`X=𝒩_{}`$. For every $`\epsilon >0`$ there is $`\delta >0`$ with the following property.
If $`x,y,z\mathrm{B}_X`$ are such that
$`\alpha {\displaystyle \frac{z}{z}}+\beta {\displaystyle \frac{x}{x}}`$ $``$ $`(1\delta )\left(|\alpha |+|\beta |\right)`$
$`\alpha {\displaystyle \frac{z}{z}}+\beta {\displaystyle \frac{y}{y}}`$ $``$ $`(1\delta )\left(|\alpha |+|\beta |\right)`$
then there are $`\stackrel{~}{x},\stackrel{~}{y},\stackrel{~}{z}X`$ of norm one such that
$`\begin{array}{c}\hfill \stackrel{~}{x}\stackrel{~}{z}\\ \hfill \stackrel{~}{y}\stackrel{~}{z}\end{array}\text{ and }\begin{array}{c}\hfill x\stackrel{~}{x}\epsilon \\ \hfill y\stackrel{~}{y}\epsilon \\ \hfill z\stackrel{~}{z}\epsilon \end{array}`$
Proof: We have already recalled the fact that two normal functionals on a von Neumann algebra are orthogonal if and only if they span $`l^1(2)`$ isometrically. A second ingredient of this proof is the fact that the ultraproduct of a family of preduals of von Neumann algebras is again a predual of a von Neumann algebra , . It is now enough to combine these two facts with a standard ultraproduct argument .
It is well known that if $`𝒩`$ is a von Neumann algebra with a finite faithful normal trace $`\tau `$ then the measure topology on the predual $`𝒩_{}`$ defined via $`\tau `$ is metrizable and makes $`𝒩_{}`$ a Hausdorff topological vector space. For the abstract measure topology $`\tau _\mu `$ on the predual of an arbitrary von Neumann algebra we only know by the following lemma (and by Lemmas 5, 6) that multiplication by scalars and addition are continuous and that on bounded sets $`\tau _\mu `$ makes the von Neumann predual a Fréchet space.
###### Proposition 23
Let $`X=𝒩_{}`$ be the predual of a von Neumann algebra $`𝒩`$. If $`X`$ is endowed with the abstract measure toplogy $`\tau _\mu `$ which it has by Theorem 4 then addition is $`\tau _\mu `$-continuous.
Proof: Let $`(x_n)`$, $`(y_n)`$ be two sequences in $`\mathrm{B}_X`$ each of which spans $`l^1`$ asymptotically. Suppose there is $`\epsilon >0`$ such that $`z_n>\epsilon `$ where $`z_n=x_n+y_n`$, suppose further that $`limx_n`$, $`limy_n`$, $`limz_n`$ exist. To show the Proposition it is enough to show that there is a subsequence $`(z_{n_k})`$ spanning $`l^1`$ asymptotically.
Let $`(\delta _n)`$ be a sequence of strictly positive numbers in $`]0,1]`$ converging to $`0`$. Set $`\eta _1=\frac{1}{3}\delta _1`$ and $`\eta _{n+1}=\frac{1}{3}\mathrm{min}(\eta _n,\delta _{n+1})`$ for $`n\mathrm{I}N`$. By induction over $`m\mathrm{I}N`$ we construct a sequence $`(n_m)`$ such that
$`{\displaystyle \underset{k=1}{\overset{m}{}}}\alpha _k{\displaystyle \frac{z_{n_k}}{z_{n_k}}}\left({\displaystyle \underset{k=1}{\overset{m}{}}}(1\delta _k)|\alpha _k|\right)+\eta _m{\displaystyle \underset{k=1}{\overset{m}{}}}|\alpha _k|`$ (18)
for all integers $`m`$ and all scalars $`\alpha _k`$. Since the construction ressembles the proofs of \[23, Th. 2\], of Lemma 6, and of Remark 18 we detail it only for $`m=2`$.
Let $`\theta _2>0`$ be arbitrary for the moment and to be determined later. Set $`n_1=1`$.
Claim: There is $`n_2\mathrm{I}N`$ such that
$`\alpha {\displaystyle \frac{z_1}{z_1}}+\beta {\displaystyle \frac{x_{n_2}}{x_{n_2}}}`$ $``$ $`(1\theta _2)\left(|\alpha |+|\beta |\right)`$ (19)
$`\alpha {\displaystyle \frac{z_1}{z_1}}+\beta {\displaystyle \frac{y_{n_2}}{y_{n_2}}}`$ $``$ $`(1\theta _2)\left(|\alpha |+|\beta |\right)`$ (20)
for all scalars $`\alpha ,\beta `$. (As usual in this paper) this is essentially due to the $`w^{}`$-lower semicontinuity of the norm by which one has $`lim\; inf_\gamma \alpha \frac{z_1}{z_1}+\beta \frac{x_{n_\gamma }}{x_{n_\gamma }}\alpha \frac{z_1}{z_1}+\beta \frac{x_s}{x_s}=|\alpha |+|\beta |`$ where $`x_sX_s`$ is the $`w^{}`$-limit of a universal net $`(x_{n_\gamma })`$ and $`x_s=lim_\gamma x_{n_\gamma }`$ (see Lemma 2); hence there are infinitely many $`n_2`$ satisfying (19). Applying the same reasoning to the corresponding subsequence of $`(y_n)`$ gives the Claim.
By Sublemma 22 there are $`\stackrel{~}{x}_2`$, $`\stackrel{~}{y}_2`$, $`\stackrel{~}{z}_1`$ in $`X`$ such that
$`x_{n_2}\stackrel{~}{x}_2\theta _2^{},\text{ }y_{n_2}\stackrel{~}{y}_2\theta _2^{},\text{ }z_1\stackrel{~}{z}_1\theta _2^{},`$
and
$`\stackrel{~}{z}_1\stackrel{~}{x}_2,\text{ }\text{ }\stackrel{~}{z}_1\stackrel{~}{y}_2`$ (21)
where $`\theta _2^{}0`$ as $`\theta _20`$. Now (21) implies that
$`\stackrel{~}{z}_1\stackrel{~}{z}_{n_2}`$
where $`\stackrel{~}{z}_{n_2}=\stackrel{~}{x}_{n_2}+\stackrel{~}{y}_{n_2}`$ which means that $`\alpha \frac{\stackrel{~}{z}_1}{\stackrel{~}{z}_1}+\beta \frac{\stackrel{~}{z}_2}{\stackrel{~}{z}_2}=|\alpha |+|\beta |`$. Hence, if $`\theta _2`$ is choosen small enough, we get (18) for $`m=2`$. It is now left to the reader to iterate the construction in order to end the induction and thus the proof.
As already noticed in the remarks concerning the questions after Corollary 9, Randrianantoanina has proved that the Kadec-Pełczyński splitting lemma holds for preduals of von Neumann algebras. With this result Komlos’ theorem follows almost immediately for von Neumann preduals.
###### Proposition 24 (Komlos)
Each bounded sequence in the predual of a von Neumann algebra admits a subsequence whose Cesaro means converge with respect to the abstract measure topology.
Proof: Let $`X=𝒩_{}`$ be the predual of a von Neumann algebra $`𝒩`$. Endow $`X`$ with its abstract measure topology $`\tau _\mu `$. Let $`(x_n)X`$ be bounded. Then by there is a subsequence $`(x_{n_k})`$ and there is a decomposition $`x_{n_k}=y_k+z_k`$ where $`(z_k)`$ $`w`$-converges and such that there is a sequence $`(\stackrel{~}{y}_k)`$ of pairwise orthogonal elements in $`X`$ with
$`\stackrel{~}{y}_ky_k0.`$ (22)
If one applies the classical theorem of Komlos to the isometric $`l^1`$-copy $`[\stackrel{~}{y}_k]`$ then one obtains, after passing to an appropriate subsequence, that the Cesaro means of any subsequence of $`(\stackrel{~}{y}_k)`$ converge with respect to the measure (=pointwise) topology of $`[\stackrel{~}{y}_k]`$ whence with respect to $`\tau _\mu `$. So do the Cesaro means $`\widehat{y}_k`$ of the $`y_k`$ by (22). It is known that von Neumann preduals have the weak Banach-Saks property . Hence, again after passing to the appropriate subsequences of $`(y_k)`$ and $`(z_k)`$, the Cesaro means $`\widehat{z}_k`$ of the $`z_k`$ converge in norm whence with respect to $`\tau _\mu `$. Since $`\widehat{x}_k=\widehat{y}_k+\widehat{z}_k`$ where $`\widehat{x}_k`$ denote the Cesaro means of the $`x_{n_k}`$, the assertion now follows from Proposition 23.
Acknowledgement I thank Dirk Werner for several helpful discussions.
Hermann Pfitzner
Université d’Orléans
BP 6759
F-45067 Orléans Cedex 2
France
e-mail: pfitzner@labomath.univ-orleans.fr |
warning/0003/hep-ph0003208.html | ar5iv | text | # Phenomenological Survey of a Minimal Superstring Standard Model aafootnote aTalk presented at PASCOS 99, Lake Tahoe, CA December 10–16 1999.
## 1 Flat Directions in Three–Generation Heterotic–String Models
Over the past decade the free fermionic formulation$`^\mathrm{?}`$ of the heterotic string has been utilized to derive the most realistic$`^{\mathrm{?},\mathrm{?}}`$ string models to date. In some of these models the observable sector gauge group directly below the string scale is a Grand Unified Theory while in others it is the (MS)SM gauge group, $`SU(3)_C\times SU(2)_L\times U(1)_Y`$, joined by a few extra $`U(1)`$ symmetries. In chiral three generation models of the latter class, one of the additional Abelian gauge groups is inevitably anomalous. That is, the trace of the $`U(1)_A`$ charge, $`\mathrm{Tr}Q^{(A)}`$, generates a Fayet–Iliopoulos (FI) term,
$`ϵ{\displaystyle \frac{g_s^2M_P^2}{192\pi ^2}}\mathrm{Tr}Q^{(A)}.`$ (1)
The FI term breaks supersymmetry near the Planck scale, and destabilizes the string vacuum. Supersymmetry is restored and the vacuum is stabilized if there exists a direction, $`\varphi =_i\alpha _i\varphi _i`$, in the scalar potential which is $`F`$–flat and also $`D`$–flat with respect to the non–anomalous gauge symmetries and in which $`_iQ_i^A|\alpha _i|^2`$ and $`ϵ`$ are of opposite sign. If such a direction exists it will acquire a vacuum expectation value (VEV) cancelling the anomalous $`D`$–term, restoring supersymmetry and stabilizing the string vacuum. Since the fields corresponding to such a flat direction typically also carry charges for the non–anomalous $`D`$–terms, a non–trivial set of constraints on the possible choices of VEVs is imposed:
$`D_A`$ $`=`$ $`{\displaystyle \underset{m}{}}Q_m^{(A)}|\phi _m|^2+ϵ=0,`$ (2)
$`D_i`$ $`=`$ $`{\displaystyle \underset{m}{}}Q_m^{(i)}|\phi _m|^2=0,`$ (3)
$`D_a^\alpha `$ $`=`$ $`{\displaystyle \underset{m}{}}\phi _m^{}T_a^\alpha \phi _m=0,`$ (4)
with $`T_a^\alpha `$ a matrix generator of the non–Abelian gauge group $`g_\alpha `$ for the representation $`\phi _m`$. These scalar VEVs will in general break some, or all, of the additional symmetries spontaneously.
Additionally one must insure that the supersymmetric vacuum is also $`F`$–flat. Each superfield $`\mathrm{\Phi }_m`$ (containing a scalar field $`\phi _m`$ and chiral spin–$`\frac{1}{2}`$ superpartner $`\psi _m`$) that appears in the superpotential imposes further constraints on the scalar VEVs. $`F`$–flatness will be broken (thereby destroying spacetime supersymmetry) at the scale of the VEVs unless,
$$F_m\frac{W}{\mathrm{\Phi }_m}=0;W=0,$$
(5)
where $`W`$ is the superpotential which contains cubic level and higher order non–renormalizable terms. The higher order terms have the generic form
$$<\mathrm{\Phi }_1^f\mathrm{\Phi }_2^f\mathrm{\Phi }_3^b\mathrm{}\mathrm{\Phi }_N^b>.$$
Some of the fields appearing in the non–renormalizable terms will in general acquire a non–vanishing VEV by the anomalous $`U(1)`$ cancellation mechanism. Thus, in this process some of the non–renormalizable terms induce effective renormalizable operators in the effective low energy field theory wherein either all fields or all fields but one are replaced with VEVs. One must insure that such terms do not violate supersymmetry at an unacceptable level.
## 2 A Minimal Superstring Standard Model
In addition to possessing an anomalous $`U(1)_A`$ symmetry, chiral three generation $`SU(3)_C\times SU(2)_L\times U(1)_Y\times _iU(1)_i`$ models generically contain numerous non–MSSM $`SU(3)_C\times SU(2)_L\times U(1)_Y`$–charged states, some of which only carry MSSM$`\times _iU(1)_i`$ charges and others of which also carry hidden sector (non)–Abelian charges. Recently, we showed that in some of these models it is actually possible to decouple all such non–MSSM states from the low energy effective field theory. For example, in the “FNY” model first presented in Ref. 3, we discovered there are several flat directions$`^{\mathrm{?},\mathrm{?}}`$ for which almost all MSSM–charged exotics receive FI masses (typically of the scale $`5\times 10^{16}`$ GeV to $`1\times 10^{17}`$ GeV) while the remaining MSSM–charged exotics (usually composed of simply a $`SU(3)_C`$ triplet/anti–triplet pair) acquire masses slightly suppressed below the FI scale by a factor of $`𝒪(\frac{1}{10}`$ to $`\frac{1}{100})`$. Some of our flat directions accomplishing this feat contain only VEVs of non–Abelian singlet fields$`^{\mathrm{?},\mathrm{?}}`$ while others of ours also contain VEVs of non–Abelian charged fields$`^{\mathrm{?},\mathrm{?}}`$. Along these directions, exactly three generations of ($`Q_i`$, $`u_i^c`$, $`d_i^c`$, $`L_i`$, $`e_i^c`$, $`N_i^c`$) fields and a pair of electroweak Higgs, $`h_u`$ and $`h_d`$, are the only MSSM–charged fields that remain massless significantly below the FI scale. The non–MSSM–charged singlet fields and hidden sector non–Abelian fields that also remain massless below the FI scale vary with the flat direction chosen.
The complete massless spectrums produced, respectively, by four singlet flat directions were presented in Ref. 6. Detailed analysis of the associated three generation mass textures was also performed therein. The leading mass terms were found to be $`Q_1u_1^c\overline{h}_1`$ and $`Q_3d_3^ch_3`$ for all four directions. The non–Abelian directions considered in Ref. 7 produced similar results. This presented a phenomenological difficulty for our first few MSSM–producing flat directions by implying that the left–handed component of the top and bottom quarks live in different multiplets. A more detailed study of the mass textures possible from non–Abelian flat directions is underway and will be presented in Ref. 9.
## 3 Phenomenogolical Implications
A string–derived three generation MSSM low energy effective field theory resulting from the decoupling of all MSSM–charged exotics via an anomaly–cancelling flat direction possesses some unique phenomenological characteristics (independent of the viability of the associated MSSM mass matrices). For instance, we found that at least one $`U_i`$ combination usually remains unbroken by a given flat direction. Generally the surviving extra Abelian symmetries could not have been embedded in $`SO(10)`$ or $`E_6`$ GUTS$`^\mathrm{?}`$. Relatedly, a common feature in the surviving $`U_i`$ combinations is a flavor non–universality. Thus, the distinctive collider signatures of a $`Z^{}`$ arising from one such symmetry will be a non–universality in the production of the different generations. An additional $`Z^{}`$ of this type has also been suggested as playing a role in suppressing proton decay in supersymmetric extensions of the Standard Model$`^\mathrm{?}`$.
## Acknowledgments
G.C. thanks the organizers of PASCOS ’99 for an excellent conference. The work discussed herein was done in collaboration with A.E. Faraggi, D.V. Nanopoulos, & J.W. Walker.
## References |
warning/0003/cond-mat0003096.html | ar5iv | text | # Why does a protein fold?
## Abstract
With the help of lattice Monte Carlo modelling of heteropolymers, we show that the necessary condition for a protein to fold on short call is to proceed through partially folded intermediates. These elementary structures are formed at an early stage in the folding process and contain, at the local level, essentially all of the amino acids found in the folding core (transition state) of the protein, providing the local guidance for its formation. The sufficient condition for the protein to fold is that the designed sequence has an energy, in the native conformation, below $`E_c`$ (the lowest energy of the structurally dissimilar compact conformations) where it has not to compete with the bulk of misfolded conformations. Sequences with energy close to $`E_c`$ can display prion–like behaviour, folding to two structurally dissimilar conformations, one of them being the native.
We wish to suggest a novel model for protein folding, where the building blocks which control the dynamics of the designed sequences are partially folded intermediates. Starting form a random coil (Fig. 1(a) ), they are formed only after $`10^2`$ steps of Monte Carlo (MC) simulations (Fig. 1(b) ), when some of the most strongly interacting amino acids establish their local contacts. They achieve $`9095\%`$ stability after $`10^5`$ MC steps, and when they assemble together after $`10^6`$ MC steps (Fig. 1(c) ) they form the (post–critical) folding core of the notional protein, from which it reaches the native conformation (Fig. 1(d) ) in a short time ($`10^3`$ MC steps), provided the energy of the system is lower than $`E_c`$. Partially folded intermediates and not the individual monomers thus take care, through local guidance and non–local long range correlations (bonding between partially folded intermediates), of the process of protein folding, as testified by the disruptive effect mutations which affect the stability of these structures have on the folding ability of the designed sequence (cf. also ).
The fast formation of few partially folded intermediates, and of their bonding, reduces, in a conspicuous way, the number of conformations that need to be searched (in case of the chain considered in Fig. 1 to $`10^{12}`$ as compared to $`10^{24}`$ for the random–coil), leading to the resolution of the Levinthal paradox. It is also a very efficient way to squeeze entropy from the system ($`50\%`$) at the very early stages of the folding process, and to repeat this feat when the partially folded intermediates come together to form the folding nucleus (cf also ref. ), at which stage the integrated decrease of entropy amounts to a large (in the case of Fig. 1, of the order of $`80\%`$) fraction of the original random–coil value.
The numbers quoted in the first paragraph were obtained using a lattice model of proteins studied earlier by us and others . The model sequences are composed of amino acids of 20 types and containing 36 monomers, which interact through contact energies obtained from a statistical analysis of real proteins , the associated standard deviation of the interaction energies between different amino acid types being $`\sigma =0.3`$. From very long Monte Carlo runs ($`10^9`$ MC steps) a sequence has been found with a sufficiently low energy in the native conformation , which in the units we use ($`RT_{room}=0.6kcal/mol`$) is equal to $`E_n=17.13`$, to be compared to the lowest energy of the structurally dissimilar part of the spectrum $`E_c=14`$ (with standard deviation $`\sigma _c=0.2`$), obtained through a set of low temperature MC samplings in conformational space. While all sequences lying below $`E_c`$ (of the order of $`10^{10}`$ ) eventually fold, in keeping with the fact that they share a (small) number of conserved contacts (folding nucleus), the folding time is correlated with the corresponding energy gap (Table 1). To state that the ability a notional protein has to fold, is connected with the presence of a small number of conserved contacts or, equivalently, of conserved amino acids , is tantamount to saying that foldability is connected with the presence of a small number of partially folded intermediates. In fact, although most of the conserved contacts found in the folding nucleus of ref. are non–local, few of them are local. These few contacts stabilize the partially folded intermediates already at the initial stage of the folding process. It is then natural that the non–local contacts of the folding nucleus arize from the assembling together of the partially folded intermediates. Because these local structures are both few and strongly interacting as they are mediated, by the few, strongly interacting, amino acids occupying ”hot” sites in the protein (cf. caption to Fig. 1 and ), they can come together both fast ($`10^510^6`$MC steps) and in a unique fashion, to form the (post–critical) folding core of the protein. Consequently, the findings displayed in Fig. 1 agree in detail with the result of ref. providing it a simple microscopic picture.
In spite of the difference in language, it also agrees with the findings of ref. . In fact, while the stability of the folding intermediates is not $`100\%`$, the corresponding contacts (cf. Fig. 2) are operative with an incidence which is much higher than that associated with the non–conserved contacts (cf. also Fig. 5 of ref. ). In keeping with the definition of the transition states (which, in the present case, are $`10^4`$) as those in which the protein has equal probability to proceed to the native conformation as it has to unfold, not all the conserved contacts are, in these states, operative with probability 1. In this sense, good folders can fold in different manners (different transition states) . On the other hand, any good folder passes, with probability 1, through the (post–critical) folding core conformation (Fig. 1(c) )en route to the native structure.
In order to investigate the dynamical behaviour of sequences with energy $`E_nEE_c`$ (that is sequences which can also be marginally stable), a database of (composition conserving) sequences of specified energy has been created making use of a Monte Carlo algorithm. The database is divided in 6 groups whose elements have energy $`17.00<E<16.50`$, $`16.50<E<15.00`$, … , $`14.50<E<14.00`$, each group containing 500 sequences. For each group the Monte Carlo selection has been performed at a temperature ($`T=0.28`$) such that the average energy lies in the associated energy interval.
Essentially all sequences ($`92\%`$) with $`E<E_c`$ fold in rather homogeneous times (cf. Table 1), a time which is much shorter than that associated with a random search in the space of compact conformations ($`10^{12}`$), let alone the full space of conformations ($`10^{24}`$). The sequences fold either to the native structure or to a unique structure with a value of the similarity parameter (defined as the ratio of native contacts of a given conformation to the total number of contacts) $`q>0.6`$. This process takes always place through partially folded intermediates, a result which seems to find strong experimental support (cf. e.g. and refs. therein). To be noted that for a given native structure, all designed sequences are characterized by a very limited choice of partially folded intermediates . For example, in the case of the native conformation chosen for the analysis (Fig. 1(d) ), these local substructures involve monomers 3–6, 11–14 and 27–30 (the only other choice for partially folded intermediates involves monomers 2–7, 11–14, 16–21). We find that in the folding process, all sequences with energy $`EE_c^{}`$ where $`E_c^{}=E_cn\sigma `$ ($`n=12`$), undergo a first order–like transition , the transition state being characterized by values of the conformational order parameter q which range, depending on the sequence, between $`0.45`$ and $`0.70`$. As expected from the definition of $`E_c`$, for sequences with energy close to $`E_c`$, the native conformation starts competing with other conformations. In fact, sequences which do not fold ($`8\%`$ of the total database) are concentrated in the energy range $`14.50<E<14.00`$. Of them, $`2.7\%`$ behave like random sequences, while $`5.3\%`$ display an unexpected prion–like behaviour, folding either to the native state or to another unique conformation. The fact that also these sequences display (according to simulations performed in the range of $`10^8`$ MC steps) a first order–like transition, suggests that a mechanism of kinetic partitioning is active, in the sense that, in the folding time scale, the unfolded and only one of the two possible folded conformations play a role in each simulation. In other words, prion–like sequences behave as if, at the very early stages of the folding process, one of the two possible folded conformations was selected (cf. also ref. ). In keeping with these results, and to the extent that lattice simulations do describe ”wild type” proteins, one could argue that the mere existence of prions testifies to the central role partially folded intermediates (only ”intelligent” structures operative at the very early stage of the compaction process) play in the folding of proteins.
We have also found that many sequences with $`EE_c`$ can still fold to the native conformation, while a consistent part of them again show prion–like behaviour. For example, among sequences with $`14.00<E<13.50`$, $`53\%`$ of them fold to the native state, $`8\%`$ fold to a unique conformation, similar to the native state ($`q>0.6`$), $`28\%`$ display prion–like behaviour, while $`11\%`$ do not fold. Among sequences with energy lying in the interval $`13.50<E<13.00`$, $`40\%`$ can still fold to the native state within $`210^8`$ MC steps, while sequences which reach a unique conformation, different from the native, drop to $`3\%`$ and those displaying prion–like behaviour become $`22\%`$.
Partially folded intermediates are also found to be present in the compaction of sequences displaying an energy, in the native conformation, much larger than $`E_c`$. Making again use of a Monte Carlo algorithm, we have investigated the average native energy of sequences at different selective temperatures $`T_s`$, together with the average energy of the partially folded intermediates and of the folding nucleus. The results are displayed in Fig. 3. Both the partially folded intermediates and the folding core undergo something resembling a first order phase transition (strongly blurred by fluctuations, in keeping with the fact that the system is small) at $`T_s0.3`$, while, in the same range of selective temperatures, the overall sequence undergoes a second order transition. The energy $`E_c`$ corresponds to a selective temperature $`T_s=0.09`$, which is far below the phase transition temperature. Consequently, partially folded intermediates and the folding core are, in average, present in all the sequences with energy as high as $`E=9`$. In spite of the fact that some of these sequences are able to fold, folding events are rare at these energies, in keeping with the fact that the folded state is immersed in a dense background of states associated with random sequences and thus of misfolded conformation.
The presence of a specific set of very favorable interactions, the folding core , which is built out of the partially folded intermediates and which depends on the geometry of the target structure, and consequently is missed by the mean field description, can lower the energy of the native state below $`E_c`$. In other words, the ground state energy of a sequence can be written as $`E=E_{core}+E_{oth}`$, where $`E_{core}`$ is the energy of the folding nucleus, an energy which, according to a first order transition interpretation of Fig. 3 can be either $`0`$ or $`J`$ (with $`J=7`$). The energy $`E_{oth}`$ of the non–core residues are distributed according to the Random Energy Model , the lowest of them being $`E_c^{\prime \prime }`$ (which is higher than $`E_c`$ because it contains less residues). Setting $`E_c^{\prime \prime }=E_c(nn_{core})/n`$ with $`n=40`$ the total number of contacts, and with $`n_{core}=9`$ the contacts belonging to the folding core (cf. Fig. 1(c) ), one obtains $`E_c^{\prime \prime }=10.9`$. Then, the energy of the lowest sequence should be $`E_{opt}=17.9`$ (to be compared to the value $`17.13`$ we obtained in MC simulations), so that the gap of the best sequence is $`\delta _{opt}=E_cE_{opt}=3.9`$ (to be compared to $`\delta _{opt}=3.13`$, the outcome of MC simulations). In keeping with this discussion, we find that there are two sets of sequences. One, which in the compaction process does not display partially folded intermediates and thus a folding core, spanning the energy interval $`E_c^{\prime \prime }<E<0`$. Another, spanning the energy interval $`E_c^{\prime \prime }+J<E<J`$, which in the compaction process form partially folded intermediates. Sequences of this type with energies $`EE_c`$ fold in times which are, within an order of magnitude, essentially the same (Table 1). Within the present model, this is a rather natural result due to the fact that the folding time is, to a large extent, determined by the time it takes for the partially folded intermediates to assemble together into the folding core, and to the result that the partially folded intermediates are the same for all sequences with $`E<E_c`$.
To the question: why does a protein fold?, the answer seems to be: because it proceeds through early formed local structures, partially folded intermediates (efficient way to squeeze entropy from the chain) carriers of the information concerning the folding core they form by assembling together, thus lowering the energy of the system below the threshold energy of random sequences, where the system has not to compete with the bulk of misfolded conformation. |
warning/0003/gr-qc0003035.html | ar5iv | text | # 1 Introduction
## 1 Introduction
It is generally agreed that gravitational energy exists, but because of the equivalence principle it cannot be localized. The notion of quasilocal energy is currently one of the most promising descriptions of energy in the context of general relativity, and can be characterized simply as follows. The total energy, including both matter and gravitational contributions, contained in a finite spatial volume $`\mathrm{\Sigma }`$ can be defined only as the integral of some energy surface density over its two-surface boundary, $`S=\mathrm{\Sigma }`$. This implies that, strictly speaking, there is no such thing as a local energy volume density, except that which arises from the small $`S`$ limit of quasilocal energy.<sup>1</sup><sup>1</sup>1A notable exception is the Tolman density, which integrates to the Komar mass . But it can be defined only when the spacetime possesses special properties, namely a timelike Killing vector field and an asymptotically flat spatial infinity, and so tells us little about the nature of energy in a general context. And even this local notion is not truly local because it cannot be integrated over a finite volume unless one is willing to ignore effects due to gravity. In short, energy is associated with closed spacelike two-surfaces in spacetime, not points.
There is also a growing consensus that the ADM and Bondi-Sachs masses are simply not enough. We need some definition of energy that is ‘more local’ than these, i.e., a quasilocal definition that does not rely on the existence of an asymptotically flat region . For example, recent proponents of this movement are Ashtekar et al , who have introduced the quasilocal idea of an ‘isolated horizon’ to describe a black hole. They articulate several reasons for this need, and it is useful to paraphrase here at least part of their argument: Let us accept that a black hole is a thermodynamic object, and so obeys the first law: $`\delta E=T\delta S+\mathrm{}`$. Now suppose that the universe is asymptotically flat in spatial directions, and contains a single black hole. Then $`E`$ in the first law is the ADM mass. But if there is anything else in the universe then $`E`$ is not the ADM mass, and the question arises, What expression is to be used for $`E`$ in the first law? In other words, we expect that we can put something else in the universe, say a galaxy somewhere, such that the black hole we started with, considered ‘by itself’, will still behave as more or less the same thermodynamic object, with the same mass, radiating at the same temperature as before, and with the same entropy equal to one quarter its area. This expectation requires the ability to compute the energy of a given system contained within a finite closed surface, rather than merely the total energy of all such systems comprising the whole universe.
Thus quasilocal energy lies between the notions of local energy density and total energy of an isolated system, in the sense that it is expected to give the energy contained in any volume, no matter how small or large. Although the equivalence principle precludes the existence of a local gravitational energy density, it does not prevent us from evaluating the (quasilocal) gravitational energy in an arbitrarily small but nonvanishing volume $`\mathrm{\Sigma }`$. This is because no matter how small $`\mathrm{\Sigma }`$ is, $`S=\mathrm{\Sigma }`$ is not a point, but rather the boundary of some neighborhood of a point, and so we are always inherently making a ‘tidal force measurement’. In this sense quasilocal energy is distinct from attempts to define a local gravitational energy density based on certain symmetries of the action, and the concept of a Nöther charge.<sup>2</sup><sup>2</sup>2A recent discussion of the connection between pseudotensor methods and the quasilocal idea can be found in Ref. . At the other extreme is the Komar mass (or the closely related ADM mass), i.e., the total energy associated with the time translation symmetry of an isolated system. As emphasized in Ref. , this gravitational conserved charge is intimately connected with a lapse function, whereas quasilocal energy need not make any reference to a lapse function. The point is, the two are conceptually distinct , even though in some circumstances one might expect their numerical values to coincide.
Currently there are several contenders for a good definition of quasilocal energy (see Refs. and the references therein). The two that interest us at the moment are Brown and York’s ‘canonical quasilocal energy’ (or CQE) , and the various definitions based on the integral over $`S`$ of the Witten-Nester two-form (the two-form used in Witten’s proof of the positive energy theorem ). The latter approach uses spinorial methods, and the different definitions are distinguished by the choice of supplementary equation the $`S`$-spinors are supposed to satisfy, for example the Sen-Witten equation , the Dougan-Mason equation , or the Ludvigsen-Vickers equation . The Brown-York definition of quasilocal energy has the form<sup>3</sup><sup>3</sup>3We use a sign convention for extrinsic curvatures opposite to that of Brown and York, hence the negative sign in front of this integral.
$$\mathrm{CQE}=\frac{1}{8\pi }_S𝑑Sk\mathrm{CQE}^{\mathrm{ref}},$$
(1.1)
in geometrized units, with $`G=c=1`$. The CQE is supposed to be the energy of the gravitational and matter fields contained in a finite spatial volume $`\mathrm{\Sigma }`$, whose boundary two-surface is $`S=\mathrm{\Sigma }`$. $`dS`$ is the induced integration measure on $`S`$, and $`k`$ is the trace of the extrinsic curvature of $`S`$ as embedded in $`\mathrm{\Sigma }`$. Thus, $`k/(8\pi )`$ is the Brown-York quasilocal energy surface density. When $`\mathrm{\Sigma }`$ is asymptotically flat the integral in Eq. (1.1) (the ‘unreferenced’ CQE) diverges as $`S`$ is taken to infinity, and a reference term, denoted $`\mathrm{CQE}^{\mathrm{ref}}`$, is required to regulate the energy. Brown and York’s prescription is to choose
$$\mathrm{CQE}^{\mathrm{ref}}=\frac{1}{8\pi }_S𝑑Sk^{\mathrm{ref}},$$
(1.2)
where $`k^{\mathrm{ref}}`$ is the trace of the extrinsic curvature of an isometric embedding of $`S`$ into some reference space, usually taken to be flat $`\text{I}\text{R}^3`$. With this choice the resulting CQE reduces to the ADM mass when $`S`$ is taken to infinity . While the CQE has a host of desirable properties, neatly summarized in Ref. , the embedding prescription needed to evaluate $`\mathrm{CQE}^{\mathrm{ref}}`$ is not entirely satisfactory, because not all two-surfaces that arise in practice can be embedded into flat $`\text{I}\text{R}^3`$. A ready example is the horizon of the Kerr black hole, which fails to be embeddable in flat $`\text{I}\text{R}^3`$ when the angular momentum exceeds the irreducible mass (but is not yet an extremal black hole), and the two-sphere develops regions with negative scalar curvature . While this is but one example, it is noteworthy that the breakdown of embeddability is in this case associated with angular momentum and negative scalar curvature. It is precisely such issues: embeddability, angular momentum, and negative scalar curvature, that will figure prominently in this paper, and will be seen to be subtly intertwined.
Although a relationship between the Brown-York quasilocal energy and the spinorial definitions based on the Witten-Nester integral is not immediately obvious, Lau has shown that spinors may always be chosen so that the resulting spinorial definition is equal to the unreferenced Brown-York quasilocal energy in Eq. (1.1). Moreover, he shows that the role of the Sen-Witten equation is to provide a definite reference point for the energy, which is not in general the same as $`\mathrm{CQE}^{\mathrm{ref}}`$ in Eq. (1.2). The point being made here is two-fold: (i) the unreferenced Brown-York quasilocal energy seems to be robust, and (ii) all of the problems lie in choosing a suitable reference energy. The various prescriptions are either not generally well defined, or they do not agree with each other.
I will now present a brief review of the Brown-York approach in a form that will be useful to us later. The classical energy-momentum tensor of matter is a local concept, associated with a spacetime point. It is defined for any field theory residing on a nondynamical background spacetime $`(M,g)`$ via the functional derivative of the (first order) matter action with respect to the metric, as follows:
$$2\delta _gI^{\mathrm{mat}}[\phi ,g]=_Md^4x\sqrt{g}T_{\mathrm{mat}}^{ab}\delta g_{ab}.$$
(1.3)
Here $`\phi `$ denotes the matter field(s) in question, and the factor of two on the left is a convention. Usually, as we will assume here, there is no boundary term arising from this variation (for ‘minimally coupled matter’), but in case there is it does not change the essence of the following argument, it just adds an interesting dimension to it. $`T_{\mathrm{mat}}^{ab}`$ so defined is covariantly conserved, as follows from the matter Euler-Lagrange equations. This prescription for learning about matter energy-momentum gives reasonable answers for all field theories, and so it is natural to try the same thing for gravity. In this case one finds, for the usual first order action ,
$$2\delta _gI^{\mathrm{grav}}[g]=_Md^4x\sqrt{g}\left(\frac{1}{8\pi }G^{ab}\right)\delta g_{ab}+_{}d^3x\sqrt{\gamma }\left(\frac{1}{8\pi }\mathrm{\Pi }^{ab}\right)\delta \gamma _{ab}.$$
(1.4)
Inspecting the bulk term one is thus tempted to define $`T_{\mathrm{grav}}^{ab}:=G^{ab}/(8\pi )`$ as the local energy-momentum tensor of the gravitational field, where $`G^{ab}`$ is the Einstein tensor. And this is perfectly reasonable: it is covariantly conserved—in this case identically so, via the contracted Bianchi identity. Moreover, its on-shell value is zero, in full accord with the equivalence principle, i.e., there is no nontrivial local energy-momentum tensor for the gravitational field.
In fact this absence of a nontrivial local energy-momentum tensor is true not only for the gravitational field, but also for any system comprised of both matter and gravity. To see this we need only make the spacetime metric dynamical, in which case the matter action in Eq. (1.3) must be augmented by the gravity action. Adding Eqs. (1.3) and (1.4) one finds for the total action
$$2\delta _gI^{\mathrm{tot}}[\phi ,g]=_Md^4x\sqrt{g}\left(T_{\mathrm{mat}}^{ab}\frac{1}{8\pi }G^{ab}\right)\delta g_{ab}+_{}d^3x\sqrt{\gamma }\left(\frac{1}{8\pi }\mathrm{\Pi }^{ab}\right)\delta \gamma _{ab}.$$
(1.5)
Thus one is led to identify $`T_{\mathrm{tot}}^{ab}:=T_{\mathrm{mat}}^{ab}G^{ab}/(8\pi )`$ as the total local energy-momentum tensor for matter plus gravity. It has the desirable property of being covariantly conserved, but turns out to be just zero by the Einstein equations. If this argument is taken seriously we learn that, as soon as we add gravity to any matter system, the notion of a nontrivial local energy-momentum tensor disappears. Furthermore, one might interpret the Einstein equations, written in the form $`T_{\mathrm{mat}}^{ab}+T_{\mathrm{grav}}^{ab}=0`$, as a ‘micro-balancing’ of local stress-energy-momentum at each spacetime point: wherever a component of matter stress-energy-momentum is positive, the corresponding component of gravitational stress-energy-momentum is negative, and vice versa, such that the total is always zero. The idea that $`G^{ab}/(8\pi )`$ is the local energy-momentum tensor of gravity is, of course, a very old idea, first put forward by Lorentz and Levi-Civita. It was rejected by Einstein, since it implies that the total energy of a closed system would always be zero, which is obviously problematic.<sup>4</sup><sup>4</sup>4See the historical discussion given on pages 176-7 in Ref. . I thank L. de Menezes for bringing this reference to my attention. It is only with hindsight that we now realize why the problem was not resolved much sooner. People then did not think about boundary terms as much as they do today. Thanks to Brown and York we now know that what comes to the rescue is the boundary term in Eq. (1.5).
In this equation the spacetime is assumed to be the topological product of a three-space $`\mathrm{\Sigma }`$ and a real line interval. The boundary component $``$ is a timelike ‘tube’, topologically the product of $`S=\mathrm{\Sigma }`$ and the real line interval. (The two spacelike ‘end-cap’ boundary components of $`M`$ have been omitted, as they play no role in this discussion.) The quantity $`\sqrt{\gamma }\mathrm{\Pi }^{ab}/(16\pi )`$, constructed in the usual way out of the extrinsic curvature of $``$, is the gravitational momentum conjugate to the three-metric $`\gamma _{ab}`$ induced on $``$. Now, in the spirit of identifying the energy-momentum tensor as the functional derivative of the action with respect to the metric, one reads off from Eq. (1.5) the energy-momentum tensor $`T_{}^{ab}:=\mathrm{\Pi }^{ab}/(8\pi )`$, which is inherently associated with the boundary $``$, rather than the bulk spacetime. Like any acceptable energy-momentum tensor, $`T_{}^{ab}`$ is covariantly conserved (as a tensor residing in $``$)—this follows from the analogue of the diffeomorphism constraint of general relativity for the three-surface $``$ (rather than $`\mathrm{\Sigma }`$). Physically, $``$ is to be thought of as the congruence of world lines of a two-parameter family of observers with four-velocity $`u^a`$, hypersurface orthogonal to a one-parameter foliation of $``$ by spacelike two-surfaces with the topology of $`S`$. Within their three-dimensional spacetime $`(,\gamma )`$, the observers measure a spatial energy density $`T_{}^{ab}u_au_b`$, which is precisely $`k/(8\pi )`$, and thus one is led to Eq. (1.1). Finally, observe that one can add to the action any covariant functional of the boundary three-metric $`\gamma _{ab}`$ without affecting the previous argument. This is the source of the reference point ambiguity $`\mathrm{CQE}^{\mathrm{ref}}`$ in Eq. (1.1). This summarizes the central idea of the Brown-York approach .
Now if energy is really quasilocal, and calculated via a surface integral involving $`T_{}^{ab}`$, one comes to the conclusion that a priori neither $`T_{\mathrm{mat}}^{ab}`$ nor $`T_{\mathrm{grav}}^{ab}`$ has anything to do with energy! While this might be unsettling at first, it is reassuring to know that a satisfactory notion of local matter energy density can be recovered from the small sphere limit of quasilocal energy. For example, in Ref. it is shown that, for a certain choice of reference term $`\mathrm{CQE}^{\mathrm{ref}}`$, the Brown-York quasilocal energy contained in an infinitesimal sphere of proper radius $`r`$ is the volume of the sphere ($`4\pi r^3/3`$) times the local matter energy density $`T_{\mathrm{mat}}^{ab}u_au_b`$ (evaluated at the center of the sphere) that would be measured by an observer with four-velocity $`u^a`$. Moreover, this is a well established property of most quasilocal energy definitions , so the result is quite robust. And at higher order in $`r`$, gravitational energy begins to appear, as will be discussed in detail later. The point is there is no contradiction between (i) the local energy-momentum tensor $`T_{\mathrm{mat}}^{ab}+T_{\mathrm{grav}}^{ab}`$ being zero, and (ii) there being nonzero stress-energy-momentum in a finite spatial volume. This is because $`T_{\mathrm{mat}}^{ab}+T_{\mathrm{grav}}^{ab}`$ is not a local energy-momentum tensor—indeed, if we accept the previous argument, there is no such thing. There is only $`T_{}^{ab}`$, associated with the fact that energy is fundamentally quasilocal.
The main purpose of this Introduction is to emphasize, firstly, that energy is fundamentally quasilocal, i.e., associated with closed spacelike two-surfaces—not points—in spacetime; and secondly, there are strong reasons to believe that $`k/(8\pi )`$ is the correct quasilocal energy surface density. The major unresolved problem is how to choose the right well defined energy reference term, $`\mathrm{CQE}^{\mathrm{ref}}`$. Rather than address this problem per se, I will begin with $`k/(8\pi )`$ as an energy surface density and construct a new definition of quasilocal energy based on analogy with the special relativity formula: $`E^2\stackrel{}{p}^{\mathrm{\hspace{0.33em}2}}=m^2`$. The new definition is both physically and geometrically natural, and lies somewhere between the Brown-York CQE and the Hawking or Hayward definitions. A reference subtraction procedure is still required, that involves a reference embedding, but this is a codimension-two embedding that is not subject to the problem that afflicts the Brown-York embedding prescription. Moreover, there is a shift in the physics: the reference embedding is not associated with determining a reference energy so much as a ‘reference angular momentum’, so to speak. Why angular momentum? Because angular momentum contributes to energy, and the new definition can be seen as a precise formulation of this fact at the quasilocal level.
The paper is organized as follows. In Sec. 2 we introduce the geometrical quantities we will use later. Sec. 3 contains the physical and geometrical motivations behind the new definition of quasilocal energy (as well as the definition itself). The reference subtraction term is discussed in Sec. 4. In Sec. 5 both the spatial and future null infinity limits of the energy are examined; the small sphere limit is considered in Sec. 6. Finally, in Sec. 7 we examine the new energy in the context of asymptotically anti-de Sitter spacetimes. Since this is a lengthy paper I have provided a summary of its results at the end, which also includes some additional discussion.
## 2 The geometry of two-dimensional spacelike submanifolds
Let $`(M,g)`$ be a four-dimensional Lorentzian geometry with signature $`+2`$, and $`S`$ be a closed two-dimensional spacelike submanifold. Let $`u^a`$ and $`n^a`$ be timelike and spacelike unit normals to $`S`$ that are orthogonal to each other: $`u^au_a=1`$, $`n^an_a=1`$, and $`u^an_a=0`$. We will assume that $`S`$ is orientable, and an open neighborhood of $`S`$ in $`M`$ is space and time orientable, so that $`u^a`$ and $`n^a`$ are globally well defined . $`u^a`$ and $`n^a`$ are fixed up to an arbitrary local boost transformation:
$`u^a`$ $`=`$ $`u^a\mathrm{cosh}\lambda +n^a\mathrm{sinh}\lambda `$
$`n^a`$ $`=`$ $`u^a\mathrm{sinh}\lambda +n^a\mathrm{cosh}\lambda .`$ (2.1)
The physical picture to keep in mind is that of a finite spatial volume $`\mathrm{\Sigma }`$, i.e., a three-dimensional spacelike submanifold, whose boundary is $`S`$. Although $`S`$ need not be connected, nor simply connected, we will often think of $`\mathrm{\Sigma }`$ as having the topology of a three-ball, and $`S`$ that of a two-sphere, and thus will sometimes refer to the direction of $`n^a`$ (assumed outward directed) as the ‘radial’ direction. Given such a three-surface $`\mathrm{\Sigma }`$ spanning $`S`$ it is natural to choose $`u^a`$ to be orthogonal to $`\mathrm{\Sigma }`$ (and future directed), in which case $`n^a`$ is tangential to $`\mathrm{\Sigma }`$. Physically, $`u^a`$ is the instantaneous four-velocity of a two-parameter family of observers on $`S`$. With $`u^a`$ thus tied to the spanning surface $`\mathrm{\Sigma }`$, a deformation of $`\mathrm{\Sigma }`$ (preserving $`S`$) will in general effect a radial boost, Eqs. (2.1).
The remainder of this section is a summary of some standard facts about the geometry of the submanifold $`S`$, as can be found, e.g., in Ref. , except here we follow a notation similar to that used in Ref. . The surface projection operator, $`𝒫_b^a`$, is a tensor defined on $`S`$ by
$$𝒫_b^a:=\delta _b^a+u^au_bn^an_b.$$
(2.2)
(All raising and lowering of indices $`a,b,c,\mathrm{}`$ will be effected with the metric $`g_{ab}`$, or its inverse, $`g^{ab}`$.) A ‘surface tensor’ is defined as a tensor on $`S`$ that is left invariant under projection of all its indices with the surface projection operator. Obviously one such tensor is the spatial two-metric
$$\sigma _{ab}:=𝒫_a^c𝒫_b^dg_{cd}=g_{ab}+u_au_bn_an_b$$
(2.3)
induced on $`S`$. Another is the corresponding volume form on $`S`$, given by
$$ϵ_{ab}:=ϵ_{abcd}u^cn^d,$$
(2.4)
where $`ϵ_{abcd}`$ is the volume form on $`M`$. The symbol $`dS`$ will be used in place of $`ϵ_{ab}`$ as the integration measure for $`(S,\sigma )`$.
If $`_a`$ denotes the Levi-Civita connection of $`(M,g)`$, then $`𝒟_a`$, the Levi-Civita connection induced on $`(S,\sigma )`$, is defined by
$$𝒟_aT_c\mathrm{}^b\mathrm{}=𝒫_a^d𝒫_e^b𝒫_c^f\mathrm{}_dT_f\mathrm{}^e\mathrm{},$$
(2.5)
where $`T_c\mathrm{}^b\mathrm{}`$ is any surface tensor. Then for any two surface vector fields $`X^a`$ and $`Y^a`$, the Gauss formula reads
$$X^a_aY^c=X^a𝒟_aY^c+h_{ab}^cX^aY^b,$$
(2.6)
where $`h_{ab}^c`$ is the second fundamental form. Its first index is normal to $`S`$, i.e., $`𝒫_c^dh_{ab}^c=0`$, whereas the remaining two are surface tensor indices, that are symmetric under interchange (as can be easily seen by interchanging $`X`$ and $`Y`$ in Eq. (2.6) and subtracting the two equations). Thus the second fundamental form can be decomposed into components along the two unit normals:
$$h_{ab}^c=u^cl_{ab}n^ck_{ab},$$
(2.7)
where the two extrinsic curvatures are (symmetric) surface tensors given by
$`l_{ab}`$ $`=`$ $`u_ch_{ab}^c=𝒫_a^c𝒫_b^d_cu_d`$
$`k_{ab}`$ $`=`$ $`n_ch_{ab}^c=𝒫_a^c𝒫_b^d_cn_d.`$ (2.8)
It is useful to decompose the extrinsic curvatures into trace and trace-free parts:
$`l_{ab}`$ $`=`$ $`{\displaystyle \frac{1}{2}}l\sigma _{ab}+\stackrel{~}{l}_{ab}`$
$`k_{ab}`$ $`=`$ $`{\displaystyle \frac{1}{2}}k\sigma _{ab}+\stackrel{~}{k}_{ab},`$ (2.9)
where $`l=\sigma ^{ab}l_{ab}`$ and $`k=\sigma ^{ab}k_{ab}`$, and a tilde appearing over any quantity in this paper will always mean ‘trace-free part of’. The mean curvature vector is then
$$H^c:=\frac{1}{2}\sigma ^{ab}h_{ab}^c=\frac{1}{2}(lu^ckn^c),$$
(2.10)
and $`HH=(k^2l^2)/4`$ is the square of the mean curvature.
Extrinsic curvature is a measure of how a unit normal vector rotates as it is parallelly propagated tangent to $`S`$ in the ambient space $`(M,g)`$. Two normal vectors means two extrinsic curvatures. However, from Eqs. (2.8) we see that $`l_{ab}`$ and $`k_{ab}`$ measure only the components of this rotation tangent to $`S`$. There is also a normal component, i.e., the component of the rotation of one normal vector along the other. Thus a complete characterization of the extrinsic geometry of $`S`$ requires also the surface one-form
$$A_a:=𝒫_a^bn^c_bu_c.$$
(2.11)
This is an $`SO(1,1)`$ connection in the normal bundle of $`S`$, and its associated curvature two-form is
$$_{ab}:=𝒟_aA_b𝒟_bA_a.$$
(2.12)
We will see later that the curvature of the normal bundle of $`S`$ plays a key role with regard to angular momentum.
While the second fundamental form (including $`H^c`$ and $`HH`$) and the curvature of the normal bundle are invariant under local radial boosts, the extrinsic curvatures and the connection on the normal bundle are not. They transform as
$`l_{ab}^{}`$ $`=`$ $`l_{ab}\mathrm{cosh}\lambda +k_{ab}\mathrm{sinh}\lambda `$
$`k_{ab}^{}`$ $`=`$ $`l_{ab}\mathrm{sinh}\lambda +k_{ab}\mathrm{cosh}\lambda `$
$`A_a^{}`$ $`=`$ $`A_a+𝒟_a\lambda .`$ (2.13)
Observe that $`A_a`$ is different from the other extrinsic curvatures in that its transformation law, which is a gauge transformation, depends on the derivative of $`\lambda `$.
Our sign conventions are such that the Riemann tensor of $`(M,g)`$ is defined by $`(_a_b_b_a)X_c=R_{abc}^dX_d`$, and similarly that of $`(S,\sigma )`$ by $`(𝒟_a𝒟_b𝒟_b𝒟_a)X_c=_{abc}^dX_d`$ ($`X_c`$ is a surface one-form in the latter case). Appropriate projections of the Riemann tensor of $`(M,g)`$ yield the Gauss equation:
$$𝒫_a^e𝒫_b^f𝒫_c^g𝒫_d^hR_{efgh}=_{abcd}+(l_{ac}l_{bd}l_{bc}l_{ad})(k_{ac}k_{bd}k_{bc}k_{ad}),$$
(2.14)
the Codazzi equations:
$`𝒫_a^e𝒫_b^f𝒫_c^gu^hR_{efgh}`$ $`=`$ $`(𝒟_al_{bc}𝒟_bl_{ac})(A_ak_{bc}A_bk_{ac}),`$
$`𝒫_a^e𝒫_b^f𝒫_c^gn^hR_{efgh}`$ $`=`$ $`(𝒟_ak_{bc}𝒟_bk_{ac})(A_al_{bc}A_bl_{ac}),`$ (2.15)
and the Ricci equation:
$$𝒫_a^e𝒫_b^fu^gn^hR_{efgh}=_{ab}+(k_a^cl_{bc}l_a^ck_{bc}).$$
(2.16)
These are the integrability conditions for the isometric embedding of $`(S,\sigma )`$ into $`(M,g)`$, and so by definition of $`S`$ are necessarily satisfied.
## 3 The invariant quasilocal energy
A physical interpretation of the various geometrical quantities introduced in the previous section can be given as follows. The expansion $`k`$ measures the fractional expansion of the area of a small element of $`S`$ when each point in the element is projected a unit distance radially outward. It will have a certain positive value if, for example, $`S`$ is a round sphere enclosing a volume of flat $`\text{I}\text{R}^3`$. (For our present purposes, imagine ‘flat $`\text{I}\text{R}^3`$’ as a $`t=\mathrm{constant}`$ surface in Minkowski space.) Now if $`S`$ is a round sphere of the same area enclosing some matter, then, according to the Einstein equations, the matter curves the space inside $`S`$ in such a way that its volume is greater than one would infer by measuring just the area of the sphere and using Euclidean geometry. Thus the expansion measured at $`S`$ must be smaller, i.e., the areas of spherical shells at larger radii will not increase as rapidly as expected. So we see that the unreferenced Brown-York quasilocal energy in Eq. (1.1) is greater (less negative) when $`S`$ contains matter, than when it does not. This is an intuitive reason why $`k`$ is a measure of the energy inside $`S`$. (It also explains the need to subtract off a reference energy of the form given in Eq. (1.2): $`k^{\mathrm{ref}}`$ is the nonzero value of $`k`$ when $`S`$ merely encloses a volume of flat $`\text{I}\text{R}^3`$, i.e., no energy.)
$`l`$ is similar to $`k`$, except that it measures the expansion of $`S`$ in time, i.e., in the direction of the observer’s four-velocity $`u^a`$. Intuitively, if the observers tend to be moving radially outward then the area of the two-surface they are on will be expanding, i.e., $`l>0`$. Conversely, a radially inward motion corresponds to $`l<0`$. Thus $`l`$ (more precisely, $`l/(8\pi )`$) can be interpreted as a ‘radial momentum surface density’ . In the case that $`(k^2l^2)`$ is positive, the observers can always make appropriate local radial boosts such that $`l=0`$ at each point of $`S`$, a situation corresponding to a ‘quasilocal rest frame’. I will comment on this notion more precisely at the end of this section.
The trace-free quantities $`\stackrel{~}{k}_{ab}`$ and $`\stackrel{~}{l}_{ab}`$ measure the shear of $`S`$, and are intimately connected with angular momentum (or at least $`\stackrel{~}{l}_{ab}`$ is). For example, consider a set of locally nonrotating observers who at coordinate time $`t`$ are on a constant $`r,t`$ sphere of the Kerr black hole in Boyer-Lindquist coordinates. Their four-velocity is given by
$$u^a=\frac{1}{N}\left(\frac{}{t}+\omega \frac{}{\varphi }\right)^a,$$
(3.1)
where $`N`$ is the lapse function, and $`\omega (r,\theta )=g_{t\varphi }/g_{\varphi \varphi }`$ is an observer’s angular velocity as measured from infinity . Starting at Eqs. (2.8) it is not difficult to show that in this case $`l=0`$, so here is an example of observers in a ‘quasilocal rest frame’ as defined above. Furthermore, one can show that the nonvanishing components of the shear in the time direction are given by
$$\stackrel{~}{l}_{\theta \varphi }=\stackrel{~}{l}_{\varphi \theta }=\frac{g_{\varphi \varphi }}{2N}\frac{\omega }{\theta }.$$
(3.2)
Physically, a nonzero $`\omega `$ reflects the frame dragging caused by the rotating black hole. The fact that the degree of frame dragging depends on $`\theta `$ is what makes the observers at different latitudes of the sphere rotate at different rates relative to ‘the distant stars’, and more to the point, relative to each other. This causes a shear effect between observers at neighboring latitudes, which obviously disappears when the angular momentum is zero.
Furthermore, let the locally nonrotating observers label themselves with coordinates $`(\theta ^{},\varphi ^{})`$, which at $`t=0`$ coincide with the Boyer-Lindquist $`(\theta ,\varphi )`$ coordinates on $`S`$. Then although the observers always measure the same two-geometry of $`S`$ as time $`t`$ goes on, the components of the two-metric $`\sigma _{ab}`$ in their $`(\theta ^{},\varphi ^{})`$ coordinates will differ from those in the $`(\theta ,\varphi )`$ coordinates by a $`t`$\- and $`\theta `$-dependent diffeomorphism along the $`\varphi `$-direction. So although one usually associates shear with a geometrical deformation, for instance a round sphere evolving into an ellipsoid, one can also have a physically meaningful shear associated with a continuous parameter family of isometric surfaces. This fact plays an important role in understanding certain embedding equations we will encounter later, and will be discussed in detail elsewhere .
The last geometrical quantity to interpret is $`A_a`$, the connection in the normal bundle. In the Brown-York analysis the quantity $`A_a/(8\pi )`$ is called the ‘momentum surface density’, and is denoted as $`j_a`$ . The momentum vector $`j^a`$ is tangential to $`S`$, corresponding to a ‘rotating two-surface’, and thus should be associated with angular momentum. Indeed, this is correct: Let $``$ denote the timelike three-surface that is the congruence of world lines belonging to the two-parameter family of observers on a two-sphere $`S`$. If $``$ admits a Killing vector field $`\varphi ^a`$, whose orbits lie in $`S`$, then one can define the angular momentum charge
$$J:=_S𝑑S\varphi ^aj_a,$$
(3.3)
which can be shown to coincide with the ADM angular momentum at infinity for asymptotically flat spacetimes . Thus we expect both the shear and the connection in the normal bundle to play a role in angular momentum at the quasilocal level, and indeed we will see that this turns out to be the case.
Now the first goal of this paper is to provide a physical motivation for the general relativistic analogue of the special relativity formula: $`E^2\stackrel{}{p}^{\mathrm{\hspace{0.33em}2}}=m^2`$. First of all, this formula applies strictly to point particles (as opposed to extended objects). One imagines determining, say, the instantaneous three-velocity of such a particle by measuring its location in space at two closely separated points in time, in some inertial reference frame. In the spirit of the quasilocal idea, the analogue of this in general relativity would be to first replace measurements at a point with measurements on a closed spacelike two-surface $`S`$. But measurements of what? It would seem that ‘measurements of the location of the point particle’ ‘in some inertial frame’ is to be replaced with ‘measurements of the two-geometry of $`S`$’ ‘in a generic spacetime’. These measurements are to be repeated at two closely separated points in time. In the point particle case this yields the three-velocity (or the three-momentum $`\stackrel{}{p}`$ if one also knows $`m`$); in the two-surface case it yields $`l_{ab}`$, the time component of the extrinsic curvature of $`S`$. Now I pointed out above that the trace of $`l_{ab}`$—more precisely $`l/(8\pi )`$—indeed has the interpretation of a momentum: it is the normal (or radial) momentum surface density . So it seems reasonable to replace $`\stackrel{}{p}`$ with $`l/(8\pi )`$. What about the trace-free part of $`l_{ab}`$? It was argued above that $`\stackrel{~}{l}_{ab}`$ is associated with angular momentum. Insofar as angular momentum is qualitatively distinct from linear momentum, its role at least at this point of the argument is not clear, and we will simply drop it for now (however, its role will become clear later). Notice that dropping $`\stackrel{~}{l}_{ab}`$ is at least roughly consistent with being interested only in the two-geometry of $`S`$, i.e., the two-metric $`\sigma _{ab}`$ modulo diffeomorphisms, since, as indicated above, $`\stackrel{~}{l}_{ab}`$ is in some cases associated with just diffeomorphisms of $`S`$.
Thus, in the expression $`E^2\stackrel{}{p}^{\mathrm{\hspace{0.33em}2}}`$ we will replace $`\stackrel{}{p}`$ with $`l/(8\pi )`$. What should replace $`E`$? Given the previous discussion, the obvious answer is the Brown-York energy surface density, $`k/(8\pi )`$. Clearly $`l/(8\pi )`$ and $`k/(8\pi )`$ are on exactly the same geometrical footing, being proportional to the timelike and spacelike components of the mean curvature vector $`H^c`$ (see Eq. (2.10)). Thus we arrive at the generalization
$$E^2\stackrel{}{p}^{\mathrm{\hspace{0.33em}2}}\frac{1}{(8\pi )^2}(k^2l^2).$$
(3.4)
Now before we accept this generalization, let us observe that there is something unexpected about it. The four-momentum $`(E,\stackrel{}{p})`$ has become a two-momentum, $`(k,l)/(8\pi )`$. What happened to the other two components of spatial momentum? $`l/(8\pi )`$ is just the radial component; shouldn’t the Brown-York momentum surface density $`j^a`$ (in our notation, $`A^a/(8\pi )`$), which is tangent to $`S`$, be the analogue of the two missing components of $`\stackrel{}{p}`$? If so, then instead of Eq. (3.4) we should have
$$E^2\stackrel{}{p}^{\mathrm{\hspace{0.33em}2}}\stackrel{\mathrm{?}}{}\frac{1}{(8\pi )^2}(k^2l^2A^aA_a).$$
(3.5)
At first sight this expression is appealing because it manifestly includes a contribution from angular momentum, and it is known that in general relativity angular momentum contributes to mass. A simple example that illustrates this phenomenon is the Kerr black hole, where the ADM mass in excess of the irreducible mass is due to rotational energy. The precise relationship is
$$M_{\mathrm{ir}}^2=M^2\left(\frac{J}{2M_{\mathrm{ir}}}\right)^2,$$
(3.6)
where $`M`$ is the ADM mass, $`M_{\mathrm{ir}}`$ the irreducible mass, and $`J`$ the angular momentum of the black hole. Comparing the right hand sides of the previous two equations suggests we conceptually identify $`|A|/(8\pi )`$ with the angular momentum term, $`J/(2M_{\mathrm{ir}})`$, which seems reasonable. This leaves $`\sqrt{k^2l^2}/(8\pi )`$ to be interpreted as an object like $`M`$, viz., a total mass, ‘total’ in the sense that it includes the contribution from angular momentum. But here then is the point: $`\sqrt{k^2l^2}/(8\pi )`$ somehow implicitly already includes the angular momentum contribution to mass. Precisely how will become clear later, but to see immediately that this is at least plausible, consider the case $`l=0`$. Then $`\sqrt{k^2l^2}/(8\pi )`$ reduces to the Brown-York energy surface density, at least when $`k`$ is nonnegative, and it is known that the (referenced) Brown-York quasilocal energy yields the ADM mass at spatial infinity, which includes the correct angular momentum contribution to mass. So we do not need the $`A^aA_a`$ term in Eq. (3.5). Besides, putting it in is counter to our goal of seeing if general relativity admits an analogue of the invariant mass, $`m`$: While the combination $`k^2l^2`$ is invariant under radial boosts,<sup>5</sup><sup>5</sup>5This was first noted in Ref. . A further discussion of boosted observers in the Brown-York framework appears in Ref. . $`A^aA_a`$ is not—see Eqs. (2.13). So from this point of view the right hand side of Eq. (3.5) is defective, not to mention the generally unsavory fact that it mixes objects with different transformation properties.<sup>6</sup><sup>6</sup>6Hayward’s definition of quasilocal energy includes an angular momentum contribution of the form $`\omega ^a\omega _a`$, analogous to the $`A^aA_a`$ term in Eq. (3.5). Hayward’s $`\omega _a`$ is a suitably normalized anholonomicity, or ‘twist’, of the pair of null normals to $`S`$, and encodes essentially the same information as $`A_a`$. The important distinction is that, unlike the connection $`A_a`$, the object $`\omega _a`$ is boost invariant, and so representing angular momentum with a term proportional to $`\omega ^a\omega _a`$, as Hayward does, is perfectly acceptable. (The relationship between $`A_a`$ and $`\omega _a`$ is discussed in Appendix B of Ref. .) However, there is no need, or even natural way for $`\omega _a`$ to enter our work here. For instance, the $`A^aA_a`$ term in Eq. (3.5) cannot simply be replaced with $`\omega ^a\omega _a`$, since the (tentative) inclusion of this term is suggested by the physical interpretation of $`A_a`$ as a momentum surface density. This interpretation arises from Brown and York’s Hamilton-Jacobi analysis of the gravitational action , and it is not clear that a similar interpretation can be given to $`\omega _a`$. The question of ‘missing momentum components’ can also be thought about as follows. A point particle has three components of spatial momentum. Likewise, each point on a two-surface $`S`$ also has three components of spatial momentum (more properly, momentum surface density): one normal, and two tangential to $`S`$. But being tangential, the latter two are associated with a ‘rotating surface’, and hence with angular momentum. In going from a point to a two-surface, two components of the linear momentum have become angular momenta. So they do not (directly at least) contribute to the expression for an invariant mass given on the right hand side of Eq. (3.4) because, as claimed, this expression already inherently includes the contribution from angular momentum.
Thus we are led to propose the following definition of an ‘invariant quasilocal energy’ (or IQE):
$$\mathrm{IQE}=\frac{1}{8\pi }_S𝑑S\sqrt{k^2l^2}\mathrm{IQE}^{\mathrm{ref}},$$
(3.7)
where $`\mathrm{IQE}^{\mathrm{ref}}`$ is a reference subtraction term that will be defined later. The word ‘invariant’ in ‘IQE’ refers to the fact just mentioned, that $`k^2l^2`$ is invariant under local radial boosts of the observers on $`S`$. And the word ‘energy’ is used instead of ‘mass’—despite our analogy between $`\sqrt{k^2l^2}/(8\pi )`$ and the mass $`m`$—because, as we will see in Sec. 7, the IQE behaves more like an energy than a mass. So the IQE can be thought of as the amount of ‘rest energy’ contained in $`S`$, a quantity independent of the motion of the observers measuring it. Notice that the unreferenced IQE is negative. Nominally the reference energy $`\mathrm{IQE}^{\mathrm{ref}}`$ is more negative, so that the referenced IQE is positive.
It is useful to express the integrand of the unreferenced IQE in two other equivalent forms. Define the pair of null normals $`\xi _\pm ^a:=u^a\pm n^a`$ on $`S`$, and the corresponding null expansions
$$\theta _\pm :=\sigma ^{ab}_a\xi _{\pm b}=l\pm k,$$
(3.8)
cf. Eqs. (2.8) and (2.9). Then we have the following three equivalent expressions:
$$\frac{1}{8\pi }\sqrt{k^2l^2}=\frac{1}{8\pi }\sqrt{\theta _+\theta _{}}=\frac{1}{4\pi }\sqrt{HH}.$$
(3.9)
For the last expression recall the definition of the mean curvature given after Eq. (2.10). Thus the unreferenced IQE in Eq. (3.7) has a very simple geometrical interpretation: up to a proportionality constant, it is just the mean curvature of $`S`$, averaged over $`S`$. Because of the square root it is defined only when $`k^2l^20`$, i.e., at each point of $`S`$ the mean curvature vector $`H^c`$ in Eq. (2.10) must be either spacelike or null, but never timelike. Roughly speaking, this means that the area of $`S`$ changes more rapidly in a radial direction, than in time. For example, this condition is satisfied for the constant $`r,t`$ two-spheres outside the horizon of a Schwarzschild black hole, but not for those inside; on the horizon the unreferenced IQE is zero.
In terms of the null expansions, recall that a future (past) trapped surface is one for which both ingoing and outgoing null expansions, $`\theta _{}`$ and $`\theta _+`$, are everywhere negative (positive) on $`S`$ . Thus, the unreferenced IQE is imaginary when $`S`$ is a future or past trapped surface. It is real only when no point on $`S`$ is ‘trapped’. Now a future trapped surface does not quite characterize a black hole, and more subtle characterizations have been proposed for a local definition of a black hole horizon . For example, Hayward has introduced the notion of a ‘future outer trapping horizon’, $`H`$, characterized by: (i) $`\theta _{}|_H<0`$ (in-going light rays converging), (ii) $`\theta _+|_H=0`$ (outgoing light rays instantaneously parallel on the horizon), (iii) $`\theta _+|_{H^+}>0`$ (outgoing light rays diverging just outside the horizon), and (iv) $`\theta _+|_H^{}<0`$ (outgoing light rays converging just inside the horizon). According to this general definition of a black hole, the unreferenced IQE is nonzero just outside the horizon, zero on the horizon, and undefined (or imaginary) just inside the horizon. In this connection see also Ref. .
Furthermore, observe that the condition for the integrand of the unreferenced IQE to be real and nonzero, namely $`k^2l^2>0`$, is precisely the same condition that ensures that the observers can always, by appropriate local boosts, go to a ‘quasilocal rest frame’ in which $`l=0`$ at each point of $`S`$. Such a two-surface is analogous to a massive particle. The case $`k^2l^2=0`$ everywhere on $`S`$, for instance when $`S`$ is a future outer trapping horizon, is analogous to a massless particle, for which no quasilocal rest frame exists. And finally, the case $`k^2l^2<0`$, say ‘inside’ a future outer trapping horizon, corresponds to a superluminal particle.<sup>7</sup><sup>7</sup>7As noted in the text, the spirit of the quasilocal idea is to replace measurements at a point (of certain aspects of a point particle, say) with measurements on a closed spacelike two-surface. If one takes seriously that Eq. (3.4) is the generalization of point particle rest mass, then one is quickly led to speculate that a closed spacelike two-surface is the generalization of a point particle. This is curiously reminiscent of string theory, except that the one-dimensional string is replaced by a two-dimensional surface.
The situation is actually more subtle than indicated in the previous two paragraphs. The conditions for $`\sqrt{\theta _+\theta _{}}`$ to be real are reminiscent of the condition $`\theta _{}0`$ required for the holomorphic case of the Dougan-Mason quasilocal energy to be nonnegative. The conditions $`\theta _+0`$ and $`\theta _{}0`$ essentially imply that the two-surface $`S`$ is ‘suitably’ convex . To emphasize that “ ‘suitably’ convex” is not a serious restriction, in particular it does not mean that $`S`$ cannot be concave, consider a two-parameter family of observers at rest in an inertial frame in flat spacetime. Suppose that at $`t=0`$ they lie on a two-sphere $`S`$ that is round except for a small indentation. Then $`l=0`$ at each point of $`S`$, and $`k`$ is positive everywhere except in a small region near the center of the indentation, where it is negative. Thus there will be a circle of points $`C`$ at which $`k=0`$. So at each point of $`S`$ we have $`k^2l^20`$, equality holding on $`C`$. One might worry that a radial boost at a point on $`C`$ will make $`l^2>0`$, and hence $`k^2l^2<0`$. But of course this will not happen: If we consider a second set of observers, boosted relative to the first, then $`k=l=0`$ on $`C`$ implies $`k^{}=l^{}=0`$ on $`C`$. So we can consider the second set of observers to be boosted radially outward in the region of the indentation, such that the indentation, and its attendant set of fixed points $`C`$, smoothly disappear as the sphere evolves in time. $`k`$ switches from negative to positive by passing through the origin of a $`k`$$`l`$ diagram. Thus we can imagine a wide class of two-surfaces, including ones with indentations, and dynamically changing in time, for which $`k^2l^20`$ everywhere on $`S`$. Moreover, bear in mind that the observers are allowed to accelerate, so there is a great deal of freedom for them to maintain a ‘physically reasonable’ $`S`$. Nevertheless, what is needed here is a careful analysis based on Raychaudhuri-like equations for a two-parameter family of accelerated timelike curves (as opposed to the more usual case of timelike geodesics). Such a detailed analysis is outside the scope set for this introductory paper.
## 4 The reference invariant quasilocal energy
As in the Brown-York case, the unreferenced IQE diverges in an asymptotically flat spacetime as the two-surface $`S`$ is taken to (spatial or null) infinity, and so must be regulated with a reference term, $`\mathrm{IQE}^{\mathrm{ref}}`$, as already anticipated in Eq. (3.7). To better understand the nature of our definition of the invariant quasilocal energy, and to help suggest a natural choice for $`\mathrm{IQE}^{\mathrm{ref}}`$, we now make use of the Gauss embedding equation given in Eq. (2.14). This equation has only one independent component. Transvecting both sides with $`\sigma ^{ac}\sigma ^{bd}`$ reduces it to the scalar equation
$$\sigma \sigma R=\frac{1}{2}(k^2l^2)+(\stackrel{~}{k}^2\stackrel{~}{l}^2),$$
(4.1)
where $`\sigma \sigma R`$ is shorthand for $`\sigma ^{ac}\sigma ^{bd}R_{abcd}`$, $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)`$ is shorthand for $`(\stackrel{~}{k}^{ab}\stackrel{~}{k}_{ab}\stackrel{~}{l}^{ab}\stackrel{~}{l}_{ab})`$, and $``$ is the scalar curvature of $`(S,\sigma )`$. Using this equation we can express the IQE given in Eq. (3.7) in the equivalent form
$$\mathrm{IQE}=\frac{1}{8\pi }_S𝑑S\sqrt{2\left[\sigma \sigma R+(\stackrel{~}{k}^2\stackrel{~}{l}^2)\right]}\mathrm{IQE}^{\mathrm{ref}}.$$
(4.2)
We remark here that $`\sigma \sigma R`$ is a natural geometrical object called the sectional curvature of $`(S,\sigma )`$ as embedded in $`(M,g)`$ . It will play an important role in what follows.
Now the definition of the unreferenced IQE is rooted in the extrinsic geometry of the submanifold $`(S,\sigma )`$, thought of as a two-surface isometrically embedded in the spacetime $`(M,g)`$. It is then natural to define the reference IQE to be of the same form as the unreferenced IQE in Eq. (4.2), i.e., to be the same geometrical object, except with $`(S,\sigma )`$ now isometrically embedded in a different spacetime—some reference spacetime, $`(M^{\mathrm{ref}},g^{\mathrm{ref}})`$. Thus $`\mathrm{IQE}^{\mathrm{ref}}`$ will be the integral in Eq. (4.2) (or (3.7)), except with all quantities referred to the reference spacetime, which we indicate with a superscript ‘ref’. Note that although the extrinsic geometry of $`S`$ will be different in $`(M^{\mathrm{ref}},g^{\mathrm{ref}})`$, its intrinsic geometry, by assumption, will not. So in the $`\mathrm{IQE}^{\mathrm{ref}}`$ integral we are constructing we can set $`dS^{\mathrm{ref}}=dS`$, $`^{\mathrm{ref}}=`$, and $`(\sigma \sigma R)^{\mathrm{ref}}=\sigma \sigma R^{\mathrm{ref}}`$. Also note that, in general, $`M^{\mathrm{ref}}M`$ (topologically). For example, $`S`$ may be a two-sphere embedded in a black hole spacetime, with $`M=R^2\times S^2`$, whereas the reference spacetime might be Minkowski space, with $`M^{\mathrm{ref}}=R^4`$. With this understanding, we define
$$\mathrm{IQE}^{\mathrm{ref}}=\frac{1}{8\pi }_S𝑑S\sqrt{(k^2l^2)^{\mathrm{ref}}}=\frac{1}{8\pi }_S𝑑S\sqrt{2\left[\sigma \sigma R^{\mathrm{ref}}+(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}\right]}.$$
(4.3)
The term $`\sigma \sigma R^{\mathrm{ref}}`$ is shorthand for $`\sigma ^{ac}\sigma ^{bd}R_{abcd}^{\mathrm{ref}}`$, where $`R_{abcd}^{\mathrm{ref}}`$ is the Riemann tensor of the reference spacetime.
Typically one is motivated to choose a reference spacetime of constant curvature, the geometrical reason being that then the Gauss, Codazzi, and Ricci embedding equations make no reference to ‘where’ $`(S,\sigma )`$ is embedded in $`(M^{\mathrm{ref}},g^{\mathrm{ref}})`$. In other words, the conditions placed on $`k_{ab}^{\mathrm{ref}}`$, $`l_{ab}^{\mathrm{ref}}`$, and $`A_a^{\mathrm{ref}}`$ by the ‘reference version’ of Eqs. (2.142.16)—which are just integrability conditions for the reference embedding—do not depend on knowing the embedding itself . This is a pleasing criterion because it keeps the reference spacetime ‘abstract’, rather than ‘concrete’. For a four-dimensional space of constant curvature we have
$$R_{abcd}^{\mathrm{ref}}=\frac{C}{12}(g_{ac}^{\mathrm{ref}}g_{bd}^{\mathrm{ref}}g_{bc}^{\mathrm{ref}}g_{ad}^{\mathrm{ref}}),$$
(4.4)
where $`C`$ is the constant value of its scalar curvature. For this choice of reference spacetime one gets
$$\sigma \sigma R^{\mathrm{ref}}=\frac{C}{6}.$$
(4.5)
For example, for Minkowski space we have $`C=0`$, and for anti-de Sitter space we have $`C=12/\mathrm{}^2`$, where $`\mathrm{}`$ is the radius of curvature of the anti-de Sitter space, and is related to the (negative) cosmological constant $`\mathrm{\Lambda }`$ by $`\mathrm{\Lambda }=3/\mathrm{}^2`$. We will return to these two examples later.
The idea of embedding $`(S,\sigma )`$ into some reference space(time) is in the same spirit as the Brown-York approach, but an important difference that arises out of using the invariant quantity $`\sqrt{k^2l^2}`$, rather than $`k`$, deserves further comment. In the Brown-York approach $`k`$ is the trace of the extrinsic curvature of $`(S,\sigma )`$ as embedded in a three-geometry, $`(\mathrm{\Sigma },h)`$, where $`\mathrm{\Sigma }`$ is a spacelike three-surface spanning $`S`$, with induced metric $`h_{ab}`$. Thus it is natural to take $`k^{\mathrm{ref}}`$ as the trace of the extrinsic curvature of $`(S,\sigma )`$ as embedded in some three-dimensional reference space, $`(\mathrm{\Sigma }^{\mathrm{ref}},h^{\mathrm{ref}})`$. So the embeddings of $`S`$, for both the unreferenced and reference CQE, inherently have a three-dimensional target space. On the other hand, $`\sqrt{k^2l^2}`$ is proportional to a geometrical invariant of $`S`$, namely its mean curvature, and makes no essential reference to a spanning three-surface $`\mathrm{\Sigma }`$ (making the IQE ‘truly quasilocal’ in the sense that it depends on $`S`$ alone<sup>8</sup><sup>8</sup>8The CQE can also be made ‘truly quasilocal’, in a slightly different sense: by relaxing the restriction that the foliation of the spacetime (i.e., $`\mathrm{\Sigma }`$) be orthogonal to the boundary $``$ , it is shown in Ref. that the resulting CQE no longer depends on $`\mathrm{\Sigma }`$, but instead just depends on the foliation of $``$.). As a consequence, the embeddings inherently have a four-dimensional target space(time).
The advantage of a three-dimensional target reference space, say flat $`\text{I}\text{R}^3`$, is that when the embedding exists, it is unique (up to translations and rotations), and so the Brown-York $`\mathrm{CQE}^{\mathrm{ref}}`$ is unique. The disadvantage, as is well known, is that such embeddings do not exist for all $`(S,\sigma )`$ of interest, and this problem is not limited to just a few isolated exceptional cases.
For a four-dimensional target reference spacetime, say Minkowski space, the situation is reversed: an embedding always exists, but it is not unique. Regarding the first half of this statement, Brinkmann has shown, by a simple explicit construction, that any $`n`$-dimensional conformally flat Riemann space can be considered as a particular cut of a light cone in $`(n+2)`$-dimensional Minkowski space. And conversely, any cut of such a light cone gives an $`n`$-dimensional conformally flat Riemann space. Now any $`n=2`$ space is, of course, conformally flat, and thus any $`(S,\sigma )`$ can always be so embedded, even if $`S`$ has regions with negative scalar curvature. In the Introduction I mentioned the example of the horizon of the Kerr black hole, which cannot be globally embedded into flat $`\text{I}\text{R}^3`$ when the angular momentum exceeds the irreducible mass, which coincides with the two-sphere developing regions with negative scalar curvature . However, it is a simple exercise to apply Brinkmann’s construction and thus globally embed the horizon into a light cone of four-dimensional Minkowski space. I will omit the details of this calculation, and just note that the embedding is valid for all angular momentum $`J`$ (up to and including the extremal black hole case), and changes smoothly with $`J`$, including at the critical point when $`J`$ equals the irreducible mass.
On the other hand, in a codimension-two (versus -one) embedding there is more ‘elbow room’, and consequently the embedding is not unique—there is a function worth of freedom (which will be discussed in detail in Ref. ). This results in an ambiguity in the reference energy, $`\mathrm{IQE}^{\mathrm{ref}}`$, which enters via the reference shear term $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$—see Eq. (4.3). This is the only term in $`\mathrm{IQE}^{\mathrm{ref}}`$ not yet determined, and the only one for which we require a reference embedding. Observe that in the Brown-York approach the undetermined quantity is $`k^{\mathrm{ref}}`$, an expansion. Here it is $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$, a shear. I will now argue that it is precisely this term that plays the key role in properly incorporating angular momentum into the IQE. The basic idea is simple, but first we will introduce some notation.
In the previous section we introduced the null normals $`\xi _\pm ^a=u^a\pm n^a`$, and corresponding null expansions $`\theta _\pm =l\pm k`$ in Eq. (3.8). Similarly, the (trace-free) shears in the two null directions are defined by $`s_{\pm ab}:=\stackrel{~}{l}_{ab}\pm \stackrel{~}{k}_{ab}`$. The curvature of the normal bundle has only one independent component, and can be written as $`_{ab}=(/2)ϵ_{ab}`$ for some scalar field $``$, where $`ϵ_{ab}`$ is the volume form on $`S`$ defined earlier in Eq. (2.4). With this notation, and assuming that the reference spacetime is one of constant curvature, i.e., Eq. (4.5) holds, the Gauss, Codazzi, and Ricci embedding equations given at the end of Sec. 2 take the form
$`{\displaystyle \frac{C}{6}}`$ $`=`$ $`+{\displaystyle \frac{1}{2}}\theta _+^{\mathrm{ref}}\theta _{}^{\mathrm{ref}}s_{+b}^{\mathrm{ref}a}s_a^{\mathrm{ref}b},`$ (4.6)
$`0`$ $`=`$ $`{\displaystyle \frac{1}{2}}(𝒟_aA_a^{\mathrm{ref}})\theta _\pm ^{\mathrm{ref}}(𝒟_bA_b^{\mathrm{ref}})s_{\pm a}^{\mathrm{ref}b},`$ (4.7)
$`0`$ $`=`$ $`^{\mathrm{ref}}+{\displaystyle \frac{1}{2}}ϵ_b^a[s_+^{\mathrm{ref}},s_{}^{\mathrm{ref}}]_a^b.`$ (4.8)
In the Ricci equation, $`[s_+^{\mathrm{ref}},s_{}^{\mathrm{ref}}]_a^b`$ denotes the commutator of the shears: $`s_{+c}^{\mathrm{ref}b}s_a^{\mathrm{ref}c}s_c^{\mathrm{ref}b}s_{+a}^{\mathrm{ref}c}`$. Notice that by using null directions, rather than $`u^a`$ and $`n^a`$, the Codazzi equations have decoupled into a ‘$`+`$’ and a ‘$``$’ set.
Our task is thus: Given $`\sigma _{ab}`$, and hence $``$, $`𝒟_a`$, and $`ϵ_{ab}`$, solve these embedding equations for the unknown quantities $`\theta _\pm ^{\mathrm{ref}}`$, $`s_{\pm ab}^{\mathrm{ref}}`$, and $`A_a^{\mathrm{ref}}`$. (Of course $`^{\mathrm{ref}}=2ϵ^{ab}𝒟_aA_b^{\mathrm{ref}}`$ is not an independent quantity.) In particular, we are interested in the solution for the boost invariant reference shear term
$$(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}s_{+b}^{\mathrm{ref}a}s_a^{\mathrm{ref}b},$$
(4.9)
appearing in the Gauss equation, which is to then be substituted into the second integral of Eq. (4.3). Or equivalently, solve for $`(k^2l^2)^{\mathrm{ref}}\theta _+^{\mathrm{ref}}\theta _{}^{\mathrm{ref}}`$ and substitute the answer into the first integral of Eq. (4.3). This is how $`\mathrm{IQE}^{\mathrm{ref}}`$ is determined.
However, as already noted, any solution we obtain is not unique. We can see this immediately by counting functional degrees of freedom. $`\theta _\pm ^{\mathrm{ref}}`$ are two functions, $`s_{\pm ab}^{\mathrm{ref}}`$ are four (the two shears are symmetric and trace-free), and $`A_a^{\mathrm{ref}}`$ are two. These eight functions are subject to six equations: Gauss is one, Codazzi are four, and Ricci is one. This leaves two arbitrary functions in the solution. But owing to the invariance of the embedding equations under a local boost transformation (see Eqs. (2.13)), one of these functions is just the boost parameter, $`\lambda `$, leaving one nontrivial arbitrary function in the solution.
The question then arises, Is there a natural way to impose one additional functional condition on the unknowns so that the embedding, subject to this additional condition, is unique, and hence $`\mathrm{IQE}^{\mathrm{ref}}`$ is unique? One of the central ideas in this paper is to impose the additional condition
$$^{\mathrm{ref}}=,$$
(4.10)
i.e., the curvature of the normal bundle of $`S`$ as embedded in the reference spacetime $`(M^{\mathrm{ref}},g^{\mathrm{ref}})`$ should equal that of $`S`$ as embedded in the original physical spacetime, $`(M,g)`$.
There are several reasons why it is geometrically natural to demand $`^{\mathrm{ref}}=`$. First, the two-surface $`S`$ has two connections: one is an $`SO(2)`$ connection on the tangent bundle of $`S`$, associated with the curvature $`_{abcd}`$ (which has only one independent component, namely $``$), and the other is an $`SO(1,1)`$ connection on the normal bundle of $`S`$, associated with the curvature $`_{ab}`$ (which also has only one independent component, namely $``$). In fact both of these connections are metric connections, associated with the metrics in the tangent and normal bundles to $`S`$, respectively . Furthermore, Szabados has considered the two-dimensional version of the Sen connection for spinors and tensors on a submanifold such as $`S`$, and has found that the two-surface spinor curvature has, essentially, imaginary part equal to $``$, and real part equal to $``$. Finally, although $`A_a`$ is a measure of extrinsic geometry, as pointed out earlier it is not really on the same footing as the extrinsic curvatures $`k_{ab}`$ and $`l_{ab}`$, since its transformation law under local radial boosts is qualitatively different—see Eqs.(2.13). It transforms like the connection that it is, and gives rise to a curvature, and so arguably has more in common with $``$ than with $`k_{ab}`$ and $`l_{ab}`$. The point is, $``$ is really on the same geometrical footing as $``$. We have already demanded that $`^{\mathrm{ref}}=`$, as a necessary condition for the embedding of $`(S,\sigma )`$ into $`(M^{\mathrm{ref}},g^{\mathrm{ref}})`$ to be isometric. So demanding also that $`^{\mathrm{ref}}=`$ is thus seen to be quite natural.
Unfortunately, implementing Eq. (4.10) seems like an intractable task. Embedding equations involving curvature of the normal bundle, i.e., codimension-two (and higher) embeddings, have, of course, been studied for a long time. With regard to solutions, although one expects to be able to express $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$ in terms of $``$, $``$, and their derivatives, I am not aware of any such general results in the literature. In fact, much of the literature on such embeddings considers the case $`=0`$, which is not the case we are particularly interested in here (a notable exception is Ref. ). One possible way to proceed is as follows. Given $`^{\mathrm{ref}}`$ ($`=`$ by Eq. (4.10)), choose $`A_a^{\mathrm{ref}}`$ such that $`^{\mathrm{ref}}=2ϵ^{ab}𝒟_aA_b^{\mathrm{ref}}`$. There may be a convenient gauge choice, such as $`𝒟A^{\mathrm{ref}}=0`$, or $`l^{\mathrm{ref}}=0`$. Then view the Codazzi equations (4.7) as a set of four linear partial differential equations for the four independent degrees of freedom in the two shears. These equations are of second order if one makes use of the fact that any trace-free symmetric tensor $`s_{ab}`$ on a two-surface $`(S,\sigma )`$ with two-sphere topology can be expressed as $`s_{ab}=𝒟_av_b+𝒟_bv_a\sigma _{ab}𝒟v`$ for some vector field $`v^a`$. Thus solve for $`s_{\pm ab}^{\mathrm{ref}}`$ in terms of the expansions, $`\theta _\pm ^{\mathrm{ref}}`$, and their derivatives. The expressions one obtains at this stage are, in general, nonlocal. Then substitute these into the Gauss and Ricci equations, (4.6) and (4.8), which are really just nonlinear algebraic constraints. But because the shears involve nonlocal operators acting on the expansions, one ends up with two nonlocal and nonlinear partial differential equations for the two expansions. Remarkably, it is almost possible to solve these equations, but in the end one encounters a certain combination of nonlocality and nonlinearity that makes the final step to a solution seem impossible. Nevertheless, it appears that the solution for $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$, if it can be found, almost certainly depends in a simple way on both $``$ and $``$, and derivatives (of a finite or possibly infinite order) of these two curvatures. I am suggesting that it is through this subtle presence of $``$ in $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$ that angular momentum is properly incorporated into the IQE.
So although a direct attack on the embedding equations has not yet yielded a solution, fortunately one can make some progress of a general nature by calculating the first and second order variations of $`^{\mathrm{ref}}`$ and $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$ under isometric deformations of a given embedding. The idea is to see how both of these quantities change under such a deformation, and thereby infer how $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$ depends on $`^{\mathrm{ref}}`$, and hence angular momentum. The results are somewhat involved, and will be given elsewhere . For now let us start by making some simple observations regarding the enigmatic object $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$.
To begin with, one might object to our argument thus far because it implies that $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$, and hence $`\mathrm{IQE}^{\mathrm{ref}}`$, depends on the extrinsic geometry of $`S`$ as embedding in the physical spacetime. In particular, through Eq. (4.10) and the reference embedding equations, $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$ depends on $``$. On the other hand, it is often stated that a reference subtraction term should be a functional of only the intrinsic geometry of $`(S,\sigma )`$. However, notice that there is no dependence on the extrinsic curvatures proper, i.e. $`l_{ab}`$ and $`k_{ab}`$, only a dependence on $``$, a quantity which I argued above is really on the same geometrical footing as the intrinsic quantity $``$. Moreover, as discussed in the Introduction (refer to Eq. (1.5)), in the Brown-York approach one is free to add to the action any functional of the boundary three-metric, $`\gamma _{ab}`$, which contains information about the two-metric $`\sigma _{ab}`$, as well as information about how $`(S,\sigma )`$ is embedded in the three-boundary $``$. For instance, one could add to the action a boundary integral of the scalar curvature of $``$, whose variation would add to $`\mathrm{\Pi }^{ab}`$ in Eq. (1.5) a term proportional to the Einstein tensor of $`\gamma _{ab}`$, as is done in Ref. . Such a term obviously depends on some extrinsic geometry of $`(S,\sigma )`$. In Brown and York’s work this fact is of course recognized, but being in a Hamiltonian framework, they restrict the form of the arbitrary boundary functional such that the energy surface density ($`k/(8\pi )`$) and momentum surface density ($`A_a/(8\pi )`$) of $`S`$ in a particular spacelike hypersurface $`\mathrm{\Sigma }`$ depend only on the canonical data on $`\mathrm{\Sigma }`$. This effectively means that their reference subtraction term can depend only on $`\sigma _{ab}`$ . But as I emphasized earlier, our approach is based on the invariant object $`\sqrt{k^2l^2}`$, and makes no essential reference to a three-surface $`\mathrm{\Sigma }`$ spanning $`S`$. The invariant quasilocal energy constructed here does not come out of a canonical analysis, so there is no reason that our subtraction term cannot depend on $``$.
So the shear term $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$ is allowed, but is it really necessary? Perhaps it is just an unsavory term resulting from a poor definition of the IQE. For instance, looking at the Gauss embedding equation (4.1) one might be tempted to write, instead of Eq. (3.4),
$$E^2\stackrel{}{p}^{\mathrm{\hspace{0.33em}2}}\stackrel{\mathrm{?}}{}\frac{1}{(8\pi )^2}[(k^2l^2)2(\stackrel{~}{k}^2\stackrel{~}{l}^2)],$$
(4.11)
where the additional shear term on the right hand side is perhaps the proper way to include angular momentum, somewhat like the $`A^aA_a`$ term we attempted in Eq. (3.5). This would have the advantage of changing $`\mathrm{IQE}^{\mathrm{ref}}`$ in Eq. (4.3) to
$$\mathrm{IQE}^{\mathrm{ref}}\stackrel{\mathrm{?}}{=}\frac{1}{8\pi }_S𝑑S\sqrt{2\left[\sigma \sigma R^{\mathrm{ref}}\right]},$$
(4.12)
which is clearly unique, and moreover, requires no reference embedding. Eq. (4.12) is a more general case of the zero point energy suggested by Lau (except that his derivation of it requires a reference embedding—we will return to this point later). But unfortunately it cannot be correct. For example, when the reference spacetime is Minkowski space, Eq. (4.5) tells us that $`\sigma \sigma R^{\mathrm{ref}}=0`$ and thus the radical reduces to $`\sqrt{2}`$, which is not defined for negative $``$. Nor is this problem properly solved by taking $`C0`$ in Eq. (4.5), since this would put an ad hoc fixed lower bound on $``$.
So the shear term $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$ is not only allowed, it is necessary (or at least its absence leads to an unsatisfactory result). In fact its role seems to be to keep nonnegative what is under the square root in Eq. (4.3). To see this more clearly, consider the special case that $`(S,\sigma )`$ is embeddable in flat $`\text{I}\text{R}^3`$. If our reference spacetime is Minkowski space, we can then choose to embed $`(S,\sigma )`$ in a $`t=\mathrm{constant}`$ slice and, within this slice, the embedding is essentially unique. In this case it is easy to see that we will have $`l_{ab}^{\mathrm{ref}}=0`$. And by assumption, $`\sigma \sigma R^{\mathrm{ref}}=0`$, so Eq. (4.3) reduces to
$$\mathrm{IQE}^{\mathrm{ref}}|_{l_{ab}^{\mathrm{ref}}=0}=\frac{1}{8\pi }_S𝑑S|k^{\mathrm{ref}}|=\frac{1}{8\pi }_S𝑑S\sqrt{2\left[+(\stackrel{~}{k}^{\mathrm{ref}})^2\right]}.$$
(4.13)
The uniqueness of the embedding means that $`k^{\mathrm{ref}}`$ and $`\stackrel{~}{k}_{ab}^{\mathrm{ref}}`$ are unique. Now clearly, no matter what the surface is, the spatial shear $`\stackrel{~}{k}_{ab}^{\mathrm{ref}}`$ must be such that what is under the square root in Eq. (4.13) is nonnegative, because $`|k^{\mathrm{ref}}|`$ is real. (More properly, one should look at the ‘reference version’ of Eq. (4.1), with $`(l^{\mathrm{ref}})^2=(\stackrel{~}{l}^{\mathrm{ref}})^2=\sigma \sigma R^{\mathrm{ref}}=0`$.) For example, consider a dumb-bell-shaped surface of revolution in flat $`\text{I}\text{R}^3`$. In a region near the throat of this surface $``$ is negative, nevertheless at every point of the surface we have $`+(\stackrel{~}{k}^{\mathrm{ref}})^20`$.
I emphasize that, even when $`(S,\sigma )`$ can be embedded in flat $`\text{I}\text{R}^3`$, its embedding in Minkowski space need not be chosen to be in a $`t=\mathrm{constant}`$ slice, as was done in the previous paragraph. One may also embed it in a light cone, or in a host of other ways—remember that there is a function-worth of freedom in our choice. I argued in the context of Eq. (4.10) that this freedom has to do with angular momentum, or more precisely, the curvature of the normal bundle. For the $`t=\mathrm{constant}`$ embedding, $`\stackrel{~}{l}_{ab}^{\mathrm{ref}}=0`$ implies $`s_{\pm ab}^{\mathrm{ref}}=\pm \stackrel{~}{k}_{ab}^{\mathrm{ref}}`$, and so inspection of Eq. (4.8) reveals that in this case $`^{\mathrm{ref}}=0`$. However, it is not hard to see that starting with such a $`t=\mathrm{constant}`$ embedding one can perform an infinitesimal isometric deformation of the embedding out of the $`t=\mathrm{constant}`$ plane, i.e., in a direction $`\phi /t`$, where $`\phi `$ is an arbitrary function. Furthermore, I show in Ref. that after such an infinitesimal deformation $`^{\mathrm{ref}}`$ is no longer zero, and can in fact be made to be essentially any infinitesimal function we like by a suitable choice of $`\phi `$. It is not hard to imagine (just hard to do!) that by integrating such isometric deformations one may be able to achieve a two-geometry $`(S,\sigma )`$ isometrically embedded with any desired curvature of the normal bundle. With this in mind, it would seem unnatural to reference-embed, e.g., a constant $`r,t`$ two-sphere of the Kerr geometry (which is easily shown to have $`0`$) as a two-sphere in a Minkowski reference spacetime with $`^{\mathrm{ref}}=0`$ (say, in a $`t=\mathrm{constant}`$ slice), when it seems possible to instead embed it with $`^{\mathrm{ref}}=`$. Note, however, that as the embedding is deformed out of the $`t=\mathrm{constant}`$ surface, $`\stackrel{~}{l}_{ab}^{\mathrm{ref}}`$ will also cease to be zero, and so will introduce a negative contribution to the quantity under the square root in Eq. (4.3). This jeopardizes the nonnegativity of this quantity. But at the same time we clearly cannot simply throw away the $`(\stackrel{~}{l}^{\mathrm{ref}})^2`$ term, since this would violate a key property of the IQE, namely its invariance under local boosts. Short of solving the embedding equations for a generic $`^{\mathrm{ref}}`$ and explicitly checking, I do not know of any guarantee of nonnegativity. Equivalently, in question here is the nonnegativity of the quantity $`(k^2l^2)^{\mathrm{ref}}`$, which is the same type of question as the nonnegativity of $`(k^2l^2)`$ discussed at the end of Sec. 3. I do not at present have a complete answer to either of these difficult questions.
Finally, one might guess that it is possible to avoid the embedding problem entirely by simply setting $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}=(\stackrel{~}{k}^2\stackrel{~}{l}^2)`$, which is in the same spirit as Eq. (4.10) in that, like $``$, the shear term has something to do with angular momentum. But this does not work. For example, it is not hard to show that, although $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)`$ is in general nonvanishing on constant $`r,t`$ spheres of the Kerr geometry, it happens to vanish on the horizon. So if we set $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}=(\stackrel{~}{k}^2\stackrel{~}{l}^2)`$, then in calculating $`\mathrm{IQE}^{\mathrm{ref}}`$ for the Kerr horizon example we would run into the same problem with negative $``$ as we did in Eq. (4.12). So such a prescription must not be valid, and we cannot avoid the embedding problem this way.
Let us conclude this section by addressing the question, What is the relationship between the IQE defined here and the Brown-York CQE? We begin by supposing that the following four conditions are satisfied: (i) $`(S,\sigma )`$ is a two-surface in the physical spacetime such that $`k^2l^2>0`$; (ii) $`k>0`$; (iii) $`(S,\sigma )`$ is such that it can be embedded in flat $`\text{I}\text{R}^3`$; and (iv) for the embedding in (iii), $`k^{\mathrm{ref}}0`$. As discussed earlier, condition (i) ensures that the unreferenced IQE is well defined—roughly speaking, $`S`$ is strictly outside of a black hole. Then it is always possible to go to a ‘quasilocal rest frame’ where $`l=0`$ on $`S`$, and the integrand in Eq. (3.7) is just $`|k|`$. Given condition (ii), the unreferenced IQE thus reduces to the unreferenced CQE in Eq. (1.1), provided the observers in the Brown-York case are in a ‘quasilocal rest frame’. Condition (iii) ensures that the Brown-York prescription is well-defined, and allows us to choose a $`t=\mathrm{constant}`$ embedding in Minkowski space, as above, and get Eq. (4.13). The first integral in this equation, together with condition (iv), shows that our $`\mathrm{IQE}^{\mathrm{ref}}`$ reduces to the Brown-York $`\mathrm{CQE}^{\mathrm{ref}}`$ in Eq. (1.2). So if these conditions hold, and we choose to use a $`t=\mathrm{constant}`$ embedding to calculate $`\mathrm{IQE}^{\mathrm{ref}}`$, then our invariant quasilocal energy is the same as the Brown-York ‘rest energy’. In most applications considered in the literature these conditions are satisfied, and the IQE will then share all of the desirable properties of the CQE. For example, it will be the thermodynamic energy that appears in the first law of black hole thermodynamics for Schwarzschild black holes, as considered in Ref. .
On the other hand, I emphasize that conditions (ii) and (iv) are easily violated. One need only think of a round sphere with a small indentation, an example discussed at the end of Sec. 3 (except here the spacetime is generic). So in general, the IQE defined here is not simply the Brown-York ‘rest energy’. Furthermore, we need not choose a $`t=\mathrm{constant}`$ embedding to calculate $`\mathrm{IQE}^{\mathrm{ref}}`$. Indeed, as I have argued, such a choice is unnatural when $`0`$. In short, the invariant quasilocal energy defined here is not quite the same object as the Brown-York quasilocal energy. Note that the aforementioned thermodynamic nature of the CQE is derived in Ref. assuming that $`S`$ is a round sphere. It would be interesting to extend this analysis to indented spheres, for instance, and determine which, if either of the CQE or IQE, is the ‘correct’ thermodynamic energy.
## 5 The large sphere limit of the IQE
Let us now assume that the spacetime $`(M,g)`$ is asymptotically flat, and evaluate the limit of the invariant quasilocal energy as $`(S,\sigma )`$ tends to a large sphere at infinity. At spatial and null infinity we might expect these limits to be the ADM and Bondi-Sachs masses, respectively. Let us see if this is so.
Let $`\tau `$ be a ‘time’ function on $`M`$ such that $`\tau =\tau _{}`$ defines a spacelike (respectively, null) hypersurface $`_\tau _{}`$ of topology $`R\times S^2`$ extending to spatial (respectively, null) infinity. Letting the parameter $`\tau _{}`$ vary over some range gives a foliation of a part of $`M`$. Let $`r`$ be a function on $`M`$ such that $`r=r_{}`$ defines a hypersurface that intersects each leaf $`_\tau _{}`$ (over the allowed range of $`\tau _{}`$) in a spacelike two-sphere, $`S_{\tau _{},r_{}}`$. The parameter $`r_{}`$ ranges to infinity, and over its range the surfaces $`S_{\tau _{},r_{}}`$ provide a foliation of $`_\tau _{}`$. We are interested in the limit $`r_{}\mathrm{}`$, with $`\tau _{}`$ arbitrary but fixed. In a rather benign abuse of notation we will refer to $`S_{\tau _{},r_{}}`$ as simply $`S`$, and ‘take the limit as $`r\mathrm{}`$ with $`\tau `$ fixed’. The metric induced on $`S`$ will, as usual, be denoted as $`\sigma _{ab}`$ (in abstract index notation).
Now assume that the functions $`\tau `$ and $`r`$ have been chosen such that $`(S,\sigma )`$ tends to a round sphere at infinity. Thus the components of its metric in spherical coordinates $`x^i=(\theta ,\varphi )`$ have an asymptotic expansion of the form
$$\sigma _{ij}=r^2\left(\begin{array}{cc}1& 0\\ 0& \mathrm{sin}^2\theta \end{array}\right)+2r\left(\begin{array}{cc}X& Y\mathrm{sin}\theta \\ Y\mathrm{sin}\theta & Z\mathrm{sin}^2\theta \end{array}\right)+O_<(r).$$
(5.1)
In this expansion, $`X`$, $`Y`$, and $`Z`$ are each arbitrary functions of $`\tau `$, $`\theta `$, and $`\varphi `$. The symbol $`O_<(r^n)`$ denotes a term that falls off faster (or grows more slowly, depending on the sign of $`n`$) than $`r^n`$, but not necessarily according to a power of $`r`$. For example, rather than $`O(1)`$, the remainder term $`O_<(r)`$ might grow as $`\mathrm{ln}r`$. The motivation for this increased generality will be explained below when we consider the large sphere limit at null infinity. Furthermore, we can choose (the function) $`r`$ to be an areal radius, in which case we may take $`\sqrt{\sigma }=r^2\mathrm{sin}\theta `$, where $`\sigma =det\sigma _{ij}`$. It is easy to see that this requires $`Z=X`$ in Eq. (5.1).
The scalar curvature of a round sphere of areal radius $`r`$ is $`2/r^2`$. Since the metric in Eq. (5.1) differs from that of a round sphere by a term one power lower in $`r`$, we immediately have that its scalar curvature $``$ is given by
$$=\frac{2}{r^2}+\frac{\mathrm{\Delta }_{}}{r^3},$$
(5.2)
where the remainder term $`\mathrm{\Delta }_{}`$ is of order one.<sup>9</sup><sup>9</sup>9 In Ref. it is shown that $`\mathrm{\Delta }_{}=𝒟v+O_<(1)`$ for some vector field $`v`$ in $`S`$. In other words, to leading order $`\mathrm{\Delta }_{}`$ is a divergence. This plays a crucial role in some of the results in Ref. . However, we will not need to use this fact, except to provide some insight into our discussion of the ‘solution’ of the Ricci embedding equation in Sec. 5.2. In our asymptotically flat spacetime the components of the Riemann tensor fall off as $`1/r^3`$, and so the same will be true of $`\sigma \sigma R`$, the sectional curvature of $`(S,\sigma )`$. In the present context the appropriate reference spacetime is Minkowski space, and so $`\sigma \sigma R^{\mathrm{ref}}=0`$. The only other terms to consider in Eqs. (4.2) and (4.3) are the shear terms, $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)`$ and $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$. We will see below that these, too, fall off at least as fast as $`1/r^3`$ in both the spatial and null infinity limits. In the large sphere limit the unreferenced IQE thus behaves as
$`\mathrm{IQE}^{\mathrm{unref}}`$ $`=`$ $`{\displaystyle \frac{1}{8\pi }}{\displaystyle _S}𝑑S\sqrt{2\left[{\displaystyle \frac{2}{r^2}}+{\displaystyle \frac{\mathrm{\Delta }_{}}{r^3}}\sigma \sigma R+(\stackrel{~}{k}^2\stackrel{~}{l}^2)\right]}`$ (5.3)
$`=`$ $`{\displaystyle \frac{1}{8\pi }}{\displaystyle _S}𝑑S{\displaystyle \frac{2}{r}}\left\{1+{\displaystyle \frac{r^2}{4}}\left[{\displaystyle \frac{\mathrm{\Delta }_{}}{r^3}}\sigma \sigma R+(\stackrel{~}{k}^2\stackrel{~}{l}^2)\right]+O(r^2)\right\}.`$
The reference IQE behaves similarly, except we have $`\sigma \sigma R^{\mathrm{ref}}=0`$. Thus
$$\mathrm{IQE}^{\mathrm{ref}}=\frac{1}{8\pi }_S𝑑S\frac{2}{r}\left\{1+\frac{r^2}{4}\left[\frac{\mathrm{\Delta }_{}}{r^3}0+(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}\right]+O(r^2)\right\}.$$
(5.4)
In forming the difference of the previous two expressions it is important to observe that not only do the divergent terms coming from the $`2/r^2`$ piece of $``$ cancel, but also the (finite) remainder terms $`\mathrm{\Delta }_{}`$ are the same in both, and thus also cancel, independent of what $`\mathrm{\Delta }_{}`$ is. Thus we find that the large sphere behavior of the (referenced) IQE is given by
$$\mathrm{IQE}=\frac{1}{16\pi }_S𝑑Sr[\sigma \sigma R(\stackrel{~}{k}^2\stackrel{~}{l}^2)+(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}+O(r^4)].$$
(5.5)
Of course if the shear terms fall off as $`1/r^4`$ or faster they get absorbed into the $`O(r^4)`$ remainder term. It is worth emphasizing that in the large sphere limit the square root in the IQE is eliminated by the fact that $``$ dominates over the other terms. The areal radius factor $`r`$ outside the brackets in Eq. (5.5) is really $`\sqrt{2/}`$.<sup>10</sup><sup>10</sup>10This mechanism works even when the sphere is not asymptotically round. In this case the shear terms contribute at a higher order, viz. $`1/r^2`$, in an effort to keep what is under the square root sign positive, as discussed earlier. In other words, the factor $`\sqrt{2/}`$ is modified in such a way that negative $``$ is likely not a problem. We will not consider this more complicated case here. In Sec. 6 we will see a similar mechanism at work in the small sphere limit. But in the intermediate regime the IQE is, in general, an integral of the difference of two radicals.
As a quick check of Eq. (5.5) let us evaluate the right hand side for the Schwarzschild geometry. In the usual Schwarzschild coordinates $`r`$ is an areal radius, and it is a simple exercise to compute the sectional curvature of a $`t,r=\mathrm{constant}`$ two-sphere. The result is: $`\sigma \sigma R=4M/r^3`$. The shear terms obviously vanish, and with $`dS=r^2d\mathrm{\Omega }`$ ($`d\mathrm{\Omega }`$ the measure on the unit round sphere) one immediately gets $`\mathrm{IQE}=M`$, the ADM mass of the black hole. Now the main task is to investigate in detail the shear terms in Eq. (5.5), which we will do separately for the spatial and null infinity limits, respectively.
### 5.1 The spatial infinity limit
Rather than proceed with complete generality, it is more instructive to consider an asymptotically flat metric that exhibits angular momentum explicitly, and then see how this angular momentum works its way into the shear terms. (We will be completely general in the more interesting null infinity limit case.) The spacetime far from any isolated stationary (nonradiating) rotating source is described asymptotically by the Kerr metric (see Secs. 19.3 and 33.3 of Ref. ), so let us take $`(M,g)`$ to be the Kerr spacetime. We choose the following basis of orthonormal one-forms:
$$e^0=Ndt,e^1=\frac{\rho }{\sqrt{\mathrm{\Delta }}}dr,e^2=\rho d\theta ,e^3=\sqrt{g_{\varphi \varphi }}(d\varphi \omega dt),$$
(5.6)
which are associated with locally nonrotating observers. The corresponding basis of orthonormal vector fields is
$$e_0=\frac{1}{N}(_t+\omega _\varphi ),e_1=\frac{\sqrt{\mathrm{\Delta }}}{\rho }_r,e_2=\frac{1}{\rho }_\theta ,e_3=\frac{1}{\sqrt{g_{\varphi \varphi }}}_\varphi .$$
(5.7)
The notation used is standard: $`x^a=(t,r,\theta ,\varphi )`$ are Boyer-Lindquist coordinates, $`\omega (r,\theta )=g_{t\varphi }/g_{\varphi \varphi }`$ is an observer’s angular velocity as measured from infinity, $`N=\sqrt{\omega ^2g_{\varphi \varphi }g_{tt}}`$ is the lapse function, etc. (see, e.g., Sec. 33.4 of Ref. ). Let $`A,B,\mathrm{}`$ be indices labeling the basis vectors and one-forms, ranging from 0 to 3, and $`I,J,\mathrm{}`$ denote the subset of these taking values 2 and 3. These indices are raised and lowered with the flat Lorentz metric $`\eta _{AB}=\eta ^{AB}=\mathrm{diagonal}(1,1,1,1)`$.
The vector fields $`e_I^a`$ are tangent to any $`r,t=\mathrm{constant}`$ two-sphere $`S`$, and so $`u^a:=e_0^a`$ and $`n^a:=e_1^a`$ are, respectively, timelike and spacelike unit vectors orthogonal to $`S`$. From Eqs. (2.8) we see that the orthonormal basis components of the extrinsic curvatures $`l_{ab}`$ and $`k_{ab}`$ are given by
$`l_{IJ}`$ $`=`$ $`e_I^ae_J^b_au_b=\omega _{0JI},`$ (5.8)
$`k_{IJ}`$ $`=`$ $`e_I^ae_J^b_an_b=\omega _{1JI}.`$ (5.9)
Here $`\omega _{CBA}=e_A^ae_B^b_ae_{Cb}`$ are Ricci rotation coefficients. Working out these coefficients<sup>11</sup><sup>11</sup>11The easiest way to do this is to recognize that, with $`Z=0`$ or 1, $`\omega _{ZJI}=\alpha _{(IJ)Z}`$, where $`\alpha _{AB}^C=i_{e_B}i_{e_A}de^C`$. Thus we need only compute the exterior derivative of $`e^I`$. I find that the trace-free parts of $`l_{IJ}`$ and $`k_{IJ}`$ are given by
$`\stackrel{~}{l}_{IJ}`$ $`=`$ $`\alpha \left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\mathrm{where}\alpha ={\displaystyle \frac{\sqrt{g_{\varphi \varphi }}}{2N\rho }}_\theta \omega ,`$ (5.12)
$`\stackrel{~}{k}_{IJ}`$ $`=`$ $`\beta \left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),\mathrm{where}\beta ={\displaystyle \frac{\sqrt{\mathrm{\Delta }}}{2\rho }}_r\mathrm{ln}{\displaystyle \frac{\rho }{\sqrt{g_{\varphi \varphi }}}}.`$ (5.15)
Geometrically, the coefficient $`\alpha `$ is just $`e_1\mathrm{\Omega }^{(\mathrm{precess})}`$, i.e., the radial component of the angular velocity vector $`\mathrm{\Omega }^{(\mathrm{precess})}`$ that measures the precession of a gyroscope carried by a locally nonrotating observer, relative to the observer’s orthonormal frame (see Eq. (33.24) of Ref. ; compare also with Eq.(3.2) above, and the discussion following it). Thus, the unreferenced shear term in Eq. (5.5) is given by
$$(\stackrel{~}{k}^2\stackrel{~}{l}^2)=2(\beta ^2\alpha ^2)=\frac{a^4}{2r^6}\mathrm{sin}^4\theta +O\left(\frac{1}{r^7}\right),$$
(5.16)
where the last expression on the right hand side is the large $`r`$ asymptotic expansion. So clearly, being of order $`1/r^6`$, $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)`$ does not contribute to the large sphere limit of the IQE at spatial infinity.
What about the reference term $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$? It is plausible that the reference term is of the same order in $`1/r`$ as the unreferenced term, viz. $`1/r^6`$, or less, and so also does not contribute. However, to be certain one needs to solve the embedding equations (4.6-4.8), subject to the condition $`^{\mathrm{ref}}=`$, as argued in Sec. 4. We will not attempt to do so here, but it is instructive to at least work out what $``$ is for the Kerr geometry. From Eq. (2.11) we see that the orthonormal basis components of the connection in the normal bundle are given by
$$A_I=e_I^bn^c_bu_c=\omega _{01I}.$$
(5.17)
Evaluating these Ricci rotation coefficients reveals that the one-form $`A=A_Ie^I`$ (pulled back to the two-sphere $`S`$) is given by
$$A=\gamma d\varphi ,\mathrm{where}\gamma =\frac{g_{\varphi \varphi }\sqrt{\mathrm{\Delta }}}{2N\rho }_r\omega =\frac{3aM}{r^2}\mathrm{sin}^2\theta +O\left(\frac{1}{r^4}\right).$$
(5.18)
Recall that $`\omega =\omega (r,\theta )`$ is a measure of the frame dragging produced by the rotating geometry. While the shear in the time direction measures the $`\theta `$ dependence of $`\omega `$ (see $`\alpha `$ in Eq. (5.12)), Eq. (5.18) shows that the connection in the normal bundle measures its $`r`$ dependence. Both are measures of angular momentum.
Calculating the exterior derivative of $`A`$ leads to the curvature in the normal bundle:
$$=\frac{1}{\sqrt{g_{\varphi \varphi }}\rho }_\theta \left(\frac{g_{\varphi \varphi }\sqrt{\mathrm{\Delta }}}{2N\rho }_r\omega \right)=\frac{6aM}{r^4}\mathrm{cos}\theta +O\left(\frac{1}{r^6}\right),$$
(5.19)
which is of order $`1/r^4`$. Inspection of the Ricci equation (4.8), or Eq. (2.16) with the left hand side set to zero, reveals that a solution to the reference embedding equations, subject to Eq. (4.10), requires that $`\stackrel{~}{l}_{IJ}^{\mathrm{ref}}`$ and $`\stackrel{~}{k}_{IJ}^{\mathrm{ref}}`$ be two matrices whose commutator is of order $`1/r^4`$. On the other hand, one might guess that the trace of the difference of their squares, $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$, might be of order $`1/r^6`$, as suggested above. It is not difficult to convince oneself that these two conditions are not incompatible, so the reference embedding equations at least do not obviously forbid the reference shear term $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$ from being of the same order of magnitude as $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)`$, such that neither contributes to the IQE.
In any case, assuming just that $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$ is at most $`O_<(r^3)`$, which is almost certainly true, we find that the large sphere limit of the IQE at spatial infinity is given by
$$\underset{r\mathrm{}}{lim}\mathrm{IQE}=\underset{r\mathrm{}}{lim}\frac{1}{16\pi }_S𝑑Sr\sigma \sigma R.$$
(5.20)
This limit of the IQE thus has a simple geometrical interpretation: apart from a factor proportional to the areal radius $`r`$, it is just the average over $`S`$ of the sectional curvature of $`(S,\sigma )`$ as embedded in the physical spacetime $`(M,g)`$.
Now let us assume that $`(M,g)`$ is vacuum ($`R_{ab}=0`$) near spatial infinity, so that there the Riemann tensor reduces to the Weyl tensor, $`C_{abcd}`$. From the definition of the two-surface metric given in Eq. (2.3), and the fact that the Weyl tensor is traceless, one immediately gets
$$\sigma \sigma R=2E_{ab}n^an^b.$$
(5.21)
Thus the sectional curvature of $`(S,\sigma )`$ is just (twice) the radial-radial (Coulomb) component of the electric part of the Weyl tensor, $`E_{ab}:=C_{acbd}u^cu^d`$. Inserting this result into Eq. (5.20) we see that in this limit the IQE is precisely the coordinate-independent expression of the ADM mass given by Ashtekar and Hansen .<sup>12</sup><sup>12</sup>12In Ref. a different definition of $`E_{ab}`$ is used, namely $`E_{ab}=C_{acbd}n^cn^d`$, but accounting for this difference in notation the two results agree.
One more remark is in order here: Hayward’s work on quasilocal energy resembles what is done here in the sense that the Gauss embedding equation plays a central role, and that the analysis is boost invariant in spirit, i.e., no reference is made to a spacelike three-surface spanning $`S`$, with its attendant preferred time direction on $`S`$, and so on. However, Hayward’s quasilocal energy is distinct from the IQE here: it does not involve a square root. Basically, Hayward starts with an integral over $`S`$ of the ‘2+2’ Hamiltonian density, which yields a dimensionless quantity, and then multiplies this quantity by the areal radius of $`S`$ to correct this ‘defect’, i.e., give it the dimensions of energy. This is in the same spirit as the areal radius factor appearing in the Hawking mass . For the large sphere limit at spatial infinity Hayward arrives at the same result given in Eq. (5.20), except with $`r`$ ‘outside the integral’, so to speak. His quasilocal energy has the very appealing feature of not requiring a reference subtraction term, at least when the sectional curvature $`\sigma \sigma R`$ falls off as $`1/r^3`$ (however it diverges if, e.g., the spacetime is asymptotically anti-de Sitter space—recall Eq. (4.5)). In our case, the square root in Eq. (4.2) ensures that the IQE has the dimensions of energy, but the price paid is that a reference subtraction term is needed. (Without the square root the large sphere limit of the unreferenced IQE would just be (negative) the Euler number of $`S`$, which is finite, but carries no information about energy.) The areal radius factor $`r`$ in Eq. (5.20) appears ‘inside the integral’: as mentioned before, it arises from the dominant scalar curvature term $``$, and is really $`\sqrt{2/}`$. I emphasize that this is a geometrically natural mechanism—$`r`$ is not put in by hand. Finally, while one might feel that there is something unattractively ad hoc about a reference subtraction term, the flexibility it affords makes it possible to deal with the wide range of boundary conditions possible in general relativity. In Sec. 7 we will consider an interesting example of this.
### 5.2 The null infinity limit
We now suppose that the physical spacetime $`(M,g)`$ is asymptotically flat at future null infinity. As in the previous subsection we begin with the generic large sphere form of the IQE given in Eq. (5.5), except now we take the large $`S`$ limit in the future null direction. More precisely, on $`M`$ we introduce the Bondi coordinates $`x^a=(w,r,\theta ,\varphi )`$, and as before, denote the subset $`(\theta ,\varphi )`$ of spherical coordinates by $`x^i`$. The retarded time $`w`$ labels a one-parameter family of outgoing null hypersurfaces, and $`r`$ is an areal radius (luminosity parameter) along the outgoing null geodesic generators of these hypersurfaces. The $`w,r=\mathrm{constant}`$ surfaces are topologically two-spheres, any one of which we denote as $`S`$. This setup is the same as discussed at the beginning of Sec. 5, where $`w`$ here is what we there called the ‘time’ coordinate, $`\tau `$. We are interested in the one-parameter family of two-spheres $`S`$ in the limit as $`r\mathrm{}`$, with $`w`$ arbitrary but fixed.
In the Bondi coordinates our asymptotically flat metric takes the standard form
$$g_{ab}dx^adx^b=UVdw^22Udwdr+\sigma _{ij}(dx^i+W^idw)(dx^j+W^jdw).$$
(5.22)
We assume the following expansions for the various terms in this metric:<sup>13</sup><sup>13</sup>13We are following closely the notation used in Ref. , as well as the spirit of the discussion in their footnote 2. The meaning of the notation $`O_<(r^n)`$ was described following Eq. (5.1). The motivation for this level of generality is that Chruściel et al have recently shown that one can allow ‘polyhomogeneous’ terms of the form $`r^n\mathrm{ln}^mr`$ in these expansions and still have a consistent framework for solving the Bondi-Sachs-type characteristic initial value problem. Allowing only expansions in powers of inverse $`r`$ is tantamount to Sachs’ ‘outgoing radiation condition’ , which they argue is overly restrictive. However, besides making the calculations ‘tighter’ as regards remainder terms, and slightly more general, we would get the same results had we assumed Sachs’ outgoing radiation condition.
$`V`$ $`=`$ $`12mr^1+O_<(r^1),`$ (5.23)
$`U`$ $`=`$ $`1{\displaystyle \frac{1}{2}}(X^2+Y^2)r^2+O_<(r^2),`$ (5.24)
$`W^\theta `$ $`=`$ $`(2X\mathrm{cot}\theta +_\theta X+\mathrm{csc}\theta _\varphi Y)r^2+O_<(r^2),`$
$`W^\varphi `$ $`=`$ $`\mathrm{csc}\theta (2Y\mathrm{cot}\theta +_\theta Y\mathrm{csc}\theta _\varphi X)r^2+O_<(r^2),`$ (5.25)
$`\sigma _{ij}`$ $`=`$ $`r^2\left(\begin{array}{cc}1& 0\\ 0& \mathrm{sin}^2\theta \end{array}\right)+2r\left(\begin{array}{cc}X& Y\mathrm{sin}\theta \\ Y\mathrm{sin}\theta & X\mathrm{sin}^2\theta \end{array}\right)+O_<(r).`$ (5.30)
The function $`V`$ contains the mass aspect, $`m(w,\theta ,\varphi )`$. Observe that the metric on $`S`$ is of the same form given in Eq. (5.1) (with $`Z=X`$ because $`r`$ is an areal radius here), except now $`X(w,\theta ,\varphi )`$ and $`Y(w,\theta ,\varphi )`$ have a significant physical interpretation: they are the real and imaginary parts of Sachs’ complex asymptotic shear $`c=X+iY`$ . Thus the scalar curvature of $`(S,\sigma )`$ will be given by Eq. (5.2), and we can begin our discussion of the IQE at Eq. (5.5). Our first task is to compute the unreferenced shear term, $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)`$.
Inspecting the metric in Eq. (5.22), we choose the following basis of one-forms:
$$e^{}=\frac{1}{2}Udw,e^+=dr+\frac{1}{2}Vdw,e^I=\gamma _i^I(dx^i+W^idw),$$
(5.31)
where indices $`I,J,\mathrm{}`$ take the values 2 and 3, and $`\gamma _i^I`$ is defined by demanding that $`\sigma _{ij}=\delta _{IJ}\gamma _i^I\gamma _j^J`$. A suitable choice for $`\gamma _i^I`$ is given by
$$\gamma _i^I=\left(\begin{array}{cc}r+X& 0\\ 2Y& (rX)\mathrm{sin}\theta \end{array}\right)+O_<(1).$$
(5.32)
In this matrix expression, $`I`$ ($`i`$) is a row (column) index. Let the indices $`A,B,\mathrm{}`$ take values in the set $`\{,+,I\}`$. Then the metric is given by $`g_{ab}=\eta _{AB}e_a^Ae_b^B`$, where
$$\eta _{AB}=\left(\begin{array}{cccc}0& 2& 0& 0\\ 2& 0& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right).$$
(5.33)
This matrix, and its inverse, $`\eta ^{AB}`$, are used to raise and lower the the basis indices. The vector fields dual to the one-forms in Eq. (5.31) are given by $`e_A^a=\eta _{AB}g^{ab}e_b^B`$, or explicitly:
$$e_{}=\frac{2}{U}\left(_w\frac{1}{2}V_rW^ii\right),e_+=_r,e_I=\gamma _I^i_i,$$
(5.34)
where $`\gamma _I^i`$ is defined by $`\gamma _I^i=\delta _{IJ}\sigma ^{ij}\gamma _j^J`$, $`\sigma ^{ij}`$ being the inverse of the matrix $`\sigma _{ij}`$.
Inspection of $`e_I`$ in Eq. (5.34) shows that these vectors are tangent to $`S`$, and so $`e_\pm `$ are two null normals to $`S`$. Since their normalization is such that $`e_+e_{}=2`$, we can set $`\xi _\pm ^a=e_\pm ^a`$, where $`\xi _\pm ^a`$ was previously defined by $`\xi _\pm ^a:=u^a\pm n^a`$ (see Eq. (3.8)). Thus, from Eqs. (2.8) we have the following result, in basis components:
$$l_{IJ}\pm k_{IJ}=e_I^ae_J^b_ae_{\pm b}=\omega _{\pm JI},$$
(5.35)
Working out the required Ricci rotation coefficients<sup>14</sup><sup>14</sup>14See the footnote on page 11, with $`Z=\pm `$. I find
$$l_{IJ}+k_{IJ}=B_{IJ},l_{IJ}k_{IJ}=\frac{2}{U}\left(A_{IJ}\frac{V}{2}B_{IJ}𝒟_{(I}W_{J)}\right),$$
(5.36)
where
$`A_{IJ}`$ $`=`$ $`\gamma _{(I}^i\dot{\gamma }_{J)i}={\displaystyle \frac{1}{r}}\left(\begin{array}{cc}\dot{X}& \dot{Y}\\ \dot{Y}& \dot{X}\end{array}\right)+O_<(r^1),`$ (5.39)
$`B_{IJ}`$ $`=`$ $`\gamma _{(I}^i\gamma _{J)i}^{}={\displaystyle \frac{1}{r}}\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right){\displaystyle \frac{1}{r^2}}\left(\begin{array}{cc}X& Y\\ Y& X\end{array}\right)+O_<(r^2).`$ (5.44)
Here we use an overdot (prime) to denote differentiation with respect to $`w`$ ($`r`$). Observe that $`𝒟_{(I}W_{J)}`$ in Eq. (5.36) is of order $`1/r^2`$. Taking the trace-free part of Eqs. (5.36) gives us the basis components of the shears in the two null directions: $`s_{\pm IJ}=\stackrel{~}{l}_{IJ}\pm \stackrel{~}{k}_{IJ}`$. Explicitly:
$$s_{+IJ}=\frac{1}{r^2}\left(\begin{array}{cc}X& Y\\ Y& X\end{array}\right)+O_<(r^2),s_{IJ}=\frac{2}{r}\left(\begin{array}{cc}\dot{X}& \dot{Y}\\ \dot{Y}& \dot{X}\end{array}\right)+O_<(r^1).$$
(5.45)
Thus the unreferenced shear term in Eq. (5.5) is given by
$$(\stackrel{~}{k}^2\stackrel{~}{l}^2)s_{+IJ}s_{}^{IJ}=\frac{4}{r^3}(X\dot{X}+Y\dot{Y})+O_<(r^3)$$
(5.46)
(cf. Eq (4.9)).
We thus learn that, in contrast to the spatial infinity limit (see Eq. (5.16)), in the null infinity limit the unreferenced shear term is of order $`1/r^3`$, and so does contribute to the IQE. We will argue below that the reference shear term $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$ is also of order $`1/r^3`$, but that it is a total derivative and therefore does not contribute. So as not to interrupt the flow or our discussion, for the moment let us assume this is true, in which case Eq. (5.5) becomes
$$\underset{r\mathrm{}}{lim}\mathrm{IQE}=\underset{r\mathrm{}}{lim}\frac{1}{16\pi }_S𝑑Sr\left[\sigma \sigma R\frac{4}{r^3}(X\dot{X}+Y\dot{Y})\right].$$
(5.47)
Because of the $`c\dot{\overline{c}}+\dot{c}\overline{c}=2(X\dot{X}+Y\dot{Y})`$ term, this result looks like it could very well be the Bondi-Sachs mass . To see that in fact it is, a straightforward calculation of the Riemann tensor of $`g_{ab}`$ projected into $`S`$ gives the following sectional curvature of $`(S,\sigma )`$:
$$\sigma \sigma R=\frac{4}{r^3}(m+X\dot{X}+Y\dot{Y})+O<(r^3).$$
(5.48)
Inserting this result into Eq. (5.47) we see that the shear terms cancel, leaving only the mass aspect, $`m`$:
$$\underset{r\mathrm{}}{lim}\mathrm{IQE}=\frac{1}{4\pi }_S𝑑\mathrm{\Omega }m(w,\theta ,\varphi ).$$
(5.49)
In obtaining this expression, recall that because $`r`$ is an areal radius we can (and did) take $`\sqrt{\sigma }=r^2\mathrm{sin}\theta `$, and so $`dS=r^2d\mathrm{\Omega }`$, where $`d\mathrm{\Omega }=\mathrm{sin}\theta d\theta d\varphi `$ is the measure on the unit round sphere.<sup>15</sup><sup>15</sup>15On this note, it might be helpful to point out an important detail in the calculation of the sectional curvature given in Eq. (5.48). One of the the terms that arises in the calculation is the trace of $`A_{IJ}`$ in Eq. (5.39), which appears to be of order $`O_<(r^1)`$. If this were so it would be problematic. But in fact it is zero, because $`A_I^I=\gamma _I^i\dot{\gamma }_i^I=(1/2)\dot{\sigma }/\sigma =0`$, since $`\sigma =r^4\mathrm{sin}^2\theta `$ does not depend on $`w`$. Thus the future null infinity limit of the IQE is the Bondi-Sachs mass .
Now there is an important lesson to be learned from this result. The unreferenced shear term $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)`$ is solely responsible for producing the all-important $`c\dot{\overline{c}}+\dot{c}\overline{c}`$ term that accounts for the mass loss due to gravitational radiation. Hence this term is necessary under the square root in Eq. (4.2), and so there is no natural way to avoid $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$ in $`\mathrm{IQE}^{\mathrm{ref}}`$, and its attendant embedding problem. Moreover, we learn that these shear terms are not only associated with angular momentum, as I have been stressing, but also encode information about gravitational radiation. We will see precisely the same phenomenon emerge in the small sphere limit in Sec. 6. Furthermore, it is emphasized in Ref. that it is easy to construct, ab initio, an integral expression involving the Riemann tensor (e.g., an integral of $`\sigma \sigma R`$ over $`S`$) that is conserved under certain circumstances. One is thus tempted to interpret such a conserved quantity as an energy. However, such attempts fail to produce, in the null infinity limit, the crucial null-surface-dependent shear terms seen in Eq. (5.47), and it is difficult to see how to modify them in a covariant way to produce these terms . The shear term $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)`$ is precisely such a covariant modification. Moreover, it arises naturally from simply replacing the Brown-York $`k`$ with the boost invariant quantity $`\sqrt{k^2l^2}`$. (Of course a similar observation can be made concerning, say, the Hawking mass , which has the same large sphere limit as in Eq. (5.5), except without the reference shear term.)
These clean results rely on our assumption that the reference shear term $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$ does not contribute to the null infinity limit of the IQE. I claimed above that this is so because it is a total derivative. To prove this would require solving the embedding equations (4.6-4.8), which we know is a very difficult task. However, I will now present a heuristic solution of the Ricci equation that leads to a substantiation of this claim. Moreover, we will see how demanding $`^{\mathrm{ref}}=`$ plays a crucial role in achieving this result, which provides our first bit of indirect but concrete evidence that this condition is required to properly account for angular momentum (and as we now see, also gravitational radiation).
To begin we need to calculate the connection in the normal bundle, $`A_a`$, and then its corresponding curvature, $``$. Proceeding as we did in the spatial infinity case (see Eq. (5.17)) we find that the basis components of $`A_a`$ are given by
$$A_I=\frac{1}{2}\omega _{+I}=\frac{1}{rU}W_I+\frac{1}{2}e_I(\mathrm{ln}U),$$
(5.50)
where $`e_I(\mathrm{ln}U)`$ denotes the derivative of $`\mathrm{ln}U`$ along the vector field $`e_I`$. This term is pure gauge. As for the other term, since we only know $`W_I`$ to leading order we can put $`U=1`$ here—see Eqs. (5.24) and (5.25). Thus, up to a gauge transformation, the connection one-form $`A`$ is just ($`1/r`$ times) the one-form $`W_idx^i`$. And the curvature is thus proportional to the curl of $`W`$:
$$=\frac{2}{r}ϵ^{IJ}𝒟_IW_J=\frac{2}{r}(𝒟_2W_3𝒟_3W_2).$$
(5.51)
Keep in mind that the numerical indices here refer to basis components, not coordinate components. It is easy to see that $``$ is of order $`1/r^3`$. It is interesting to compare $``$ with the scalar curvature $``$, whose form was given in Eq. (5.2). Using the metric in Eq. (5.30) it is not difficult to evaluate the remainder term, $`\mathrm{\Delta }_{}`$. The net result is :
$$=\frac{2}{r^2}+\frac{2}{r}𝒟W+O_<(r^3),$$
(5.52)
Note from Eq. (5.25) that $`W^i`$ is of order $`1/r^2`$, so the term above involving $`W`$ is, indeed, of order $`1/r^3`$, as it should be. Thus we see that $``$ is associated with the divergence of $`W`$, and $``$ with its curl. This is an explicit example of a point made earlier, namely that both curvatures are on the same geometrical footing: To capture the two pieces of information in $`W`$—its divergence and its curl—requires precisely both $``$ and $``$.
Let us now turn to the null shears of $`S`$ as they appear in the physical spacetime, Eq. (5.45). The form of these shears suggests we make the following ansatz for the null shears of $`S`$ in the reference (Minkowski) spacetime:
$$s_{+IJ}^{\mathrm{ref}}=\frac{1}{r^2}\left(\begin{array}{cc}\alpha & \beta \\ \beta & \alpha \end{array}\right)+O_<(r^2),s_{IJ}^{\mathrm{ref}}=\frac{1}{r}\left(\begin{array}{cc}\gamma & \delta \\ \delta & \gamma \end{array}\right)+O_<(r^1).$$
(5.53)
Observe that we might expect the pair $`(\alpha ,\beta )`$ to play a role distinct from the pair $`(\gamma ,\delta )`$. Comparing Eqs. (5.53) and (5.45) suggests that $`\alpha `$ and $`\beta `$ will be like $`X`$ and $`Y`$ in that they have something to do with the intrinsic geometry of $`(S,\sigma )`$. The ‘more important’ terms will be $`\gamma `$ and $`\delta `$, because they occur at the dominant power of inverse $`r`$. Also, we expect them to be related to the extrinsic geometry of $`S`$, since their counterparts, $`\dot{X}`$ and $`\dot{Y}`$, measure how $`\sigma _{ij}`$ changes as a function of the retarded time $`w`$—they are the two ‘news’ functions .
With these observations in mind, we will now ‘solve’ the Ricci embedding equation, (4.8). We first compute (in basis components)
$$\frac{1}{2}ϵ_J^I[s_+^{\mathrm{ref}},s_{}^{\mathrm{ref}}]_I^J=\frac{2}{r^3}(\alpha \delta \beta \gamma )=\frac{2}{r^3}det\left(\begin{array}{cc}\alpha & \beta \\ \gamma & \delta \end{array}\right).$$
(5.54)
Next we impose the condition $`^{\mathrm{ref}}=`$, and observe that $``$ in Eq. (5.51) can also be expressed as a ‘determinant’, i.e.,
$$^{\mathrm{ref}}==\frac{2}{r}det\left(\begin{array}{cc}𝒟_2& 𝒟_3\\ W_2& W_3\end{array}\right)$$
(5.55)
The Ricci embedding equation instructs us to equate the two previous determinant expressions. One solution is to make the identifications $`\alpha r𝒟_2`$ and $`\beta r𝒟_3`$ (which are consistent with our expectation that $`\alpha `$ and $`\beta `$ be associated with intrinsic geometry), together with $`\gamma rW_2`$ and $`\delta rW_3`$ (which are consistent with $`\gamma `$ and $`\delta `$ being associated with extrinsic geometry, since $`W`$ is proportional to the connection in the normal bundle—a measure of extrinsic geometry). Notice that this means $`s_{+IJ}^{\mathrm{ref}}`$ is a derivative operator! To make this more palatable one may go to a ‘Fourier transform space’, where the derivative operators $`𝒟_I`$ become momenta, $`𝒦_I`$. Accepting these heuristic identifications, and recalling Eq. (4.9), the key point now is to observe that
$$(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}s_{+J}^{\mathrm{ref}I}s_I^{\mathrm{ref}J}=\frac{2}{r^3}(\alpha \gamma +\beta \delta )=\frac{2}{r}(𝒟_2W_2+𝒟_3W_3)=\frac{2}{r}𝒟W,$$
(5.56)
the result we desired: being a total derivative, $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$ does not contribute to the IQE.
Now recall that the purpose of solving the embedding equations is to establish a relationship between $`^{\mathrm{ref}}`$ and $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$. It is via this relationship that we envision the angular momentum information in $`^{\mathrm{ref}}`$ ($`=`$) to enter $`\mathrm{IQE}^{\mathrm{ref}}`$, which sees only $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$. The previous result implies the following relationship: $`^{\mathrm{ref}}`$ is the curl of $`W`$, and $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$ is the divergence of $`W`$. Hence, no relationship! Is this a problem? Certainly not. We require the relationship in question only when different isometric embeddings of the same $`(S,\sigma )`$ result in a different $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$, making the IQE ambiguous. Such a relationship can be used to resolve this ambiguity by setting $`^{\mathrm{ref}}=`$. But accepting the argument that $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$ is always a total derivative, any such ambiguity would be of no consequence here. Furthermore, in this simple case there is no need for $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$ to carry any information about angular momentum (or gravitational radiation) coming from $``$, because all of the relevant information is already carried in the unreferenced shear term, $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)`$. This is not to say that the condition $`^{\mathrm{ref}}=`$ is not important here. According to our heuristic argument, it is precisely this condition which ensures that $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$ is a total derivative, and so carries no information.
The null infinity limit is a simple case which only minimally exercises the consequences of the condition $`^{\mathrm{ref}}=`$. For $`(S,\sigma )`$ a finite two-surface the situation is much richer: it is highly nonlinear, since the square roots do not disappear, and almost certainly requires effectively a one-to-one relationship between $`^{\mathrm{ref}}`$ and $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$ to remove the embedding ambiguity inherent in the IQE.
## 6 The small sphere limit of the IQE
Having considered the large sphere case, we now turn our attention to evaluating the IQE when $`(S,\sigma )`$ is a small sphere. The large and small sphere limits are similar in that in both cases $`S`$ approaches an asymptotically flat region of $`(M,g)`$. In the latter case, the asymptotically flat region is the infinitesimal neighborhood of a generic spacetime point $`pM`$, which is the ‘center’ of our shrinking sphere. For simplicity we will suppose that $`(S,\sigma )`$ is asymptotically round. Another feature in common with the large sphere limit is that in this codimension-two setting, the limit can be approached from different directions, either spatial or null. More precisely, fix a set of Riemann normal coordinates $`(t,x^i)`$ about the point $`p`$, set $`r^2:=\delta _{ij}x^ix^j`$, and define $`S_{}`$ by the condition $`(r,t)=(r_{},\alpha r_{})`$, where $`\alpha `$ is a direction parameter. Then consider the limit of $`S_{}`$ as $`r_{}0`$. As before, we will henceforth omit the subscript ‘$``$’. The case $`\alpha =0`$ is a spatial limit, since then $`S`$ always lies entirely in the $`t=0`$ spacelike three-surface containing $`p`$. $`\alpha =\pm 1`$ is the null limit, in which $`S`$ lies in the future/past light cone of the point $`p`$. The latter case was considered by Horowitz and Schmidt in their classic work on the small sphere limit of the Hawking mass. Brown, Lau, and York also consider this same limit of the Brown-York quasilocal energy. We will be borrowing some results from these two references.
Explicitly, for a given value of the parameter $`r`$, $`S`$ is defined as a submanifold of $`(M,g)`$ by embedding a topological two-sphere with coordinates $`\theta `$ and $`\varphi `$ into the Riemann normal coordinate system, as follows: $`t=\alpha r`$, $`x^1=r\mathrm{sin}\theta \mathrm{cos}\varphi `$, $`x^2=r\mathrm{sin}\theta \mathrm{sin}\varphi `$, $`x^3=r\mathrm{cos}\theta `$. Since (with $`t=\alpha r`$) the deviation from the flat metric in Riemann normal coordinates is $`O(r^2)`$, the induced metric $`\sigma _{ab}`$ on $`S`$ will differ from that of the round sphere to this same order, and so the scalar curvature of $`(S,\sigma )`$ will have an expansion in $`r`$ of the form
$$=\frac{2}{r^2}+^{(0)}+r^{(1)}+r^2^{(2)}+O(r^3),$$
(6.1)
where each of the coefficients $`^{(n)}`$ is a function of $`\theta `$, $`\varphi `$, and the parameter $`\alpha `$. To evaluate the IQE we will also need similar expansions for the other quantities appearing in Eqs. (4.2) and (4.3). These are written as follows:
$`\sigma \sigma R`$ $`=`$ $`\sigma \sigma R^{(0)}+r\sigma \sigma R^{(1)}+r^2\sigma \sigma R^{(2)}+O(r^3),`$ (6.2)
$`(\stackrel{~}{k}^2\stackrel{~}{l}^2)`$ $`=`$ $`r^2(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{(2)}+O(r^3),`$ (6.3)
$`(\stackrel{~}{k}^2\stackrel{~}{l}^2)`$ $`=`$ $`r^2(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{(2)\mathrm{ref}}+O(r^3),`$ (6.4)
where each of the coefficients on the right hand side is similarly a function of $`\theta `$, $`\varphi `$, and $`\alpha `$. Since the appropriate reference spacetime in this case is Minkowski space, we have $`\sigma \sigma R^{\mathrm{ref}}=0`$. Substituting these expansions into Eq. (4.2) we find that in the small sphere limit the unreferenced IQE behaves as
$`\mathrm{IQE}^{\mathrm{unref}}={\displaystyle \frac{1}{8\pi }}{\displaystyle _S}`$ $`dS{\displaystyle \frac{2}{r}}\{1+{\displaystyle \frac{r^2}{4}}[(^{(0)}\sigma \sigma R^{(0)})+r(^{(1)}\sigma \sigma R^{(1)})`$
$`+r^2(^{(2)}\sigma \sigma R^{(2)}+(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{(2)}{\displaystyle \frac{1}{8}}(^{(0)}\sigma \sigma R^{(0)})^2)]+O(r^5)\}.`$
Similarly, the reference IQE behaves as
$`\mathrm{IQE}^{\mathrm{ref}}={\displaystyle \frac{1}{8\pi }}{\displaystyle _S}𝑑S{\displaystyle \frac{2}{r}}`$ $`\{1+{\displaystyle \frac{r^2}{4}}[^{(0)}+r^{(1)}`$ (6.6)
$`+r^2(^{(2)}+(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{(2)\mathrm{ref}}{\displaystyle \frac{1}{8}}(^{(0)})^2)]+O(r^5)\}.`$
Notice that, unlike in the large sphere limit, neither the unreferenced nor reference energies diverge as $`r0`$. Nevertheless, the reference subtraction procedure is still necessary to eliminate the leading $`O(r)`$ term in $`\mathrm{IQE}^{\mathrm{unref}}`$, which has nothing to do with energy. Thus, forming the difference of the previous two expressions we find that the small sphere behavior of the (referenced) IQE is given by
$$\mathrm{IQE}=\frac{1}{16\pi }_S𝑑Sr\left[\sigma \sigma R(\stackrel{~}{k}^2\stackrel{~}{l}^2)+(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}+\frac{1}{8}r^2\sigma \sigma R^{(0)}(\sigma \sigma R^{(0)}2^{(0)})+O(r^3)\right].$$
(6.7)
The sectional curvature and shear terms have been resummed according to Eqs. (6.2-6.4), and the resulting expression is valid to the order indicated. Notice that, as in the large sphere limit, the scalar curvature $``$ dominates the other terms under the square root, allowing us to expand the radical about $`\sqrt{4/r^2}=2/r`$.<sup>16</sup><sup>16</sup>16As remarked in the footnote on page 10, if the sphere is not asymptotically round the shear terms will contribute at order $`1/r^2`$, and the $`2/r`$ term outside the braces in Eqs. (6) and (6.6) will be modified accordingly. We will not consider this more complicated case here. And after the reference subtraction is performed what remains again is the sectional curvature of $`(S,\sigma )`$ as the dominant term contributing to the energy. Comparing Eqs. (6.7) and (5.5) we observe that both the small and large sphere limits of the IQE are very nearly formally identical.
We will split the IQE into three pieces, each to be discussed separately: $`\mathrm{IQE}=\mathrm{IQE}_1+\mathrm{IQE}_2+\mathrm{IQE}_3+O(r^6)`$, where
$`\mathrm{IQE}_1`$ $`=`$ $`{\displaystyle \frac{1}{16\pi }}{\displaystyle _S}𝑑Sr\left[\sigma \sigma R(\stackrel{~}{k}^2\stackrel{~}{l}^2)\right],`$ (6.8)
$`\mathrm{IQE}_2`$ $`=`$ $`{\displaystyle \frac{1}{128\pi }}{\displaystyle _S}𝑑Sr^3\left[\sigma \sigma R^{(0)}(\sigma \sigma R^{(0)}2^{(0)})\right],`$ (6.9)
$`\mathrm{IQE}_3`$ $`=`$ $`{\displaystyle \frac{1}{16\pi }}{\displaystyle _S}𝑑Sr\left[(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}\right].`$ (6.10)
We begin with $`\mathrm{IQE}_1`$, and show that this piece is essentially the Hawking mass . To see this, we combine Eqs. (4.1) and (2.10) to get
$$\sigma \sigma R(\stackrel{~}{k}^2\stackrel{~}{l}^2)=\frac{1}{2}(k^2l^2)=2HH.$$
(6.11)
Replacing the integrand of $`\mathrm{IQE}_1`$ with the last expression, and using the Gauss-Bonnet theorem to integrate the $``$ term, we find
$$\mathrm{IQE}_1=\frac{1}{4\pi }\sqrt{\frac{A}{4\pi }}\left[2\pi \frac{1}{2}_S𝑑SHH\right],$$
(6.12)
where we pulled $`r`$ outside the integral and replaced it with $`\sqrt{A/(4\pi )}`$, where $`A`$ is the area of $`(S,\sigma )`$. This form of $`\mathrm{IQE}_1`$ is precisely the expression of the Hawking mass given in Ref. . (Our $`H^c`$ in Eq. (2.10) is their $`N^c/2`$, and the sign of their metric signature is opposite to ours.) Comparing Eq. (6.12) with Eqs. (3.7) and (3.9) we observe that, while the unreferenced IQE involves the mean curvature itself, $`\sqrt{HH}`$, the Hawking mass is constructed from the square of the mean curvature. As mentioned above, the square root in $`\sqrt{HH}`$ effectively disappears in the small (and large) sphere limits due to the presence of the dominant scalar curvature term, and consequently the leading order contribution to the IQE reduces to essentially the Hawking mass.
There are two subtleties worth mentioning: (i) Replacing $`r`$ with $`\sqrt{A/(4\pi )}`$ is in general not valid because it requires that $`r`$ be an areal radius which, in general, it is not. However, it certainly is to lowest order in $`r`$, which will be sufficient for our purposes here. But to higher order, $`\mathrm{IQE}_1`$ and the Hawking mass will in general give different results. (ii) It is well known that the Hawking mass runs into difficulties when $`(S,\sigma )`$ is not a round sphere , a problem that was addressed by Hayward in Ref. . It might be that this problem is a result of having to insert by hand the factor $`\sqrt{A/(4\pi )}`$ outside the integral, versus having $`r`$ ‘inside the integral’ generated automatically by $`\sqrt{2/}`$. A related issue was discussed at the end of Sec. 5.1 in connection with Hayward’s definition of quasilocal energy.
The connection between $`\mathrm{IQE}_1`$ and the Hawking mass allows us to borrow some results from Ref. , which are calculated for the null limit case ($`\alpha =1`$). When matter is present Horowitz and Schmidt find that, to lowest order in $`r`$, the Hawking mass is
$$\mathrm{IQE}_1=\left(\frac{4}{3}\pi r^3\right)T_{ab}^{\mathrm{mat}}u^au^b|_p+O(r^4).$$
(6.13)
Here $`T_{ab}^{\mathrm{mat}}`$ is the energy-momentum tensor of matter, and the expression is to be evaluated at the point $`p`$, where the unit timelike vector $`u^a`$ is just $`(/t)^a`$ in our Riemann normal coordinates. This is a standard result in the literature on quasilocal energy , and a very significant one. As emphasized in the Introduction, the quasilocal idea asserts that the time-time component of the energy-momentum tensor of matter a priori has nothing to do with energy. It is only from the small sphere limit of the quasilocal energy that we learn this quantity is an energy volume density, i.e., multiplying it by the volume factor $`4\pi r^3/3`$ gives the energy in an infinitesimal sphere of proper radius $`r`$. However, integrating this energy volume density over a finite volume to determine the total energy inside is not, in general, valid unless one wishes to ignore gravitational effects which, as we will see in moment, come at higher order in $`r`$.<sup>17</sup><sup>17</sup>17A well established example of this phenomenon is the Tolman density, which integrates to the Komar mass, and is defined in the special case that the spacetime is stationary and asymptotically flat . It is noteworthy that it is not $`T_{ab}`$ that appears in the Tolman density, but rather the combination $`T_{ab}(1/2)Tg_{ab}`$. The extra term involving the trace of $`T_{ab}`$ is associated with gravitational effects—see Ref. , and problems 4 and 5 in Chapter 11 of Ref. . It is in this sense that the quasilocal idea implies that even energy due to matter is not localizable in the context of general relativity.
Now let us assume the spacetime is vacuum in the neighborhood of $`p`$. Then the leading order contribution to the Hawking mass is
$$\mathrm{IQE}_1=\frac{1}{90}r^5T_{abcd}u^au^bu^cu^d|_p+O(r^6),$$
(6.14)
where $`T_{abcd}`$ is the Bel-Robinson tensor . Thus gravitational energy begins to appear at $`O(r^5)`$. This same result is obtained for the Brown-York quasilocal energy for a suitable choice of reference embedding . However, this is not a universal result in the literature . For example, Hayward’s quasilocal mass gives a similar result as above, but with the numerical factor $`1/90`$ replaced with $`2/45`$ . Given that gravitational energy is such a difficult problem it is not surprising that a consensus has not yet been reached.
We now turn our attention to the second contribution to the IQE, namely $`\mathrm{IQE}_2`$ given in Eq. (6.9). This quantity represents a deviation from the Hawking mass due to the fact that the IQE is roughly the square root of the former. Actually, it is clearer to compare with the Brown-York quasilocal energy, since in constructing the IQE we simply replaced the Brown-York $`k`$ with $`\sqrt{k^2l^2}`$. In the context of our generalization given in Eq. (3.4), we therefore have the heuristic comparison:
$$m=\sqrt{E^2\stackrel{}{p}^{\mathrm{\hspace{0.33em}2}}}=E\frac{\stackrel{}{p}^{\mathrm{\hspace{0.33em}2}}}{2E}\mathrm{}\frac{1}{8\pi }\sqrt{k^2l^2}=\frac{1}{8\pi }\left(k\frac{l^2}{2k}\mathrm{}\right).$$
(6.15)
So $`\mathrm{IQE}_2`$ might be thought of as the analogue of the term $`l^2/(2k)`$, and as such would be expected to reduce the magnitude of the IQE from the result given in Eq. (6.14).
In order to calculate $`\mathrm{IQE}_2`$ we need to evaluate the quantities $`^{(0)}`$ and $`\sigma \sigma R^{(0)}`$ in Eq. (6.9). To do so we appeal to the Gauss equation, (4.1). Up to zeroth order in $`r`$, this equation reads
$$\sigma \sigma R^{(0)}=\frac{2}{r^2}+^{(0)}\frac{1}{2}(k^2l^2)+O(r),$$
(6.16)
where we made use of Eq. (6.1). We will show later that
$$(k^2l^2)=\frac{4}{r^2}+\frac{4}{3}(1+2\alpha ^2)E_{ab}n^an^b+O(r),$$
(6.17)
where $`\alpha `$ is the direction parameter introduced at the beginning of this section. $`E_{ab}n^an^b`$ is the radial-radial component of the electric part of the Weyl tensor, which we saw earlier in Eq. (5.21). This quantity ($`E_{ab}n^an^b`$) is to be evaluated at the point $`p`$, where in our Riemann normal coordinates the radial unit vector $`n^a`$ has only spatial components, given by $`n^i=x^i/r`$. We also need $`dS`$ to lowest order, which is just $`r^2d\mathrm{\Omega }`$, $`d\mathrm{\Omega }`$ being the measure on the unit sphere. Putting these results together we have
$$\mathrm{IQE}_2=\frac{1}{96\pi }(5+4\alpha ^2)r^5E_{ij}E_{kl}𝑑\mathrm{\Omega }n^in^jn^kn^l.$$
(6.18)
By symmetry, the integral over the product of radial vectors must be proportional to $`\delta ^{(ij}\delta ^{kl)}`$. Transvecting both this term and the integral in question with $`\delta _{ij}\delta _{kl}`$, we easily obtain that the proportionality constant is $`4\pi /5`$. Using the fact that $`E_{ab}`$ is symmetric, trace-free, and orthogonal to $`u^a`$, we find that
$$\mathrm{IQE}_2=\frac{1}{180}(5+4\alpha ^2)r^5E_{ab}E^{ab}.$$
(6.19)
So $`\mathrm{IQE}_2`$ is negative, as expected, and this negative contribution is to be added to $`\mathrm{IQE}_1`$ in Eq. (6.14). Of course we can only consider the $`\alpha =1`$ case, since this is the case assumed in Eq. (6.14). Recall that the time component of the Bel-Robinson tensor can be expressed in terms of the electric and magnetic parts of the Weyl tensor : $`T_{abcd}u^au^bu^cu^d=E_{ab}E^{ab}+B_{ab}B^{ab}`$, so $`\mathrm{IQE}_1`$ is nonnegative. Inspection of Eq. (6.19) shows that adding to $`\mathrm{IQE}_1`$ the $`\alpha =1`$ value of $`\mathrm{IQE}_2`$ makes the energy have indefinite sign. It is positive (negative) if the magnetic (electric) part dominates. This seems like a strange result, but it is only an intermediate result. We have not yet considered the last contribution, $`\mathrm{IQE}_3`$, involving the reference shear term. But unfortunately at present I do not know how to solve the embedding equations to determine this term.
Now one can construct a heuristic argument much like the one given at the end of Sec. 5.2, which suggests that $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$ is a total derivative, and so does not contribute. However, the argument is much less believable in this case. In contrast to Eq. (5.53) it turns out that, because $``$ is $`O(1)`$ in $`r`$ (as we shall see later), we must expand the reference shears $`s_\pm ^{\mathrm{ref}}`$ over three orders of magnitude in $`r`$, from $`O(1)`$ to $`O(r^2)`$. One might trust a heuristic argument working to leading order, but believing the higher order corrections is less palatable. In short, I do not know what energy prediction the IQE gives at $`O(r^5)`$, and until a solution to the embedding equations is found there is no sense in speculating.
However, before leaving this section I will provide an intriguing interpretation of how a definition of quasilocal energy such as the Hawking mass (or the IQE) provides a measure of the gravitational energy contained inside a small sphere.
In the $`\alpha =1`$ null limit case, the lowest order contribution to the gravitational energy, namely $`(1/90)r^5T_{abcd}u^au^bu^cu^d|_p`$ in Eq. (6.14), originates in the $`O(r^2)`$ terms inside the brackets of the integrand of $`\mathrm{IQE}_1`$ in Eq. (6.8). In the terminology of Eqs. (6.2) and (6.3), this means we are interested in the coefficients $`\sigma \sigma R^{(2)}`$ and $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{(2)}`$. Inspecting the Appendix of Ref. reveals that these two coefficients differ only by a numerical factor. They are both proportional to $`\mathrm{\Psi }_0\overline{\mathrm{\Psi }}_0|_p`$ (in Newman-Penrose notation), and the two numerical factors conspire to produce the $`1/90`$ factor in the final result. Thus, to understand how the integrand of $`\mathrm{IQE}_1`$ encodes information about gravitational energy it suffices to study the shear term, $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{(2)}`$. We will now compute this term for arbitrary $`\alpha [1,1]`$ to see how it behaves in both the spatial and null limit cases of the small sphere.
Denoting our Riemann normal coordinates $`(t,x^i)`$ collectively as $`x^a`$, the metric in these coordinates take the form
$$g_{ab}=\eta _{ab}\frac{1}{3}J_{abcd}x^cx^d+O(x^3),$$
(6.20)
where $`J_{abcd}=(R_{acbd}+R_{adbc})/2`$ is the Jacobi curvature tensor . We first construct a pair of mutually orthogonal unit normal vector fields $`u^a`$ and $`n^a`$ on $`S`$, with $`u^a`$ normal to the $`t=\mathrm{constant}`$ surface passing through $`S`$. These are given by
$`u^0`$ $`=`$ $`{\displaystyle \frac{1}{N}},u^i={\displaystyle \frac{1}{3}}r^2\beta ^{i0}+O(r^3),u_0=N,u_i=0,`$ (6.21)
$`n^0`$ $`=`$ $`0,n^i=\rho \left[{\displaystyle \frac{x^i}{r}}+{\displaystyle \frac{1}{3}}r^2\beta ^{ij}{\displaystyle \frac{x_j}{r}}+O(r^3)\right],n_0={\displaystyle \frac{1}{3}}r^2\beta _{0j}{\displaystyle \frac{x^j}{r}}+O(r^3),n_i=\rho {\displaystyle \frac{x_i}{r}},`$ (6.22)
where
$`N`$ $`=`$ $`1+{\displaystyle \frac{1}{6}}r^2\beta _{00}+O(r^3),`$ (6.23)
$`\rho `$ $`=`$ $`1{\displaystyle \frac{1}{6}}r^2\beta _{ij}{\displaystyle \frac{x^ix^j}{r^2}}+O(r^3),`$ (6.24)
$`\beta _{ab}`$ $`=`$ $`\alpha ^2J_{ab00}+2\alpha J_{ab0j}{\displaystyle \frac{x^j}{r}}+J_{abij}{\displaystyle \frac{x^ix^j}{r^2}}.`$ (6.25)
Since the Jacobi tensor in Eq (6.20) is evaluated at the coordinate origin $`p`$, its indices, and thus those of $`\beta _{ab}`$, are raised and lowered with the flat spacetime metric $`\eta _{ab}=\eta ^{ab}=\mathrm{diagonal}(1,1,1,1)`$. Similarly, $`x_i:=\delta _{ij}x^j`$.
Now define on $`S`$ a pair of mutually orthogonal unit tangent vector fields $`e_I^a`$, where indices $`I,J,\mathrm{}`$ take the values 2 and 3. The set $`\{e_0^a:=u^a,e_1^a:=n^a,e_I^a\}`$ thus comprises an orthonormal basis adapted to $`S`$. Let basis indices $`A,B,\mathrm{}`$ run from 0 to 3, and $`\alpha ,\beta ,\mathrm{}`$ from 1 to 3. Beginning with this setup it is straightforward to compute the basis components of the extrinsic curvatures defined in Eqs. (2.8). I find:
$`l_{IJ}`$ $`=`$ $`e_I^ae_J^b_au_b={\displaystyle \frac{2}{3}}r[\alpha J_{00IJ}+J_{01IJ}]+O(r^2),`$ (6.26)
$`k_{IJ}`$ $`=`$ $`e_I^ae_J^b_an_b={\displaystyle \frac{1}{r}}\delta _{IJ}{\displaystyle \frac{1}{3}}r\left[J_{11IJ}\alpha ^2\left(J_{00IJ}{\displaystyle \frac{1}{2}}J_{0011}\delta _{IJ}\right)\right]+O(r^2).`$ (6.27)
In these equations a quantity such as $`J_{01IJ}`$ means $`[e_0^ae_1^be_I^ce_J^dJ_{abcd}]|_p`$, which is a function of only the angles $`\theta `$ and $`\varphi `$ on $`S`$.
Since we are interested in purely gravitational energy we shall restrict ourselves to the vacuum case. In our basis components the electric and magnetic parts of the Weyl tensor are defined by :
$$E_{\alpha \beta }=C_{0\alpha 0\beta }\mathrm{and}B_{\alpha \beta }=C_{0\alpha 0\beta },$$
(6.28)
where $`C_{ABCD}=(1/2)ϵ_{AB}^{EF}C_{EFCD}`$. These are symmetric trace-free three-dimensional tensors associated with $`t=\mathrm{constant}`$ spacelike hypersurfaces. As such, each has five independent components, which together comprise the ten independent components of the Weyl tensor. In terms of these fields, the components of the Jacobi curvature tensor relevant to Eqs. (6.26-6.27) read:
$$J_{0011}=E_{11},J_{00IJ}=E_{IJ},J_{01IJ}=\stackrel{~}{B}_{IJ},J_{11IJ}=E_{IJ}E_{11}\delta _{IJ}.$$
(6.29)
Here $`\stackrel{~}{B}_{IJ}`$ is the trace-free part of $`B_{IJ}`$, and $`\stackrel{~}{B}_{IJ}=ϵ_I^K\stackrel{~}{B}_{KJ}`$ is its dual in $`(S,\sigma )`$, which is also trace-free. The trace of $`E_{IJ}`$ is $`\delta ^{IJ}E_{IJ}=\delta ^{\alpha \beta }E_{\alpha \beta }E_{11}=E_{11}`$, since $`E_{\alpha \beta }`$ is trace-free. Thus, the trace of the extrinsic curvatures in Eqs. (6.26-6.27) is found to be
$`l`$ $`=`$ $`{\displaystyle \frac{2}{3}}r\alpha E_{11}+O(r^2),`$ (6.30)
$`k`$ $`=`$ $`{\displaystyle \frac{2}{r}}+{\displaystyle \frac{1}{3}}r(1+2\alpha ^2)E_{11}+O(r^2).`$ (6.31)
Squaring these and forming their difference leads to Eq. (6.17) written earlier. The trace-free parts are
$`\stackrel{~}{l}_{IJ}`$ $`=`$ $`{\displaystyle \frac{2}{3}}r(\alpha \stackrel{~}{E}_{IJ}+\stackrel{~}{B}_{IJ})+O(r^2),`$ (6.32)
$`\stackrel{~}{k}_{IJ}`$ $`=`$ $`{\displaystyle \frac{1}{3}}r(1\alpha ^2)\stackrel{~}{E}_{IJ}+O(r^2).`$ (6.33)
Now $`\stackrel{~}{E}_{IJ}`$ and $`\stackrel{~}{B}_{IJ}`$ each have two independent components, and it is useful to introduce the notation
$`\stackrel{~}{E}_{IJ}`$ $`=`$ $`\left(\begin{array}{cc}\frac{1}{2}(E_{22}E_{33})& E_{23}\\ E_{23}& \frac{1}{2}(E_{22}E_{33})\end{array}\right)=:\left(\begin{array}{cc}e^2& e^3\\ e^3& e^2\end{array}\right),`$ (6.38)
$`\stackrel{~}{B}_{IJ}`$ $`=`$ $`\left(\begin{array}{cc}\frac{1}{2}(B_{22}B_{33})& B_{23}\\ B_{23}& \frac{1}{2}(B_{22}B_{33})\end{array}\right)=:\left(\begin{array}{cc}b^2& b^3\\ b^3& b^2\end{array}\right),`$ (6.43)
and thus define a pair of two-vectors $`\stackrel{}{e}:=e^Ie_I^a`$ and $`\stackrel{}{b}:=b^Ie_I^a`$ tangent to $`S`$. Since $`\stackrel{}{e}`$ and $`\stackrel{}{b}`$ represent the pullbacks to $`S`$ of $`E_{\alpha \beta }`$ and $`B_{\alpha \beta }`$, respectively, they are to be thought of as ‘electric and magnetic fields’ induced on $`S`$ by the Weyl curvature that $`S`$ is embedded in. It is easy to see that under a rotation of the basis vectors $`e_I^a`$ through an angle $`\gamma `$, the components of $`\stackrel{}{e}`$ and $`\stackrel{}{b}`$ rotate through an angle $`2\gamma `$, so $`\stackrel{}{e}`$ and $`\stackrel{}{b}`$ are not true (spin one) vectors, but rather spin two objects, as one would expect.
Thus we arrive at the results we are interested in. To lowest order in $`r`$ we have
$`\stackrel{~}{l}^2`$ $`=`$ $`{\displaystyle \frac{4}{9}}r^2(\alpha \stackrel{~}{E}+\stackrel{~}{B})^2={\displaystyle \frac{8}{9}}r^2(\alpha ^2\stackrel{}{e}\stackrel{}{e}+\stackrel{}{b}\stackrel{}{b}2\alpha \stackrel{}{e}\times \stackrel{}{b}),`$ (6.44)
$`\stackrel{~}{k}^2`$ $`=`$ $`{\displaystyle \frac{1}{9}}r^2(1\alpha ^2)^2\stackrel{~}{E}^2={\displaystyle \frac{2}{9}}r^2(1\alpha ^2)^2\stackrel{}{e}\stackrel{}{e}.`$ (6.45)
The difference of these two gives the $`O(r^2)`$ piece of the unreferenced shear term appearing in the integrand of $`\mathrm{IQE}_1`$ in Eq. (6.8). Notice that it is the appearance of $`\stackrel{~}{B}_{IJ}`$ (rather than $`\stackrel{~}{B}_{IJ}`$) that gives rise to the cross product term $`\stackrel{}{e}\times \stackrel{}{b}=e^2b^3e^3b^2`$ in Eq. (6.44).
We first consider the case $`\alpha =1`$, in which $`S`$ lies in the future light cone of the point $`p`$. Then $`\stackrel{~}{k}^2=0`$ and so the shear term $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)`$ is proportional to $`r^2(\stackrel{~}{E}+\stackrel{~}{B})^2`$. As a quick check, it is easy to verify that $`(\stackrel{~}{E}+\stackrel{~}{B})^2`$, in turn, is proportional to $`\mathrm{\Psi }_0\overline{\mathrm{\Psi }}_0|_p`$, in agreement with Ref. . I mentioned above that in this case the $`\sigma \sigma R`$ term in Eq. (6.8) also contributes a term proportional to $`r^2\mathrm{\Psi }_0\overline{\mathrm{\Psi }}_0|_p`$ . Putting in the numerical factors I find that
$$\mathrm{IQE}_1=_S𝑑S\frac{r^3}{9}\left[\frac{1}{8\pi }(\stackrel{}{e}\stackrel{}{e}+\stackrel{}{b}\stackrel{}{b})\frac{1}{4\pi }\stackrel{}{e}\times \stackrel{}{b}\right]+O(r^6).$$
(6.46)
Now $`(\stackrel{}{e}\stackrel{}{e}+\stackrel{}{b}\stackrel{}{b})/(8\pi )`$ looks like the energy surface density of the ‘electromagnetic field’, but we must be careful about its dimension. $`(\stackrel{}{e}\stackrel{}{e}+\stackrel{}{b}\stackrel{}{b})/(8\pi )`$ has dimension $`L^4`$, where $`L`$ means ‘length’, which is not correct. However, the additional factor of $`r^3/9`$ in the integrand suggests that it is really $`:=r^3(\stackrel{}{e}\stackrel{}{e}+\stackrel{}{b}\stackrel{}{b})/(72\pi )`$ that is the proper energy density. $``$ has dimension $`L^1`$, consistent with it being interpreted as the ‘electromagnetic’ energy per unit area of $`S`$. Besides giving the right dimension, the additional $`r^3`$ factor forces $`_S𝑑S`$ to go to zero as $`r^5`$, consistent with the fact that there can be no gravitational energy at order $`r^3`$. We interpret $`_S𝑑S`$ as the total ‘electromagnetic’ energy that was in $`S`$ at $`t=0`$, or equivalently, the total gravitational energy that was in the small volume spanning $`S`$ at $`t=0`$.
To further justify this interpretation we now turn to the radiation term in Eq. (6.46). Clearly $`\stackrel{}{e}\times \stackrel{}{b}/(4\pi )`$ might be thought of as the gravitational analogue of the electromagnetic Poynting flux, directed radially outward from $`S`$. But again, its dimension is wrong. Of course the factor of $`r^3/9`$ will fix this problem, as before, but the situation is more interesting this time. Multiplying by $`r^2/9`$ we get the proper Poynting flux, $`𝒫:=r^2\stackrel{}{e}\times \stackrel{}{b}/(36\pi )`$. $`𝒫`$ has dimension $`L^2`$, consistent with interpreting it as the ‘electromagnetic’ energy per unit time per unit area. So $`_S𝑑S𝒫`$ gives the ‘electromagnetic’ energy per unit time radiating from (or through) the surface $`S`$. The factor of $`r^2`$ indicates that the efficiency of a small volume to radiate gravitationally grows in proportion to its surface area, in analogy with an electromagnetic antenna. But there is one more factor of $`r`$, which one might imagine is the $`r`$ outside the brackets in Eq. (6.8), i.e., $`\sqrt{2/}`$. This distinction between $`r`$s is suggested by the close analogy between $`𝒫`$ and the shear term responsible for radiation in the null infinity limit—see Eqs. (5.46-5.47). This additional factor of $`r`$ has the interpretation of a time lapse, i.e., $`r_S𝑑S𝒫`$ is the amount of ‘electromagnetic’ energy radiated from $`S`$ between time $`t=0`$ and $`t=r`$. Thus, the following picture has emerged regarding Eq. (6.46). The ‘electromagnetic’ energy in $`S`$ (or equivalently, the gravitational energy in the volume spanning $`S`$) at time $`t=r`$ is the energy at $`t=0`$ minus the amount of energy radiated during this time interval. (Keep in mind that $`\stackrel{}{e}`$ and $`\stackrel{}{b}`$ are evaluated at $`p`$, and hence at $`t=0`$.)
The case $`\alpha =1`$ is similar, except now $`S`$ lies in the past light cone of the point $`p`$. Inspection of Eq. (6.44) reveals that the radiation term in Eq. (6.46) now appears with the opposite sign. The fact that this sign change comes out correctly is reassurance that our picture is correct: The energy at time $`t=r`$ is the energy at $`t=0`$ plus the amount of energy that will be radiated from the sphere during the interval from $`t=r`$ to $`t=0`$.<sup>18</sup><sup>18</sup>18The reader may have noticed that $`\stackrel{}{b}`$ in Eq. (6.43) was defined with an awkward minus sign. This sign was chosen to give the picture just described. Reversing the sign is equivalent to replacing $`\alpha `$ with $`\alpha `$. Insofar as $`\stackrel{}{e}`$ and $`\stackrel{}{b}`$ (like $`\stackrel{}{E}`$ and $`\stackrel{}{B}`$ in electromagnetism) are defined by their physical interpretation, choosing the sign of $`\stackrel{}{b}`$ to give a result with the correct interpretation is legitimate. But this assumes we know what the correct interpretation is, and it is not certain we do. For example, I mentioned above that at $`O(r^5)`$ Hayward’s quasilocal energy gives a negative gravitational energy . If this is correct, then we should replace the definition of $`\stackrel{}{b}`$ with $`\stackrel{}{b}`$.
Finally, we consider the spatial limit case, $`\alpha =0`$. According to our discussion above we would expect $`\mathrm{IQE}_1`$ to be the same as in Eq. (6.46), except with the radiation term absent. Inspection of Eqs. (6.44-6.45) reveals that this is not the case. However, it is only when $`\alpha =1`$ (and presumably also when $`\alpha =1`$) that we know that $`\sigma \sigma R^{(2)}`$ in Eq. (6.8) is proportional to $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{(2)}`$, in which case it suffices to consider only the shear term. Unfortunately, it is not possible to compute $`\sigma \sigma R`$ to $`O(r^2)`$ within the framework of our $`O(r^2)`$ Riemann normal coordinates, so we cannot learn if this simple proportionality between the two persists when $`|\alpha |<1`$. One might guess that it almost certainly does not, but I will leave this question for future work. Nevertheless, since we expect $`\stackrel{~}{l}^2`$ to play the key role with regards to radiation, Eq. (6.44) is still of some qualitative value when $`|\alpha |<1`$. From this equation we see that the radiation term is zero when $`\alpha `$ is zero, and turns on in proportion to $`\alpha `$, precisely as it should since the time lapse is now $`\alpha r`$, instead of $`r`$.
As satisfying as it might seem, the picture just given is not truly quasilocal, in the sense that $`\stackrel{}{e}`$ and $`\stackrel{}{b}`$ are evaluated at the point $`p`$. To be truly quasilocal we need ‘electromagnetic’ fields, call them $`\stackrel{}{E}`$ and $`\stackrel{}{B}`$, evaluated on $`S`$. This is where observers reside, and measurements are made, according to the quasilocal idea. Such a quasilocal picture is achieved very naturally as follows. The basic idea is: $`\stackrel{}{e}`$ and $`\stackrel{}{b}`$ are certain components of the Weyl tensor evaluated at $`p`$ ($`r=0`$). But this information is contained in the $`O(r)`$ piece of certain connection coefficients evaluated on $`S`$ ($`r>0`$). Thus we expect the desired $`\stackrel{}{E}`$ and $`\stackrel{}{B}`$ fields to be associated with connection coefficients.
For simplicity we will restrict ourselves to the case $`\alpha =1`$, which also allows us to borrow some results from Ref. . We first observe that
$$\mathrm{\Psi }_0|_p=2[(e^2b^3)+i(e^3+b^2)].$$
(6.47)
On the left hand side is a component of the Weyl tensor in Newman-Penrose (NP) notation. Using Eqs. (6.28) and (6.38-6.43), $`\mathrm{\Psi }_0`$ is easily converted to the expression given on the right hand side. From Eq. (B5b) of Ref. we have
$$\sigma =\frac{r}{3}\mathrm{\Psi }_0|_p+O(r^2),$$
(6.48)
where $`\sigma `$ is one of the NP spin coefficients. Thus we find that
$$\frac{1}{4}\sigma \overline{\sigma }=\frac{r^2}{9}[\stackrel{}{e}\stackrel{}{e}+\stackrel{}{b}\stackrel{}{b}2\stackrel{}{e}\times \stackrel{}{b}]+O(r^3).$$
(6.49)
Comparing this with the integrand in Eq. (6.46) we are led to define
$$\stackrel{}{E}:=\frac{r}{3}\stackrel{}{e}+O(r^2),\stackrel{}{B}:=\frac{r}{3}\stackrel{}{b}+O(r^2),$$
(6.50)
or in other words,
$$\sigma =2[(E^2B^3)+i(E^3+B^2)]+O(r^2).$$
(6.51)
Thus, to $`O(r)`$, $`\stackrel{}{E}`$ and $`\stackrel{}{B}`$ are related to connection coefficients, as we expected.
Substituting Eq. (6.50) into Eq. (6.46) we have
$$\mathrm{IQE}_1=_S𝑑Sr\left[\frac{1}{8\pi }(\stackrel{}{E}\stackrel{}{E}+\stackrel{}{B}\stackrel{}{B})\frac{1}{4\pi }\stackrel{}{E}\times \stackrel{}{B}\right]+O(r^6).$$
(6.52)
Observe that the mysterious $`r^2/9`$ factor has disappeared, and the analogy with electromagnetism is improved: $`\stackrel{}{E}`$ and $`\stackrel{}{B}`$ now have their usual dimension ($`L^1`$), as do the energy density and Poynting flux terms. I emphasize again that, in contrast to $`\stackrel{}{e}`$ and $`\stackrel{}{b}`$ in Eq. (6.46), $`\stackrel{}{E}`$ and $`\stackrel{}{B}`$ are fields measured by observers residing in $`S`$, in the true quasilocal spirit. $`\stackrel{}{E}`$ is clearly associated with tidal forces tangential to $`S`$, and $`\stackrel{}{B}`$ is a measure of frame dragging effects. Notice that $`\stackrel{}{E}`$ and $`\stackrel{}{B}`$ vanish as $`r0`$, in accord with the equivalence principle.
To conclude this section we consider the connection in the normal bundle and its associated curvature. I find that
$$A_J=\frac{2}{3}r\left[\alpha E_{1J}+\frac{1}{2}B_{1J}\right]+O(r^2),$$
(6.53)
where $`B_{1J}=ϵ_J^KB_{1K}`$, and
$$=2B_{11}+O(r).$$
(6.54)
So to leading order the curvature of the normal bundle is (twice) the radial-radial component of the magnetic part of the Weyl tensor, and is thus associated with ‘gravitational magnetic charge’. There are both local and global dimensions to this result. Locally, $``$ is associated with frame dragging, a ready example being $``$ for the Kerr black hole given in Eq. (5.19), which is proportional to the angular momentum. Globally, it is known that in exact analogy with the scalar curvature $``$, the integral of $``$ over $`S`$ is proportional to the Euler number of the normal bundle . For a Euclidean-signature spacetime the normal bundle is an $`SO(2)`$ (rather than $`SO(1,1)`$) bundle, and there can be a nontrivial winding number, corresponding to a gravitational magnetic monopole. In the case of the Kerr spacetime there is no monopole present since, as is obvious from inspection of Eq. (5.19), the integral of $``$ is zero. It might be interesting to explore topologically nontrivial cases in the context of the IQE.
The result in Eq. (6.54) can actually be obtained immediately by inspection of Eq. (2.16), assuming that the shear terms are higher order in $`r`$ than $``$ is. To lowest order in $`r`$ we then see that $`=2R_{0123}`$. But $`R_{0123}=C_{0123}`$ is identically true, and thus we are led to Eq. (6.54). So this equation is true whether or not matter is present, and is also independent of $`\alpha `$. It is instructive to compare this result with the sectional curvature of $`S`$ in vacuo:
$$\sigma \sigma R=2E_{11}+O(r).$$
(6.55)
This is essentially the same as the spatial infinity limit result given in Eq. (5.21), and is derived similarly. Comparing the previous two equations we see a striking electric/magnetic duality between the sectional curvature of $`S`$ (electric), and the curvature of its normal bundle (magnetic). When matter is present, the right hand side of the equation above acquires an additional term, and it is precisely this term that is responsible for the $`O(r^3)`$ matter contribution seen in Eq. (6.13). So the sectional curvature is the dominant term in the energy that encodes information about the matter content of the spacetime. It seems reasonable to expect inertial effects (frame dragging) produced by this matter to also play a role in the energy. But consideration of such effects is subtle, because the magnetic part of the Weyl tensor has no Newtonian gravity analogue. I have argued that the procedure suggested in Sec. 4 is a geometrically natural way to incorporate such inertial effects into the energy: one demands $`^{\mathrm{ref}}=`$, and then solves the embedding equations for the reference shear term, present in the reference energy. In this way the inertia information contained in $``$ makes its presence felt in the energy. Moreover, by inspection of the purely spatial ($`\alpha =0`$) case of Eq. (6.53), one observes that the set of magnetic quantities: $``$, $`A`$, and $`\stackrel{}{b}`$, precisely encode the five independent components of the magnetic part of the Weyl tensor. It seems likely that the phenomenon of gravitational energy is subtle enough to be sensitive to this full set. Out of this set, in this section we have seen only the role of $`\stackrel{}{b}`$. To see whether or not the other components play a role (via $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$) will have to wait until a solution to the embedding equations is found.
## 7 Asymptotically anti-de Sitter spacetimes
In this last section we will explore the significance of the $`\sigma \sigma R^{\mathrm{ref}}`$ term in Eq. (4.3). Suppose our physical spacetime $`(M,g)`$ is asymptotically anti-de Sitter space. The $`\sigma \sigma R^{\mathrm{ref}}`$ term in $`\mathrm{IQE}^{\mathrm{ref}}`$ gives us the freedom to specify the Riemann tensor of a reference spacetime, which in this case is naturally the Riemann tensor of anti-de Sitter space. Thus, according to Eq. (4.5) we have $`\sigma \sigma R^{\mathrm{ref}}=2/\mathrm{}^2`$, and so
$$\mathrm{IQE}^{\mathrm{ref}}=\frac{1}{8\pi }_S𝑑S\sqrt{2\left[\frac{2}{\mathrm{}^2}++(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}\right]}.$$
(7.1)
In the large sphere limit it is clear that the cosmological constant term will dominate, rather than $``$, and the behavior of the IQE is qualitatively different from that for asymptotically flat spacetimes.
Let us now specialize to the case that $`(M,g)`$ is the $`\mathrm{AdS}_4`$-Schwarzschild spacetime, so that our main argument is not obscured by consideration of the shear terms, which will obviously be just zero. The line element in this case is given by
$$ds^2=N^2dt^2+\frac{1}{f^2}dr^2+r^2d\mathrm{\Omega }^2,\mathrm{where}N(r)=f(r)=\left(\frac{r^2}{l^2}+1\frac{2M}{r}\right)^{\frac{1}{2}},$$
(7.2)
and $`d\mathrm{\Omega }^2=d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2`$ is the line element on the unit round sphere. Let $`S`$ be a $`t,r=\mathrm{constant}`$ two-sphere. Its scalar curvature is $`=2/r^2`$, and a simple calculation shows that its sectional curvature is given by
$$\sigma \sigma R=\frac{2}{\mathrm{}^2}+\frac{4M}{r^3},$$
(7.3)
the dominant term coming from the anti-de Sitter ‘background’. Substituting these results into Eqs. (4.2) and (7.1) we find
$$\mathrm{IQE}=\frac{1}{8\pi }_S𝑑S\sqrt{2\left[\frac{2}{\mathrm{}^2}+\frac{2}{r^2}\frac{4M}{r^3}\right]}+\frac{1}{8\pi }_S𝑑S\sqrt{2\left[\frac{2}{\mathrm{}^2}+\frac{2}{r^2}\right]}=\frac{M\mathrm{}}{r}+O\left(\frac{1}{r^3}\right).$$
(7.4)
The divergent terms due to the cosmological constant cancel, so the limit of the IQE as $`r\mathrm{}`$ exists, and this limit is zero. This would be the expected result if the IQE had the interpretation of an energy, which should be red shifted to zero by the cosmological horizon. In contrast, we do not expect an invariant mass to be redshifted. This is why the terminology ‘invariant quasilocal energy’ was chosen rather than ‘invariant quasilocal mass’, even though the IQE is the analogue of the mass $`m`$ in the formula: $`m=\sqrt{E^2\stackrel{}{p}^{\mathrm{\hspace{0.33em}2}}}`$.
However, one can easily modify the definition of the IQE—to give it the interpretation of mass—by multiplying the right hand side of Eq. (3.4) by a lapse function. Thus one replaces Eq. (4.2) with
$$\mathrm{IQE}[N_{}]=\frac{1}{8\pi }_S𝑑SN_{}\sqrt{2\left[\sigma \sigma R\right]}+\frac{1}{8\pi }_S𝑑SN_{}^{\mathrm{ref}}\sqrt{2\left[\sigma \sigma R^{\mathrm{ref}}\right]}$$
(7.5)
(ignoring the shear terms). Here the smearing function, $`N_{}`$, is the lapse function in the timelike three-boundary, $``$, swept out by the two-parameter family of observers (cf. Eqs. (11) and (13) in Ref. ). In the $`\mathrm{AdS}_4`$-Schwarzschild example $``$ is an $`r=\mathrm{constant}`$ surface, and $`N_{}=N`$. The question arises, What are we to put for $`N_{}^{\mathrm{ref}}`$? The answer that works is $`N_{}^{\mathrm{ref}}=N_{}`$, which is intuitively justified as follows: we are already isometrically embedding $`(S,\sigma )`$ into $`(M^{\mathrm{ref}},g^{\mathrm{ref}})`$, and the condition $`N_{}^{\mathrm{ref}}=N_{}`$ represents the next would-be step towards an isometric embedding of $`(,\gamma )`$ into $`(M^{\mathrm{ref}},g^{\mathrm{ref}})`$, where $`\gamma _{ab}`$ is the three-metric in $``$. (‘Would-be’ in the sense that, while the lapse carries some information about $`\gamma _{ab}`$, we still only need to embed $`(S,\sigma )`$ into $`(M^{\mathrm{ref}},g^{\mathrm{ref}})`$, not $`(,\gamma )`$ into $`(M^{\mathrm{ref}},g^{\mathrm{ref}})`$.) By comparing Eq. (7.5) with (7.4), and using the fact that the lapse function goes as $`r/\mathrm{}`$ for large $`r`$, it is easy to see that with $`N_{}^{\mathrm{ref}}=N_{}`$ we get $`lim_r\mathrm{}\mathrm{IQE}[N_{}]=M`$ for the $`\mathrm{AdS}_4`$-Schwarzschild case. Thus $`\mathrm{IQE}[N_{}]`$ has the interpretation of a mass, as claimed. Unlike the original IQE, it is not red shifted to zero, and is thus a different physical quantity. Regarding the comment at the end of the previous paragraph, since special relativity does not know about lapse functions, the generalization given in Eq. (3.4) is open to this ambiguity: one can define both an invariant quasilocal energy and an ‘invariant’ quasilocal mass.
It is instructive to evaluate $`\mathrm{IQE}[N_{}]`$ also at the horizon, $`r=r_+`$ (where $`N_{}(r_+)=0`$), and compare with what one gets using the unsmeared IQE. The following results for the $`\mathrm{AdS}_4`$-Schwarzschild example are easily established:
$`\mathrm{IQE}`$ $`=`$ $`\{\begin{array}{cc}\sqrt{2Mr_+}\hfill & \text{at }r=r_+\hfill \\ 0\hfill & \text{at }r=\mathrm{}\hfill \end{array}`$ (7.8)
$`\mathrm{IQE}[N_{}]`$ $`=`$ $`\{\begin{array}{cc}0\hfill & \text{at }r=r_+\hfill \\ M\hfill & \text{at }r=\mathrm{}\hfill \end{array}`$ (7.11)
We thus learn that the IQE decreases with $`r`$, which can be interpreted as the result of negative binding energy—another reason to think of the IQE as an energy. On the other hand, $`\mathrm{IQE}[N_{}]`$ increases with $`r`$, which might be interpreted as saying that, for larger $`r`$, ‘more mass is enclosed’; it starts from zero at the horizon (no mass in the interior of the black hole) and accumulates to $`M`$ at infinity. This is reminiscent of the old notion that the substance of mass is nothing but the curvature of spacetime itself. A similar behavior is observed for the usual Schwarzschild case:
$`\mathrm{IQE}`$ $`=`$ $`\{\begin{array}{cc}2M\hfill & \text{at }r=2M\hfill \\ M\hfill & \text{at }r=\mathrm{}\hfill \end{array}`$ (7.14)
$`\mathrm{IQE}[N_{}]`$ $`=`$ $`\{\begin{array}{cc}0\hfill & \text{at }r=2M\hfill \\ M\hfill & \text{at }r=\mathrm{}\hfill \end{array}`$ (7.17)
The only qualitative difference occurs at $`r=\mathrm{}`$, where $`\mathrm{IQE}[N_{}]=\mathrm{IQE}`$ in the Schwarzschild case because, of course, the lapse function goes to one in this limit. There is no cosmological horizon.
Despite these appealing features of $`\mathrm{IQE}[N_{}]`$, it is unsatisfactory from the point of view taken here because the presence of the lapse function means it depends on a choice of three-surface passing through $`S`$ (hence the use of inverted commas in the terminology ‘invariant’ quasilocal mass above). The situation might be improved by replacing $`k`$ with $`N_{}k`$ and $`l`$ with $`N_\mathrm{\Sigma }l`$ in Eq. (3.4), and then proceeding as before. Here $`N_{}`$ and $`N_\mathrm{\Sigma }`$ are time and radial lapse functions, respectively (equal to $`N`$ and $`1/f`$ in the $`\mathrm{AdS}_4`$-Schwarzschild example). Admittedly such a procedure is ad hoc, and unless it can be improved upon we are not particularly interested in $`\mathrm{IQE}[N_{}]`$. It was introduced simply to illustrate the distinction between mass and energy, but for the remainder of this section we will return to the original definition of the IQE.
The main point of this section is to draw attention to a remarkable similarity between the reference subtraction term given in Eq. (7.1), and a certain counterterm action recently suggested in the context of the conjectured AdS/CFT correspondence. We begin by observing that when the reference shear term vanishes, Eq. (7.1) reduces to precisely the same reference subtraction term suggested by Lau , in the context of the Brown-York quasilocal energy (except that Lau’s expression has a lapse function present in the manner discussed above). However, our derivations of this expression are different. Lau employs a ‘light cone reference’ embedding of $`(S,\sigma )`$, together with a ‘rest frame’ assumption, $`l^{\mathrm{ref}}=0`$, to derive an expression for $`k^{\mathrm{ref}}`$, which is then used to construct his reference subtraction term. In our case we get the same end result, but we get it without any recourse to a reference embedding! This is because the cosmological constant term in Eq. (7.1) comes directly from the explicit dependence of $`\mathrm{IQE}^{\mathrm{ref}}`$ on the Riemann tensor of the reference spacetime, i.e., the term $`\sigma \sigma R^{\mathrm{ref}}`$ in Eq. (4.3). Our reference embedding of $`(S,\sigma )`$ into $`(M^{\mathrm{ref}},g^{\mathrm{ref}})`$ is required only to evaluate the reference shear term, $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$. This ‘higher order correction’—which I have argued accounts for angular momentum—is not present in Lau’s reference subtraction term. Also, his additional ‘rest frame’ assumption is not required here because the IQE is already naturally a ‘rest frame energy’.
Let us now return to Eq. (1.5). We know that when space is non-compact the boundary (or quasilocal) energy-momentum tensor $`T_{}^{ab}=\mathrm{\Pi }^{ab}/(8\pi )`$ diverges in general as $``$ is taken to infinity.<sup>19</sup><sup>19</sup>19As elsewhere in this paper, we use the symbol $``$ loosely to refer to either a timelike three-surface in the interior of $`M`$, bounding a finite spatial region, or the boundary at infinity. It’s meaning should be clear from the context in which it is used. To render it finite, Brown and York suggest the use of a reference subtraction term that involves an isometric embedding of $`(,\gamma )`$ into a suitable reference spacetime. However, like their prescription to embed $`(S,\sigma )`$ into a suitable three-dimensional reference space, this prescription suffers from the drawback that such a codimension-one embedding does not always exist. Recently, Balasubramanian and Kraus have proposed an alternative procedure: Since it is always possible to add to the action a local functional of the intrinsic geometry of the boundary without affecting the equations of motion or the symmetries (but of course this alters $`T_{}^{ab}`$), their idea is to choose this functional such that its divergences cancel those of the original $`T_{}^{ab}`$, rendering the improved boundary energy-momentum tensor finite as $``$ is taken to infinity. No recourse is made to a reference embedding. This procedure was first applied to spacetimes that are asymptotically AdS space, in which case the required ‘counterterms’ amounted to a simple finite polynomial in the curvature invariants of $``$ . This idea is exactly analogous to the standard prescription for removing ultraviolet divergences in quantum field theory by adding to the Lagrangian a finite polynomial in the fields. Moreover, the conjectured AdS/CFT correspondence implies that the two procedures are not merely analogous, they are one and the same .
Now since flat spacetime is recovered from AdS space by taking $`\mathrm{}`$ to infinity, one might expect that in this same limit the counterterms found by Balasubramanian and Kraus would produce counterterms suitable for asymptotically flat spacetimes. It is not obvious that this is so . However, Mann has suggested the following generalization of their counterterm action:
$$I_{\mathrm{ct}}=\frac{1}{8\pi }_{_{\mathrm{}}}d^3x\sqrt{\gamma }\sqrt{2\left[\frac{2}{\mathrm{}^2}+R(\gamma )\right]},$$
(7.18)
where $`R(\gamma )`$ is the scalar curvature of $`(,\gamma )`$, and $`_{\mathrm{}}`$ indicates that we are to take the limit as $``$ goes to infinity. For small $`\mathrm{}`$ Mann’s formula reduces to the one given by Balasubramanian and Kraus, but in addition it has a smooth flat spacetime limit as $`\mathrm{}\mathrm{}`$. Moreover, Mann showed that in many explicit examples it leads to a cancellation of all divergences, and the remaining finite part agrees with that obtained using the reference spacetime procedure . While a counterterm action and a reference mass are not the same thing, the resemblance between the expressions in Eqs. (7.18) and (7.1) is nevertheless striking.<sup>20</sup><sup>20</sup>20I am indebted to R.B. Mann for pointing out to me the significance of the $`\sigma \sigma R^{\mathrm{ref}}`$ term in $`\mathrm{IQE}^{\mathrm{ref}}`$, and emphasizing that it provides at least some measure of geometrical motivation for his expression in Eq. (7.18).
To see that the connection between $`\mathrm{IQE}^{\mathrm{ref}}`$ in Eq. (7.1) and the AdS/CFT-inspired counterterm action is probably much deeper than mere resemblance, we now turn to recent work done by Kraus et al . Besides providing an independent derivation of Mann’s formula, and its generalization to higher dimensions, of most interest to us here is their geometrical argument suggesting what the counterterm for $`\mathrm{\Pi }^{ab}`$ in Eq. (1.5) should be in order to cancel divergences. Their result is an expansion in powers of $`\mathrm{}`$. Denoting their counterterm ($`\stackrel{~}{\mathrm{\Pi }}_{ab}`$) as $`\mathrm{\Pi }_{ab}^{\mathrm{ct}}`$, and specializing their result to a three-dimensional boundary $``$, they find:
$`\mathrm{\Pi }_{ab}^{\mathrm{ct}}`$ $`=`$ $`{\displaystyle \frac{2}{\mathrm{}}}\gamma _{ab}+\mathrm{}(R_{ab}{\displaystyle \frac{1}{2}}\gamma _{ab}R)+\mathrm{}^3\{{\displaystyle \frac{1}{2}}\gamma _{ab}(R_{cd}R^{cd}{\displaystyle \frac{3}{8}}R^2)`$ (7.19)
$`+{\displaystyle \frac{3}{4}}RR_{ab}2R^{cd}R_{acbd}+{\displaystyle \frac{1}{4}}D_aD_bR\mathrm{}R_{ab}+{\displaystyle \frac{1}{4}}\gamma _{ab}\mathrm{}R\}+O(\mathrm{}^5).`$
The curvatures and covariant derivatives in this expression all refer to the induced timelike three-metric $`\gamma _{ab}`$ on $``$.
Now consider our usual two-surface $`(S,\sigma )`$ in the physical spacetime $`(M,g)`$, which is an asymptotically AdS space. As we did at the end of Sec. 4, suppose that $`S`$ is such that $`(k^2l^2)>0`$ and $`k>0`$. Then we can always find a timelike unit vector $`u^a`$ normal to $`S`$ such that $`l=0`$, and so $`\sqrt{k^2l^2}=k=\mathrm{\Pi }_{ab}u^au^b`$ is the Brown-York energy surface density (modulo the factor of $`1/(8\pi )`$). In other words, our unreferenced IQE reduces to the unreferenced Brown-York CQE, which would be called the ‘unrenormalized’ energy in Ref. . The counterterm required to renormalize the energy surface density is thus $`\mathrm{\Pi }_{ab}^{\mathrm{ct}}u^au^b`$, which we will denote at $`E_{\mathrm{ct}}`$. Hence, our task is to compare $`E_{\mathrm{ct}}:=\mathrm{\Pi }_{ab}^{\mathrm{ct}}u^au^b`$ with the integrand of $`\mathrm{IQE}^{\mathrm{ref}}`$ in Eq. (7.1); we expect to see at least some measure of agreement between the two.
This comparison will not be straightforward, however, because on the one hand we expect the integrand of $`\mathrm{IQE}^{\mathrm{ref}}`$ to depend on $``$, $``$, and their derivatives in $`S`$, as discussed previously, whereas on the other hand $`E_{\mathrm{ct}}`$ depends on the three-metric $`\gamma _{ab}`$. Nevertheless, let us see how far we can get. Let $``$ be a three-surface in $`(M,g)`$ passing through $`S`$ in a direction tangent to $`u^a`$ on $`S`$. Different choices of $``$ satisfying these conditions will lead to different induced metrics $`\gamma _{ab}`$, but this ambiguity will not affect our considerations. At least $`\gamma _{ab}`$ on $`S`$ is uniquely determined, and some information about $`\gamma _{ab}`$ in the neighborhood of $`S`$ is determined by the condition $`l=0`$. Our choice of $``$ means that $`l_{ab}`$ defined in Eqs. (2.8) is the extrinsic curvature of $`(S,\sigma )`$ as embedded in $`(,\gamma )`$, and so the corresponding codimension-one Gauss embedding equation reads
$$𝒫_a^e𝒫_b^f𝒫_c^g𝒫_d^hR_{efgh}=_{abcd}+(l_{ac}l_{bd}l_{bc}l_{ad}).$$
(7.20)
This is just a truncated version of Eq. (2.14), except here $`R_{efgh}`$ is the Riemann tensor of $`(,\gamma )`$, not $`(M,g)`$.
Now let $`E_{\mathrm{ct}}^{(n)}`$ denote the term in $`E_{\mathrm{ct}}`$ of order $`\mathrm{}^n`$. Inspection of Eq. (7.19) shows that $`E_{\mathrm{ct}}^{(1)}=2/\mathrm{}`$. The term $`E_{\mathrm{ct}}^{(1)}`$ can be written in terms of $`G_{ab}`$, the Einstein tensor of $`(,\gamma )`$, and we have
$$E_{\mathrm{ct}}^{(1)}=\mathrm{}G_{ab}u^au^b=\frac{\mathrm{}}{2}\sigma ^{ac}\sigma ^{bd}R_{abcd}=\frac{\mathrm{}}{2}(\stackrel{~}{l}^2).$$
(7.21)
The second equality is an easily derived identity (valid in any codimension-one setting) relating the $`uu`$ component of the Einstein tensor to the sectional curvature of the hypersurface orthogonal to $`u^a`$ (in this case the hypersurface $`S`$ in $``$). The third equality follows from contracting Eq. (7.20) with $`\sigma ^{ac}\sigma ^{bd}`$, and using the fact that $`l=0`$ by our choice of $`u^a`$. (And as usual, $`\stackrel{~}{l}^2`$ is shorthand for $`\stackrel{~}{l}_{ab}\stackrel{~}{l}^{ab}`$.) Thus, to order $`\mathrm{}`$ we have
$$E_{\mathrm{ct}}=\frac{2}{\mathrm{}}+\frac{\mathrm{}}{2}(\stackrel{~}{l}^2)+O(\mathrm{}^3)=\sqrt{2\left[\frac{2}{\mathrm{}^2}+\stackrel{~}{l}^2+O(\mathrm{}^2)\right]}.$$
(7.22)
Comparing the last expression with Eq. (7.1) suggests the correspondence:
$$(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}\stackrel{~}{l}^2+O(\mathrm{}^2).$$
(7.23)
Immediately we see something odd: we are identifying a boost invariant quantity with one that is not, i.e., it seems that a $`\stackrel{~}{k}^2`$ is missing from the right hand side. I will comment on this shortly. Let us assume for the moment that the right hand side reads $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)`$, in which case Eq. (7.23) seems reasonable: it suggests that, if we solve the embedding equations (4.6-4.8) for $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$ we will find that, to lowest order in $`\mathrm{}`$, the reference shear term is the same as the unreferenced shear term, the difference to be seen at a higher order in $`\mathrm{}`$. On the other hand, this seems like a problem: Would it not mean, e.g., that the shear terms in Eq. (5.5) basically cancel, thus ruining the Bondi-Sachs mass result in Eq.(5.49), which depends so crucially on $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)`$? The answer is No, because Eq. (5.5) is valid in the asymptotically flat case, not the asymptotically AdS case. To make a statement that is valid in the asymptotically flat case ($`\mathrm{}\mathrm{}`$) we need to know $`E_{\mathrm{ct}}`$ to all orders in $`\mathrm{}`$, then sum the infinite series, and finally take the limit $`\mathrm{}\mathrm{}`$. So being at the ‘other end’ of the series, Eq. (7.23) has nothing to say about the asymptotically flat case. But we also expected $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$ to depend on $``$, $``$, and their derivatives. Why do we not see these quantities on the right hand side of Eq. (7.23)? The answer is, We will—we just have to calculate $`E_{\mathrm{ct}}`$ to the next order in $`\mathrm{}`$.
But before doing so I will comment on the ‘missing’ $`\stackrel{~}{k}^2`$ in Eq. (7.23). Kraus et al have devised an algorithm to compute the extrinsic geometrical quantity $`\mathrm{\Pi }_{ab}^{\mathrm{ct}}`$ from the intrinsic geometry of $``$. Insofar as $`l_{ab}`$ (and thus $`\stackrel{~}{l}^2`$) depends only on the metric $`\gamma _{ab}`$, there is no doubt that the $`\stackrel{~}{l}^2`$ term in Eq. (7.23) is ‘correct’. On the other hand, $`k_{ab}`$ (and thus $`\stackrel{~}{k}^2`$) depends on the extrinsic geometry of $``$, being just a certain projection of $`\mathrm{\Pi }_{ab}`$ into $`S`$. The algorithm of Kraus et al relies on the fact that the divergent part of the derivative of $`\mathrm{\Pi }_{ab}`$ in the direction normal to $``$ can be expressed in terms of just the intrinsic geometry of $``$. In essence, their algorithm is designed precisely to compute the $`\mathrm{𝑑𝑖𝑣𝑒𝑟𝑔𝑒𝑛𝑡}`$ part of $`\mathrm{\Pi }_{ab}`$. The ‘correctness’ of the accompanying finite part is a subtle issue. In a slightly different context, they discuss two different counterterm actions that both properly cancel divergences, but that lead to different finite terms in the action. Furthermore, they point out that their algorithm, when carried to all orders in $`\mathrm{}`$, might imply singularities in the bulk spacetime, but that this is of no concern because they truncate their counterterm expressions to a finite number of terms, enough at least to cancel the divergences. In our case we have the quasilocal idea in mind, and so are interested in all of the finite terms—it matters what happens in the bulk. But going further with this discussion will take us beyond the scope set for this simple comparison. I will just conclude by saying that, insofar as the shear terms almost certainly represent a finite contribution to the energy, we do not necessarily expect the algorithm of Kraus et al to produce a $`\stackrel{~}{k}^2`$ term on the right hand side of Eq. (7.23). Our goals are slightly different, and it is too much to expect exact agreement between $`E_{\mathrm{ct}}`$ and the integrand of $`\mathrm{IQE}^{\mathrm{ref}}`$.
Nevertheless, it is still instructive to proceed with the comparison to the next order in $`\mathrm{}`$. In light of my previous remarks, we will make the simplifying assumption that the metric on $``$ has a product structure: $`\gamma _{ab}dx^adx^b=N^2dt^2+\sigma _{ij}(x)dx^idx^j`$, where $`x^a=(t,x^i)`$ are local coordinates on $``$, $`N`$ is a constant lapse, and $`\sigma _{ij}(x)`$ is the metric on any $`t=\mathrm{constant}`$ two-surface $`S`$. The idea is that $`E_{\mathrm{ct}}`$ and the integrand of $`\mathrm{IQE}^{\mathrm{ref}}`$ should agree at least in their dependence on the intrinsic geometry of $`(S,\sigma )`$. Assuming such a product structure for $`\gamma _{ab}`$ is a convenient was to isolate this dependence, and ignore everything else.
In this case the only nonvanishing components of the Riemann tensor of $`\gamma _{ab}`$ are $`R_{ijkl}=_{ijkl}`$, the Riemann tensor of $`\sigma _{ij}`$. And clearly $`l_{ab}=0`$. Thus $`E_{\mathrm{ct}}^{(1)}`$ in Eq. (7.21) reduces to $`\mathrm{}/2`$, and it is a simple exercise to work out $`E_{\mathrm{ct}}^{(3)}`$. The net result is
$$E_{\mathrm{ct}}=\frac{2}{\mathrm{}}+\frac{\mathrm{}}{2}\frac{\mathrm{}^3}{16}\left(^24\mathrm{}\right)+O(\mathrm{}^5)=\sqrt{2\left[\frac{2}{\mathrm{}^2}+\frac{\mathrm{}^2}{2}\mathrm{}+O(\mathrm{}^4)\right]},$$
(7.24)
where $`\mathrm{}`$ is the Laplacian in $`(S,\sigma )`$. Comparing the last expression with Eq. (7.1) we now have the higher order correspondence:
$$(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}\frac{\mathrm{}^2}{2}\mathrm{}+O(\mathrm{}^4).$$
(7.25)
Thus we begin to see how a solution to our embedding equations might yield an expression for $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$ in terms of $``$, $``$, and their derivatives, as we have expected all along.
To conclude this section we make two general observations. First, it is especially clear from the higher order expression in Eq. (7.24) that the AdS/CFT-inspired counterterm energy is, in fact, the square root of some quantity. This is not surprising, since the algorithm of Kraus et al is a means of solving a Gauss embedding equation for $`\mathrm{\Pi }_{ab}^{\mathrm{ct}}`$, and this equation is quadratic in $`\mathrm{\Pi }_{ab}^{\mathrm{ct}}`$. But it is significant. Beginning simply with the definition of the quasilocal energy-momentum tensor as the functional derivative of the action with respect to the boundary metric , in which there is no square root in sight, the counterterm energy required to cancel divergences unmistakably involves a square root. Moreover, it concurs with the square root introduced here, in the context of the IQE, as the general relativistic analogue of the special relativistic formula: $`m=\sqrt{E^2\stackrel{}{p}^{\mathrm{\hspace{0.33em}2}}}`$. I believe it is unlikely this is a mere coincidence. Given that it is nonanalytic, a square root is too unusual an object to occur without good reason.
Secondly, under the square root (in our case) is $`2/\mathrm{}^2++(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$. In the case of the AdS/CFT-inspired counterterm energy , it is $`2/\mathrm{}^2++X`$, where $`X`$ is an infinite series in increasing powers of $`\mathrm{}`$. That $`X`$ is clearly not zero lends strong support for our additional term $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$, which is thus seen to be a necessary generalization of Lau’s suggestion . I have argued that its necessity is closely linked to the proper inclusion of angular momentum in the energy. Given that angular momentum is a subtle notion in general relativity, especially so at the quasilocal level we envision here, it is not surprising that our biggest difficulty lies in evaluating $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$. In light of the algorithm given by Kraus et al , work is currently in progress to try to apply similar techniques to solve the embedding equations (4.6-4.8). Since these embedding equations are manifestly boost invariant, I expect at least to recover the ‘missing’ $`\stackrel{~}{k}^2`$ term in Eq. (7.23), and hopefully the entire series.
## 8 Summary and discussion
In this paper I have introduced a new definition of quasilocal energy that is a simple modification of the Brown-York quasilocal energy. I just replace their energy surface density $`k`$ with $`\sqrt{k^2l^2}`$, where $`l`$ is the radial momentum surface density. (For ease of exposition here I will omit the $`1/(8\pi )`$ factors.) The principle motivation for doing this stems from an analogy with the formula: $`m=\sqrt{E^2\stackrel{}{p}^{\mathrm{\hspace{0.33em}2}}}`$ in special relativity. Identifying $`E`$ with $`k`$ (which are both energies), and $`\stackrel{}{p}`$ with $`l`$ (both momenta), identifies $`m`$ with $`\sqrt{k^2l^2}`$. Like $`m`$, $`\sqrt{k^2l^2}`$ is a boost invariant quantity, and hence the integral of $`\sqrt{k^2l^2}`$ over a spacelike two-surface $`S`$ gives rise to an ‘invariant quasilocal energy’, or IQE. In what follows I will refer to the Brown-York quasilocal energy as the CQE—canonical quasilocal energy.
There are several important consequences of replacing $`k`$ with $`\sqrt{k^2l^2}`$:
* While $`k`$ is always well defined for any spacelike two-surface $`S`$, $`\sqrt{k^2l^2}`$ is not. Roughly speaking, it is real when $`S`$ lies in the exterior region of a black hole, zero when it is on the horizon, and imaginary in the black hole interior. Thus (again roughly speaking) the IQE asserts that energy is real only outside of a black hole.
* Both the CQE and the IQE require a reference energy subtraction procedure. Since $`k`$ is associated with a spacelike three-surface spanning $`S`$, the reference space into which $`S`$ is to be isometrically embedded is inherently three-dimensional. Such a codimension-one embedding does not always exist, but when it does, it is essentially unique. This means the CQE, when it is defined, is unique. In contrast, $`\sqrt{k^2l^2}`$ makes no reference to a three-surface spanning $`S`$, and so the reference space(time) is inherently four-dimensional. Such codimension-two embeddings (at least of a generic non-round sphere into Minkowski space) always exist , but are not unique. However, in this situation there are two curvatures associated with $`S`$: the scalar curvature $``$, and the curvature of the normal bundle, $``$. A necessary condition for an isometric embedding is that $`^{\mathrm{ref}}=`$. I argued that demanding also $`^{\mathrm{ref}}=`$ is both a means to make the embedding essentially unique, and at the same time, a geometrically natural way to properly incorporate angular momentum into energy at the quasilocal level. Indeed, since angular momentum is associated with rotational kinetic energy, it should contribute to the energy in some way.
* While $`\mathrm{CQE}^{\mathrm{ref}}`$ is associated with a reference energy density $`k^{\mathrm{ref}}`$, $`\mathrm{IQE}^{\mathrm{ref}}`$ is concerned with a reference shear term, $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$. ($`\stackrel{~}{k}_{ab}`$ and $`\stackrel{~}{l}_{ab}`$ are the trace-free parts of the two extrinsic curvatures of $`S`$.) In a certain sense, the IQE already inherently contains the correct reference energy, without recourse to a reference embedding. The reference embedding is required only to determine the reference shear term, which is a higher order correction to the energy associated with angular momentum.
* The CQE is sensitive to the sign of $`k`$, whereas since it involves $`\sqrt{k^2l^2}`$, the IQE is not. Thus one can easily construct simple examples for which the two energies give different results, even when $`l=0`$. Thus the IQE is not simply the ‘rest energy’ version of the CQE. Note: the IQE naturally assigns zero energy to any two-surface in flat spacetime. This is because the natural reference spacetime in this case is the very same spacetime, namely flat spacetime. So obviously one can always reference-embed the two-surface identically (up to Poincaré transformations) to the way it is embedded in the physical spacetime, and get $`\mathrm{IQE}=0`$. The only subtlety that may arise is if the two conditions: ‘isometric embedding’ and ‘$`^{\mathrm{ref}}=`$’ do not uniquely determine $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$. Then the flat spacetime result ($`\mathrm{IQE}=0`$) may be reduced to a choice, rather than a necessary fact. To properly address this subtlety requires an in depth understanding of the embedding equations. But in any case, the fact that $`\mathrm{IQE}=0`$ in flat spacetime is independent of the motion of the observers. In contrast, moving observers in flat spacetime could measure nonzero energy in the Brown-York approach . This is because under a radial boost the Brown-York energy surface density dilates by a Lorentz factor, as in special relativity, whereas the reference energy surface density does not. According to Ref. the latter depends only on the intrinsic geometry of $`S`$, and therefore does not know about the time derivative of this geometry.<sup>21</sup><sup>21</sup>21I thank R.B. Mann and I.S. Booth for this remark on the Brown-York case.
We examined both the large and small sphere limits of the IQE, taking $`S`$ to be asymptotically round for simplicity. In an asymptotically flat spacetime, the large sphere limit of the IQE in a spatial direction yields the ADM mass. In the future null direction it reduces to the Bondi-Sachs mass, provided the reference shear term is a total divergence. Short of solving the embedding equations, I gave a heuristic argument which shows that is is. It is significant that this argument relies on the condition $`^{\mathrm{ref}}=`$, since this provides evidence that the curvature of the normal bundle is involved in quasilocal energy, albeit its involvement in this simple example is minimal.
The quantity $`\sqrt{k^2l^2}`$ is proportional to the mean curvature of $`S`$ as a two-surface embedded in the physical spacetime, and so the IQE is a natural geometrical invariant of $`S`$. Since the Hawking mass is constructed using $`(k^2l^2)`$, the IQE can be thought of roughly as the square root of the Hawking mass. In the small sphere limit the square root disappears, and to leading order the IQE reduces to the Hawking mass (but differs from it at higher order). Thus, when matter is present, the lowest order contribution to the IQE gives the standard result: $`(4\pi r^3/3)T_{ab}^{\mathrm{mat}}u^au^b`$, i.e., the expected matter energy contained in a small sphere of proper radius $`r`$. Note that $`u^a`$ here is not necessarily the four-velocity of any observer on $`S`$, since the IQE is boost invariant, and so independent of the observers’ velocities on $`S`$. Rather, $`u^a`$ is the four-velocity that observers would have if they were in the rest frame determined by $`S`$. More precisely, in the small sphere limit we considered, namely a $`t,r=\mathrm{constant}`$ two-sphere in Riemann normal coordinates (with $`tr`$), $`u^a=(/t)^a`$ evaluated at the center of the sphere. In the limit $`r0`$, the four-velocity $`(/t)^a`$ corresponds to observers who at each point on $`S`$ have zero radial momentum, i.e., $`l=0`$. In general, since the IQE is an energy rather than a mass, the question arises, In whose rest frame is the energy measured?<sup>22</sup><sup>22</sup>22I thank A. Ashtekar for posing this question. The answer is, The ‘quasilocal rest frame’ determined by the condition $`l=0`$ at each point on $`S`$. Whenever $`(k^2l^2)>0`$, observers on $`S`$ can always achieve this state of motion by appropriate local radial boosts. This (or more precisely, $`(k^2l^2)0`$) is the same condition required for the unreferenced IQE to be well defined in the first place—refer to the text in the first bullet on page 8.
Returning to the small sphere case, in vacuo the leading order contribution due to gravitational energy occurs at order $`r^5`$. At this order the IQE results are inconclusive because it is expected that the reference shear term will play a significant role, and without a solution to the embedding equations (which is an extremely difficult problem) this term cannot be determined. Nevertheless, it was possible to show that in the small sphere limit, the Hawking mass, which in this case is closely related to the IQE, can be understood as a measure of the gravitational energy contained in $`S`$ by considering certain tangential ‘electromagnetic’ fields $`\stackrel{}{E}`$ and $`\stackrel{}{B}`$ induced on $`S`$ by the Weyl curvature $`S`$ is embedded in. In terms of $`\stackrel{}{E}`$ and $`\stackrel{}{B}`$, gravitational energy and radiation are essentially identical in nature to their counterparts in electromagnetism, except for one crucial difference: the density $`r(\stackrel{}{E}\stackrel{}{E}+\stackrel{}{B}\stackrel{}{B})/(8\pi )`$ is integrated over the surface $`S`$ to determine the gravitational energy contained in the spatial volume that $`S`$ encloses. Here $`r`$ can be thought of as the areal radius of $`S`$, so the measurement is truly quasilocal.
The IQE was analyzed in the context of asymptotically anti-de Sitter spacetimes. The fact that $`\mathrm{IQE}^{\mathrm{ref}}`$ depends explicitly on the Riemann tensor of the reference spacetime (naturally taken to be anti-de Sitter space) was seen to play a significant role. A connection was established between $`\mathrm{IQE}^{\mathrm{ref}}`$ and a certain counterterm energy that has recently been proposed in the context of the conjectured AdS/CFT correspondence. Two similarities are striking: (i) Both energies involve a square root, and (ii) the two leading terms under the square root match. The remaining term under the square root in our case is the reference shear term, $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$; in the case of Ref. it is an infinite series in increasing powers of $`\mathrm{}`$, where $`\mathrm{}`$ is the radius of curvature of the AdS space. It was shown that the first two nontrivial terms of this series (i.e., to the highest order given in Ref. ) can plausibly be identified with $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$. This agreement is impressive because $`\mathrm{IQE}^{\mathrm{ref}}`$ and the AdS/CFT-inspired counterterm energy are independently motivated, and derived quite differently. It might be possible to use techniques developed in Ref. to solve our embedding equations for $`(\stackrel{~}{k}^2\stackrel{~}{l}^2)^{\mathrm{ref}}`$. The present lack of a solution to these equations is the main outstanding obstacle to further understanding the nature of the IQE.
A final remark is in order. Most definitions of quasilocal energy, including the IQE, assume that energy is associated with a closed spacelike two-surface, $`S`$. Given such a two-surface one can always find a timelike unit normal vector field $`u^a`$, which at each point on $`S`$ is supposed to correspond to an observer’s instantaneous four-velocity. But this may not be a general enough setting. While a two-parameter family of observers will always sweep out a timelike three-surface $``$, the two-surface elements orthogonal to their world lines in $``$ are not, in general, integrable. Thus a shift in emphasis from $`S`$ to $``$, i.e., from Eulerian to Lorentzian observers , might lead to a deeper understanding of quasilocal energy, in particular of gravitational radiation at the quasilocal level. This shift would also bring the quasilocal energy idea closer in line with the conjectured AdS/CFT correspondence. Whether or not this is the right direction, the results in Sec. 7 strongly suggest that this is the direction the IQE is pointing in.
Acknowledgments
I would like to thank Joseph Samuel, Madhavan Varadarajan, and Sukanya Sinha at the Raman Research Institute for their continued encouragement. I am also indebted to Robert Mann, Ivan Booth, and Luis de Menezes at the University of Waterloo for many enlightening discussions, and Joseph Samuel and Robert Mann for providing critiques of the manuscript. I thank Gabor Kunstatter and Steven Carlip for input during the early stages of this work. This work was financially supported in part by the Natural Sciences and Engineering Research Council of Canada, and the Department of Science and Technology, Government of India. |
warning/0003/gr-qc0003062.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Recently the existence of the cosmological constant becomes quite probable from the observation of the deep galaxy survey . While the dark matter is necessary in various observations such as the rotational curve of the spiral galaxy or the missing of the ordinary matter in the cosmological scale.
Then we start from the theory of general relativity with the dark matter and the cosmological constant in order to study the standard scenario of the astrophysics, that is, the physics of the cosmological, the galactic or solar scale. Though the neutrino is the promising candidate of the dark matter, there is no established direct observation of the dark matter as the ordinary matter. There is another attempt to explain the rotation curves in the theory of the Brans-Dicke theory , where the Newtonian force is modified by the effect of the scalar field. In this paper, we consider the scalar field as a candidate of the dark matter together with the cosmological constant and study their effects on time development of the scale factor of the universe in the cosmological scale and the gravitational potential in the galactic or solar scale.
The famous scalar-tensor theory is the Brans-Dicke theory, but we adopt the Einstein theory with minimally coupled scalar field instead of the Brans-Dicke theory. Our principle of the choice of the theory is the following. For the scalar-tensor gravity theory, we can transform one from the Jordan frame to the Einstein frame by the conformal transformation . We prefer to adopt the Einstein frame because the post-Newtonian test of the general relativity such as the radar echo delay is quite stringent . Also we prefer the Einstein theory with minimally coupled scalar field because the scalar field as the dark matter means that the scalar field has no direct interaction of the gravitational field with the ordinary matter.
## 2 Classical Solution with Minimally Coupled Scalar and Cosmological Constant
### 2.1 Einstein Theory vs. Brans-Dicke Theory
The Brans-Dicke theory is the typical scalar-tensor theory of the gravity. The action of the Brans-Dicke theory with the cosmological term is given by
$`I_{\mathrm{BD}}={\displaystyle d^4x\sqrt{g}\left[\frac{1}{16\pi G}\left(\xi \varphi ^2R2\mathrm{\Lambda }\varphi ^n\right)\frac{1}{2}g^{\mu \nu }_\mu \varphi _\nu \varphi +_{\mathrm{ordinary}\mathrm{matter}}\right]}.`$
Uehara-Kim found the general solution for $`n=2`$ and matter dominant case, and Fujii has found the special solution for general $`n`$.
Putting $`g_{\mu \nu }(x)=\mathrm{\Omega }^2(x)g_{\mu \nu }(x)`$ with $`\mathrm{\Omega }(x)=\sqrt{\xi }\varphi (x)`$, we obtain the following action
$`I_{\mathrm{BD}}={\displaystyle d^4x\sqrt{g}\left[\frac{1}{16\pi G}\left(R_{}2\mathrm{\Lambda }e^{(n4)\zeta \varphi _{}}\right)\frac{1}{2}g_{}^{\mu \nu }_\mu \varphi _{}_\nu \varphi _{}+\xi ^2e^{4\zeta \varphi _{}}_{\mathrm{ordinary}\mathrm{matter}}\right]},`$
where $`\varphi =\mathrm{exp}(\zeta \varphi _{})`$ with $`\zeta ^1=\sqrt{1/\xi +3/4\pi G}`$ and $`_{\mathrm{ordinary}\mathrm{matter}}`$ is obtained from $`_{\mathrm{ordinary}\mathrm{matter}}`$ by replacing the metric part in the form $`g_{\mu \nu }g_{\mu \nu }\xi ^1\mathrm{exp}(\zeta \varphi _{})`$.
Our philosophy to fix the theory comes from the following two principles: i) the kinetic part of the gravity is of the standard Einstein form, because of the stringent constraint of the post-Newtonian test such as the delay of the radar echo experiment, ii) the scalar field has no direct coupling to the ordinary matter nor gives the effect on the geodesic equation of the particle. From these principles, we do not adopt the Brans-Dicke theory. In the following, we adopt the Einstein theory with standard cosmological term and the minimally coupled scalar field.
### 2.2 Einstein Theory with Minimally Coupled Scalar Field
By considering the minimally coupled scalar field as some kind of dark matter, we study the prototype of the time development of the scale factor of the universe in the cosmological scale and the gravitational potential in the galactic or solar scale.
We use Misner-Thorne-Wheeler notation and consider Einstein action with the cosmological constant, the minimally coupled scalar field and the ordinary matter
$`I={\displaystyle d^4x\sqrt{g}\left[\frac{1}{16\pi G}\left(R2\mathrm{\Lambda }\right)\frac{1}{2}g^{\mu \nu }_\mu \varphi _\nu \varphi +_{\mathrm{ordinary}\mathrm{matter}}\right]},`$ (2.1)
where $`G`$ is the gravitational constant, $`R`$ is the scalar curvature and $`\varphi `$ is the minimally coupled scalar field. The equations of motion in this system are given by
$`R_{\mu \nu }{\displaystyle \frac{1}{2}}g_{\mu \nu }R+\mathrm{\Lambda }g_{\mu \nu }`$ $`=`$ $`8\pi G(T_{\mu \nu }^\varphi +T_{\mu \nu }),`$ (2.2)
$`_\mu (\sqrt{g}g^{\mu \nu }_\nu \varphi )`$ $`=`$ $`0,`$ (2.3)
where $`T_{\mu \nu }^\varphi =_\mu \varphi _\nu \varphi \frac{1}{2}g_{\mu \nu }g^{\rho \sigma }_\rho \varphi _\sigma \varphi `$ and $`T_{\mu \nu }`$ is the energy-momentum tensor of the ordinary matter.
## 3 Cosmological Exact Solution
In order to study classical solutions in cosmology, we substitute the homogeneous, isotropic and flat metric
$$ds^2=dt^2+a(t)^2\left[dr^2+r^2(d\theta ^2+\mathrm{sin}\theta ^2d\phi ^2)\right],$$
(3.1)
and the perfect fluid expression of the ordinary matter $`T_{\mu \nu }=(\rho +p)u_\mu u_\nu +pg_{\mu \nu }`$ into equations of motion. We denote $`\rho `$ and $`p`$ as the density and the pressure of the perfect fluid respectively and we can take $`u_\mu =(1,0,0,0)`$ in the co-moving system. Then equations of motion to be solved become
$`\left({\displaystyle \frac{\dot{a}}{a}}\right)^2{\displaystyle \frac{\mathrm{\Lambda }}{3}}={\displaystyle \frac{8\pi G}{3}}\left(\rho +{\displaystyle \frac{\dot{\varphi }^2}{2}}\right),`$ (3.2)
$`\left({\displaystyle \frac{\dot{a}}{a}}\right)^2+2{\displaystyle \frac{\ddot{a}}{a}}\mathrm{\Lambda }=8\pi G\left(p+{\displaystyle \frac{\dot{\varphi }^2}{2}}\right),`$ (3.3)
$`{\displaystyle \frac{\ddot{\varphi }}{\dot{\varphi }}}+3{\displaystyle \frac{\dot{a}}{a}}=0.`$ (3.4)
We consider the perfect fluid characterized by $`p=\gamma \rho `$, and we obtain the conservation law of the ordinary matter density by taking the linear combination of Eqs.(3.2), (3.3) and (3.4). From Eq.(3.4), we have another conservation law. Then we have the following two conservation laws
$`\rho =\rho _0a^{3(1+\gamma )},`$ (3.5)
$`\dot{\varphi }={\displaystyle \frac{k}{a^3}},`$ (3.6)
where $`\rho _0`$ and $`k`$ are integration constants. The equation to be solved becomes Eq.(3.2) with the conditions Eqs.(3.5) and (3.6). Substituting Eqs.(3.5) and (3.6) into Eq.(3.2), we have the equation of the form
$`\left({\displaystyle \frac{\dot{a}}{a}}\right)^2{\displaystyle \frac{\mathrm{\Lambda }}{3}}={\displaystyle \frac{8\pi G}{3}}\left({\displaystyle \frac{\rho _0}{a^{3(1+\gamma )}}}+{\displaystyle \frac{k^2}{2a^6}}\right).`$ (3.7)
As the mathematical problem, we can solve exactly in the $`\gamma =1`$ and $`\gamma =0`$ case, but $`\gamma =1`$ case is unphysical. Then we consider only the $`\gamma =0`$ case, that is, the matter dominant case. In this case, we obtain
$`\left({\displaystyle \frac{\dot{a}}{a}}\right)^2{\displaystyle \frac{4\pi G}{3}}\left(\dot{\varphi }+{\displaystyle \frac{\rho _0}{k}}\right)^2={\displaystyle \frac{\mathrm{\Lambda }}{3}}{\displaystyle \frac{4\pi G\rho _{0}^{}{}_{}{}^{2}}{3k^2}}`$ (3.8)
by using Eqs.(3.2), (3.5) and (3.6).
### 3.1 $`k^2\mathrm{\Lambda }>4\pi G\rho _{0}^{}{}_{}{}^{2}`$ case
In this case the cosmological and/or the scalar term are dominant, and we parametrize
$`{\displaystyle \frac{\dot{a}}{a}}=\sqrt{{\displaystyle \frac{\mathrm{\Lambda }}{3}}{\displaystyle \frac{4\pi G\rho _{0}^{}{}_{}{}^{2}}{3k^2}}}\mathrm{cosh}\mathrm{\Theta },`$ (3.9)
$`\dot{\varphi }={\displaystyle \frac{\rho _0}{k}}+\sqrt{{\displaystyle \frac{\mathrm{\Lambda }}{4\pi G}}{\displaystyle \frac{\rho _{0}^{}{}_{}{}^{2}}{k^2}}}\mathrm{sinh}\mathrm{\Theta }.`$ (3.10)
Substituting this parametrization into Eq.(3.4), we have
$`\dot{\mathrm{\Theta }}+\sqrt{{\displaystyle \frac{3\mathrm{\Lambda }}{1+A^2}}}(\mathrm{sinh}\mathrm{\Theta }A)=0,`$ (3.11)
where $`A^1=\sqrt{k^2\mathrm{\Lambda }/4\pi G\rho _{0}^{}{}_{}{}^{2}1}`$. Then we obtain
$`{\displaystyle \frac{1}{\sqrt{1+A^2}}}\mathrm{log}\left|{\displaystyle \frac{A\mathrm{tanh}(\mathrm{\Theta }/2)1+\sqrt{1+A^2}}{A\mathrm{tanh}(\mathrm{\Theta }/2)1\sqrt{1+A^2}}}\right|=\sqrt{{\displaystyle \frac{3\mathrm{\Lambda }}{1+A^2}}}(tt_0)`$ (3.12)
by using the formula
$`{\displaystyle \frac{d\mathrm{\Theta }}{\mathrm{sinh}\mathrm{\Theta }A}}={\displaystyle \frac{1}{\sqrt{1+A^2}}}\mathrm{log}\left|{\displaystyle \frac{A\mathrm{tanh}(\mathrm{\Theta }/2)1+\sqrt{1+A^2}}{A\mathrm{tanh}(\mathrm{\Theta }/2)1\sqrt{1+A^2}}}\right|.`$ (3.13)
And then we have the relation
$`{\displaystyle \frac{A\mathrm{tanh}(\mathrm{\Theta }/2)1+\sqrt{1+A^2}}{A\mathrm{tanh}(\mathrm{\Theta }/2)1\sqrt{1+A^2}}}=\mathrm{exp}\left(\sqrt{3\mathrm{\Lambda }}\left(tt_0\right)\right),`$ (3.14)
which gives the relation
$`\mathrm{tanh}(\mathrm{\Theta }/2)={\displaystyle \frac{1}{A}}+{\displaystyle \frac{\sqrt{1+A^2}}{A\mathrm{tanh}\left(\sqrt{3\mathrm{\Lambda }}\left(tt_0\right)/2\right)}}.`$ (3.15)
Using this relation, we can write $`\mathrm{cosh}\mathrm{\Theta }`$ and $`\mathrm{sinh}\mathrm{\Theta }`$ in the form
$`\mathrm{cosh}\mathrm{\Theta }={\displaystyle \frac{1+\mathrm{tanh}^2(\mathrm{\Theta }/2)}{1\mathrm{tanh}^2(\mathrm{\Theta }/2)}}={\displaystyle \frac{(1+A^2)\mathrm{cosh}\left(TT_0\right)\sqrt{1+A^2}\mathrm{sinh}\left(TT_0\right)}{A^2\mathrm{cosh}\left(TT_0\right)+\sqrt{1+A^2}\mathrm{sinh}\left(TT_0\right)}},`$ (3.16)
$`\mathrm{sinh}\mathrm{\Theta }={\displaystyle \frac{2\mathrm{tanh}(\mathrm{\Theta }/2)}{1\mathrm{tanh}^2(\mathrm{\Theta }/2)}}={\displaystyle \frac{A\left(1\mathrm{cosh}\left(TT_0\right)+\sqrt{1+A^2}\mathrm{sinh}\left(TT_0\right)\right)}{A^2\mathrm{cosh}\left(TT_0\right)+\sqrt{1+A^2}\mathrm{sinh}\left(TT_0\right)}},`$ (3.17)
where $`T=\sqrt{3\mathrm{\Lambda }}t`$ and $`T_0=\sqrt{3\mathrm{\Lambda }}t_0`$. Introducing $`\mathrm{\Theta }_0`$ through the relation $`\mathrm{cosh}\mathrm{\Theta }_0=\sqrt{1+A^2}/A`$, $`\mathrm{sinh}\mathrm{\Theta }_0=1/A`$, we can simplify the above expression in the form
$`\mathrm{cosh}\mathrm{\Theta }={\displaystyle \frac{\sqrt{1+A^2}\mathrm{cosh}(TT_0\mathrm{\Theta }_0)}{\mathrm{sinh}(TT_0\mathrm{\Theta }_0)A}},`$ (3.18)
$`\mathrm{sinh}\mathrm{\Theta }=A\left(1+{\displaystyle \frac{(A+1/A)}{\mathrm{sinh}(TT_0\mathrm{\Theta }_0)A}}\right).`$ (3.19)
Using Eqs.(3.9) and (3.18), we have
$`\mathrm{log}a`$ $`=`$ $`{\displaystyle \frac{da}{a}}=\sqrt{{\displaystyle \frac{\mathrm{\Lambda }}{3}}{\displaystyle \frac{4\pi G\rho _{0}^{}{}_{}{}^{2}}{3k^2}}}{\displaystyle 𝑑t\mathrm{cosh}\mathrm{\Theta }}`$ (3.20)
$`=`$ $`{\displaystyle \frac{1}{3}}\sqrt{1{\displaystyle \frac{4\pi G\rho _{0}^{}{}_{}{}^{2}}{k^2\mathrm{\Lambda }}}}\sqrt{1+A^2}{\displaystyle 𝑑T\frac{\mathrm{cosh}\left(TT_0\mathrm{\Theta }_0\right)}{\mathrm{sinh}\left(TT_0\mathrm{\Theta }_0\right)A}}`$
$`=`$ $`{\displaystyle \frac{1}{3}}\mathrm{log}\left|\mathrm{sinh}(TT_0\mathrm{\Theta }_0)A\right|+\mathrm{const}.,`$
where we use the relation $`4\pi G\rho _{0}^{}{}_{}{}^{2}/k^2\mathrm{\Lambda }=A^2/(1+A^2)`$. Therefore we have
$`a(t)=a_0\left(\mathrm{sinh}(TT_0\mathrm{\Theta }_0)A\right)^{1/3},`$ (3.21)
where $`a_0`$ is the constant.
Similarly, from Eqs.(3.10) and (3.19), we have
$`\varphi `$ $`=`$ $`{\displaystyle 𝑑t\left(\frac{\rho _0}{k}+\sqrt{\frac{\mathrm{\Lambda }}{4\pi G}\frac{\rho _{0}^{}{}_{}{}^{2}}{k^2}}\mathrm{sinh}\mathrm{\Theta }\right)}={\displaystyle \frac{\rho _0}{k\sqrt{3\mathrm{\Lambda }}}}{\displaystyle 𝑑T\left(1\frac{\mathrm{sinh}\mathrm{\Theta }}{A}\right)}`$ (3.22)
$`=`$ $`{\displaystyle \frac{\rho _0(1+A^2)}{kA\sqrt{3\mathrm{\Lambda }}}}{\displaystyle \frac{dT}{\mathrm{sinh}(TT_0\mathrm{\Theta }_0)A}}`$
$`=`$ $`\varphi _0+{\displaystyle \frac{1}{\sqrt{12\pi G}}}\mathrm{log}\left|{\displaystyle \frac{A\mathrm{tanh}\left((TT_0\mathrm{\Theta }_0)/2\right)+1\sqrt{1+A^2}}{A\mathrm{tanh}\left((TT_0\mathrm{\Theta }_0)/2\right)+1+\sqrt{1+A^2}}}\right|`$
$`=`$ $`\varphi _1+{\displaystyle \frac{1}{\sqrt{12\pi G}}}\mathrm{log}\left|{\displaystyle \frac{\mathrm{exp}(TT_0\mathrm{\Theta }_0)A\sqrt{1+A^2}}{\mathrm{exp}(TT_0\mathrm{\Theta }_0)A+\sqrt{1+A^2}}}\right|,`$
where $`\varphi _0`$ is the constant and $`\varphi _1`$ is given by $`\varphi _1=\varphi _0+{\displaystyle \frac{1}{\sqrt{12\pi G}}}\mathrm{log}\left|{\displaystyle \frac{1+A\sqrt{1+A^2}}{1+A+\sqrt{1+A^2}}}\right|`$.
The integration constant $`a_0`$ is not the independent integration constant but it can be expressed by $`k`$ and $`\rho _0`$. From Eqs.(3.6), (3.21) and (3.22), we have
$`\dot{\varphi }`$ $`=`$ $`{\displaystyle \frac{\rho _0(1+A^2)}{kA\left(\mathrm{sinh}(TT_0\mathrm{\Theta }_0)A\right)}}`$ (3.23)
$`=`$ $`{\displaystyle \frac{\rho _0(1+A^2)a_{0}^{}{}_{}{}^{3}}{kAa^3}}={\displaystyle \frac{k}{a^3}}.`$
which gives the relation $`a_0^3=k^2A/\rho _0(1+A^2)`$. This can be written in the form
$`k^2\mathrm{\Lambda }4\pi G\rho _{0}^{}{}_{}{}^{2}={\displaystyle \frac{\mathrm{\Lambda }^2a_{0}^{}{}_{}{}^{6}}{4\pi G}}.`$ (3.24)
Then we can obtain
$`a_0=\left(4\pi G({\displaystyle \frac{k^2}{\mathrm{\Lambda }}}{\displaystyle \frac{4\pi G\rho _0^2}{\mathrm{\Lambda }^2}})\right)^{1/6}.`$ (3.25)
Therefore we have the exact solution in the form
$`a(t)`$ $`=`$ $`\left(4\pi G({\displaystyle \frac{k^2}{\mathrm{\Lambda }}}{\displaystyle \frac{4\pi G\rho _0^2}{\mathrm{\Lambda }^2}})\right)^{1/6}\left(\mathrm{sinh}(TT_1)A\right)^{1/3},`$ (3.26)
$`\varphi (t)`$ $`=`$ $`\varphi _1+{\displaystyle \frac{1}{\sqrt{12\pi G}}}\mathrm{log}\left|{\displaystyle \frac{\mathrm{exp}(TT_1)A\sqrt{1+A^2}}{\mathrm{exp}(TT_1)A+\sqrt{1+A^2}}}\right|,`$ (3.27)
where
$`T=\sqrt{3\mathrm{\Lambda }}t,T_1=T_0\mathrm{\Theta }_0=\sqrt{3\mathrm{\Lambda }}t_1=\mathrm{const}.,`$
$`A^1=\sqrt{{\displaystyle \frac{k^2\mathrm{\Lambda }}{4\pi G\rho _0^2}}1},\varphi _1=\mathrm{const}..`$
### 3.2 $`4\pi G\rho _{0}^{}{}_{}{}^{2}>k^2\mathrm{\Lambda }`$ case
In this case the ordinary matter is dominant, and we parametrize
$`{\displaystyle \frac{\dot{a}}{a}}=\sqrt{{\displaystyle \frac{4\pi G\rho _{0}^{}{}_{}{}^{2}}{3k^2}}{\displaystyle \frac{\mathrm{\Lambda }}{3}}}\mathrm{sinh}\mathrm{\Theta },`$ (3.28)
$`\dot{\varphi }={\displaystyle \frac{\rho _0}{k}}+\sqrt{{\displaystyle \frac{\rho _{0}^{}{}_{}{}^{2}}{k^2}}{\displaystyle \frac{\mathrm{\Lambda }}{4\pi G}}}\mathrm{cosh}\mathrm{\Theta }.`$ (3.29)
Substituting this parametrization into Eq.(3.4), we have
$`\dot{\mathrm{\Theta }}+\sqrt{{\displaystyle \frac{3\mathrm{\Lambda }}{B^21}}}(\mathrm{cosh}\mathrm{\Theta }B)=0,`$ (3.30)
where $`B^1=\sqrt{1k^2\mathrm{\Lambda }/4\pi G\rho _{0}^{}{}_{}{}^{2}}`$. Then we obtain
$`{\displaystyle \frac{1}{\sqrt{B^21}}}\mathrm{log}\left|{\displaystyle \frac{1B+\sqrt{B^21}\mathrm{tanh}(\mathrm{\Theta }/2)}{1B\sqrt{B^21}\mathrm{tanh}(\mathrm{\Theta }/2)}}\right|=\sqrt{{\displaystyle \frac{3\mathrm{\Lambda }}{B^21}}}(tt_0),`$ (3.31)
by using the formula
$`{\displaystyle \frac{d\mathrm{\Theta }}{\mathrm{cosh}\mathrm{\Theta }B}}={\displaystyle \frac{1}{\sqrt{B^21}}}\mathrm{log}\left|{\displaystyle \frac{1B+\sqrt{B^21}\mathrm{tanh}(\mathrm{\Theta }/2)}{1B\sqrt{B^21}\mathrm{tanh}(\mathrm{\Theta }/2)}}\right|.`$
Then we have the relation
$`{\displaystyle \frac{1B+\sqrt{B^21}\mathrm{tanh}(\mathrm{\Theta }/2)}{1B\sqrt{B^21}\mathrm{tanh}(\mathrm{\Theta }/2)}}=\mathrm{exp}\left(\sqrt{3\mathrm{\Lambda }}(tt_0)\right),`$ (3.32)
where we take the branch of the logarithm in such a way as the scale factor of the universe behaves as the power law in time at the very early age of the universe. Then we have the relation
$`\mathrm{tanh}(\mathrm{\Theta }/2)=\sqrt{{\displaystyle \frac{B1}{B+1}}}\mathrm{tanh}\left(\sqrt{3\mathrm{\Lambda }}(tt_0)/2\right).`$ (3.33)
Using this relation, we can write $`\mathrm{sinh}\mathrm{\Theta }`$ and $`\mathrm{cosh}\mathrm{\Theta }`$ in the form
$`\mathrm{sinh}\mathrm{\Theta }`$ $`=`$ $`{\displaystyle \frac{2\mathrm{tanh}(\mathrm{\Theta }/2)}{1\mathrm{tanh}^2(\mathrm{\Theta }/2)}}={\displaystyle \frac{\sqrt{B^21}\mathrm{sinh}\left(TT_0\right)}{\mathrm{cosh}\left(TT_0\right)B}},`$ (3.34)
$`\mathrm{cosh}\mathrm{\Theta }`$ $`=`$ $`{\displaystyle \frac{1+\mathrm{tanh}^2(\mathrm{\Theta }/2)}{1\mathrm{tanh}^2(\mathrm{\Theta }/2)}}={\displaystyle \frac{B\mathrm{cosh}\left(TT_0\right)1}{\mathrm{cosh}\left(TT_0\right)B}},`$ (3.35)
where $`T=\sqrt{3\mathrm{\Lambda }}t`$ and $`T_0=\sqrt{3\mathrm{\Lambda }}t_0`$.
Using Eqs.(3.28) and (3.34), we have
$`\mathrm{log}a`$ $`=`$ $`{\displaystyle \frac{da}{a}}=\sqrt{{\displaystyle \frac{4\pi G\rho _{0}^{}{}_{}{}^{2}}{3k^2}}{\displaystyle \frac{\mathrm{\Lambda }}{3}}}{\displaystyle 𝑑t\mathrm{sinh}\mathrm{\Theta }}`$ (3.36)
$`=`$ $`{\displaystyle \frac{1}{3}}\sqrt{{\displaystyle \frac{4\pi G\rho _{0}^{}{}_{}{}^{2}}{k^2\mathrm{\Lambda }}}1}{\displaystyle 𝑑T\frac{\sqrt{B^21}\mathrm{sinh}\left(TT_0\right)}{\mathrm{cosh}\left(TT_0\right)B}}`$
$`=`$ $`{\displaystyle \frac{1}{3}}\mathrm{log}\left|\mathrm{cosh}\left(TT_0\right)B\right|+\mathrm{const}.,`$
where we use the relation $`4\pi G\rho _{0}^{}{}_{}{}^{2}/k^2\mathrm{\Lambda }=B^2/(B^21)`$. Therefore we have
$`a(t)=a_0\left(\mathrm{cosh}(TT_0)B\right)^{1/3},`$ (3.37)
where $`a_0`$ is the constant and $`(TT_0)=\sqrt{3\mathrm{\Lambda }}(tt_0)`$.
Similarly, from Eqs.(3.29) and (3.35), we have
$`\varphi `$ $`=`$ $`{\displaystyle 𝑑t\left(\sqrt{1\frac{k^2\mathrm{\Lambda }}{4\pi G\rho _{0}^{}{}_{}{}^{2}}}\mathrm{cosh}\mathrm{\Theta }1\right)}={\displaystyle \frac{\rho _0}{k\sqrt{3\mathrm{\Lambda }}}}{\displaystyle 𝑑T\left(\frac{\mathrm{cosh}\mathrm{\Theta }}{B}1\right)}`$ (3.38)
$`=`$ $`{\displaystyle \frac{\rho _0(B^21)}{kB\sqrt{3\mathrm{\Lambda }}}}{\displaystyle \frac{dT}{\mathrm{cosh}\left(TT_0\right)B}}`$
$`=`$ $`\varphi _0+{\displaystyle \frac{\rho _0\sqrt{B^21}}{kB\sqrt{3\mathrm{\Lambda }}}}\mathrm{log}\left|{\displaystyle \frac{1B+\sqrt{B^21}\mathrm{tanh}\left((TT_0)/2\right)}{1B\sqrt{B^21}\mathrm{tanh}\left((TT_0)/2\right)}}\right|`$
$`=`$ $`\varphi _1+{\displaystyle \frac{1}{\sqrt{12\pi G}}}\mathrm{log}\left|{\displaystyle \frac{\mathrm{exp}\left(TT_0\right)B\sqrt{B^21}}{\mathrm{exp}\left(TT_0\right)B+\sqrt{B^21}}}\right|,`$
where $`\varphi _0`$ is the constant and $`\varphi _1`$ is given by $`\varphi _1=\varphi _0+{\displaystyle \frac{1}{\sqrt{12\pi G}}}\mathrm{log}\left|{\displaystyle \frac{1+B\sqrt{B^21}}{1+B+\sqrt{B^21}}}\right|`$.
The integration constant $`a_0`$ is not the independent integration constant but it can be expressed by $`k`$ and $`\rho _0`$. From Eqs.(3.6), (3.37) and (3.38), we have
$`\dot{\varphi }={\displaystyle \frac{(B^21)\rho _0}{kB\left(\mathrm{cosh}(TT_0)B\right)}}`$
$`={\displaystyle \frac{(B^21)\rho _0a_{0}^{}{}_{}{}^{3}}{kBa^3}}={\displaystyle \frac{k}{a^3}},`$ (3.39)
which gives the relation $`k^2=(B^21)\rho _0a_{0}^{}{}_{}{}^{3}/B`$. This can be written in the form
$`4\pi G\rho _{0}^{}{}_{}{}^{2}k^2\mathrm{\Lambda }={\displaystyle \frac{\mathrm{\Lambda }^2a_{0}^{}{}_{}{}^{6}}{4\pi G}}.`$ (3.40)
Then we can obtain
$`a_0=\left(4\pi G({\displaystyle \frac{4\pi G\rho _{0}^{}{}_{}{}^{2}}{\mathrm{\Lambda }^2}}{\displaystyle \frac{k^2}{\mathrm{\Lambda }}})\right)^{1/6}.`$ (3.41)
Therefore we have the exact solution in the form
$`a(t)`$ $`=`$ $`\left(4\pi G({\displaystyle \frac{4\pi G\rho _{0}^{}{}_{}{}^{2}}{\mathrm{\Lambda }^2}}{\displaystyle \frac{k^2}{\mathrm{\Lambda }}})\right)^{1/6}\left(\mathrm{cosh}(TT_0)B\right)^{1/3},`$ (3.42)
$`\varphi (t)`$ $`=`$ $`\varphi _0+{\displaystyle \frac{1}{\sqrt{12\pi G}}}\mathrm{log}\left|{\displaystyle \frac{\mathrm{exp}\left(TT_0\right)B\sqrt{B^21}}{\mathrm{exp}\left(TT_0\right)B+\sqrt{B^21}}}\right|,`$ (3.43)
where
$`T=\sqrt{3\mathrm{\Lambda }}t,T_0=\sqrt{3\mathrm{\Lambda }}t_0=\mathrm{const}.,`$
$`B^1=\sqrt{1{\displaystyle \frac{k^2\mathrm{\Lambda }}{4\pi G\rho _{0}^{}{}_{}{}^{2}}}},\varphi _0=\mathrm{const}..`$
### 3.3 $`4\pi G\rho _{0}^{}{}_{}{}^{2}=k^2\mathrm{\Lambda }`$ case
For the completeness of the solution, we give the exact solution in this case. As the method to solve the equation is similar, we give only the result. The solution is given by
$`a(t)`$ $`=`$ $`\left({\displaystyle \frac{4\pi G\rho _0}{\mathrm{\Lambda }}}\right)^{1/3}\left(\mathrm{exp}(TT_0)1\right)^{1/3},`$ (3.44)
$`\varphi (t)`$ $`=`$ $`\varphi _0+{\displaystyle \frac{1}{\sqrt{12\pi G}}}\mathrm{log}\left|1\mathrm{exp}\left((TT_0)\right)\right|,`$ (3.45)
where
$`T=\sqrt{3\mathrm{\Lambda }}t,T_0=\sqrt{3\mathrm{\Lambda }}t_0=\mathrm{const}.,\varphi _0=\mathrm{const}..`$
### 3.4 Special Limiting Case
i) $`\rho _0=0`$ case (no ordinary matter)
In case there is no ordinary matter $`\rho _00`$, which corresponds to $`A\sqrt{4\pi G\rho _{0}^{}{}_{}{}^{2}/k^2\mathrm{\Lambda }}`$, we have the expression
$`a(t)`$ $`=`$ $`\left({\displaystyle \frac{4\pi Gk^2}{\mathrm{\Lambda }}}\right)^{1/6}\left(\mathrm{sinh}(TT_1)\right)^{1/3},`$ (3.46)
$`\varphi (t)`$ $`=`$ $`\varphi _1+{\displaystyle \frac{1}{\sqrt{12\pi G}}}\mathrm{log}\left|\mathrm{tanh}\left((TT_1)/2\right)\right|.`$ (3.47)
ii) $`k=0`$ case (no scalar matter): Lemaître universe
In case there is no scalar matter $`k0`$, which corresponds to $`B1`$, we have the expression
$`a(t)`$ $`=`$ $`\left({\displaystyle \frac{4\pi G\rho _0}{\mathrm{\Lambda }}}\right)^{1/3}\left(\mathrm{cosh}(TT_1)1\right)^{1/3},`$ (3.48)
$`\varphi (t)`$ $`=`$ $`\varphi _1.`$ (3.49)
iii) $`\mathrm{\Lambda }=0`$ case (no cosmological constant)
In case there is no cosmological term $`\mathrm{\Lambda }0`$, which corresponds to $`B1+k^2\mathrm{\Lambda }/8\pi G\rho _{0}^{}{}_{}{}^{2}`$, we have the expression
$`a(t)`$ $`=`$ $`\underset{\mathrm{\Lambda }0}{lim}\left({\displaystyle \frac{4\pi G\rho _0}{\mathrm{\Lambda }}}\right)^{1/3}\left({\displaystyle \frac{3\mathrm{\Lambda }\left(tt_0\right)^2}{2}}{\displaystyle \frac{k^2\mathrm{\Lambda }}{8\pi G\rho _{0}^{}{}_{}{}^{2}}}\right)^{1/3}`$ (3.50)
$`=`$ $`\left(6\pi G\rho _0\right)^{1/3}\left(\left(tt_0\right)^2{\displaystyle \frac{k^2}{12\pi G\rho _{0}^{}{}_{}{}^{2}}}\right)^{1/3}`$
$`\varphi (t)`$ $`=`$ $`\underset{\mathrm{\Lambda }0}{lim}\left\{\varphi _0+{\displaystyle \frac{1}{\sqrt{12\pi G}}}\mathrm{log}\left|{\displaystyle \frac{\sqrt{3\mathrm{\Lambda }}(tt_0)k\sqrt{\mathrm{\Lambda }}/\sqrt{4\pi G\rho _0^2}}{\sqrt{3\mathrm{\Lambda }}(tt_0)+k\sqrt{\mathrm{\Lambda }}/\sqrt{4\pi G\rho _0^2}}}\right|\right\}`$ (3.51)
$`=`$ $`\varphi _0+{\displaystyle \frac{1}{\sqrt{12\pi G}}}\mathrm{log}\left|{\displaystyle \frac{(tt_0)k/\sqrt{12\pi G\rho _0^2}}{(tt_0)+k/\sqrt{12\pi G\rho _0^2}}}\right|.`$
## 4 Effect on the Gravitational Potential
In this section, we calculate the effect of the cosmological constant and the scalar matter to the gravitational potential. For the special case of i) no scalar matter or ii) no cosmological term, the exact solutions are well-known.
### 4.1 Exact solution for special cases
i) No scalar matter case
In this case, we take the standard metric in the form
$$ds^2=h(r)^2dt^2+f(r)^2dr^2+\left[r^2(d\theta ^2+\mathrm{sin}\theta ^2d\phi ^2)\right],$$
(4.1)
and the exact solution is given by
$`h(r)^2=1r_0/r\mathrm{\Lambda }r^2/3,`$ (4.2)
$`f(r)^2={\displaystyle \frac{1}{1r_0/r\mathrm{\Lambda }r^2/3}}.`$ (4.3)
ii) No cosmological term case
In this case, we take the isotropic metric in the form
$$ds^2=h_1(r)^2dt^2+f_1(r)^2\left[dr^2+r^2(d\theta ^2+\mathrm{sin}\theta ^2d\phi ^2)\right],$$
(4.4)
and the exact solution is given by
$`\varphi (r)`$ $`=`$ $`\varphi _0\mathrm{log}\left({\displaystyle \frac{rr_0}{r+r_0}}\right),`$ (4.5)
$`h_1(r)^2`$ $`=`$ $`\left({\displaystyle \frac{rr_0}{r+r_0}}\right)^{2C},`$ (4.6)
$`f_1(r)^2`$ $`=`$ $`\left(1{\displaystyle \frac{r_{0}^{}{}_{}{}^{2}}{r^2}}\right)^2\left({\displaystyle \frac{r+r_0}{rr_0}}\right)^{2C},`$ (4.7)
where
$`\varphi _0`$ $`=`$ $`\sqrt{{\displaystyle \frac{2(1C^2)}{8\pi G}}},C=\mathrm{const}..`$
### 4.2 Cosmological term and scalar matter co-existing case
When the cosmological term and scalar matter co-exist, we cannot solve analytically, and we calculate the effect on the gravitational potential approximately. For this purpose, we take the standard metric Eq.(4.1) and the equations of motion Eqs.(2.2) and (2.3) are given by
$`{\displaystyle \frac{4rf^{^{}}}{f^2}}+2f{\displaystyle \frac{2}{f}}2\mathrm{\Lambda }r^2f`$ $`=`$ $`{\displaystyle \frac{8\pi r^2G\varphi _{}^{^{}}{}_{}{}^{2}}{f}},`$ (4.8)
$`{\displaystyle \frac{4rh^{^{}}}{f^2}}{\displaystyle \frac{2h}{f^2}}+2h2\mathrm{\Lambda }r^2h`$ $`=`$ $`{\displaystyle \frac{8\pi r^2Gh\varphi _{}^{^{}}{}_{}{}^{2}}{f^2}},`$ (4.9)
$`\left({\displaystyle \frac{r^2h\varphi ^{^{}}}{f}}\right)^{^{}}`$ $`=`$ $`0.`$ (4.10)
These can be rewritten into the form
$`{\displaystyle \frac{f^{^{}}}{f}}{\displaystyle \frac{h^{^{}}}{h}}+{\displaystyle \frac{f^2}{r}}{\displaystyle \frac{1}{r}}`$ $`=`$ $`\mathrm{\Lambda }rf^2,`$ (4.11)
$`{\displaystyle \frac{f^{^{}}}{f}}+{\displaystyle \frac{h^{^{}}}{h}}`$ $`=`$ $`4\pi Gr\varphi _{}^{^{}}{}_{}{}^{2},`$ (4.12)
$`\left({\displaystyle \frac{r^2h\varphi ^{^{}}}{f}}\right)^{^{}}`$ $`=`$ $`0.`$ (4.13)
From Eq.(4.13), we have $`\varphi ^{^{}}=\alpha f/r^2h`$ with constant $`\alpha `$. Then we have
$`\left(\mathrm{log}f/h\right)^{^{}}={\displaystyle \frac{f^{^{}}}{f}}{\displaystyle \frac{h^{^{}}}{h}}`$ $`=`$ $`{\displaystyle \frac{f^2}{r}}+{\displaystyle \frac{1}{r}}+\mathrm{\Lambda }rf^2,`$ (4.14)
$`\left(\mathrm{log}fh\right)^{^{}}={\displaystyle \frac{f^{^{}}}{f}}+{\displaystyle \frac{h^{^{}}}{h}}`$ $`=`$ $`{\displaystyle \frac{4\pi G\alpha ^2f^2}{r^3h^2}}.`$ (4.15)
We introduce the new variables $`X`$, $`Y`$ in the form $`\mathrm{exp}(X)=fh`$, $`\mathrm{exp}(Y)=f/h`$, then the above equation becomes in the form
$`Y^{^{}}={\displaystyle \frac{\mathrm{exp}(X+Y)}{r}}+{\displaystyle \frac{1}{r}}+\mathrm{\Lambda }r\mathrm{exp}(X+Y),`$ (4.16)
$`X^{^{}}={\displaystyle \frac{4\pi G\alpha ^2\mathrm{exp}(2Y)}{r^3}}.`$ (4.17)
For $`\alpha =0`$ case (no scalar matter case), we have the solution
$`X=0,\mathrm{exp}(Y)={\displaystyle \frac{1}{1r_0/r\mathrm{\Lambda }r^2/3}},`$ (4.18)
which is the exact solution to Eqs.(4.2) and (4.3). Then we calculate the gravitational potential by considering the region of $`r`$ where $`r_0/r,\mathrm{\Lambda }r^2,4\pi G\alpha ^2/r^21`$. In this approximation, we have
$`Y`$ $``$ $`{\displaystyle \frac{r_0}{r}}+{\displaystyle \frac{\mathrm{\Lambda }r^2}{3}},`$ (4.19)
$`X`$ $``$ $`2\pi G\alpha ^2\left({\displaystyle \frac{1}{r_{1}^{}{}_{}{}^{2}}}{\displaystyle \frac{1}{r^2}}\right)`$ (4.20)
from Eqs.(4.17) and (4.18) where $`r_1`$ is constant.
In order to find the solution for $`\alpha 0`$, we put $`r_0(=\mathrm{const}.)r_0(r)(\mathrm{function}\mathrm{of}r)`$. Then Eq.(4.16) becomes
$`Y^{^{}}`$ $`=`$ $`\left({\displaystyle \frac{r_0(r)}{r}}+{\displaystyle \frac{\mathrm{\Lambda }r^2}{3}}\right)^{^{}}={\displaystyle \frac{r_0}{r^2}}+{\displaystyle \frac{r_{0}^{}{}_{}{}^{^{}}}{r}}+{\displaystyle \frac{2\mathrm{\Lambda }r}{3}}`$ (4.21)
$``$ $`{\displaystyle \frac{r_0}{r^2}}+{\displaystyle \frac{2\mathrm{\Lambda }r}{3}}{\displaystyle \frac{X}{r}}`$
which gives $`r_{0}^{}{}_{}{}^{^{}}=X`$. Using Eq.(4.20), we have
$`r_0(r)=r_22\pi G\alpha ^2\left({\displaystyle \frac{r}{r_{1}^{}{}_{}{}^{2}}}+{\displaystyle \frac{1}{r}}\right),`$ (4.22)
where $`r_2`$ is constant.
The gravitational potential $`\mathrm{\Phi }`$ is given by
$`g_{00}`$ $`=`$ $`(1+2\mathrm{\Phi })=\mathrm{exp}(XY)`$ (4.23)
$``$ $`(1{\displaystyle \frac{r_0(r)}{r}}{\displaystyle \frac{\mathrm{\Lambda }r^2}{3}}+X)`$
$``$ $`(1+{\displaystyle \frac{4\pi G\alpha ^2}{r_1^2}}{\displaystyle \frac{r_2}{r}}{\displaystyle \frac{\mathrm{\Lambda }r^2}{3}}),`$
which gives $`\mathrm{\Phi }=2\pi G\alpha ^2/r_1^2r_2/2r\mathrm{\Lambda }r^2/6`$. Therefore the scalar matter does not contribute to the gravitational force $`F_r=\mathrm{\Phi }/r=r_2/2r^2+\mathrm{\Lambda }r/3`$ within our approximation. The cosmological term contribute to the repulsive force within the approximation.
## 5 Summary and Discussion
We consider the scalar field as the candidate of the dark matter. Then, in order to give the standard scenario of the astrophysics, we study the Einstein theory with minimally coupled scalar field and the cosmological constant. We have studied various classical solutions with minimally coupled scalar and the cosmological term in the cosmological, the galactic or solar scale. We obtained the exact solution in the cosmology scale, where the scale factor expand in the power law in the first beginning and then expand exponentially. In the galactic or solar scale, we cannot find the exact solution, and examine the contribution from the scalar field to the gravitational potential and find that the scalar field does not contribute to the gravitational force within our approximation. In this way, in the cosmological scale, the scalar field play the role of the dark matter in some sense. While, in the galactic or solar scale, the scalar field does not pay the role of the dark matter.
For the ordinary matter, we first start from the classical Lagrangian and quantize the field and treat it as the classical smeared matter and make the perfect fluid approximation. While, in our approach, we treat the scalar field as the classical field in the same level as the classical gravitational field. If the metric is homogeneous, it may give the same effect whether we treat the scalar field as the classical field or the quantized and classically smeared matter. But if the metric is not homogeneous and is space-dependent, there is the quite big gap in the step of the quantization and the treatment of the classically smeared matter. In this sense, the scalar field may give the contribution to the gravitational force if we treat the scalar field as the quantized matter field.
Acknowledgement:
Two of us (K.S. and K.U.) are grateful to the academic research funds of Tezukayama University. |
warning/0003/nlin0003056.html | ar5iv | text | # Phase Dynamics of Nearly Stationary Patterns in Activator-Inhibitor Systems
## I Introduction
Studies of stationary patterns in activator-inhibitor systems have focused primarily on localized structures such as pulses and spots in excitable and bistable media TyKe:88 ; KeOs:89 ; meron:92 ; KM:94 ; LeSw:95 ; MuOs:96 ; WSBOP:96 ; SOBP:97 , and periodic patterns near a Turing bifurcation CDBD:90 ; OuSw:91n ; AAP:97 . Localized structures have instabilities to traveling patterns, breathing motion, and transverse deformations KeOs:89 ; OMK:89 ; HaMe:94a ; GMP:96 . Periodic patterns have been analyzed near the onset of a Turing instability and also near the codimension-two point of a Turing instability and a Hopf bifurcation RoMe:92 ; HBP:93 ; PDDK:93 ; DLDB:96 ; OB:98 . But very few studies have explored instabilities of periodic (nonlocal) stationary patterns in excitable and bistable media, or of periodic stationary patterns far beyond the Turing instability DK:89 ; Osipov:96 . The latter case includes pattern formation studies on the CIMA chemical reaction PDDK:93 ; OuSw:91 ; DP:94 .
In this paper we study instabilities of stationary periodic patterns by deriving a Cross-Newell phase equation KBBC:82 ; CN:84 ; NPBEI:96 . The derivation is not restricted to the immediate neighborhood of a Turing instability and applies to periodic patterns with space-scale separation that arise far from onset or in excitable and bistable media. The Cross-Newell equation was originally derived in the context of fluid dynamics and has recently been applied in a laser system LMN:96 .
We choose to study the FitzHugh-Nagumo (FHN) equations, a canonical model for activator-inhibitor systems,
$`{\displaystyle \frac{u}{t}}`$ $`=`$ $`uu^3v+^2u,`$ (1)
$`{\displaystyle \frac{v}{t}}`$ $`=`$ $`ϵ(ua_1va_0)+\delta ^2v.`$
Here, $`u`$ is the activator and $`v`$ the inhibitor. The parameters $`a_0`$ and $`a_1`$ can be chosen so that the FHN model (1) represents an excitable medium, a bistable medium, or a system with a Turing instability HaMe:94a . All three cases support stationary periodic solutions for $`\delta `$ sufficiently large.
In Section II we derive a phase equation describing weak modulations of periodic stripe pattern in the FHN model. In Section III we evaluate stability thresholds for Eckhaus and zig-zag instabilities and for a transition from stationary to traveling patterns. These thresholds suggest a number of spatial or spatio-temporal behaviors which we test in Section IV with numerical solutions of Eqs. (1).
## II The Phase Equation
Let $`u_0(\theta ;k)=u_0(\theta +2\pi ;k)`$, $`v_0(\theta ;k)=v_0(\theta +2\pi ;k)`$ be a stationary periodic solution of Eqs. (1) with phase $`\theta `$ and wavenumber $`k`$. We consider weak spatial modulations of this periodic pattern and assume that those modulations have a length scale $`L`$ that is much larger than the wavelength $`1/k`$. The ratio of the length scales $`\lambda =1/(kL)`$ can then be used as a small parameter to write modulated solutions as an asymptotic expansion about the periodic solution
$`u(\theta ,𝐑,T)=u_0(\theta ;k)+\lambda u_1(\theta ,𝐑,T)+\lambda ^2u_2(\theta ,𝐑,T)+\mathrm{}`$
$`v(\theta ,𝐑,T)=v_0(\theta ;k)+\lambda v_1(\theta ,𝐑,T)+\lambda ^2v_2(\theta ,𝐑,T)+\mathrm{}`$ (2)
where $`𝐑=\lambda 𝐫`$ and $`T=\lambda ^2t`$ are slow space and time variables. The phase $`\theta `$ in Eq. (2) is an undetermined function of space and time and $`k=|𝐤|=|\theta |`$ is the local wavenumber. Our objective is to derive an equation for the slow phase
$$\mathrm{\Theta }(𝐑,T):=\lambda \theta (𝐫,𝐑,T).$$
In terms of this phase the local wavevector is
$$𝐤(𝐑,T)=_R\mathrm{\Theta }.$$
Inserting the expansions (2) in Eqs. (1) we find at order unity
$`u_0u_0^3v_0+k^2{\displaystyle \frac{^2u_0}{\theta ^2}}`$ $`=`$ $`0,`$ (3a)
$`ϵ(u_0a_1v_0a_0)+\delta k^2{\displaystyle \frac{^2v_0}{\theta ^2}}`$ $`=`$ $`0,`$ (3b)
where $`k^2=𝐤𝐤`$. At order $`\lambda `$
$`\left(k^2{\displaystyle \frac{^2}{\theta ^2}}+13u_0^2\right)u_1v_1`$ $`=`$ $`𝒟{\displaystyle \frac{u_0}{\theta }},`$ (4a)
$`ϵu_1+\left(\delta k^2{\displaystyle \frac{^2}{\theta ^2}}ϵa_1\right)v_1`$ $`=`$ $`𝒟{\displaystyle \frac{v_0}{\theta }},`$ (4b)
where
$`𝒟={\displaystyle \frac{\mathrm{\Theta }}{T}}_R𝐤2𝐤_R.`$ (5)
Projecting the right hand side of (4) onto $`(_\theta u_0,ϵ^1_\theta v_0)`$, the solution of the adjoint problem, produces the phase equation
$$\tau \frac{\mathrm{\Theta }}{T}=_R(𝐤B),$$
where
$`\tau `$ $`=`$ $`<(_\theta u_0)^2>ϵ^1<(_\theta v_0)^2>,`$ (6)
$`B`$ $`=`$ $`<(_\theta u_0)^2>+\delta ϵ^1<(_\theta v_0)^2>.`$ (7)
In these equations $`<(.)>:=\frac{1}{2\pi }_0^{2\pi }(.)d\theta `$.
The quantities $`B`$ and $`\tau `$ contain information about various instabilities of the periodic stripe pattern. The condition $`\frac{d}{dk}[kB(k)]=0`$ implies the onset of an Eckhaus instability and the condition $`B=0`$ the onset of a zigzag instability CN:84 . In Appendix A we show that the condition $`\tau =0`$ indicates a transition to traveling waves.
To implement these conditions we need to solve Eqs. (3) for the periodic solution $`(u_0,v_0)`$. For parameter values that satisfy $`ϵ/\delta :=\mu 1`$ an approximate solution can be computed as shown in Appendix B. Using this solution to calculate $`\tau `$ and $`B`$, as shown in Appendix C, gives the following expressions:
$`\tau `$ $`=`$ $`{\displaystyle \frac{2\sqrt{2}}{3\pi k}}{\displaystyle \frac{v_{}}{q\pi k\eta }}\beta (\mathrm{\Lambda }_{})\gamma (\mathrm{\Lambda }_{},\mathrm{\Lambda }_+),`$
$`B`$ $`=`$ $`{\displaystyle \frac{2\sqrt{2}}{3\pi k}}{\displaystyle \frac{v_{}}{q\pi k\sqrt{\mu }}}\beta (\mathrm{\Lambda }_{})\gamma (\mathrm{\Lambda }_{},\mathrm{\Lambda }_+),`$ (8)
$$\mathrm{\Lambda }_{}+\mathrm{\Lambda }_+=\frac{2\pi \sqrt{\mu }}{k},$$
(9)
$$v_+\beta (\mathrm{\Lambda }_+)+v_{}\beta (\mathrm{\Lambda }_{})=0,$$
(10)
where $`\mu =ϵ/\delta `$, $`\eta =\sqrt{ϵ\delta }`$, $`v_\pm =(\pm 1a_0)/q^2`$, $`q^2=a_1+1/2`$,
$$\beta (x)=\mathrm{coth}qx\mathrm{csch}qx,$$
(11)
and
$`\gamma (\mathrm{\Lambda }_{},\mathrm{\Lambda }_+)`$ $`=`$ $`1+{\displaystyle \frac{1}{2}}(1+a_0)q\mathrm{\Lambda }_{}\mathrm{csch}q\mathrm{\Lambda }_{}`$ (12)
$`+{\displaystyle \frac{1}{2}}(1a_0)q\mathrm{\Lambda }_+\mathrm{csch}q\mathrm{\Lambda }_+.`$
The quantities $`\mathrm{\Lambda }_+`$ and $`\mathrm{\Lambda }_{}`$ denote the widths of domains with high and low values of $`u`$ and $`v`$, respectively. The width is measured with respect to the spatial coordinate $`z=\frac{\sqrt{\mu }}{k}\theta `$ (see Appendix B). Given $`k`$, Eqs. (9) and (10) can be solved for $`\mathrm{\Lambda }_+(k)`$ and $`\mathrm{\Lambda }_{}(k)`$. Using these solutions in Eq. (8) graphs of $`\tau `$ and $`kB`$ as functions of $`k`$ can be produced.
## III Stability Thresholds
Explicit forms for $`\tau (k)`$ and $`B(k)`$ are available in the symmetric case, $`a_0=0`$, where $`\mathrm{\Lambda }_+=\mathrm{\Lambda }_{}=\frac{\pi \sqrt{\mu }}{k}`$:
$`\tau (k)={\displaystyle \frac{1}{\pi k\eta _cq^3}}\left[1{\displaystyle \frac{\eta _c}{\eta }}f(\pi q\sqrt{\mu }/k)\right],`$
$`B(k)={\displaystyle \frac{1}{\pi k\eta _cq^3}}\left[1{\displaystyle \frac{\eta _c}{\sqrt{\mu }}}f(\pi q\sqrt{\mu }/k)\right],`$ (13)
where $`\eta _c=\frac{3}{2\sqrt{2}q^3}`$ and
$$f(x)=(1x\mathrm{csch}x)(\mathrm{coth}x\mathrm{csch}x).$$
(14)
Figure 1 shows graphs of $`\tau (k)`$ and $`kB(k)`$ for a bistable medium obtained with Eqs. (III) (thick lines) and with Eqs. (6) and (7) using numerically calculated solutions $`u_0,v_0`$ (circles). A very good agreement is obtained within the validity range of the analysis, $`k𝒪(\sqrt{\mu })1`$. For $`k𝒪(1)`$ the deviations become large. In particular the minimum of $`kB(k)`$ which designates the Eckhaus instability threshold, is not reproduced by the analytical form (13).
The instability to traveling waves occurs at $`\tau =0`$ or at
$$ϵ=\eta _c^2f^2(\pi q\sqrt{\mu }/k)\delta ^1.$$
(15)
The zigzag instability occurs at $`B=0`$ or at
$$ϵ=\eta _c^2f^2(\pi q\sqrt{\mu }/k)\delta .$$
(16)
The condition $`\frac{d}{dk}(kB)=0`$ for the Eckhaus instability becomes
$$\frac{df}{dx}|_{x=\pi q\sqrt{\mu }/k}=0.$$
Consider first the limit $`k0`$ in which the periodic pattern approaches an array of isolated front structures. In this limit $`f(\pi q\sqrt{\mu }/k)1`$ and the condition for the onset of traveling waves becomes $`ϵ=\eta _c^2\delta ^1`$. This is precisely the nonequilibrium Ising-Bloch (NIB) bifurcation point, where a stationary front loses stability to a pair of counter-propagating fronts. The condition for the zigzag instability becomes $`ϵ=\eta _c^2\delta `$. This is the threshold for the transverse front instability HaMe:94c .
The neutral stability curves for a bistable medium corresponding to Eqs. (15) and (16) are shown in Figs. 2a and 2b for fixed $`\delta `$ and $`ϵ`$, respectively. They imply that high wavenumber stationary planar patterns are stabilized against zigzag and traveling wave instabilities. Notice that for $`\delta =1`$ the neutral stability curves $`\tau =0`$ and $`B=0`$ coincide (see Eqs. (15) and (16) or Fig. 2b). For $`\delta >1`$, upon decreasing the wavenumber at constant $`ϵ`$, a high wavenumber pattern is destabilized to a zigzag pattern, whereas for $`\delta <1`$ the destabilization is to traveling waves. Similar neutral stability curves are found for the nonsymmetric case, $`a_00`$, for excitable media and for systems (far) beyond the Turing instability.
## IV Comparisons with Numerical Solutions
We have computed numerical solutions of Eqs. (1) to test the stabilization of zigzag and traveling-wave instabilities at high wavenumbers. Figure 3 shows a low wavenumber zigzag pattern and a high wavenumber planar pattern computed for the same parameter values. This behavior is well known in other contexts CrHo:93 . The zigzag instability is a mechanism by which the system locally increases the wavenumber. Figure 4 shows coexistence of a low wavenumber traveling wave and a high wavenumber stationary pattern. These numerical results are for a bistable system but similar results are found for excitable and Turing unstable systems. Coexistence of stationary and traveling waves has been found in experiments on the CIMA reaction PDDK:93 ; DP:94 . and analyzed using different theoretical approaches DK:89 ; IO:92 ; KO:95 ; Osipov:96 .
We have also tested the condition for the Eckhaus instability in a bistable system using numerical computations of $`\tau `$ and $`B`$. Choosing wavenumbers $`k>k_c`$ where $`k_c`$ corresponds to the minimum of $`kB`$, we found initial periodic patterns either collapse to uniform states or to a lower wavenumber pattern through phase slips. Fig. 5 demonstrates these two cases. Similar conclusions hold for excitable systems. An unstable Turing pattern, on the other hand, always converges to a lower wavenumber pattern since the single uniform state is unstable.
## V Conclusion
We have shown that the Cross-Newell phase equation provides a powerful tool for studying instabilities of stationary periodic patterns in activator-inhibitor systems. The equation contains information not only on the Eckhaus and zigzag instabilities, but also on the destabilization of stationary periodic patterns to traveling waves. The same equation applies to bistable, excitable, and Turing unstable systems. Equations of that kind should prove useful in identifying parameters and initial conditions where zigzag and Eckhaus instabilities couple to traveling wave modes. Such coupling may lead to complex spatiotemporal behavior analogous to the coupling of the NIB front bifurcation to a transverse front instability HaMe:94c ; HaMe:94b
###### Acknowledgements.
This study was supported in part by grant No 95-00112 from the US-Israel Binational Science Foundation (BSF) and by the Department of Energy, under contract W-7405-ENG-36.
## Appendix A The Meaning of $`\tau =0`$
We show here that the condition $`\tau =0`$ defines the critical value of $`ϵ`$ at which traveling solutions bifurcate from the stationary solution. We look for traveling solutions $`u(\theta ),v(\theta )`$ of Eqs. (1), where $`\theta =kx\omega t`$, that bifurcate from the stationary solution branch $`\omega =0`$ at some $`ϵ=ϵ_c`$. Near the bifurcation where $`\omega 1`$ we can expand the traveling solutions as power series in $`\omega `$ around the stationary solution $`u_0,v_0`$:
$`u(\theta )=u_0(\theta )+\omega u_1(\theta )+\mathrm{},`$
$`v(\theta )=v_0(\theta )+\omega v_1(\theta )+\mathrm{},`$ (17)
Expanding $`ϵ`$ as
$$ϵ=ϵ_c+ϵ_1\omega +\mathrm{},$$
(18)
and using these expansions in Eqs. (1) we find at order $`\omega `$
$`\left(k^2{\displaystyle \frac{^2}{\theta ^2}}+13u_0^2\right)u_1v_1`$ $`=`$ $`u_{0}^{}{}_{}{}^{},`$
$`ϵ_cu_1+\left(\delta k^2{\displaystyle \frac{^2}{\theta ^2}}ϵa_1\right)v_1`$ $`=`$ $`v_0^{}ϵ_1(u_0a_1v_0a_0).`$
Projecting the right hand side onto $`(u_0^{},ϵ_c^1v_0^{})`$ gives
$$ϵ_c=\frac{<v_{0}^{}{}_{}{}^{2}>}{<u_{0}^{}{}_{}{}^{2}>},$$
(19)
where we used Eq. (3b) and switched to the notation of a prime for the derivative with respect to the single argument $`\theta `$. Using the definition (6) of $`\tau `$ and (19) we find
$$\tau =\left(1\frac{ϵ_c}{ϵ}\right)<u_{0}^{}{}_{}{}^{2}>.$$
(20)
Thus, $`\tau =0`$ implies $`ϵ=ϵ_c`$ or the onset of traveling solutions.
## Appendix B Approximate Stationary Solution
For $`\mu =ϵ/\delta 1`$ a singular perturbation approach can be used to approximate the stationary solution $`u_0(\theta ),v_0(\theta )`$. Rescaling the space coordinate as $`z=\frac{\sqrt{\mu }}{k}\theta `$, Eqs. (3) become
$`u_0u_0^3v_0+\mu u_0^{\prime \prime }`$ $`=`$ $`0,`$
$`u_0a_1v_0a_0+v_0^{\prime \prime }`$ $`=`$ $`0,`$
where the prime denotes now the derivative with respect to $`z`$. Since the small parameter $`\mu `$ multiplies the second derivative term $`u_0^{\prime \prime }`$, two types of spatial regions can be distinguished. Outer regions where $`u_0(z)`$ varies on a scale of order unity and the term $`\mu u_0^{\prime \prime }`$ is negligible, and inner regions where $`u_0(z)`$ varies on a very short scale of order $`\sqrt{\mu }`$ and the term $`\mu u_0^{\prime \prime }`$ cannot be neglected. In these regions, however, $`v_0`$ hardly changes.
The analysis of the inner regions leads to the solutions
$$u_0=\pm \mathrm{tanh}\frac{\theta }{\sqrt{2}k},v_0=0.$$
(21)
These solution represent front structures separating two types of outer regions: domains of high activator values, $`u=u_+(v_0)`$ (”up state”), and domains of low activator values $`u=u_{}(v_0)`$ (”down state”), where $`u_\pm (v_0)`$ are the extreme roots of $`u_0u_0^3v_0=0`$. We look for periodic stationary solutions with wavelength $`\mathrm{\Lambda }=\mathrm{\Lambda }_{}+\mathrm{\Lambda }_+`$, where $`\mathrm{\Lambda }_+`$ and $`\mathrm{\Lambda }_{}`$ are the widths of up and down states respectively. Consider now a down state spanning the spatial range $`\mathrm{\Lambda }_{}<z<0`$ followed by an up state spanning the range $`0<z<\mathrm{\Lambda }_+`$. The equations for $`v`$ at these outer regions are
$$v_0^{\prime \prime }q^2(v_0v_{})=0,\mathrm{\Lambda }_{}<z<0,$$
(22)
with the boundary conditions $`v_0(\mathrm{\Lambda }_{})=v_0(0)=0`$, and
$$v_0^{\prime \prime }q^2(v_0v_+)=0,0<z<\mathrm{\Lambda }_+,$$
(23)
with the boundary conditions $`v_0(0)=v_0(\mathrm{\Lambda }_+)=0`$. In obtaining these equations we approximated $`u_\pm (v_0)=\pm 1v_0/2`$. This approximation is particularly good for bistable media with $`a_0`$ small and $`a_1`$ relatively large. These values restrict $`v_0`$ to a small range around $`v_0=0`$. For excitable media and systems undergoing Turing instability, a large value of $`\delta `$ might be needed to keep $`v_0`$ small.
The solutions to Eqs. (22) and (23) are
$$v_0=\frac{v_{}}{\mathrm{sinh}q\mathrm{\Lambda }_{}}\left[\mathrm{sinh}qz\mathrm{sinh}q(z+\mathrm{\Lambda }_{})\right]+v_{},$$
(24)
for $`\mathrm{\Lambda }_{}<z<0`$, and
$$v_0=\frac{v_+}{\mathrm{sinh}q\mathrm{\Lambda }_+}\left[\mathrm{sinh}q(z\mathrm{\Lambda }_+)\mathrm{sinh}qz\right]+v_+,$$
(25)
for $`0<z<\mathrm{\Lambda }_+`$. To determine $`\mathrm{\Lambda }_\pm `$ for a given $`\mathrm{\Lambda }`$ we match the derivatives of $`v_0`$ at the front positions
$$v_0^{}(0^+)=v_0^{}(0^{}),v_0^{}(\mathrm{\Lambda }_+)=v_0^{}(\mathrm{\Lambda }_{}).$$
This leads to the relation
$$v_+\beta (q\mathrm{\Lambda }_+)+v_{}\beta (q\mathrm{\Lambda }_{})=0,$$
where $`\beta (x)`$ is given by Eq. (11).
## Appendix C Calculation of $`\tau `$ and $`B`$
The quantities $`\tau `$ and $`B`$ are given by Eqs. (6) and (7). Consider first the integral
$$<u_0^{}(\theta )^2>=\frac{1}{2\pi }_0^{2\pi }u_0^{}(\theta )^2𝑑\theta .$$
It has a contribution from two inner regions at $`z=0`$ and $`z=\mathrm{\Lambda }_+`$ where $`u_0`$ is given by Eq. (21), and a contribution from two outer regions, $`\mathrm{\Lambda }_{}<z<0`$ and $`0<z<\mathrm{\Lambda }_+`$ where $`u_0=1v_0/2`$ and $`u_0=1v_0/2`$ with $`v_0`$ given by Eq. (24) and Eq. (25), respectively. (Recall that $`z=\frac{\sqrt{\mu }}{k}\theta `$).
The contribution from the two inner regions is
$`<u_0^{}(\theta )^2>_{inner}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi k^2}}{\displaystyle _{inner}}\mathrm{sech}^4\left({\displaystyle \frac{\theta }{\sqrt{2}k}}\right)𝑑\theta ,`$
$``$ $`{\displaystyle \frac{1}{\sqrt{2}\pi k}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{sech}^4x𝑑x={\displaystyle \frac{2\sqrt{2}}{3\pi k}}.`$
We have used here the fact that $`k𝒪(\sqrt{\mu })1`$. The integral over a narrow inner region is transformed into an integral over a wide region after stretching the $`\theta `$ variable to the $`x=\frac{\theta }{\sqrt{2}k}`$ variable. The contribution from the two outer regions is
$`<u_0^{}(\theta )^2>_{outer}`$ $`=`$ $`{\displaystyle \frac{\sqrt{\mu }}{8\pi k}}`$
$`\times `$ $`\left[{\displaystyle _\mathrm{\Lambda }_{}^0}v_0(z)_{}^{}{}_{}{}^{2}𝑑z+{\displaystyle _0^{\mathrm{\Lambda }_+}}v_0(z)_{}^{}{}_{}{}^{2}𝑑z\right],`$
where we used in the two outer regions $`u_0(z)^{}=\frac{1}{2}v_0(z)^{}`$. Altogether,
$`<u_0^{}(\theta )^2>`$ $`=`$ $`{\displaystyle \frac{2\sqrt{2}}{3\pi k}}+{\displaystyle \frac{\sqrt{\mu }}{8\pi k}}`$
$`\times `$ $`\left[{\displaystyle _\mathrm{\Lambda }_{}^0}v_0(z)_{}^{}{}_{}{}^{2}𝑑z+{\displaystyle _0^{\mathrm{\Lambda }_+}}v_0(z)_{}^{}{}_{}{}^{2}𝑑z\right].`$
The second term on the right hand side of Eq. (C) is small (since $`\sqrt{\mu }1`$) and will not contribute to the leading order forms of $`\tau `$ and $`B`$.
Consider now the integral
$$<v_0^{}(\theta )^2>=\frac{1}{2\pi }_0^{2\pi }v_0^{}(\theta )^2𝑑\theta .$$
The contribution from the inner regions to this integral is negligible for $`\mu 1`$. Thus
$`<v_0^{}(\theta )^2>`$ $`=`$ $`{\displaystyle \frac{\sqrt{\mu }}{2\pi k}}`$ (27)
$`\times `$ $`\left[{\displaystyle _\mathrm{\Lambda }_{}^0}v_0(z)_{}^{}{}_{}{}^{2}𝑑z+{\displaystyle _0^{\mathrm{\Lambda }_+}}v_0(z)_{}^{}{}_{}{}^{2}𝑑z\right].`$
Using the solutions (24) and (25) in the integrals (C) and (27) and using the expressions for $`\tau `$ and $`B`$ we obtain the expressions (8). |
warning/0003/hep-ph0003004.html | ar5iv | text | # SPbU-IP-00-06 Nucleus-nucleus scattering in perturbative QCD with 𝑁_𝑐→∞
## 1 Introduction
In the framework of the colour dipole model of A.H.Mueller it follows that in the high-colour limit $`N_c\mathrm{}`$ the scattering on a heavy nucleus is exactly described by the sum of fan diagrams constructed of BFKL pomerons, each of them splitting into two . The equation for the sum of BFKL fan diagrams was first written by I.Balitsky in his original operator expansion formalism. Then it was rederived by Yu.Kovchegov in the colour dipole framework and by the author by directly summing the BFKL fan diagrams . The perturbative solution of this equation in the region of small non-linearity (outside the saturation region) was studied in . Asymptotic estimates of the solution were presented in . Finally in the exact solution of the equation was obtained by direct numerical methods. The main physical results following from these studies are that, first, at high rapidities $`Y`$ the hA total cross-section saturates to its geometrical limit $`2\pi R_A^2`$ and, second, the gluon density in the nucleus aquires a form of a soliton in $`Y\mathrm{ln}k`$ space moving towards higher momenta with nearly a constant velocity as $`Y`$ increases. This last property supports the applicability of the perturbative treatment, since the well-know diffusion of the BFFKL pomeron towards small momenta results to be stopped.
In the present paper we attempt to generalize these results to nucleus-nucleus (AB) scattering. In this case, in the $`N_c\mathrm{}`$ limit the total amplitude is given by the sum of all tree diagrams constructed of BFKL pomerons and the triple pomeron vertex. In contrast to the hA case, the vertex now describes not only splitting of a pomeron into two but also fusion of two pomerons into one. The diagrams for the amplitude accordingly result much more complicated than the fan diagrams relevant for the hA case. However, using the effective field theory methods developed for summing such diagrams long ago , one can construct an equation which describes the AB amplitude. Naturally this equation (in fact a pair of equations) results much more complicated than for the case of hA scattering. Its exact (numerical) solution does not seem realistic. However simple asymptotic estimates, analogous to the ones made in \[7,8 \], show that the total AB cross-section tends to its geometrical limit at high rapidities similar to the hA case.
Unfortunately the gluon density in the overlapping area cannot be found from these estimates, but rather requires knowledge of the solution in more detail. We leave this problem for future studies.
## 2 Effective field theory for AB scattering
At fixed overall impact parameter $`b`$ the AB amplitude $`𝒜(Y,b)`$ can be presented as an exponential of its connected part:
$$𝒜(Y,b)=2is\left(1e^{T(Y,b)}\right)$$
(1)
The dimensionless $`T`$ is an integral over two impact parameters $`b_A`$ and $`b_B`$ of the collision point relative to the centers of the nuclei A and B:
$$T(Y,b)=d^2b_Ad^2b_B\delta ^2(bb_A+b_B)T(Y,b_A,b_B)$$
(2)
As mentioned in the Introduction, in the perturbative QCD with $`N_c\mathrm{}`$ the amplitude $`T(Y,b_A,b_B)`$ is given by a sum of all connected tree diagrams constructed of BFKL pomerons and the triple pomeron vertex. More concretely, in these diagrams a line (”propagator”) connecting two points $`y_1,r_1`$ and $`y_2,r_2`$ corresponds to one half of the forward BFKL Green function :
$$G(y_1y_2,r_1,r_2)=\frac{r_1r_2}{32\pi ^2}\theta (y_1y_2)\underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}e^{in(\varphi _1\varphi _2)}\frac{d\nu e^{(y_1y_2)\omega (\nu )}(r_1/r_2)^{2i\nu }}{\left(\nu ^2+\frac{(n1)^2}{4}\right)\left(\nu ^2+\frac{(n+1)^2}{4}\right)},$$
(3)
where $`\varphi _{1,2}`$ and are the azimuthal angles and
$$\omega (\nu )=\frac{\alpha _sN_c}{\pi }(\psi (1)\mathrm{Re}\psi (1/2+i\nu ))$$
(4)
are the BFKL levels. Due to the azimuthal symmetry of the projectile and target colour densities one may retain only the term with zero orbital momenta $`n=0`$ in (3).
The interaction between the pomerons is realized via the triple pomeron vertex. It is non-local and not symmetric in the incoming and outgoing pomerons. Its form for the splitting of a pomeron into two was established in . At $`N_c\mathrm{}`$ for the transition $`12+3`$ the three BFKL Green functions are connected by it as follows (see Fig. 1a)
$$\frac{4\alpha _s^2N_c}{\pi }\frac{d^2r_1d^2r_2d^2r_3}{r_1^2r_2^2r_3^2}\delta ^2(r_1+r_2+r_3)G(y_1^{}y,r_1^{},r_1)_1^4r_1^4G(yy_2^{},r_2,r_2^{})G(yy_3^{},r_3,r_3^{})$$
(5)
Here it is assumed that the operator $`_1`$ acts on the left. The form of the vertex for the fusion of two pomerons into one is actually not known. However, the symmetry between target and projectile prompts us to assume that for the inverse process $`2+31`$ the BFKL functions are to be joined as (Fig. 1b)
$$\frac{4\alpha _s^2N_c}{\pi }\frac{d^2r_1d^2r_2d^2r_3}{r_1^2r_2^2r_3^2}\delta ^2(r_1+r_2+r_3)G(y_2^{}y,r_2^{},r_2)G(y_3^{}y,r_3^{},r_3)_1^4r_1^4G(yy_1,r_1,r_1^{})$$
(6)
Finally we have to describe the interaction of the pomerons with the two nuclei. The BFKL Green functions corresponding to the external legs of the diagrams are to be integrated with the colour density of each nucleus. We take the target nucleus at rest, that is, at rapidity zero. Then each outgoing external BFKL Green function is to be transformed into
$$g^2AT_A(b_A)d^2r^{}G(y,r,r^{})\rho _N(r^{})𝑑y^{}d^2r^{}G(yy^{},r,r^{})\tau _A(y^{},r^{})$$
(7)
where $`\rho _N`$ is the colour density of the nucleon, $`T_A`$ is the profile function of the nucleus A and we define
$$\tau _A(y,r)=g^2AT_A(b_A)\rho _N(r)\delta (y)$$
(8)
(with dependence on $`b_A`$ implicit). Similarly each ingoing BFKL external Green function is transformed into
$$𝑑y^{}d^2r^{}G(y^{}y,r^{},r)\tau _B(y^{},r^{})$$
(9)
where
$$\tau _B(y,r)=g^2BT_B(b_B)\rho _N(r)\delta (yY)$$
(10)
To find the amplitude one has to sum over all connected diagrams with $`M`$ ingoing and $`N`$ outgoing lines, corresponding to $`M`$ interactions with the projectile and $`N`$ interactions with the target, divided by $`M!N!`$.
It is trivial to see that this sum exactly corresponds to the sum of tree diagrams generated by an effective quantum theory of two pomeronic fields $`\mathrm{\Phi }(y,r)`$ and $`\mathrm{\Phi }^{}(y,r)`$ with action
$$S=S_0+S_I+S_E$$
(11)
consisting of three terms, which correspond to free pomerons, their mutual interaction and their interaction with external sourses (nuclei) respectively.
To give the correct propagators $`S_0`$ has to be chosen as
$$S_0=2𝑑y_1d^2r_1𝑑y_2d^2r_2\mathrm{\Phi }(y_1,r_1)G^1(y_1y_2,r_1,r_2)\mathrm{\Phi }^{}(y_2,r_2)2\mathrm{\Phi }|G^1|\mathrm{\Phi }^{}$$
(12)
where $`|`$ means the integration over $`y,r`$. Note that the sign of $`S_0`$ corresponds to the substitution of the conventionally defined field variables $`\mathrm{\Phi }`$ and $`\mathrm{\Phi }^{}`$:
$$\mathrm{\Phi }i\mathrm{\Phi },\mathrm{\Phi }^{}i\mathrm{\Phi }^{}$$
(13)
which allows to make all terms of the action real.
According to (5), (6) the interaction term $`S_I`$ is local in rapidity
$$S_I=\frac{4\alpha _s^2N_c}{\pi }dy\frac{d^2r_1d^2r_2d^2r_3}{r_1^2r_2^2r_3^2}\delta ^2(r_1+r_2+r_3)[_1^4r_1^4\mathrm{\Phi }^{}(y,r_1)\mathrm{\Phi }(y,r_2)\mathrm{\Phi }(y,r_3)+c.c.]$$
(14)
The overall sign combines the initial factor $`i`$ and $`i^3`$ from the substitution (13).
Finally the interaction with the nuclei is local both in rapidity and coordinates:
$$S_E=𝑑yd^2r\left[\mathrm{\Phi }(y,r)\tau _A(y,r)+\mathrm{\Phi }^{}(y,r)\tau _B(y,r)\right]$$
(15)
The minus sign comes from the initial $`i`$ and the substitution (13).
The amplitude $`T(Y,b_A,b_B)`$ is then expressed through a functional integral
$$Z=D\mathrm{\Phi }D\mathrm{\Phi }^{}e^{S/\mu ^2}$$
(16)
where $`\mu `$ is an arbitrary mass scale necessary to adjust the dimensions of various parts of the action. Keeping only the connected diagrams one finds
$$T(Y,b_A,b_b)=\mu ^2\mathrm{ln}\frac{Z}{Z_0}$$
(17)
where $`Z_0`$ is the value of $`Z`$ for $`S_E=0`$. Functional integral $`Z`$ is to be calculated in the classical approximation to retain only the tree diagrams. This gives
$$T(Y,b_A,b_B)=S_E\{\mathrm{\Phi },\mathrm{\Phi }^{}\}=d^2r\left[\mathrm{\Phi }(0,r)\widehat{\tau }_A(r)+\mathrm{\Phi }^{}(Y,r)\widehat{\tau }_B(r)\right]$$
(18)
where $`\mathrm{\Phi }`$ and $`\mathrm{\Phi }^{}`$ are the solutions of the classical equation of motion and $`\widehat{\tau }`$’s are (8) and (10) with the $`\delta `$ functions of rapidity dropped. We see that the result is independent of the scale $`\mu `$, as it should be.
## 3 Equations of motion
Before writing out the classical equation of motion, we transform the functional integral (16) to new variables in which the non-locality of the Lagrangian becomes substantially reduced. We put
$$\mathrm{\Phi }(y,r)=r^2\varphi (y,r),\mathrm{\Phi }^{}(y,r)=r^2\varphi ^{}(r,y)$$
(19)
In these variables the interaction term becomes
$$S_I=\frac{4\alpha _s^2N_c}{\pi }dyd^2r_1d^2r_2d^2r_3\delta ^2(r_1+r_2+r_3)[K_1\varphi ^{}(y,r_1)\varphi (y,r_2)\varphi (y,r_3)+c.c.]$$
(20)
where (dimensionles) operator $`K`$ has the form
$$K=r^2_r^4r^2$$
(21)
Now we transform (20) to the momentum space to obtain
$$S_I=\frac{4\alpha _s^2N_c}{\pi }dy\frac{d^2q}{(2\pi )^2}[K\varphi ^{}(y,q)\varphi (y,q)\varphi (y,q)+c.c.]$$
(22)
with operator $`K`$ local in the momentum space
$$K=_q^2q^4_q^2$$
(23)
As we observe, the interaction has become local in the momentum space.
Now we turn to the free part $`S_0`$. In new variables it takes the form
$$S_0=2\varphi |r^2G^1r^2|\varphi ^{}$$
(24)
It was shown in that
$$r_1^2_1^4G(y,r_1,r_2)=g(y,r_1,r_2)r_2^2$$
(25)
where $`g`$ is the Green function of the BFKL equation for the so-called semi-amputated wave function:
$$\left(\frac{}{y}+H_1\right)g(y,r_1,r_2)=\delta (y)\delta ^2(r_1r_2)$$
(26)
Here $`H_1`$ is the BFKL Hamiltonian acting on $`r_1`$. We rewrite (26) in the operatorial form
$$Kr^2Gr^2=\left(\frac{}{y}+H\right)^1$$
(27)
wherefrom
$$r^2G^1r^2K^1=\frac{}{y}+H$$
(28)
and finally
$$r^2G^1r^2=\left(\frac{}{y}+H\right)K$$
(29)
Since both $`G(y,r_1,r_2)`$ and $`g(y,r_1,r_2)`$ are symmetric in $`r_1,r_2`$ we also find in the same manner
$$r^2G^1r^2=K\left(\frac{}{y}+H\right)$$
(30)
so that $`K`$ commutes with $`H`$. Using (29) and (30) in (24) we see that the free part has become local in rapidity and can be expressed via the BFKL Hamiltonian $`H`$:
$$S_0=2\varphi |K\left(\frac{}{y}+H\right)|\varphi ^{}$$
(31)
This part remains non-local both in the coordinate and momentum spaces due to the non-locality of $`H`$.
The interaction with the nucleus part in the new variables takes the form
$$S_E=w_A|\varphi \varphi ^{}|w_B$$
(32)
where, in the coordinate space,
$$w_{A,B}(y,r)=r^2\tau _{A,B}(y,r)$$
(33)
Now, with the action nearly completely local (except for the Hamiltonian term) we can write out the equation of motion. We find
$$\frac{\delta S}{\delta \varphi (y,q)}=2K(\frac{}{y}+H)\varphi ^{}(y,q)+\frac{4\alpha _s^2N_c}{\pi }(K\varphi _{}^{}{}_{}{}^{2}(y,q)+2\varphi (y,q)K\varphi ^{}(y.q))w_A(y,q)=0$$
(34)
and
$$\frac{\delta S}{\delta \varphi ^{}(y,q)}=2K(\frac{}{y}+H)\varphi (y,q)+\frac{4\alpha _s^2N_c}{\pi }(K\varphi ^2(y,q)+2\varphi ^{}(y,q)K\varphi (y.q))w_B(y,q)=0$$
(35)
Applying operator $`(1/2)K^1`$ from the left we find our final equations of motion
$$(\frac{}{y}+H)\varphi ^{}(y,q)+\frac{2\alpha _s^2N_c}{\pi }(\varphi _{}^{}{}_{}{}^{2}(y,q)+2K^1[\varphi (y,q)K\varphi ^{}(y.q)])\frac{1}{2}K^1w_A(y,q)=0$$
(36)
and
$$(\frac{}{y}+H)\varphi (y,q)+\frac{2\alpha _s^2N_c}{\pi }(\varphi ^2(y,q)+2K^1[\varphi ^{}(y,q)K\varphi (y.q)])\frac{1}{2}K^1w_B(y,q)=0$$
(37)
As we see the resulting equations are rather complicated, since they involve nonlocal terms, bilinear in $`\varphi `$ and $`\varphi ^{}`$, which interconnect the two equations. Summing fan diagrams in the hA case in fact leads to the the same equations, in which however $`w_B=0`$. Then one immediately finds that $`\varphi =0`$ identically, which converts the first equation into
$$\left(\frac{}{y}+H\right)\varphi ^{}(y,q)+\frac{2\alpha _s^2N_c}{\pi }\varphi _{}^{}{}_{}{}^{2}(y,q)\frac{1}{2}K^1w_A(y,q)=0$$
(38)
This local equation is just the one studied in \[4-6\].
To conclude this section we present the non-local terms in Eqs. (36) and (37) in a more explicit way. To this end we calculate the kernel of the operator $`K^1`$ in the momentum space. We have in the coordinate space
$$K^1=r^2_r^4r^2$$
(39)
Using the identity
$$_1^4G(0,r_1,r_2)=\delta ^2(r_1r_2)$$
(40)
we can write the kernel of $`K^1`$ in the coordinate space as
$$K^1(r_1,r_2)=r_1^2G(0,r_1,r_2)r_2^2$$
(41)
Fourier transforming (41) and using (3) we find the kernel in the momentum space:
$$K^1(q_1,q_2)=\frac{d^2r_1d^2r_2}{r_1^2r_2^2}e^{iq_2r_2iq_1r_1}G(0,r_1,r_2)=\frac{1}{8}\frac{d\nu }{(\nu ^2+1/4)^2}I(\nu ,q_1)I^{}(\nu ,q_2)$$
(42)
where
$$I(\nu ,q)=_0^{\mathrm{}}𝑑rr^{2i\nu }J_0(qr)=2^{2i\nu }q^{1+2i\nu }\frac{\mathrm{\Gamma }(1/2i\nu )}{\mathrm{\Gamma }(1/2+i\nu )}$$
(43)
Doing the integral over $`\nu `$ we finally find
$$K^1(q_1,q_2)=\frac{\pi }{2}\frac{1}{q_>^2}\left(\mathrm{ln}\frac{q_>}{q_<}+1\right)$$
(44)
where $`q_{>(<)}=\mathrm{max}(\mathrm{min})\{q_1,q_2\}`$.
Using (44) we can rewrite the nonlocal term in Eq. (36) as
$$4\alpha _s^2N_c\frac{d^2q_1}{(2\pi )^2q_>^2}\left(\mathrm{ln}\frac{q_>}{q_<}+1\right)\varphi (y,q_1)_1^2q_1^4_1^2\varphi ^{}(y,q_1)$$
(45)
where $`q_{>(<)}=\mathrm{max}(\mathrm{min})\{q,q_1\}`$. The nonlocal term in Eq. (37) is obtained by complex conjugation.
## 4 The total cross-section
The obtained equations which determine the classical fields $`\varphi `$ and $`\varphi ^{}`$ are very difficult to solve even numerically. The trouble lies not so in the their explicit non-locality, but in the appearence of two different sources at two different rapidities. Due to conditions $`\varphi =0`$ for $`y>Y`$ and $`\varphi ^{}=0`$ for $`y<0`$ and the $`\delta `$-like dependence of the sources on rapidities, one can drop the sources in Eq. (36) and (37) substituting them by conditions
$$\varphi ^{}(y,q)_{y=0}=K^1\widehat{w}_A(q),\varphi (y,q)_{y=Y}=K^1\widehat{w}_B(q)$$
(46)
In the hA case one has only the first of these conditions, which converts Eq. (36) into an evolution equation in rapidity, relatively easily solved by conventional methods. As mentioned the non-local term is zero in this case, but its presence would only slightly complicate the solution. After all the BFKL Hamiltonian is also non-local (although linear).
For the nucleus-nucleus scattering we have to solve homogeneous Eqs (36) and (37) with both conditions (46) imposed upon the solution. The Cauchy problem is thus transformed into an essentially more difficult boundary problem. A possible method of the solution is to transform Eqs. (36) and (37) into a system of two non-linear integral equations in combined rapidity-momentum space, which one may try to solve by iterations. Our experience in the hA problem shows that for the solution to have a reasonable precision one requires at least 800 points in the momentum and 400 points for 5 units of rapidity. Thus to find the amplitude for say $`Y=15`$ one has to perform $`1200\times 800^310^{10}`$ operations per iteration. On top of that the convergence properties of the iteration procedure is unknown.
Here we shall not attempt at solving Eqs. (36) and (37) with any reasonable degree of precision at all values of rapidity and momentum. Instead we shall again use our experience with the case of hA scattering (fan diagrams), where at any fixed momentum and $`y\mathrm{}`$ the solution $`\varphi ^{}(y,q)`$ aquires a simple form, independent of the target properties
$$\varphi ^{}(y,q)_y\mathrm{}=\frac{2\pi }{g^2}\mathrm{ln}\frac{Q(y)}{q}$$
(47)
with $`\mathrm{ln}Q(y)2.34(\alpha _sN_c/\pi )y`$. Note that (47) is not the solution at all $`y`$ and $`q`$. In particular (47) is not valid at $`qQ`$, which is just the region which determines the gluon density. However (47) is sufficient to establish that the hA scattering cross-section tends to its geometric limit at high $`y`$ .
Our guess is that also in the nucleus-nucleus case function $`\varphi ^{}(y,q)`$ aquires the form (47) at large rapidities $`yY`$ and $`\varphi (y,q)`$ aquires the same form with $`yYy`$. To support this behaviour we are going to demonstrate that in these limits the non-local terms in Eqs. (36) and (37) can be neglected, so that the equations decouple and become similar to the hA case.
Actually the demonstration is quite simple. Take Eq. (36) at large $`y`$ and put the conjectured asymptotics (47) into it. According to (45) the mixing non-local term will then be given by
$$\alpha _sN_c\frac{d^2q_1}{(2\pi )^2q_>^2}\left(\mathrm{ln}\frac{q_>}{q_<}+1\right)\varphi (y,q_1)_1^2q_1^4_1^2\mathrm{ln}\frac{Q(y)}{q_1}$$
(48)
Note that $`\varphi (y,q_1)`$ enters at large $`y`$ and small $`Yy`$. Its exact form is unknown but we can safely assume that it rapidly falls with $`q_1`$ similar to the inhomogeneous term $`K^1w_B`$ from which it is separated by a relatively small distance in rapidity. Action of the operator $`K=^2q^4^2`$ on the asymptotic form of $`\varphi ^{}`$ however gives zero. In fact
$$_1^2\mathrm{ln}\frac{Q(y)}{q_1}=2\pi \delta ^2(q_1)$$
Subsequent integration over $`q_1`$ gives zero at any finite $`q`$ due to factor $`q_1^4`$. Thus the non-local term is zero in Eq. (36) at $`yY\mathrm{}`$. Therefore the equation aquires the same form as for hA scattering (which implies that fan diagrams going from top to bottom dominate). This means that the asymptotical behaviour (47) is indeed true. The same result is found for the non-local term in Eq. (37) at small $`y`$ and $`Yy\mathrm{}`$ assuming the asymptotic form (47) for $`\varphi (Yy,q)`$. Its meaning is that fan diagrams going from bottom to top dominate in this limit.
Functions $`\mathrm{\Phi }(y,r)`$ and $`\mathrm{\Phi }^{}(y,r)`$ which actually determine the amplitude according to (18) are related to $`\varphi (y,q)`$ and $`\varphi ^{}(y,q)`$ by
$$\mathrm{\Phi }(y,r)=\frac{d^2q}{(2\pi )^2}e^{iqr}_q^2\varphi (y,q)$$
(49)
and similarly for the conjugated function. Using the asymptotical expression (47) we then get at high $`Y`$
$$\mathrm{\Phi }(0,r)=\frac{1}{g^2}\theta (R_Bb_B),\mathrm{\Phi }^{}(Y,r)=\frac{1}{g^2}\theta (R_Ab_A)$$
(50)
The two $`\theta `$ functions appear because according to Eqs. (36) and (37) $`\varphi =0`$ ($`\varphi ^{}=0`$) when $`w_B=0`$ ($`w_B=0`$), that is for $`b_B>R_B`$ ($`b_A>R_A`$).
Putting (50) in (18) we obtain the connected part of the amplitude at $`Y>>1`$ as
$$T(Y,b_A,b_B)=AT_A(b_A)\theta (R_Bb_B)+T_B(b_B)\theta (R_Ab_A)$$
(51)
It results independent of $`Y`$. After the integration over $`b_B`$ and $`b_B`$ we find
$$T(Y,b)=d^2b_Ad^2b_B\delta ^2(bb_A+b_B)\theta (R_Ab_A)\theta (R_Bb_B)[AT_A(b_A)+BT_B(b_B)]$$
(52)
According to (1) the total AB cross-section is given by
$$\sigma ^{tot}(Y)=2d^2b\left(1e^{T(Y,b)}\right)$$
(53)
¿From (52) and (53) one concludes that at large $`Y`$ the cross-section does not depend on $`Y`$. It saturates at a value which is purely geometrical and for $`A>>1`$ or/and $`B>>1`$. approaches the black disk limit in the overlap area.
## 5 Conclusions
We have derived a pair of equations which describe the nucleus-nucleus scattering in the perturbative QCD with a large number of colours (or, alternatively in the quasi-classical limit, or, in the limit $`A,B\mathrm{}`$). The equations contain mixing terms which are both non-linear and non-local. In contrast to the hA case the equations are to be solved with given boundary conditions at rapidities both of the projectile and target, which complicates the solution enormously.
However the asymptotical form of the solution at fixed momentum and large rapidities is shown to be the same as for the hA case. This allows to demonstrate that at large rapidities the total AB cross-section becomes independent of energy and given by purely geometric considerations. At large $`A`$ or/and $`B`$ it corresponds to the scattering on the black disc in the overlap area.
Going to particle production in AB collisions, the situation at central rapidities seems to be rather simple. The inclusive cross-section will be described by diagrams like shown in Fig. 2, with the target and projectile parts joined by a single pomeron, from which the observed particle is emitted. Evidently this contribution is just a convolution of the production vertex with two gluon densities of the projectile and target, each one corresponding to fan diagrams and found in . At rapiditiy distances from the target or projectile $`\delta y1/\mathrm{\Delta }`$ where $`\mathrm{\Delta }`$ is the pomeron intercept the problem does not look so simple, since the AGK rules are rather complicated in this region (see e. g. ) and the usual cancellation of all diagrams except of the structure shown in Fig. 2 is not at all obvious. We leave this problem for future studies.
## 6 References
1. A.Mueller, Nucl. Phys.,B415 (1994) 373.
2. A.Mueller and B.Patel, Nucl. Phys.,B425 (1994) 471.
3. M.A.Braun and G.P.Vacca, Eur. Phys. J C6 (1999) 147.
4. I.Balitsky, hep-ph/9706411; Nucl. Phys. B463 (1996) 99.
5. Yu. Kovchegov, Phys. Rev D60 (1999) 034008.
6. M.Braun, preprint LU TP 00-06 (hep-ph/0001268)
7. Yu. Kovchegov, preprint CERN-TH/99-166 (hep-ph/9905214).
8 .E.Levin and K.Tuchin, preprint DESY 99-108, TAUP 2592-99 (hep-ph/9908317).
9. A.Schwimmer, Nucl. Phys. B94 (1975)445.
10. D.Amati, L.Caneshi and R.Jengo, Nucl. Phys. B101 (1975) 397.
11. L.N.Lipatov in: ”Perturbative QCD”, Ed. A.H.Mueller, World Sci., Singapore (1989) 411.
12. J.Bartels and M.Wuesthoff, Z.Phys., C66 (1995) 157.
13. M.A.Braun, Eur. Phys. J C6 (1999) 321.
14. M.Ciafaloni et al., Nucl. Phys. B98 (1975) 493.
## 7 Figure captions
Fig. 1. The triple pomeron vertex for the splitting of a pomeron into two (a) and fusion of two pomerons into one (b).
Fig. 2. The generic diagram for the inclusive particle production in the central region in AB collisions |
warning/0003/cond-mat0003406.html | ar5iv | text | # Photoemission and the Origin of High Temperature Superconductivity
## Abstract
The condensation energy can be shown to be a moment of the change in the occupied part of the spectral function when going from the normal to the superconducting state. As a consequence, there is a one to one correspondence between the energy gain associated with forming the superconducting ground state, and the dramatic changes seen in angle resolved photoemission spectra. Some implications this observation has are offered.
In 1956, Chester published an interesting paper which dealt with the the difference in energy between the normal and superconducting states (the condensation energy). For an isotope coefficient of one-half, he demonstrated that the condensation energy could be equated to the change in ion kinetic energy. Historically, the paper did not play a major role, since in the same year, the BCS theory of superconductivity was being developed, solving the problem of classical superconductors.
In this millenial year, though, we are faced with the unsolved problem of high $`T_c`$ superconductivity. Because of this, several of us have invoked the name of Chester. The hope is that by directly focusing on the condensation energy, some light might be shed on the solution to the high $`T_c`$ problem. This is particularly relevant if, as most of us suspect, the origin of high $`T_c`$ is associated directly with electron-electron interactions. In such a case, treating the pair glue as an external object, as in the electron-phonon problem, could be misleading. In essence, the electrons are gluing themselves together. This indicates that new ways of thinking may be important.
This has led a number of authors to concentrate on the change in various response functions when going into the superconducting state. This is illustrated as follows:
$`\mathrm{\Delta }G,FRF_NF_S\mathrm{\Delta }`$ (1)
The idea is to assume a non-zero superconducting order parameter, $`\mathrm{\Delta }`$. This leads to changes in the normal Greens function, $`G`$, and to the creation of an anomalous Greens function, the Gor’kov $`F`$ function. This in turn causes changes in various two-particle response functions, $`R`$: the dynamic spin susceptibility, the dielectric function, the optical conductivity, etc. This in turn leads to a change in the free energy between the normal ($`F_N`$) and superconducting ($`F_S`$) states. If the free energy is lowered, then a non-zero superconducting order parameter is self-consistently stabilized. Note that nowhere in this argument does the question of the pair glue arise. That is, superconductivity is generated simply if the response function change is such as to lower the free energy.
A number of theories with this philosophy have been advocated, each focusing on a different response function. One due to Anderson and co-workers suggests that the c-axis kinetic energy is lowered in the superconducting state . This leads to a change in the c-axis optical conductivity, which has received some experimental support . A different suggestion has been made in regards to the planar conductivity .
Turning to “potential energy” explanations, Leggett has advocated that the energy savings comes from the density-density response function. Perhaps better known is the work of Scalapino and White , where a lowering of the exchange energy is suggested based on a change in the spin-spin response function. This idea was then connected to the appearance below $`T_c`$ of the neutron resonance mode by Demler and Zhang , which has also received experimental support . In all of these cases, a particular part of the free energy is being singled out, with the connection being made via a two particle correlation function.
Here, a different approach is advocated . This is illustrated as follows:
$`\mathrm{\Delta }GF_NF_S\mathrm{\Delta }`$ (2)
Note the simplified nature of this diagram relative to the first one. This argumentation is based on the following relation:
$`U_NU_S=`$
$`{\displaystyle \underset{𝐤}{}}{\displaystyle 𝑑\omega (\omega +ϵ_k)f(\omega )\left[A_N(𝐤,\omega )A_S(𝐤,\omega )\right]}`$ (3)
where $`U_N`$ ($`U_S`$) is the internal energy of the normal (superconducting) state, $`A(𝐤,\omega )`$ the single-particle spectral function, $`f(\omega )`$ the Fermi function, and $`ϵ_k`$ the bare energy dispersion. Eq. 3 is based on a reduced (single-band) Hamiltonian with two particle interactions, and in principle can be generalized to the multi-band case by replacing the scalar quantities in this equation by matrices in reciprocal lattice space. It is easily demonstrated that this equation using the BCS reduced Hamiltonian generates the BCS condensation energy, $`\frac{1}{2}N(0)\mathrm{\Delta }^2`$.
Note that the right hand side of Eq. 3 is a moment of the occupied part of the single-particle spectral function $`A`$ ($`A^{}`$). There are strong arguments that $`A^{}`$ is being measured by angle-resolved photoemission (ARPES) measurements in quasi-2D systems . As a consequence, we see that the high $`T_c`$ phenomenon is intimally connected with the dramatic change in the photoemission lineshape when going below $`T_c`$.
A useful decomposition, especially in regards to the various theories mentioned above, is to break the right hand side of Eq. 3 up into separate kinetic and potential energy pieces. This is easily implemented by rewriting $`(\omega +ϵ_k)`$ as $`(2ϵ_k)+(\omega ϵ_k)`$, the first term being the kinetic energy, the second the potential energy. Therefore, changes in the kinetic energy are associated with changes in the momentum distribution function (which is related to the integrated ARPES spectral weight ), whereas on the Fermi surface, the potential energy contribution reduces to the first moment of $`A^{}`$.
This is easily illustrated for the case of BCS theory . In this case, the potential energy is lowered by $`\mathrm{\Delta }^2/V`$, where $`V`$ is the pair potential, and the kinetic energy increased by $`\mathrm{\Delta }^2/V\frac{1}{2}N(0)\mathrm{\Delta }^2`$, the sum being the BCS condensation energy. The kinetic energy change is a consequence of the broadening of the momentum distribution function by the BCS coherence factors. The potential energy change is easily explained as well. It is due to the difference in $`ϵ`$ and $`E=\sqrt{ϵ^2+\mathrm{\Delta }^2}`$. On the Fermi surface, this difference is maximal, leading to a potential energy lowering of $`\mathrm{\Delta }/2`$ (the factor of $`\frac{1}{2}`$ coming from the coherence factors). When integrated over $`ϵ`$, one then obtains $`\mathrm{\Delta }^2/V`$.
There are two interesting points about the BCS example. First, the transition is potential energy driven (physically, this occurs because the ion terms which actually drive the transition are absorbed into the effective potential of the reduced Hamiltonian). Second, the condensation energy is confined to the vicinity of the Fermi surface by the coherence factors. Note that although the ultraviolet cut-off (the Debye energy) enters the individual kinetic and potential energy terms, it drops out of the net term.
There are several reasons to believe that this BCS analogy may be misleading in the high $`T_c`$ problem. First, it assumes the existence of quasiparticles. This can be contrasted with the high $`T_c`$ case, where although quasiparticle peaks exist below $`T_c`$, they do not exist above . That is, even though the superconducting state is almost certainly a (superfluid) Fermi liquid, the normal state appears to be a non Fermi liquid . As a Fermi liquid does a better job of diagonalizing the kinetic energy than a non Fermi liquid, then one might conjecture that the kinetic energy is indeed lowered in the superconducting state despite the coherence factors. Note that this does not violate the considerations of Chester . That is, although Chester demonstrated that the potential energy of the electrons must be lowered and the kinetic energy raised in the superconducting state, these refer to the potential and kinetic energy terms of the total Hamiltonian, not those of the reduced one. Therefore, there is nothing that prevents the kinetic energy of the reduced Hamiltonian from being lowered.
Second, the BCS condensation energy is confined to a narrow shell around the Fermi surface. For the electron-electron case, though, we can anticipate that the whole Brillouin zone could be affected. This is corroborated by ARPES spectra, which do show changes in the spectra even well away from the Fermi surface. Moreover, ARPES data are characterized by regions of the zone where the dispersion is weak. These same regions of the zone are associated with the anomalous pseudogap seen in underdoped cuprates. Thinking about these regions of the zone, even in the superconducting state, in terms of standard Fermi surface based concepts may be misleading.
Eq. 3 suggests that the best way to get a handle on these issues is by a detailed study of the change in the ARPES lineshape throughout the zone. One objection is that the number being sought is small. For instance, Loram, based on specific heat data, estimates the condensation energy to be only 3K per plane for optimal doped YBCO . On the other hand, the relative contribution from a given $`𝐤`$ point is a different matter. As noted above, in BCS theory, k points on the Fermi surface yield a contribution to the condensation energy of $`\mathrm{\Delta }/2`$, which is a substantial number.
A quick look at the data is sufficient to indicate potential pitfalls, along with suggestions about what may be going on. The $`(\pi ,0)`$ point is singled out since it exhibits the most dramatic lineshape change. In the normal state, one has an extremely broad spectral peak, with a width of order the entire bandwidth. In the superconducting state, this gets dramatically rearranged into a sharp spectral peak, followed at higher binding energy by a dip and hump. By comparing these two spectra, we get some idea about contributions of each spectral feature to the condensation energy.
We begin by looking at the first moment ($`\omega `$) part of Eq. 3. There will be a positive contribution from the sharp spectral peak, which will be followed by a negative contribution from the spectral dip. Although the latter is smaller in weight than the former, it is enhanced because of the $`\omega `$ weighting. This is then followed by the hump and subsequent tail region, and therein lies the rub. Since this region is weighted by $`\omega `$, it is sensitive to how the data are normalized. Typically, ARPES data are normalized in such a way that the high energy tails match, and thus to first approximation there is no tail contribution. But, if one assumes (on the Fermi surface) that the data are normalized by having equal integrated weight, then typically the tails do not quite match. The resulting tail contribution to Eq. 3 can be quite large. Similar considerations enter for the kinetic ($`ϵ_k`$) part of Eq. 3, where one deals with the change in integrated area. The message here is that since there are varying positive and negative contributions to the condensation energy from any given spectrum, then the normalization issue will have to be resolved before we can gain insight from Eq. 3 based on ARPES data.
Still, a qualitative statement can be made concerning the doping dependence of the condensation energy. It is now well known that the weight in the quasiparticle peak is dramatically suppressed as the doping is reduced . In the context of Eq. 3, this implies a decrease in the condensation energy with underdoping. Moreover, since the normal state in the underdoped case is gapped, one expects a further reduction. These two observations go a long way in explaining the dramatic reduction in the condensation energy inferred from specific heat data .
To gain further insight, we have studied the so-called mode model , developed to explain the lineshape change noted above. The idea is that the superconducting lineshape is very similar to that expected for electrons interacting with a collective mode. Detailed analysis of the data vs. doping has verified that the mode is almost certainly the resonance mode observed by neutron scattering . The model consists of a constant scattering rate ($`\mathrm{\Gamma }`$) in the normal state, which becomes gapped in the superconducting state (this gap, which defines the spectral dip, is equal to $`\mathrm{\Delta }`$+$`\mathrm{\Omega }_0`$ where $`\mathrm{\Omega }_0`$ is the mode energy). Of particular note, the quasiparticle peak is a consequence of a non-zero $`\mathrm{\Omega }_0`$.
The results for the condensation energy are surprising . Near the Fermi surface, the kinetic energy is lowered, and the potential energy raised. The former occurs because the formation of quasiparticle peaks has a larger effect on sharpening the momentum distribution than the coherence factors have on broadening it. The potential energy lowering is a more subtle matter. The gap in Im$`\mathrm{\Sigma }`$ causes Re$`\mathrm{\Sigma }`$ by the Kramers-Kronig relation to have a logarithmic behavior. Because of this, the normal and superconducting state spectra only asymptotically approach one another. This leads to a large negative tail contribution to the first moment in Eq. 3, much larger than the positive contribution from the quasiparticle peak. As a result, the potential energy is raised. The kinetic energy driven nature of the transition is surprising, given the expectation that the neutron resonance mode should lead to a lowering of the exchange energy . This behavior, though, is sensitive to the size of $`\mathrm{\Gamma }`$. As $`\mathrm{\Gamma }`$ is reduced (the doping is increased), the normal state becomes more Fermi-liquid like, and one crosses over to the BCS limit, where the transition becomes potential energy driven.
This behavior is reminiscent of a phase diagram for the cuprates recently suggested by Phil Anderson . The potential energy is lowered below a temperature $`T^{}`$ due to the formation of the pseudogap. This line merges with $`T_c`$ on the overdoped side, and so the superconducting transition is potential energy driven on this side. On the underdoped side, though, the transition is kinetic energy driven, since the potential energy savings already occurs at $`T^{}`$, and the additional kinetic energy savings is driven by quasiparticle formation. It is again important to note that despite the presence of a large spectral gap in the underdoped case, quasiparticle peak formation is still associated with $`T_c`$ . Also, this conjecture is consistent with Basov’s results , in that the sum rule violation he reports (indicating a lowering of the kinetic energy) only occurs on the underdoped side of the phase diagram.
In conclusion, we suggest that Eq. 3 will be very useful when thinking about the origin of high $`T_c`$, and that ARPES data will play a major role in this endeavor.
This work is supported by the U. S. Dept. of Energy, Basic Energy Sciences, under Contract W-31-109-ENG-38, the National Science Foundation DMR 9624048, and (M. R. ) the Indian DST through a Swarnajayanti fellowship. |
warning/0003/cond-mat0003336.html | ar5iv | text | # A conjectured scenario for order-parameter fluctuations in spin glasses
## I INTRODUCTION
It is well known that mean-field spin glasses are characterized by strong (non vanishing in the thermodynamic limit) sample to sample fluctuations of the order parameter . Despite the fact that extensive thermodynamic quantities (such as free energy and all its finite order derivatives) are self-averaging in the thermodynamic limit (i.e. their intensive part does not depend on the realization of the quenched randomness) the same result cannot be extended to order parameter fluctuations. It is widely believed that absence of self-averageness of the order parameter is strongly related to replica symmetry breaking, i.e. the existence of several ergodic components not related by any symmetry of the Hamiltonian.
Recently, Guerra suggested that sample to sample fluctuations of the order parameter (hereafter referred to as OPF) verify some sum rules which are generally valid in any disordered system. This claim assumes that the system is stochastically stable in the presence of a mean-field perturbation, a property which may strongly depend on the nature of the equilibrium state. A system is stochastically stable if its properties (static or dynamic) smoothly change in the presence of a small random perturbation. These sum rules have been recently used to define a new dimensionless parameter (hereafter called $`G`$) which takes into account sample to sample fluctuations . This parameter has been shown to provide an alternative and powerful way to locate phase transition points in disordered systems. The advantage of $`G`$ respect to more canonical ones (such as the Binder cumulant ratio used in ordered systems) relies on the fact that it works very well also in the absence of time-reversal symmetry in the Hamiltonian or other more complex disordered systems. In particular, the method has been recently applied for Ising spin glasses , Migdal-Kadanoff spin glasses , Potts glasses , Heisenberg spin glasses, which display a chiral phase transition as well as some protein folding models .
The purpose of this paper is to show, by using general arguments, analytic computations for simple models and numerical simulations, that indeed this new parameter is the appropiate tool to investigate phase transitions in disordered systems much like the Binder cumulant is for ordered systems. We conjecture and prove that this parameter $`G`$ takes the universal value 1/3 at zero temperature for any disordered system (finite or infinite) with the only condition of uniqueness of the ground state and absence of a zero-temperature gap in the local field distribution. This condition is satisfied by all spin-glass models with continuous distribution of couplings and no gap at zero coupling. At finite temperature $`G`$ certainly depends on the system size. We claim that due to the property of replica equivalence, for models in which OPF are finite, $`G`$ converges in the infinite-volume limit to zero if the system is in a paramagnetic phase and to the same zero-temperature value 1/3 if the system is in the spin-glass phase. When OPF vanish this does not necessarily hold and we discuss in what conditions the universal value 1/3 may be recovered.
The paper is organized as follows. Section II is a reminder of the definition of the $`G`$ parameter as well as some other useful ones. Section III presents a detailed computation on a simple disordered model which serves as an illustrative example of the main results. Section IV proves the zero-temperature conjecture under some general conditions for any disordered system. Section V presents detailed calculations on the one-dimensional Ising spin-glass model using the transfer matrix approach. Section VI addresses the validity of the conjecture at finite temperature by studying the Sherrington-Kirkpatrick spherical spin glass, a model where OPF vanish. Finally we discuss the results and present the conclusions.
## II The G parameter and replica equivalence
The definition of the $`G`$ parameter is based on some exact relations obtained for the sample to sample fluctuations of the order parameter in the Sherrington-Kirkpatrick (SK) model . The SK model is defined by the disordered mean-field Hamiltonian,
$$_{SK}=\underset{i<j}{}J_{ij}\sigma _i\sigma _j,$$
(1)
where the $`J_{ij}`$ are quenched Gaussian variables with zero average and variance $`1/N`$ where $`N`$ is the number of sites. The SK model presents a second order phase transition at $`T_c=1`$ below which replica symmetry breaks down and ergodicity is broken. The spin-glass phase is described by an order parameter function $`P_J(q_{12})`$ where $`q_{12}=_{i=1}^N\sigma _i^1\sigma _i^2`$ is the replica overlap and the subindex $`J`$ stands for the realization of the quenched randomness. $`P_J(q)`$ is a simple object in the paramagnetic phase above $`T_c`$ ($`P_J(q)=\delta (q)`$) but develops strong sample to sample fluctuations below $`T_c`$ inside the spin-glass phase. Fluctuations in the order parameter were originally derived by Bray, Moore and Young using the property of replica equivalence . This property states that the sum of all elements contained in a given row (or column) in the replica matrix $`Q_{ab}`$ is independent of the row (or column). As shown by Parisi this is a necessary condition for the replicated free energy to be proportional to the number of replicas $`n`$ and have a well defined free energy in the limit $`n0`$. Fluctuations are then described by the following exact relation in the $`N\mathrm{}`$ limit ,
$$\overline{P_J(q_{12},q_{34})}=\frac{1}{3}\overline{P_J(q_{12})}\delta (q_{12}q_{34})+\frac{2}{3}\overline{P_J(q_{12})}\overline{P_J(q_{34})},$$
(2)
where $`\overline{(.)}`$ stands for disorder average and $`1,2,3,4`$ denote replica indices. Therefore,
$$\overline{P_J(q_{12},q_{34})}\overline{P_J(q_{12})}\overline{P_J(q_{34})},$$
(3)
so $`P_J`$ fluctuates with $`J_{ij}`$ in a non-trivial way. Multiplying both sides of eq.(2) by $`q_{12}^2`$ and $`q_{34}^2`$ and integrating over all possible values of the overlaps $`q_{12},q_{34}`$ one obtains the following sum rule ,
$$\overline{q_{12}^2^2}=\frac{1}{3}\overline{q_{12}^4}+\frac{2}{3}\overline{q_{12}}^2.$$
(4)
where $`\mathrm{}`$ stands for thermal average. This relationship has been also rederived by Guerra using general arguments based on self-averaging properties of the internal energy as well as its finite derivatives . Now let us define the following ratio,
$$G=\frac{\overline{q^2^2}\overline{q^2}^2}{\overline{q^4}\overline{q^2}^2}$$
(5)
Note that the numerator in (5) corresponds (except for the absence of a multiplicative constant $`N^2`$) to the sample fluctuations of the spin-glass susceptibility. For the SK model, because of the sum rule (5), it is possible to show that $`G`$ takes only two values. $`G`$ is equal to 1/3 in the replica symmetry broken phase and vanishes above $`T_c`$,
$$G=\frac{1}{3}\mathrm{\Theta }(T_cT).$$
(6)
The generality of the replica-equivalence property suggests that (6) will hold in any system (even beyond mean-field) if OPF do not vanish in the limit $`V\mathrm{}`$. But may well happen that OPF vanish. Then both numerator and denominator in (5) vanish in the $`V\mathrm{}`$ limit. In this case replica equivalence is not enough to decide what the value of $`G`$ is. The value of $`G`$ is then determined by the form of the finite-size corrections to the order parameter (and in particular its prefactors), which in principle could not satisfy sum rules such as (4). Despite this uncertainty, in this paper we propose three possible scenarios for the parameter $`G`$,
* 1. OPF remain finite in the thermodynamic limit. This is the general situation encountered in mean-field models with a replica broken phase. So both numerator and denominator in (5) are finite in the infinite volume limit. The property of replica equivalence and also stochastic stability indicate that the same should be valid for any finite-dimensional disordered system (assuming that for those systems OPF are finite) leading to $`G=1/3`$ in the spin-glass phase.
* 2. OPF vanish in the large volume limit like $`1/V`$. This is the situation typically encountered in the paramagnetic phase. The ratio may then be zero or finite depending on the particular case.
* 3. OPF vanish in the large volume limit slower than $`1/V`$ (for instance, like $`\frac{1}{V^\alpha }`$ with $`\alpha <1`$). This situation is typical of disordered systems with a marginally stable replica symmetric phase. Both numerator and denominator in (5) vanish, the ratio $`G`$ is finite but may be different from 1/3 at finite temperature. In this case the property of replica equivalence cannot be used for the reason discussed before and stochastic stability may not hold. Actually the property of stochastic stability may breakdown if the equilibrium phase is drastically changed in the presence of a mean-field perturbation. This situation may be found in spin-glass models without OPF such as hierarchichal lattices (i.e. spin glasses in the Migdal-Kadanoff approximation), the Sherrington-Kirkpatrick spherical spin glass (see section VI) or finite-dimensional models described by a unique low-temperature state such as the droplet model.
Despite the main hypothesis of stochastic stability remains still to be proven all previous three cases seem quite reasonable and we do not know of non-trivial counterexamples. Note that there is not any direct relationship between OPF and the value of $`G`$ in the low-temperature phase. Actually, the previous possibilities 1 and 3 may yield the same value of $`G`$ although the physical description of the low-temperature phase is much different. As has been observed in the non vanishing of $`G`$ should not be taken as direct evidence for non-vanishing OPF or replica symmetry breaking. In order to better evidenciate whether OPF survive in the infinite-volume limit, it is necessary to consider another dimensionless parameter, which has not the ambiguity of the ratio of two vanishing quantities. For instance one may define the $`A`$ parameter ,
$$A=\frac{\overline{q^2^2}\overline{q^2}^2}{\overline{q^2}^2},$$
(7)
which is nothing else than the numerator of (5) appropriately normalized. We will see later that the nice properties of $`G`$ are not present in the parameter $`A`$ and the former is much more convenient to locate phase transitions. Generally, one expects $`A`$ to be a non trivial function of both volume and temperature vanishing (in the $`V\mathrm{}`$ limit) only when OPF vanish (for instance, in a paramagnetic phase). If OPF are finite $`A`$ may take a finite value but an identity such as (6) for $`A`$ does not hold.
In this paper we will show examples for all three behaviors, by explicit analytic computations and some numerical calculations. Furthermore, we will show that for models with a unique ground state and without gap in the ground-state local field distribution,
$$\underset{T0}{lim}G(V,T)=\frac{1}{3},$$
(8)
so the $`G`$ parameter is 1/3 at $`T=0`$ for any finite volume $`V`$. This is not anymore true at finite temperature where the parameter $`G`$ may take the value $`1/3`$ only in the infinite volume limit.
Before finishing this section let us remind that in references another quantities similar to (5) and (7) have been proposed for systems without time-reversal symmetry. These are defined by considering the connected overlaps,
$$G_c=\frac{\overline{(qq)^2^2}\overline{(qq)^2}^2}{\overline{(qq)^4}\overline{(qq)^2}^2},$$
(9)
$$A_c=\frac{\overline{(qq)^2^2}\overline{(qq)^2}^2}{\overline{(qq)^2}^2}.$$
(10)
We will see that a result like (8) applies also to the parameter $`G_c`$ and our result reads:
$$\underset{T0}{lim}G_c(V,T)=\frac{13}{31}.$$
(11)
For the SK model the quantity $`G_c`$ is defined by restricting thermal averages to the $`q>0`$ part of the $`P(q)`$. $`G_c`$ does not satisfy the identity (6) so this is not the best quantity to look at in numerical simulations despite the fact that both $`G_c`$ and $`G`$ (and also $`A_c`$ and $`A`$) may take similar values in the vicinity of the critical region. This explains why similar results were obtained for both sets of quantities in numerical simulations.
## III An instructive example
Here we analyze in detail a solvable case which will be useful to illustrate the main contents of the paper and how disorder expectation values of the overlaps are computed. Moreover, the analysis of this section will prove to be useful for a constructive proof of the zero-temperature results (8) and (11) to be presented later on. Consider the following Hamiltonian,
$$=\underset{i=1}{\overset{V}{}}h_i\sigma _i,$$
(12)
where the spins may take the values $`\pm 1`$ and the fields $`h_i`$ are uncorrelated and randomly taken from a distribution $`P(h)`$ with finite weight at zero field (i.e. $`P(0)`$ finite). In principle $`P(h)`$ may be any function
$$P(h)=\stackrel{~}{P}(h)+\underset{k}{}c_k\delta (hh_k)$$
(13)
with $`\stackrel{~}{P}(h)`$ any continuous function, all $`h_k0`$ and $`\stackrel{~}{P}(0)0`$. This condition is enough to ensure the non-degeneracy of the ground state because there is a single configuration which minimizes the energy $`\sigma _i^{}=\mathrm{sign}(h_i)`$. Note that if a finite fraction of the fields $`h_i`$ were zero then the ground state would be degenerate. With this very general condition we may exactly compute the parameters $`G`$ and $`A`$ introduced in the previous section. Everything reduces to compute the three overlap quantities: $`\overline{q^2},\overline{q^4}`$ and $`\overline{q^2^2}`$. The computations are quite elementary and here we present the final results. For the numerator and denominator of eq.(5) we get,
$$Numerator\frac{2(V1)}{V^3}\overline{R^2}^2+\frac{4(V1)(V2)}{V^3}\overline{R^2}\overline{R}^2\frac{2(2V^25V+3)}{V^3}\overline{R}^4,$$
(14)
$$Denominator\frac{2}{V^2}\frac{2}{V^3}+\frac{4(V1)(V2)}{V^3}\overline{R}^2\frac{2(2V^25V+3)}{V^3}\overline{R}^4,$$
(15)
where
$$\overline{R}=_{\mathrm{}}^{\mathrm{}}𝑑hP(h)\mathrm{tanh}^2(\beta h),$$
(16)
$$\overline{R^2}=_{\mathrm{}}^{\mathrm{}}𝑑hP(h)\mathrm{tanh}^4(\beta h),$$
(17)
and $`P(h)`$ is the generic distribution (13). The expressions for the parameters $`G`$ and $`A`$ may be further simplified yielding,
$$G=\frac{(\overline{R^2}\overline{R}^2)(\overline{R^2}+(2V3)\overline{R}^2)}{(1\overline{R}^2)(1+(2V3)\overline{R}^2)},$$
(18)
and
$$A=\frac{2(V1)(\overline{R^2}\overline{R}^2)(\overline{R^2}+(2V3)\overline{R}^2)}{V(1+(V1)\overline{R}^2)^2}.$$
(19)
Note that in the limit $`V\mathrm{}`$ both numerator and denominators in (14) and (15) vanish. The quantity $`A`$ also vanishes like $`1/V`$ but the ratio $`G`$ stays finite,
$$\underset{V\mathrm{}}{lim}G(V,T)=\frac{\overline{R^2}\overline{R}^2}{1\overline{R}^2}.$$
(20)
The finite volume quantity $`G(T,V)`$ in (18) satisfies the conjecture (8). A simple integration by parts reveals that the asymptotic low-temperature behavior of $`\overline{R}`$ and $`\overline{R^2}`$ is given by,
$$\overline{R}=1TD+O(T^2),\overline{R^2}=\overline{R}\frac{T}{3}D+O(T^2),$$
(21)
where $`D`$ is a positive constant given by
$$D=2P(0).$$
(22)
Substituting the asymptotic behavior (21) in (18) we obtain $`G(V,T=0)=1/3`$. Note that the same result is obtained substituting (21) in (20) because in this simple example the two limits $`T0`$ and $`V\mathrm{}`$ may be interchanged. This is not generally true: in particular when a phase transition takes place at $`T=0`$ the two limits may not be interchanged anymore.
For the parameters $`G_c`$ and $`A_c`$ introduce in (9)(10) we get:
$$G_c=\frac{\overline{R^4}(\overline{R^2})^2}{2V24(V2)\overline{R^2}+(2V3)(\overline{R^2})^23\overline{R^4}}$$
(23)
$$A_c=\frac{\overline{R^4}(\overline{R^2})^2}{V(1\overline{R^2})^2}$$
(24)
We observe that $`G_c`$ behaves in a different way. It tends to zero for $`T`$ finite and is again independent of the volume for $`T=0`$ but takes the value $`13/31`$. For $`G_c`$ the two limits ($`V\mathrm{}`$ and $`T0`$) now do not commute. Figures 1234 show the behavior of $`G,G_c,A,A_c`$ as functions of temperature for different values of $`V`$ for the case of a Gaussian fields distribution $`P(h)=(2\mathrm{\Pi })^{1/2}\mathrm{exp}(h^2/2)`$.
In case of a gap of amplitude $`\mathrm{\Delta }`$ in the field distribution one finds that both $`A`$ and $`G`$ vanish exponentially with that gap $`GT\mathrm{exp}(\beta \mathrm{\Delta })`$ and the conjecture does not hold anymore.
We will now prove that, under some general conditions, the conjectured zero-temperature values for $`G`$ and $`G_c`$ hold for any disordered system.
## IV A proof of the conjecture
To generally prove (8) and (11) we start by considering a general Hamiltonian $`(\{\sigma \})`$ where the $`\{\sigma _i;i=1,..,V\}`$ are Ising variables which may take the values $`\pm 1`$ . This Hamiltonian may be written in terms of the local fields,
$$=\underset{i=1}{\overset{V}{}}h_i\sigma _i,$$
(25)
where the $`h_i`$ are local fields proportional to
$$h_i\frac{}{\sigma _i},$$
(26)
which depend on the configuration $`\{\sigma \}`$. Suppose now that the Hamiltonian $``$ may only take continuous values so there is no ground state degeneracy (apart from a global symmetry in the Hamiltonian such as time-reversal symmetry; this case will be discussed below). In particular, no local field $`h_i`$ vanishes. Let us denote by $`\{\sigma ^{}\}`$ the (unique) ground state configuration. The ground state is stable with respect all possible number of spin flips, so that the value of the energy in that configuration $`(\{\sigma ^{}\})`$ is an absolute minimum. In particular the ground state is stable respect to single spin flips and the local fields evaluated at the ground state satisfy the property,
$$\sigma _i^{}=sign(h_i^{})$$
(27)
where $`h_i^{}`$ are evaluated at $`\sigma ^{}`$. To prove the conjecture we need to prove the following statement:
* Statement: Excitations which involve the reversal of a single spin yield the dominant contribution to the low-temperature behavior for all the quantities $`(\overline{q^k})^l`$ for any positive integers $`k,l`$ and by extension, to the numerator and denominator in (5,7,9,10).
This statement somehow allows to map the most probable excitations in (25) with those of the instructive example presented before (12). Nevertheless, we must emphasize two points. The first one is that the ground-state local field distribution in the previous example (12) was taken uncorrelated for different sites and also the same distribution was taken for each spin $`i`$. In general this is not true. Local fields at different sites may be correlated and the distribution on a given site may depend on the site. For instance, in models with open boundaries the local field distribution for the sites located on the surface is certainly different from the distribution of those in the bulk. The second observation is that, in general, the lowest excitations in (25) may involve groups of several spins (and not a single spin flip like in the simple case (12)). So in order to prove the conjecture we must show that excitations in (25) which involve the reversal of any number of spins larger than one always yield subdominant low-temperature corrections to the single-spin excitation case.
In what follows we present a constructive proof of the previous statement without need to refer to the results of the instructive example which had some restrictive assumptions. We start from the general Hamiltonian (25) and analyze the low-temperature behavior of the order parameters $`\overline{q^2},\overline{q^4}`$ and $`\overline{q^2^2}`$. We will first consider the case of one spin excitations and later on the more general one of higher-order excitations.
* One-spin excitations
The calculation proceeds as follows. Consider the ground state $`\{\sigma ^{}\}`$ of (25) as unique and one-spin excitations which involves the reversal of a single spin. If we consider the ground state plus this class of $`V`$ possible excitations we can compute the correlation function $`\sigma _i\sigma _j`$ ($`ij`$), obtaining the result:
$$\sigma _i\sigma _j=\sigma _i^{}\sigma _j^{}\left(12\frac{\mathrm{exp}(2\beta h_i^{}\sigma _i^{})+\mathrm{exp}(2\beta h_j^{}\sigma _j^{})}{1+_{l=1}^V\mathrm{exp}(2\beta h_l^{}\sigma _l^{})}\right)$$
(28)
Defining $`x_i=\mathrm{exp}(2\beta h_i^{}\sigma _i^{})`$, we get in the $`\beta \mathrm{}`$ limit,
$$\sigma _i\sigma _j=\sigma _i^{}\sigma _j^{}(12(x_i+x_j))(ij),$$
(29)
where we have aproximated by 1 the term in the denominator of the ratio in (28). Such an approximation is allowed provided one performs the limit $`\beta \mathrm{}`$ before the infinite volume limit. Note that, in that denominator, each exponential contributes to the sum at most with a term proportional to the temperature (see below). Because there are $`V`$ terms of that type, at most that term is of order $`VT`$. Hence, in the limit $`TV1`$, that denominator equals 1. The result (29) is the only quantity we need in order to evaluate $`q^2`$ and $`q^2^2`$. In terms of the variable $`T_{ij}=\sigma _i\sigma _j^2`$, these are given by,
$$q^2=\frac{1}{V}+\frac{1}{V^2}\underset{ij}{}T_{ij},$$
(30)
$$q^2^2=\frac{1}{V^2}+\frac{2}{V^3}\underset{ij}{}T_{ij}+\frac{2}{V^4}\underset{ij}{}T_{ij}^2+\frac{4}{V^4}\underset{(ijk)}{}T_{ij}T_{ik}+\frac{1}{V^4}\underset{(ijkl)}{}T_{ij}T_{kl},$$
(31)
where the indexes in the sums run from $`1`$ to $`V`$ and correspond to different sites. To average (30),(31) over the disorder we need to compute disorder averages of terms of the type $`x_i^mx_j^n`$ where $`i,j`$ denote sites and $`m,n`$ positive integers. It is easy to show that, in the absence of gap in the ground-state local-field distribution, the terms with $`i=j`$ yield the dominant low-temperature corrections and vanish linearly with $`T`$. Terms with $`ij`$ yield higher-order $`O(T^2)`$ contributions. Suppose $`P(\{h_i^{}\})`$ stands for the ground-state local-field probability distribution. For the terms $`x_i^mx_j^n`$ ($`ij`$), we have
$$\overline{x_i^mx_j^n}=_{\mathrm{}}^{\mathrm{}}\mathrm{exp}(2m\beta h_i^{}\sigma _i^{})\mathrm{exp}(2n\beta h_j^{}\sigma _j^{})P(h_1^{},..,h_V^{})dh_1^{}..dh_V^{}.$$
(32)
The field variables $`h_k^{}(ki,j)`$ may be integrated out, yielding the following expression
$$\overline{x_i^mx_j^n}=_{\mathrm{}}^{\mathrm{}}\mathrm{exp}(2m\beta h_i^{}\sigma _i^{})\mathrm{exp}(2n\beta h_j^{}\sigma _j^{})\widehat{P}_{ij}(h_i^{},h_j^{})𝑑h_i^{}𝑑h_j^{},$$
(33)
$$\widehat{P}_{ij}(h_i^{},h_j^{})=_{\mathrm{}}^{\mathrm{}}P(h_1^{},..,h_V^{})\underset{k(i,j)}{}dh_k^{}.$$
(34)
If the local field distribution $`P(\{h_i^{}\})`$ has finite weight at the point $`h_i=0i`$, then the same holds for the two-sites probability $`\widehat{P}_{ij}(0,0)`$ so that we may expand this term around $`(0,0)`$ in (33) obtaining thereby,
$$\overline{x_i^mx_j^n}=_{\mathrm{}}^{\mathrm{}}\mathrm{exp}(2m\beta h_i^{}\sigma _i^{})\mathrm{exp}(2n\beta h_j^{}\sigma _j^{})\left(\widehat{P}_{ij}(0,0)+\left(\frac{\widehat{P}_{ij}}{h_i^{}}\right)_{(0,0)}h_i^{}+\left(\frac{\widehat{P}_{ij}}{h_j^{}}\right)_{(0,0)}h_j^{}+O(h_i^{}h_j^{})\right)𝑑h_i^{}𝑑h_j^{},$$
(35)
where $`O(h_i^{}h_j^{})`$ denote higher-order terms at least quadratic in the fields. A simple saddle-point calculation (in the $`\beta \mathrm{}`$ limit) gives then,
$$\overline{x_i^mx_j^n}=\frac{T^2}{mn}\widehat{P}_{ij}(0,0)+O(T^3).$$
(36)
The dominant terms in the limit $`T0`$ correspond to terms of the type $`\overline{x_i^n}`$, which give
$$\overline{x_i^n}=\frac{T\widehat{P}_i(0)}{n},$$
(37)
where $`\widehat{P}_i(0)`$ is the value of the single-site probability distribution on the site $`i`$ evaluated at $`h=0`$,
$$\widehat{P}_i(h^{})=_{\mathrm{}}^{\mathrm{}}P(h_1,..,h_V)\delta (h_ih^{})\underset{ki}{}dh_k.$$
(38)
This probability is not independent of the spin $`i`$, as our Hamiltonian can contain terms which introduce asymmetry between different sites. This is an important difference with respect to the computation of the previous section where the local field distribution (13) was site independent. Actually, this independency was necessary in the “instructive example” to fully carry out the analytic computation of $`G`$ and $`G_c`$. The key point is that, at the level of one-spin excitations, low-temperature corrections to overlap averages are linear in $`T`$ and $`\widehat{P}_i(0)`$. According to expressions (30), (31) all sites are equivalent (inequivalence of sites enters only through the value of $`\widehat{P}_i(0)`$), so the only invariant term linear in $`P`$ is $`_i\widehat{P}_i(0)`$. The numerator in (5) yields,
$$\overline{q^2^2}\overline{q^2}^2=\frac{16T_iP_i(0)}{3V^4}(V1)^2+O(T^2).$$
(39)
To compute the overlap $`\overline{q^4}`$ we use the expression,
$$\overline{q^4}=\frac{1}{V^4}\left(3V^22V+(6V8)\underset{ij}{}T_{ij}+\underset{(i,j,k,l)}{}T_{ijkl}\right),$$
(40)
where $`T_{ijkl}=\sigma _i\sigma _j\sigma _k\sigma _l^2`$. Similarly as for the two-point correlation function (29) we obtain
$$\sigma _i\sigma _j\sigma _k\sigma _l=\sigma _i^{}\sigma _j^{}\sigma _k^{}\sigma _l^{}(12(x_i+x_j+x_k+x_l))(i,j,k,l\text{all different}).$$
(41)
With the same assumptions as for the two-points function we obtain for the denominator in (4)
$$\overline{q^4}\overline{q^2}^2=\frac{16T{}_{i}{}^{}P_{i}^{}(0)}{V^4}(V1)^2+O(T^2),$$
(42)
which finally yields,
$$G=\frac{1}{3}+O(T).$$
(43)
A similar calculation for $`G_c`$ yields $`G_c=\frac{13}{31}+O(T)`$.
* Two-spin excitations
Let us consider now excitations which involve only two different spins in the lattice ($`V(V1)/2`$ different type of excitations). In this calculation one-spin excitations are not included. It is easy to check that these excitations yield $`O(T^2)`$ corrections to the two-spin and four-spin correlations. Under the same conditions as before these are given by,
$$\sigma _i\sigma _j=\sigma _i^{}\sigma _j^{}(14(x_i+x_j)\underset{li,j}{}x_l)(ij)$$
(44)
$$\sigma _i\sigma _j\sigma _k\sigma _l=\sigma _i^{}\sigma _j^{}\sigma _k^{}\sigma _l^{}(14(x_i+x_j+x_k+x_l)\underset{mi,j,k,l}{}x_m)(\mathrm{i},\mathrm{j},\mathrm{k},\mathrm{l}\mathrm{all}\mathrm{different})$$
(45)
A saddle point calculation shows that corrections to the ground-state correlation functions are quadratic in $`T`$. Finite $`T`$ corrections now depend on both $`\overline{x_i}\overline{x_j}`$ and $`\overline{x_ix_j}`$ for $`ij`$. Now, for the quantity $`G`$ we expect a dependence of both numerator and denominator on terms of the type $`\widehat{P}_i(0)\widehat{P}_j(0)`$ as well as $`\widehat{P}_{ij}(0,0)`$. They can enter in different forms, for instance $`_{ij}\widehat{P}_{ij}(0,0)`$, $`(_i\widehat{P}_i(0))^2`$ or $`(_i\widehat{P}_i(0)^2)`$. A universal value for $`G`$ is not guaranteed anymore. In particular, supposing uncorrelated local fields (which in principle may not be true) and independency of the one-site probability distribution $`\widehat{P}_i(0)`$ on the site $`i`$ we obtain, after a simple but lengthy calculation,
$$\overline{q^2}^2\overline{q^2}^2=\frac{128T^2P(0)^2}{9V^3}(V2)^2(V1)+O(T^3)$$
(46)
$$\overline{q^4}\overline{q^2}^2=\frac{64T^2P(0)^2}{V^3}(V2)^2(V1)+O(T^3)$$
(47)
and their ratio yields $`G=\frac{2}{9}+O(T)`$ which is different than before. We stress again that the result $`2/9`$ is not universal and will certainly not hold in the most general case. This calculation has been shown to stress how the $`1/3`$ value is a fingerprint of the dominancy in the limit $`T0`$ of the one-spin excitations.
* Higher-order excitations
The generalisation to the most general case of $`K`$-spin excitations is straightforward. Including only this class of excitations we obtain $`O(T^K)`$ corrections to correlations which involve any finite number of spins. This can be easily seen from the fact that any possible excitation of this type will involve the reversal of $`K`$ different spins, each spin $`i`$ contributing by a factor $`x_i=\mathrm{exp}(2\beta h_i^{})`$ to the correction. The simultaneous effect of all spins yields a product type $`_{i=1}^Kx_i`$ which immediately gives (in the limit $`\beta \mathrm{}`$) the $`T^K`$ term. The numerator and denominator in $`G`$ are of order $`T^K`$ with $`O(T^{K+1})`$ corrections. The final result for $`G`$ for any value of $`K`$ is not easy to compute and, as previously discussed, will depend on a larger number of invariants which involve different combinations of the terms $`\widehat{P}_i(0),\widehat{P}_{i_1i_2}(0,0),..,\widehat{P}_{i_1,..i_K}(0,0)`$.
When all possible excitations are treated together the calculation proceeds as before. The dominant contribution for OPF will always come from samples whose lowest excitations are one-spin excitations. Consequently, in the zero-temperature limit (for $`V`$ finite) one-spin excitations dominate the correction to correlation functions proving our conjecture. Note that the result we are stating here is quite natural. OPF at very low temperatures are always dominated by those rare samples characterised by local fields $`\beta h<<1`$ where one spin-excitations yield the largest contribution. From a numerical point of view this implies that more samples are needed to compute with a reasonable precision the values of $`G`$ and $`G_c`$ as $`T`$ goes down. This is because for $`T0`$ the effect from rare samples on OPF becomes more and more important. Let us stress again that the present derivation assumed that $`TV1`$. In the opposite limit or in an intermediate regime the result obviously does not hold. In that case, it may well happen that dominant contributions in OPF involve the reversal of a large number of spins (domain excitations) which, in the limit $`TV1`$, may also involve the whole system .
The hypothesis of a unique-ground state is apparently in contradiction with the case in which there is time-reversal symmetry. Indeed all spin-correlations computed in this section are invariant under time-reversal symmetry and the present conclusions remain unchanged. The situation is certainly different in disordered systems with non-trivially degenerate ground states (for instance, finite-dimensional spin glasses with discrete couplings) where we expect that $`G(V,T)`$ vanishes exponentially with $`1/T`$ like in the instructive example of the previous section. Again, in the other limit (finite temperature and $`V\mathrm{}`$) the behavior of these degenerate models may completely change and $`G`$ could be finite again .
## V The 1D Ising spin glass
In this section we present an analysis of the 1D-Ising spin-glass model with free boundary conditions. We consider the folllowing hamiltonian:
$$=\underset{i=1}{\overset{N1}{}}J_i\sigma _i\sigma _{i+1},$$
(48)
where the couplings are randomly distributed according to the probability distribution $`P(J)`$. Our aim is to obtain an analytic expression for $`G`$ and $`A`$ eqs. (5) and (7). As this model has the transition at $`T=0`$, we expect that in the large volume limit $`G`$ will go to zero except at $`T=0`$, where $`G=1/3`$. Moreover, we show that at zero temperature $`G=1/3`$ for any finite system, although here the two limits ($`V\mathrm{}`$ and $`T0`$) do not commute. In order to obtain an expression for the moment of the order parameter $`q`$ we have computed the following object:
$$\overline{e^{yq}^m},$$
(49)
where $`m`$ is a positive integer and $`q`$ is the overlap between two different configurations of spins, which is the generator of the moments of the overlap $`\overline{q^p^s}`$. Once obtained this expression, by partial derivation respect to $`y`$ we will obtain expressions for the expectation values of all the moments of $`q`$, such as:
$$\overline{q^n}=\frac{^n\overline{e^{yq}}}{y^n}|_{y=0}.$$
(50)
In our computation we are only interested on the quantities: $`\overline{q^2},\overline{q^4}`$ and $`\overline{q^2^2}`$. Consequently we only need to compute (49) for $`m=1,2`$. The former can be easily computed by (50). By doing some more work we can obtain an expression for $`\overline{q^2^2}`$:
$$\overline{q^2^2}=\frac{1}{3}\left[\frac{^4\overline{e^{yq}}^2}{y^4}\frac{^4\overline{e^{yq}}}{y^4}\right]_{y=0},$$
(51)
where we have used the fact that in this model $`q=\mathrm{\hspace{0.33em}0}`$.
### A The transfer matrix method
For general $`m`$, (49) can be computed through the transference matrix method . We have to compute:
$$\underset{\alpha =1}{\overset{m}{}}\frac{_{\{\sigma ^\alpha \}\{\tau ^\alpha \}}\mathrm{exp}\left(y_{i=1,N}\frac{\sigma _i^\alpha \tau _i^\alpha }{N}+\beta _{i=1}^{N1}J_i(\sigma _i^\alpha \sigma _{i+1}^\alpha +\tau _i^\alpha \tau _{i+1}^\alpha )\right)}{𝒵^2}$$
(52)
where $`𝒵`$$`=2_i2\mathrm{cosh}(\beta J_i)`$ is the partition function of a 1D chain, $`\alpha `$ is the index for each pair of replicas and we have $`m`$ systems of two replicas.
In order to perform the average over the disorder, we are interested in considering the transfer matrix associated to each point $`i`$, so that it contains all the dependence of the $`J_i`$. For a single pair of replicas this matrix reads:
$$V_iV(\sigma _i,\tau _i;\sigma _{i+1},\tau _{i+1})=\frac{exp\left(y\frac{\sigma _i\tau _i+\sigma _{i+1}\tau _{i+1}}{2N}+\beta J_i(\sigma _i\sigma _{i+1}+\tau _i\tau _{i+1})\right)}{(2\mathrm{cosh}(\beta J_i))^2}.$$
(53)
For general $`m`$ our matrix associated to each point consists of the tensorial product of $`m`$ matrices $`V_i`$. At this stage we are ready to perform the average over the disorder and for any $`i`$ we have:
$$\overline{T}=\overline{T_i}=\overline{\underset{1}{\overset{m}{}}V_i}$$
(54)
Then our calculation is reduced to:
$$\overline{e^{yq}^m}=\frac{1}{4}e^{y\frac{_\alpha \sigma _1^\alpha \tau _1^\alpha }{2N}}\overline{T}^{N1}e^{y\frac{_\alpha \sigma _N^\alpha \tau _N^\alpha }{2N}},$$
(55)
so we must compute the trace of the product
$$\overline{T}^{N1}B,$$
(56)
where $`A`$ is a $`4m\times 4m`$ matrix, which is the tensorial product of $`m`$ matrices, which contain the terms of the two edges which had fallen out in the symmetrization process,
$$B=\underset{\alpha }{}\frac{1}{2^2}e^{y\frac{\sigma _1^\alpha \tau _1^\alpha +\sigma _N^\alpha \tau _N^\alpha }{2N}}.$$
(57)
The rest of the calculation is straightforward: In first place we have to diagonalise $`\overline{T}`$, and obtain the set of eigenvalues and eigenvectors, so that in this new base we have
$$\overline{T}_{\lambda }^{}{}_{}{}^{N1}=\left(\begin{array}{ccc}\lambda _1^{N1}& \mathrm{}\mathrm{}& \mathrm{}\mathrm{}.\\ \mathrm{}..& \lambda _2^{N1}& \mathrm{}\mathrm{}.\\ \mathrm{}..& \mathrm{}\mathrm{}..& \mathrm{}\mathrm{}.\\ \mathrm{}..& \mathrm{}\mathrm{}& \lambda _{2^{2m}}^{N1}\end{array}\right),$$
(58)
where the subindex $`\lambda `$ stands for the diagonalised matrix. We then have to obtain the change of base matrix $`M`$ which expresses the new set of eigenvectors $`\{\lambda ^i\}`$ in terms of the old base $`\{\sigma ^\alpha \}`$. We finally obtain:
$$\overline{e^{yq}^m}=TrM\overline{T}_{\lambda }^{}{}_{}{}^{N1}M^TB.$$
(59)
We have to point out that the case $`m=1`$ is easy to solve. However, the case $`m=2`$ turns up to be more difficult as the diagonalization of $`\overline{V}`$ is not trivial. To compute $`\overline{q^2^2}`$ one can always use the traditional method by using the fact that
$$\sigma _i\sigma _j=\underset{p=i,j1}{}\mathrm{tanh}\beta J_pij.$$
(60)
### B Results
Here we report on the obtained results in the low-temperature limit and in the infinite-volume limit. The relevant quantities $`\overline{q^2},\overline{q^4}`$ and $`\overline{q^2^2}`$ only depend on $`N`$, $`\overline{R}`$ and $`\overline{R^2}`$ which have been introduced in section III, and whose low-temperature behavior is given by (21). At finite temperature, where $`\overline{R}^N`$ and $`\overline{R^2}^N`$ $`1`$ we obtain for the numerator and denominator in (5):
$$numerator=\frac{4(1+\overline{R})(\overline{R}^2\overline{R^2})}{N^3(\overline{R}1)^3(\overline{R^2}1)}+𝒪(\frac{1}{N^2}),$$
(61)
$$denominator=\frac{4(1+\overline{R})(\overline{R}^2\overline{R^2})}{N^2(\overline{R}1)^3(\overline{R^2}1)}+𝒪(\frac{1}{N^3}),$$
(62)
where we have kept the lowest orders in $`1/N`$ and we have made the following aproximations $`lim_N\mathrm{}\overline{R}^N,\overline{R^2}^N0`$. We see that in this limit $`G`$ goes to zero as $`1/N`$. However, if we take the low temperature limit (21), where $`\overline{A},\overline{A^2}1`$ then we get the expressions
$$numerator=\frac{4D(N^41)T}{45N^3}+𝒪(T^2),$$
(63)
$$denominator=\frac{4D(N^41)T}{15N^3}+𝒪(T^2),$$
(64)
where $`D`$ is given by $`D=2P(0)`$. This yields $`G=\frac{1}{3}+𝒪(T)`$, independently of the size of the system. A detailed computation up to second order in $`T`$ gives us that in the large volume limit: $`G=\frac{1}{3}BTN`$, $`B`$ being a constant. In fact for the parameter $`A`$, we get in the limit $`T0`$:
$$A=\frac{4D(N^41)T}{45N^3}+𝒪(T^2).$$
(65)
In figures (5) and (6) we show $`G`$ and $`A`$ as a function of the temperature for a Gaussian distribution odf couplings $`P(J)=(2\mathrm{\Pi })^{1/2}exp(J^2/2)`$. Note that the low-temperature correction to $`G`$ and (65) scale as $`TN`$ when $`N\mathrm{}`$ reflecting the fact that as we get close to the transition point $`T=0`$, the correlation length diverges as $`1/T`$. We recover the desired result at $`T=0`$, however we have to stress out that in this model both limits $`T0`$ and $`N\mathrm{}`$ do not commute.
## VI The spherical Sherrington-Kirkpatrick spin glass
In this section we present some numerical simulations for the values of $`G`$ and $`A`$ in the Sherrington-Kirkpatrick (SK) spherical spin glass. This case is quite interesting because its low-temperature behavior corresponds to the second possibility mentioned in section II where OPF vanish (in the $`V\mathrm{}`$ limit) much slower than the paramagnetic example studied in the previous section. Correspondingly the study of OPF in this model turns out to be very complicated because the equilibrium solution is marginally stable. The model is defined by
$$=\underset{i<j}{}J_{ij}\sigma _i\sigma _j,$$
(66)
where $`\mathrm{}<\sigma _i<\mathrm{}`$ and the values of $`\sigma _i`$ satisfy the spherical global constraint $`_{i=1}^N\sigma _i^2=N`$. The couplings have average zero and variance $`1/N`$. The statics of this model can be solved with and without replicas . In the former case one finds a transition at a temperature $`T_c=1`$ where the Edwards-Anderson parameter is different from zero and equal to $`1T`$. In the latter case the transition corresponds to a macroscopic condensation of spin configurations onto the eigenvector corresponding to the largest eigenvalue. In the replica framework it has been shown that the replica symmetric solution is the only possible one within the Parisi scheme. Since OPF vanish, the computation of $`G`$ requires the knowledge of finite-size corrections in the numerator and denominator in (5). A simple calculation reveals that the replica symmetric solution is marginally stable (the replicon eigenvalue vanishes everywhere below $`T_c`$) so the spin-glass susceptibility diverges. The situation is similar to what happens in the usual Sherrington-Kirkpatrick model with Ising spins. There the spin-glass susceptibility diverges proportionally to the volume while now the divergence is much slower (like $`N^{1/3}`$). This is so because in the present model OPF vanish like $`N^{2/3}`$ while in the original SK model OPF are finite.
Again, to compute $`G`$ we need to know the precise value of the amplitudes entering in the finite-size corrections in the parameters $`\overline{q^2},\overline{q^2^2},\overline{q^4}`$. It is well known that analytic calculations of finite-size corrections in spin glasses are extremely difficult, specially for the amplitudes which are the quantities we are interested in. For the SK model these amplitudes are partially known only for some quantities . For the present case we will use theoretical considerations and numerical simulations to estimate the asymptotic behavior of the different overlaps.
We have simulated model (66) with a Monte Carlo dynamics where a change of a randomly chosen spin is proposed $`\sigma _i\sigma _i+\delta r_i`$ where $`\delta `$ is a constant number typically of order 1 and $`r_i`$ is a random number uniformly distributed between $`1/2`$ and $`1/2`$. The value of $`\delta `$ is chosen to have a reasonable acceptance rate. The value of all other spins is recalculated in order to satisfy the global spherical constraint. Moves are accepted according to the Glauber algorithm. Note that although we need to recalculate the value of all spins (changing them by multiplying by a normalization constant) the change in the energy can be simply calculated in a finite number of operations independent of $`N`$ and simulations are as fast as with Ising spins. Our investigation has focused on small sizes, which reveal how $`G`$ is a powerful tool to investigate phase transitions. The number of samples simulated are typically several thousands for very small sizes ($`N=4,6,8,12,16`$) and several hundreds for larger ones ($`N=24,32,40,48,64`$). Overlaps have been computed by collecting statistics over a large time window (tipically of order $`10^5`$ Monte Carlo steps for each sample). We have evaluated, $`\overline{q^2}^2,\overline{q^2^2},\overline{q^4}`$ for different sizes and temperatures.
Figure 7 shows the results for $`G`$. Note that already for the smallest sizes there is a crossing of the different curves. The crossing appears for values of $`T`$ well above $`T_c=1`$ for the smallest sizes and moves to lower temperatures as the size increases converging to the expected value $`T_c=1`$. It is quite surprising that already for very small sizes the transition can be clearly seen. The crossing moves down in temperature as the sizes increase and already for several tens of spins converges to the correct value $`T=1`$. As a comparison we show in figure 8 the behavior of the usual Binder parameter defined as
$$B=\frac{1}{2}\left(3\frac{\overline{q^4}}{\overline{q^2}^2}\right)$$
(67)
In this case the crossing point appears at low temperatures for small sizes and moves up very slowly as the size increases. But already for the largest sizes the crossing is still at $`T0.8`$ quite far from $`T=1`$. A similar effect has been observed in simulations of the Sherrington-Kirkpatrick model with Ising spins . These results indicate that a numerical study of the parameter $`G`$ can be extremely useful to locate phase transitions in disordered systems by studying very small sizes .
To analyze better the behavior of $`G`$ at low temperatures we have tried to extrapolate $`G`$ to the large $`N`$ limit. Below $`T_c`$ we expect for all three quantities $`\overline{q^2},\overline{q^2^2},\overline{q^4}`$ the following finite-size corrections,
$$\overline{q^2}^2,\overline{q^2^2},\overline{q^4}=q_{EA}^4+\frac{a}{N^{2/3}}+\frac{b}{N}+\frac{c}{N^{4/3}}+\frac{d}{N^{5/3}},$$
(68)
with $`q_{EA}=1T`$. From these expressions we expect for $`G`$ the following behavior,
$$G=G_{\mathrm{}}+\frac{A}{N^{1/3}}+\frac{B}{N^{2/3}}+O(1/N).$$
(69)
We have fitted the values of $`G`$ to this expression with $`G_{\mathrm{}},A,B`$ as fitting parameters. The results and the fits are shown in figure 9. The extrapolated values for the lowest temperatures $`T=0.6,0.7`$ are $`G_{\mathrm{}}(T=0.6)=0.34\pm 0.2,(A(T=0.6)=0.71\pm 0.1`$ and $`B(T=0.6)=0.49\pm 0.13)`$, $`G_{\mathrm{}}(T=0.7)=0.29\pm 0.2,(A(T=0.7)=0.66\pm 0.1`$ and $`B(T=0.7)=\mathrm{0..49}\pm 0.12)`$. Within errors these are compatible with the value $`1/3`$. Trying to have an estimate of $`G_{\mathrm{}}`$ at higher temperatures is very difficult because critical effects are strong.
We must conclude that for this model the universal value $`1/3`$ is well compatible with the data suggesting that this may be a generic result for a spin-glass phase. Still we should do more extensive simulations to reach a final conclusion. Although going to larger sizes at the lowest temperatures may be factible this will require much longer computational time.
## VII Outlook and discussion
In this paper we have investigated order parameter fluctuations (OPF) in spin glasses. In particular we have considered four different parameters: $`G,A`$ for disconnected thermal averages and $`G_c,A_c`$ for connected thermal averages. It has been recently shown that these models can be very useful to investigate phase transitions in disordered systems and several recent numerical works indeed support this conclusion. In this work we have concentrated our attention to obtain general results and to apply them to certain solvable cases where these can be explictly checked.
We have demonstrated that for models with a unique ground state and no gap in the ground-state local field distribution (for instance, all discrete models with continuous disordered couplings taken from a distribution without gap) $`G`$ and $`G_c`$ take the respective universal values $`G=1/3`$, $`G_c=13/31`$ at zero temperature for any finite volume. This is consequence of the dominancy of one-spin excitations in OPF. For infinite volume this result still holds only in the regime where the limit $`T0`$ is taken before the limit $`V\mathrm{}`$ and fast enough such that $`TV0`$. This result has then been checked calculating OPF in an instructive example without many body interactions and for the case of the one-dimensional Ising spin glass where explicit computations can be done using the transfer matrix method. All these good properties suggest that both parameters $`G,G_c`$ are ideal candidates to investigate phase transitions in disordered systems much alike the Binder cumulant is for ordered systems.
The extension of this result to the other limit where $`V\mathrm{}`$ is taken before $`T0`$ or, more generally, the limit $`V\mathrm{}`$ for $`T`$ finite is far from trivial. In this last case, $`G(V,T)`$ is not volume independent anymore. So the question is whether $`G(V,T)`$ converges in the large $`V`$ limit to the universal temperature independent value 1/3. At finite temperatures there are different possible scenarios for the value of $`G`$. In case OPF are finite in the $`V\mathrm{}`$ limit stochastic stability arguments and replica equivalence suggest that $`G`$ should be 1/3 everywhere in the spin-glass phase. Replica equivalence is a very generic property which, to our knowledge, has not been emphasized before in the present context and implies that the free energy of a replicated disordered system must be proportional to the number of replicas. Note that at zero temperature replica equivalence cannot be used because the limits $`V\mathrm{}`$ and $`T0`$ may not commute in that case. Actually, as we proved in section IV only for models with a unique ground state and absence of gap in the fields distribution, $`G`$ takes the universal value $`1/3`$ but vanishes (exponentially fast with $`1/T`$) in the presence of a finite gap in that distribution.
The other interesting case is when OPF vanish. And here we can offer only more speculative arguments. A possible scenario is that which distinguishes two possibilities depending whether, in the infinite-volume limit, OPF vanish like $`1/V`$ or slower like $`1/V^\alpha `$ with $`\alpha <1`$. If OPF vanish like $`1/V`$, $`G`$ may take the value 0 typical of a paramagnetic phase (for instance the case of the one-dimensional spin-glass model) or a temperature dependent value (the instructive example of section III). For these two solvable cases the parameter $`G`$ is quite different. In the one-dimensional Ising spin glass we find $`G=\frac{1}{3}\delta _{T,0}`$ while in the instructive example we find $`G=\widehat{G}(T)`$ with $`\widehat{G}`$ a monotonous decreasing function of $`T`$ with $`\widehat{G}(0)=1/3`$. The reason for these two different behaviors in a disordered phase may be adscribed to the fact that, in the first case, there is a critical point at $`T=0`$ while in the second there is no critical point at all. So $`G`$ is a good indicator for a phase transition. But this observation must be taken with caution because the parameter $`G_c`$ shows a different behavior for the instructive example $`G_c=\frac{13}{31}\delta _{T,0}`$ similar to the behavior of $`G`$ in the one-dimensional spin glass. We expect the interesting behavior to be present in models where OPF vanish like $`1/V^\alpha `$ with $`\alpha <1`$. This class of models includes disordered systems where the replica symmetric solution is marginally stable and eventually finite-dimensional spin glasses if replica symmetry is not broken, a question still unsolved. This case is much more subtle because replica equivalence cannot be used (nor probably the stochastic stability property) and finite-size corrections must be known. To address this question we have done a numerical study of the spherical Sherrington-Kirkpatrick spin glass. Two are the main outcomes: 1) The parameter $`G`$ is an excellent tool to locate the spin-glass transition already for very small sizes (more precise than the usual Binder parameter) and 2) An infinite-volume numerical extrapolation (compatible with the expected form for the finite-size corrections) of the value of $`G`$ in the spin-glass phase is well compatible with the value 1/3.
shares the same p
Before concluding we want to stress that, apart from their applicability to the study of spin-glass transitions, OPF are interesting quantities which deserve further investigation. The outcome of the proof in section IV is that OPF are much sensitive and rely completely on the effect of rare samples. Actually, rare samples are those which induce the largest OPF and fix the value of $`G`$ to $`1/3`$. A comprehensive study of rare events in disordered systems is still missing. Averaging of extensive quantities such as the replicated free energy in standard renormalization group approaches may wipe out a large number of effects such as those discussed here. Certainly more detailed investigations are needed to clarify the situation. Although a final theorem which resolves this problem may be at hand we think that the search for non-trivial counterexamples of the different possibilities discussed in this paper could be very useful.
Acknowledgments. We acknowledge discussions with M. Picco and A. J. Bray. We are indebted to A. A. Garriga and A. Rocco for a careful reading of the manuscript. F.R is supported by the Ministerio de Educación y Ciencia in Spain (PB97-0971). M. S. is supported by the Ministerio de Educación y Ciencia of Spain, grant AP-98 36523875. |
warning/0003/hep-ph0003051.html | ar5iv | text | # Unifying the Strengths of Forces in Higher Dimensions
\[
## Abstract
We consider the embedding of the Standard Model fields in a $`(4+d)`$-dimensional theory while gravitons may propagate in $`d^{}`$ extra, compact dimensions. We study the modification of strengths of the gravitational and gauge interactions and, for various values of $`d`$ and $`d^{}`$, we determine the energy scale at which these strengths are unified. Special cases where the unification of strengths is characterized by the absence of any hierarchy problem are also presented.
\]
It is widely believed that any fundamental theory capable of describing our world at higher energy scales always predicts the existence of extra, spatial, compact dimensions. Being motivated by attempts to lower the string scale at the gauge unification scale or the TeV scale , the concept of large or small extra dimensions, that are being felt only by gravitons, has been used in order to attack the hierarchy problem. On the other hand, the existence of extra dimensions, where gauge bosons can propagate, was proposed in an attempt to explain the size of the supersymmetry breaking scale or to lower the unification scale of the gauge interactions . In both cases, the Kaluza-Klein excitations of either gravitons or Standard Model particles appear in the framework of the 4-dimensional, effective theory and, in principle, may modify the low-energy physics. The accuracy with which the electroweak and strong interactions have been probed at low energies demands that any contributions to SM processes coming from the KK excitations be extremely small which, in turn, places a lower bound on the energy scale associated with the extra dimensions.
In this letter, we assume the existence of two different sets of extra dimensions: one being felt only by gravitons and one where gauge bosons can propagate. The presence of the extra dimensions modifies simultaneously the strength of both the gravitational and gauge interactions. We search for specific combinations of the numbers and sizes of the extra dimensions that leads to the unification of strengths of forces in the framework of the higher-dimensional theory.
We start our analysis by considering a higher-dimensional formulation of an $`SU(N)`$ Yang-Mills theory with matter. We assume that the vector and scalar particles, i.e. the gauge bosons, $`\widehat{A}_M^\alpha `$, and the Higgs field, $`\widehat{\varphi }`$, may propagate in $`4+d`$ dimensions, while all the fermionic particles, $`\mathrm{\Psi }`$, are localized on the 4-dimensional boundary<sup>*</sup><sup>*</sup>*Our results apply also in the case where the fermions live in the bulk.. The $`(4+d)`$-dimensional action functional, that describes the above theory and preserves Lorentz invariance, may be written as
$`S_{4+d}={\displaystyle }d^4xd^dz\{`$ (1)
$``$ $`{\displaystyle \frac{1}{4}}\left(_M\widehat{A}_N^\alpha _N\widehat{A}_M^\alpha +{\displaystyle \frac{\widehat{g}}{\sqrt{\mathrm{\Lambda }^d}}}C_{\alpha \beta \gamma }\widehat{A}_M^\beta \widehat{A}_N^\gamma \right)^2`$ (2)
$``$ $`\overline{\mathrm{\Psi }}_{L,R}\gamma ^\mu \left(_\mu {\displaystyle \frac{i\widehat{g}}{\sqrt{\mathrm{\Lambda }^d}}}t^\alpha \widehat{A}_\mu ^\alpha \right)\mathrm{\Psi }_{L,R}\delta (\stackrel{}{z})`$ (3)
$``$ $`\left|\left(_M{\displaystyle \frac{i\widehat{g}}{\sqrt{\mathrm{\Lambda }^d}}}t^\alpha \widehat{A}_M^\alpha \right)\widehat{\varphi }\right|^2\widehat{\mu }^2\widehat{\varphi }^{}\widehat{\varphi }{\displaystyle \frac{\widehat{\lambda }}{2\mathrm{\Lambda }^d}}(\widehat{\varphi }^{}\widehat{\varphi })^2`$ (4)
$``$ $`({\displaystyle \frac{\widehat{Y}_1}{\sqrt{\mathrm{\Lambda }^d}}}\overline{\mathrm{\Psi }}_L\widehat{\varphi }\mathrm{\Psi }_R+{\displaystyle \frac{\widehat{Y}_2}{\sqrt{\mathrm{\Lambda }^d}}}\overline{\mathrm{\Psi }}_L\widehat{\varphi }_c\mathrm{\Psi }_R+h.c.\left)\delta (\stackrel{}{z})\right\},`$ (5)
where $`\widehat{g}`$, $`C_{\alpha \beta \gamma }`$ and $`t^a`$ are the coupling constant, the structure constants and the generators, respectively, of the $`SU(N)`$ gauge group while $`\widehat{\mu }^2`$ and $`\widehat{\lambda }`$ are the mass and coupling constant of the Higgs field. In order to render the coupling constants of the theory dimensionless in $`(4+d)`$-dimensions, an arbitrary energy scale $`\mathrm{\Lambda }`$ has been introduced. Note that, in the above, $`M,N=\{t,x_1,x_2,x_3,z_1,z_2,\mathrm{},z_d\}`$, $`\mu ,\nu =\{t,x_1,x_2,x_3\}`$ and the hat denotes $`(4+d)`$-dimensional quantities.
Next, we assume that the extra $`d`$ dimensions are compactified over an internal manifold with the size of every compact dimension being $`2L`$. Then, we can Fourier expand the $`(4+d)`$-dimensional vector and scalar fields along the compact dimensions in the following way
$$\widehat{\mathrm{\Phi }}(x,z)=\widehat{\mathrm{\Phi }}^{(0)}(x)+\underset{\stackrel{}{n}=1}{\overset{\mathrm{}}{}}\frac{\widehat{\mathrm{\Phi }}^{(\stackrel{}{n})}(x)}{\sqrt{2}}\left(e^{i\frac{\pi \stackrel{}{n}}{L}\stackrel{}{z}}+e^{i\frac{\pi \stackrel{}{n}}{L}\stackrel{}{z}}\right),$$
(6)
where $`\stackrel{}{n}=\{n_1,n_2,\mathrm{},n_d\}`$. By performing a Kaluza-Klein compactification, i.e by using the above field expansion and integrating over the extra dimensions, the action (5) reduces to an effective 4-dimensional theory. A prominent feature of this effective theory is its complexity due to extra terms involving the massive Kaluza-Klein (KK) excitations, $`\widehat{\mathrm{\Phi }}^{(\stackrel{}{n})}(x)`$, of all the fields propagating in the extra dimensions, apart from the usual, massless zero-modes, $`\widehat{\mathrm{\Phi }}^{(0)}(x)`$. Here, we are only interested in the part of the effective theory that contains the zero modes of the various fields and, more specifically, in the relations that hold between the $`(4+d)`$ and 4-dimensional couplings and masses. These are found to be
$$g=\frac{\widehat{g}}{(2L\mathrm{\Lambda })^{d/2}},Y_{1,2}=\frac{\widehat{Y}_{1,2}}{(2L\mathrm{\Lambda })^{d/2}},$$
(7)
$$\mu ^2=\widehat{\mu }^2,\lambda =\frac{\widehat{\lambda }}{(2L\mathrm{\Lambda })^d},$$
(8)
where we have used the following field redefinitions
$$A_\mu ^a=(2L)^{d/2}\widehat{A}_\mu ^a,\varphi =(2L)^{d/2}\widehat{\varphi }$$
(9)
in order to obtain canonical kinetic terms in 4-dimensions. Note that the fermionic fields being always localized on the 4-dimensional boundary remain unchanged.
In the framework of the 4-dimensional effective theory, the gauge bosons and fermions acquire mass through the spontaneous symmetry breaking of the gauge group. Their masses are given in terms of the vacuum expectation value of the Higgs field which is also redefined according to eq. (9). The question, then, arises whether the mass spectrum of the 4-dimensional effective theory corresponds to a different one in the framework of the original (4+d)-dimensional theory. To answer this question, we consider the masses of the fermions which are given by the following expression
$$m_{\mathrm{\Psi }_{L,R}}=Y_{1,2}\frac{\upsilon }{\sqrt{2}}=\frac{\widehat{Y}_{1,2}}{\mathrm{\Lambda }^{d/2}}\frac{\widehat{\upsilon }}{\sqrt{2}}=\widehat{m}_{\mathrm{\Psi }_{L,R}}.$$
(10)
Similar relations can also be written for the masses of the gauge bosons and the Higgs field. In all cases, the tree-level masses, that are generated via the Higgs mechanism in the 4-dimensional theory, remain unaltered when one embeds this theory in a higher-dimensional one.
However, this is not the case with the coupling constant $`\widehat{g}`$ that eventually determines the strength of the gauge interactions of the theory. The $`SU(N)`$ gauge group of the $`(4+d)`$-dimensional theory (5) could be replaced by the $`U(1)\times SU(2)\times SU(3)`$ group with the gauge field $`\widehat{A}_M^\alpha `$, the generators $`t^\alpha `$ and the coupling constant $`\widehat{g}`$ standing for each one of the corresponding quantities of the Standard Model group, i.e. $`(\widehat{B}_M,\widehat{W}_M^a,\widehat{G}_M^\alpha )`$, $`(Y/2,\tau ^a/2,\lambda ^\alpha /2)`$ and $`(\widehat{g}_1,\widehat{g}_2,\widehat{g}_3)`$, respectively. In that case, each coupling constant and gauge boson of the SM group is redefined according to eqs. (7) and (9), respectively.
By making use of the redefinition (7) for the gauge coupling constants, we find that the electric charge $`e`$ changes as follows
$$\widehat{e}=\widehat{g}_1\mathrm{cos}\widehat{\theta }_W=(2L\mathrm{\Lambda })^{d/2}g_1\mathrm{cos}\theta _W=(2L\mathrm{\Lambda })^{d/2}e,$$
(11)
where we have used the fact that the Weinberg angle $`\theta _W`$ remains unaltered since its tangent is given by the ratio of the gauge couplings $`g_1`$ and $`g_2`$. The above rescaling of the electric charge inevitably affects the strength of the electromagnetic interactions in higher dimensions. The strength of a force is classically defined as the potential energy of the corresponding interaction between two identical particles, separated by a distance equal to the particle’s Compton wavelength, compared to the energy of the rest mass of each particle. As is well known, in 4-dimensions and for two particles with mass $`m`$ and charge $`e`$, the above rule gives the result
$$\alpha _{EM}=\frac{E_{int}}{E_m}=\frac{1}{mc^2}\frac{e^2}{4\pi ϵ_0(\mathrm{}/mc)}=\frac{e^2}{4\pi ϵ_0}\frac{1}{137},$$
(12)
where we have set $`c=\mathrm{}=1`$. Generalizing the above definition of the strength of a force in $`(4+d)`$ dimensions, we obtain
$$\widehat{\alpha }_{EM}=\frac{\widehat{E}_{int}}{\widehat{E}_m}=\frac{1}{\widehat{m}}\frac{(\widehat{e}^2/\mathrm{\Lambda }^d)}{4\pi ϵ_0(1/\widehat{m})^{1+d}}=\alpha _{EM}\left(\frac{m}{M_x}\right)^d,$$
(13)
where we have used the fact that the masses of the particles do not change, according to (10), and where we have defined $`M_x(2L)^1`$. The above result clearly reveals that the strength of the electromagnetic force changes as the gauge bosons start “feeling” the extra dimensions. Moreover, the strength of the force strongly depends not only on the number and size of extra dimensions but also on the mass of the test particle, i.e. on the energy scale where the measurement takes place. It is also worth noting that the auxiliary energy scale $`\mathrm{\Lambda }`$ does not appear in the electromagnetic strength formula (13).
So far, we have not considered the gravitational interactions. We now assume that gravitons may propagate, apart from the usual 4 dimensions, in $`d^{}=\delta +d`$ extra dimensions, where $`\delta `$ is the number of transverse dimensions, with size $`2L^{}`$, felt only by gravitons. In that case, the $`(4+d^{})`$-dimensional action functional
$$S_{4+d^{}}=\frac{M_{GR}^{2+d^{}}}{16\pi }d^4xd^d^{}z\sqrt{G_{4+d^{}}}R_{4+d^{}}$$
(14)
reduces to an effective, 4-dimensional Einstein’s theory of gravity only when the following relation between the sizes of the extra dimensions and the energy scales of gravity, $`M_P`$ and $`M_{GR}`$, in 4 and $`(4+d^{})`$ dimensions, respectively,
$$(2L^{})^\delta (2L)^dM_{GR}^{2+d^{}}=M_P^2$$
(15)
holds. The strength of the gravitational interactions also changes when one introduces extra dimensions for the gravitons. Using the same rule as above, the strength of the gravitational interaction between two particles with mass $`m`$, in 4 dimensions, is given by the expression
$$\alpha _{GR}=\frac{E_{int}}{E_m}=\frac{1}{mc^2}\frac{G_Nm^2}{(\mathrm{}/mc)}=\left(\frac{m}{M_P}\right)^2,$$
(16)
where $`G_NM_P^2`$ and natural units, $`\mathrm{}=c=1`$, have been used again. The corresponding expression in $`4+d^{}`$ dimensions takes the form
$$\widehat{\alpha }_{GR}=\frac{\widehat{E}_{int}}{\widehat{E}_m}=\frac{1}{\widehat{m}}\frac{\widehat{G}_N\widehat{m}^2}{(1/\widehat{m})^{1+d^{}}}=\left(\frac{m}{M_{GR}}\right)^{2+d^{}},$$
(17)
where, now, $`\widehat{G}_N=1/M_{GR}^{2+d^{}}`$. By choosing appropriately the sizes of the transverse and longitudinal, extra dimensions, the higher-dimensional gravity scale, $`M_{GR}`$, could be much lower than the 4-dimensional one, a fact which, subsequently, changes the strength of gravity.
The question we would like to address is the following: can the gravitational and electromagnetic forcesUnder the assumption that the electromagnetic, weak and strong forces “feel” the same number $`d`$ of extra dimensions, we may assume that their strengths remain comparable at every scale even in $`(4+d)`$-dimensions. have comparable strengths in a world where both gravitons and gauge bosons feel extra dimensions? If so,
$$\widehat{\alpha }_{EM}=\widehat{\alpha }_{GR}\alpha _{EM}(\frac{m_U}{M_x})^d=\left(\frac{m_U}{M_{GR}}\right)^{2+d^{}},$$
(18)
where $`m_U`$ denotes the scale where the unification of strengths takes place. Hereafter, for simplicity, we will take $`\alpha _{EM}10^2`$. Regarding the values $`\delta `$ and $`d`$ of the extra dimensions, we may form four distinct categories :
(i) $`\delta =d=0`$ . In this case, $`M_{GR}M_P`$ and only at energies $`m_U10^{18}`$ GeV, i.e. close to the Planck scale, the strengths of the forces become equal.
(ii) $`\delta 0`$, $`d=0`$ . Here, we assume that only the gravitational fields can feel extra $`\delta `$ dimensions. Then, from eq. (18), we may easily conclude that the unification of strengths takes place at a scale $`m_U=(0.20.6)M_{GR}`$, for $`1\delta 7`$Motivated by M-theory, we assume that $`\delta `$, $`d7`$.. This means that the strength of the gravitational force becomes comparable to that of the other forces only when the measurement takes place near the gravity scale, independently of where $`M_{GR}`$ lies. The contact with M-theory is made for $`\delta =1`$ and $`M_c(2L^{})^1=10^{12}`$ GeV. Then, the unification takes place at $`m_UM_{GUT}M_{GR}10^{16}`$ GeV. In the case where $`M_{GR}1`$ TeV , the hierarchy problem is removed by bringing the gravitational scale down to the electroweak one, $`M_W`$. However, the energy gap between $`M_W`$ and the compactification scale $`M_c`$, which, for $`\delta =2`$, amounts to 15 orders of magnitude, remains unexplained.
One could resolve the above problem by increasing the number $`\delta `$ of extra dimensions accompanied with a small increase in the value of the higher-dimensional gravity scale $`M_{GR}`$. For example, allowing the compactification scale, $`M_c`$, to be close to the electroweak one, i.e. $`M_c=10^{\pm 1}M_W`$, the gravity scale lies in the range $`M_{GR}=(10^510^7)`$ GeV, for $`\delta 6`$. In this case, there is no energy gap between $`M_c`$ and $`M_W`$ while the gravity scale is still close to the electroweak scale (the ratio $`M_{GR}/M_W`$, in this case, is the same as the one of the top quark mass over the down quark mass, $`m_t/m_d10^4`$).
(iii) $`\delta =0`$ , $`d0`$ . In this case, it is the $`d`$ extra, longitudinal dimensions that “open up” for the gauge and scalar fields, as well as for gravitons, at the scale $`M_x`$. Substituting eq. (15) into eq. (18) and solving for $`m_U`$, we find that the unification of strengths takes place only near the string scale, i.e. $`m_U=M_s10^{18}`$ GeV, independently of the number $`d`$ of extra dimensions. If we take the new gravity scale, $`M_{GR}`$, to be the string scale, then, from eq. (15), we are led to the following relation between $`M_x`$ and $`M_{GR}`$:
$$M_x=10^{2/d}M_{GR}=(0.010.5)M_{GR},$$
(19)
for $`1d7`$. In this case, the compactification scale $`M_x`$ is 1 or 2 orders of magnitude smaller than the gravity scale. A quite interesting result arises for $`d=1`$. Then, eq. (19) gives the result $`M_x10^{16}`$ GeV, which coincides with the scale of the unification of gauge couplings, $`M_{GUT}`$, in the minimal supersymmetric standard model (according to the power-law mechanism for the unification of gauge couplings , the presence of the extra dimensions does not affect the unification scale $`M_{GUT}`$ when $`M_xM_{GUT}`$). The compactification scale, $`M_x`$, being very close to $`M_s`$ does not introduce any new energy scale in the theory and is sufficiently high to suppress the “dangerous” baryon number violating operators. From that point of view, the above arrangement of energy scales resembles the one that arises in the framework of string unification where the compactification scale, $`M_x`$, is very close to the string scale, $`M_s`$. This arrangement is indeed necessary in order to prevent the gauge coupling $`\widehat{\alpha }_{EM}`$ from acquiring an unacceptable small value near the unification scale, a goal which is also accomplished in our case. Although the hierarchy between the electroweak scale, $`M_W`$, and the gravity scale, $`M_{GR}`$, still remains, the existence of an extra, longitudinal dimension that “opens up” just above $`M_{GUT}`$ may provide a natural explanation for the energy gap between $`M_{GUT}`$ and $`M_s`$.
(iv) $`\delta 0`$ , $`d0`$ . This is the most general case where both gravitational and gauge fields feel a number of extra dimensions. According to eq. (18), the unification of strengths, now, takes place at the energy scale
$$m_U^{2+\delta }=\alpha _{EM}M_P^2M_c^\delta .$$
(20)
We make the natural assumption that the ultimate unification of strengths occurs near the new gravity scale, i.e. $`m_U=M_{GR}`$. Then, the combination of eqs. (15) and (20) leads to the same relation (19) between the compactification scale $`M_x`$ and the gravity scale $`M_{GR}`$. On the other hand, the scale $`M_c`$, associated with the size of the transverse dimensions felt only by gravitons, can be found from the following expression:
$$M_c^\delta =\alpha _{EM}^1\frac{M_{GR}^{2+\delta }}{M_P^2}.$$
(21)
Note that it is only the number of transverse dimensions $`\delta `$ that determines the compactification scale $`M_c`$, in terms of $`M_P`$ and $`M_{GR}`$, exactly in the same way that it is only the number of longitudinal dimensions $`d`$ that determines the corresponding expression of $`M_x`$. If we, now, choose $`M_{GR}=M_s10^{18}`$ GeV and $`d=1`$, the interesting result $`M_xM_{GUT}`$ of case (iii) arises once again. However, in order for this picture to be viable in the existence of extra, transverse dimensions for gravitons, the compactification scale $`M_c`$ should lie very close to the string scale for every value of $`\delta `$. Since $`M_c,M_xM_{GUT}`$, the unification pattern of gauge couplings is not affected by the presence of the extra dimensions, either transverse or longitudinal, and the ultimate unification of all forces takes place exactly at the string scale. Moreover, the compactification scale $`M_c`$, being very close to $`M_s`$, does not introduce any new energy scale in the theory. The only other case where this problem is also resolved is when $`M_c=10^{\pm 1}M_W`$, according to the analysis presented in case (ii). For $`\delta =6`$ and $`M_c=10`$ GeV, the gravity scale lies at the energy scale $`2\times 10^5`$ GeV. Then, for $`d=1`$, eq. (19) shows that the unification of strengths takes place only if $`M_x`$ is exactly two orders of magnitude smaller than $`M_{GR}`$, i.e. $`M_x=2`$ TeV. Remarkably, this is in perfect agreement with the proposal of the existence of a single extra dimension for the gauge bosons with size at the TeV scale necessary for the explanation of the supersymmetry breaking scale . The existence of an extra dimension of this size would lead to the unification of gauge couplings, through the power-law mechanism , at an energy scale which is, approximately, one order of magnitude larger than $`M_x`$ . The extra hypothesis, that gravitons may feel 6 extra, compact dimensions, brings the scale of gravity down to $`200`$ TeV, thus, completing the picture of unification of forces and removing the hierarchy problem (by changing appropriately the values of $`\delta `$ and $`M_c`$, we may easily recover the case where $`M_{GUT}=M_{GR}=M_s=10`$ TeV ).
Finally, let us stress that all the results presented above can be consistently embedded in a Type I string theory framework. Combining and rearranging eqs. (15) and (18), we can write our unification constraint in the form
$$(\frac{m_U}{M_P})^2=\widehat{\alpha }_{EM}(m_U\mathrm{\hspace{0.17em}2}L^{})^\delta (m_U\mathrm{\hspace{0.17em}2}L)^d,$$
(22)
which reduces to the Type I string constraint $`M_s^2(\alpha _{gauge}M_P\sqrt{V})^1`$ if we identify $`m_U`$ with $`M_s`$, and $`\widehat{\alpha }_{EM}`$ with $`\alpha _{gauge}^2`$. The total compactified volume is now $`V=(2L)^d(2L^{})^\delta `$ and a T-duality transformation $`m_Ur(m_Ur)^1`$ needs to be performed, where $`r2L,2L^{}`$.
In conclusion, in this letter, we have demonstrated that the addition of extra, compact dimensions for gravitons and gauge bosons modify the strengths of gravitational and gauge interactions and consequently the scale of their unification. When the gravitons propagate in $`4+\delta `$ dimensions, the unification takes place only near the higher dimensional gravity scale, $`M_{GR}`$. If we choose $`M_cM_W`$, no new scale is introduced in the theory while $`M_{GR}`$ turns out to be a few orders of magnitude larger than $`M_W`$. In the case where the gauge bosons feel $`d`$ extra dimensions, the unification occurs only near the string scale, $`M_s`$. If $`M_{GR}M_s`$, $`M_x`$ turns out to be, for $`d=1`$, of the order of the gauge coupling unification scale. Thus, the energy gap between $`M_s`$ and $`M_{GUT}`$ is attributed to the existence of one extra dimension for the gauge bosons which “opens up” just above $`M_{GUT}`$. An even more remarkable result arises in the case where both $`\delta `$ and $`d`$ are non-zero under the assumption that the gravitons can feel 6 extra dimensions with $`M_c=10`$ GeV. Then, for $`d=1`$, $`M_x`$ has exactly the right value to explain the size of the supersymmetry breaking scale with the gravity scale being only 2 orders of magnitude larger than $`M_x`$. In the framework of this 11-dimensional theory, the unification of all forces takes place at $`200`$ TeV and the hierarchy problem is completely removed. It is worth noting that all the results derived from the above analysis, based on a non-renormalizable, effective field theory for gauge and gravitational interactions, can be consistently embedded in the framework of string theory unification.
The authors would like to thank K.R. Dienes, H.P. Nilles, R.G. Roberts, G.G. Ross and K. Tamvakis for fruitful discussions and comments. A.D is supported from Marie Curie Research Training Grants ERB-FMBI-CT98-3438. P.K. is grateful to the Particle Physics Group at Rutherford Appleton Laboratory for the warm hospitality and financial support. |
warning/0003/astro-ph0003126.html | ar5iv | text | # Broad-Band X-Ray Study of a Transient Pulsar RX J0059.2-7138
## 1. Introduction
Accretion-powered X-ray pulsars are binary systems consisting of a neutron star and a stellar companion (here, X-ray binary pulsars or XBPs in short). The gravitational energy of accreting matter is converted to X-ray radiation, hence its luminosity is generally variable, depending on the mass-accretion rate. Transient behavior is also observed from many XBPs, most of which have a Be star companion with an eccentric orbit (e.g. Stella et al. 1982; Bildsten et al. 1997).
XBPs have been observed mainly in the hard X-ray band ($`2`$–40 keV) with non-imaging satellites (e.g. Nagase 1989). Their spectra are generally described by a power-law with a high-energy exponential cut-off (hereafter, ECUT power-law) with a photon index of $`\mathrm{\Gamma }1`$ and a cut-off energy at around 10–20 keV. A fluorescent emission line at 6.4 keV from neutral iron atoms has been observed in many bright XBPs, which serves as a probe of the circumstellar matter of the binary system. In the soft X-ray band (below $`2`$ keV), however, their spectra have not been well studied due to a large interstellar absorption, because most of the XBPs are located in the galactic plane (e.g. Bildsten et al. 1997).
In order to study soft X-rays from XBPs, therefore, the Magellanic Cloud sources have great advantages, owing to the low interstellar absorption and well-calibrated distances. Woo et al. (1995, 1996) made broad-band spectroscopic studies in the 0.1–40 keV band of SMC X-1 and LMC X-4, respectively, and found a “soft excess” below $`2`$ keV, in addition to the usual hard spectrum given by an ECUT power-law model. Such a soft excess was also discovered from a transient XBP XTE J0111.2$``$7317 in the Small Magellanic Cloud (SMC), by an ASCA observation in the 0.5–10.0 keV band (Yokogawa et al. 2000). However, since such a broad-band spectroscopy has been performed on only a few XBPs, whether or not the soft component is common among XBPs is still unclear, especially for transient XBPs.
Hughes (1994) serendipitously discovered a new transient XBP, RX J0059.2$``$7138, in the SMC with the soft X-ray band observation of ROSAT, and found that the spectrum is unusually soft, which is composed of a blackbody with a temperature of $`kT35`$ eV and a steep power-law with a photon index of $`\mathrm{\Gamma }2.4`$. A proposed candidate for the optical counterpart was revealed to be a Be star by Southwell and Charles (1996).
In order to examine the nature of the unusually soft spectrum of RX J0059.2$``$7138, observations in the higher energy band are essential. RX J0059.2$``$7138 was also detected in a simultaneous observation of the SNR E 0102$``$72.2 with the ASCA satellite, which is sensitive in the $`0.7`$–10.0 keV band. A preliminary short report on the ASCA results can be found in Kylafis (1996).
In this paper, we combine the ROSAT and ASCA data, and perform a broad-band timing spectroscopy, covering $`0.1`$–10.0 keV. We show that the spectrum of RX J0059.2$``$7138 has a hard component, which resembles the spectra of usual XBPs well, and that the pulsations are due to the hard component. In addition, a non-pulsating soft component is present in the spectrum. We assume the distance to RX J0059.2$``$7138 to be 60 kpc, the nominal value to the SMC (Mathewson 1985).
## 2. Observations
The transient XBP RX J0059.2$``$7138 was serendipitously detected on May 12–13 in 1993, in the simultaneous ASCA and ROSAT observations of E 0102$``$72.2, the brightest supernova remnant in the SMC. ASCA and ROSAT observations spanned about 99 ks (MJD 49119.38–MJD 49120.52) and 20 ks (MJD 49119.95–MJD 49120.18), respectively.
ASCA carries four X-ray Telescopes (XRT) with two Gas Imaging Spectrometers (GIS 2 and GIS 3) and two Solid-state Imaging Spectrometers (SIS 0 and SIS 1) on the focal planes (Tanaka et al. 1994; Serlemitsos et al. 1995; Ohashi et al. 1996; Burke et al. 1994). Since RX J0059.2$``$7138 was outside of the SISs’ field of view and was located near the calibration source of GIS 3, we used only the GIS 2 data in this paper. The GISs were operated in the normal PH mode with a time resolution of 0.0625 s (high bit rate) or 0.5 s (medium bit rate). We rejected any data obtained when the satellite was in the SAA (South Atlantic Anomaly) region, when the elevation angle of the Earth limb was less than $`5^{}`$ or when the cut-off rigidity was less than 4 GV, and also removed particle events using the rise-time discrimination method. The total exposure time of GIS 2 was $``$35 ks after screening.
ROSAT carries a soft X-ray telescope with a High Resolution Imager (HRI) and a Position Sensitive Proportional Counter (PSPC) with one of the two on the focal plane (Trümper 1983). In this observation, PSPC-B was on the focus and operated in the pointing mode. We took the screened data from the HEASARC archive, of which the total exposure time was $`5`$ ks.
## 3. Analyses and Results
For the GIS data, X-ray photons were collected from an elliptical region around RX J0059.2$``$7138 (figure 1). The background data were extracted from off-source areas on the edge of GIS 2 at the same off-axis angle as the source (figure 1), because both the diffuse X-rays and non-X-ray background depend on the off-axis angle.
As for the ROSAT PSPC, we selected a circular region of $`3^{}`$-radius around the source, while the data from an annulus of inner and outer radii $`3^{}`$ and $`6^{}`$ were used as the background. Since a super-soft source (SSS) 1E 0056.8$``$7154 was located at $`4^{}`$ northwest of RX J0059.2$``$7138 and the X-ray flux of this SSS was fairly large in the ROSAT band below $`0.5`$ keV, we excluded a circular region of a $`3^{}`$-radius around this SSS from the analyses of the ROSAT PSPC.
### 3.1. Timing Analysis
In total, the numbers of counts from the source regions are 21847 in the 0.7–10.0 keV band (ASCA GIS 2) and 23250 in the 0.07–2.4 keV band (ROSAT PSPC). After the barycentric correction of arrival time for each event, we carried out timing analyses.
The pulse period was already determined to be 2.7632$`\pm 0.0002`$ s with the ROSAT data alone (Hughes 1994). We used the ASCA data to obtain a more precise period, because the duration of the ASCA observation is nearly five-times longer than the ROSAT observation. We performed epoch-folding on all of the ASCA GIS 2 data, and determined a trial pulse period to be 2.76322 s. We then divided the ASCA observation into three segments, each having a duration of $`33`$ ks, and folded each of the light curves with a trial period of 2.76322 s.
To the folded pulse profiles of the first and the last segments, we made a cross correlation, and determined the average apparent barycentric pulse period during the observation to be 2.763221$`\pm `$0.000004 s. This value is consistent with, and more precise than, the previous result of 2.7632$`\pm 0.0002`$ s (Hughes 1994).
The cross correlation between the middle and the first/last segments shows the upper limit of the phase difference to be $`0.04`$. Based on the assumption of a constant period change during the observation, this upper limit of the phase difference is converted to the upper limit of the period derivative of $`|\dot{p}|<6\times 10^{10}`$ s s<sup>-1</sup>.
Figure 2 shows the pulse profiles in the six energy bands (0.07–0.4 keV, 0.4–1.0 keV, and 1.0–2.4 keV in ROSAT PSPC, and 0.7–2.0 keV, 2.0–4.0 keV, and 4.0–10.0 keV in ASCA GIS) folded with a pulse period of 2.763221 s. The pulse fractions, defined as $`(I_{\mathrm{max}}I_{\mathrm{min}})/2I_{\mathrm{average}}`$, in these energy bands are 0.07$`\pm `$0.06, 0.14$`\pm `$0.04, 0.26$`\pm `$0.04 (0.07–0.4 keV, 0.4–1.0 keV, and 1.0–2.4 keV in ROSAT PSPC) and 0.24$`\pm `$0.05, 0.37$`\pm `$0.05, and 0.36$`\pm `$0.05 (0.7–2.0 keV, 2.0–4.0 keV, and 4.0–10.0 keV in ASCA GIS), respectively. The pulse shape is nearly sinusoidal and independent of the X-ray energy. The apparent decrease of the pulse fraction with decreasing energy, which was already found in the ROSAT data (Hughes 1994), is also seen in the ASCA data below $`2`$ keV.
On a longer timescale (50–100 s), Hughes (1994) reported a flickering variability of about 30% of the mean rate in the ROSAT data. Such variability was also seen in the ASCA data with no energy dependence, while no larger variability, like a burst, was detected.
### 3.2. Spectral Analysis
Figure 3 shows the phase-averaged spectra after background subtraction. Since the pulse profile is essentially energy independent while the pulse fraction increases with energy below $`2`$ keV, we can naturally assume that the spectrum of RX J0059.2$``$7138 consists of two components: a hard component dominant above $`2`$ keV, which has an energy-independent pulsation, and a soft component dominant below $`2`$ keV with no pulsation.
With this working hypothesis, we first determined the hard component using the ASCA spectrum above $`2`$ keV. Since no significant emission line was detected, we fitted the spectrum with a single power-law model. However, because systematic negative residuals remained above $`7`$ keV, we introduced a high-energy exponential cut-off to the single power-law model (ECUT power-law) and fitted the spectrum again. We then obtained an acceptable fit to an ECUT power-law model with a photon index of $`\mathrm{\Gamma }0.4`$, a cut-off energy of $`E_\mathrm{c}6.6`$ keV and a folding energy of $`E_\mathrm{f}8.8`$ keV.
To investigate the soft component, we then simultaneously fitted the ROSAT and ASCA spectra in the 0.1–10.0 keV band with various two-component models, in which an ECUT power-law model was always adopted to the hard component. For the additional soft component, we adopted a blackbody, disk blackbody, bremsstrahlung, thin-thermal plasma (Mewe et al. 1985), and broken power-law model. In these two-component model fits, a common absorbing gas column with solar abundance was adopted for both the hard and soft components. All models were statistically rejected at the 90% confidence level with wavy residuals below $`1`$ keV. In particular, an edge-like structure was seen at $`0.5`$ keV, which may correspond to the absorption edge of neutral oxygen atoms. Allowing the oxygen abundance in the absorption column to be free, we obtained statistically acceptable fits only when the soft component was either a thin-thermal plasma or a broken power-law. The best-fit parameters for these models are separately listed in tables 1 and 2, while the best-fit model is shown in figure 3.
We separately made phase-resolved spectra from phases 0–0.5 (‘on-pulse’) and 0.5–1 (‘off-pulse’). We fit these spectra with the two accepted models, and found that only the normalization of the hard component varies with the pulse phase. The best-fit parameters of these models are also listed in tables 1 and 2, while the best-fit model is shown in figure 4. We also extracted the “pulsed spectrum” by subtracting the off-pulse spectrum from that of on-pulse. We found that the pulsed spectrum is fitted with an ECUT power-law with absorption, of which the best-fit parameters are consistent with those of the hard component of the phase-averaged spectrum. These facts indicate that the pulsating X-rays contribute only to the hard component, consistent with the initial working hypothesis.
## 4. Discussion
### 4.1. Comparison to the Previous Results
Hughes (1994) analyzed only the ROSAT data, and showed that the spectrum is fitted by a model having an extremely soft spectrum: a blackbody ($`kT35`$ eV) plus a steep power-law ($`\mathrm{\Gamma }2.4`$). However, we have shown that the spectrum in the ASCA band (the hard component) is as hard as those of usual XBPs, and the ROSAT spectrum consists of mostly the “soft excess” below $`2`$ keV.
For a consistency check, we first fitted our ROSAT spectrum with Hughes’ model, and obtained best-fit parameters consistent with Hughes (1994). However, as can be also seen in figure 2 in Hughes (1994), systematic positive residuals were found above $`1.5`$ keV, which implies the existence of a hard component. We then fitted the ASCA and ROSAT spectra simultaneously with a three-component model, consisting of the hard ECUT power-law to the Hughes’ two-component model. In this model, the best-fit blackbody temperature became much lower than the original result. After corrections for absorption and detector efficiency, the bolometric luminosity of the soft black body has an unrealistic value of $`10^{45}\mathrm{erg}\mathrm{s}^1`$.
On the other hand, as we found, the combined ASCA and ROSAT spectra can be fitted with two different components, with a more reasonable luminosity of $`10^{38}\mathrm{erg}\mathrm{s}^1`$ for both components. Also, the energy-dependent pulse fraction and energy-independent pulse profile are naturally explained by our two-component model. We thus infer that the present model is more probable, hence the spectrum of RX J0059.2$``$7138 is not unusual for an XBP, at least in the hard energy band.
### 4.2. Oxygen Over-Abundance
We found an edge-like structure at 0.5 keV, which would be a hint of over-abundance of oxygen in the absorbing matter. One may argue, however, that the ROSAT spectrum fitted by Hughes (1994) does not show any structure at 0.5 keV. We should note that the blackbody and power-law components in the Hughes’s model cross at around 0.5 keV, which may produce an artificial dip, and hence would compensate for the edge-like structure.
To check whether the structure is real or not, we fitted a narrow-band spectrum (0.2–0.8 keV) with a power-law absorbed by cool matter of solar abundance. Because the structure remained near 0.5 keV, we claim that the edge-like structure is real.
Since the oxygen abundance of the SMC interstellar matter is less than the solar value (Russell, Dopita 1992), we can infer that there exists oxygen-rich matter around the binary system of RX J0059.2$``$7138. The companion star, which should have ejected the oxygen-rich matter, may be in an evolved stage where the stellar surface becomes abundant only in oxygen, although such stars (e.g. WO stars) are very rare.
On the other hand, metals such as carbon, nitrogen, and oxygen could be brought to the surface of a massive star by turbulent diffusion, due to the rapid rotation or tides in a binary system (Maeder 1987). We thus fitted the spectra again, based on the assumption that abundances of carbon and nitrogen in the absorbing matter are the same as that of oxygen. The abundance was then determined to be $`6`$ solar, which leads us to a more comfortable possibility that the surface of the companion star is enriched in CNO, which may be caused by tidal effects in the close binary system.
We found that a thin-thermal plasma with very low abundance of 0.02 solar can fit the soft X-ray spectrum. As we discuss in the next subsection, the soft component is likely to originate from a large region surrounding the binary system. This creates a dilemma that the X-ray emitting circumstellar medium is extremely under-abundant, while the X-ray absorbing circumstellar gas is CNO over-abundant. Although we have no clear idea of how to solve this dilemma, we suggest that the thin-thermal plasma is not a physical model, but is a phenomenological model like a broken power-law for the soft component.
### 4.3. Origin of the Hard and Soft Components
Since the pulsation was found to be mostly due to the hard component, we can conclude that the hard component originates from a small region, probably near to the polar cap of the neutron star. Hence, the origin would be the same as those of other XBPs. However, the cut-off energy is lower than that of the usual XBPs (10–20 keV; e.g. Mihara 1995a). According to the correlation between the cut-off energy and magnetic field strength (Mihara 1995a, 1995b), the magnetic field of RX J0059.2$``$7138 is estimated to be $`6\times 10^{11}<B<1\times 10^{12}`$ G, which is one of the weakest magnetic fields in XBPs.
A puzzle is the origin of the soft component. Since it shows no pulsation, it is likely that the soft component originates from a relatively large region, covering a significant fraction of the full binary system. In fact, when we adopt the thin-thermal plasma model as the soft component, the emission measure is very large: $`n^2V10^{61}\mathrm{cm}^3`$ (table 1), where $`n`$ and $`V`$ are number density and volume of the plasma. If we assume that the plasma is distributed spherically, its radius should be larger than 0.07 AU under the optically thin condition, which does not conflict with the thin thermal scenario for the soft component.
At present, only a few XBPs are known to have a soft excess. Her X-1 has a soft component which is described by a blackbody ($`kT0.09`$ keV; Dal Fiume et al. 1998), and is pulsating. McCray et al. (1982) found the phase shift between pulse profiles of the soft and the hard components, and argues that the soft component is produced by reprocessing of the hard component in the inner accretion disk. Another XBP, LMC X-4 has a soft component of a blackbody ($`kT0.03`$ keV) plus a bremsstrahlung ($`kT0.4`$ keV) (Woo et al. 1996). The blackbody component is almost constant during the pulse phase, while the bremsstrahlung component is pulsating with some phase shift from the hard component. Woo et al. (1996) argued that the blackbody component is emitted from the accretion disk, and the bremsstrahlung component is from somewhere near the neutron star, collimated as a fan-beam.
The soft component of RX J0059.2$``$7138 is, however, different from these XBPs; it is not pulsating, it is not a blackbody spectrum and it has no emission line. Hence, the origin of the soft X-rays would be different from those proposed by the above-mentioned authors. As a matter of fact, RX J0059.2$``$7138 is a unique source among the XBPs with soft excess; because it probably has a Be star companion (Southwell, Charles 1996), it is different from either low-mass companion systems (Her X-1) or from high-mass supergiant companion systems (LMC X-4).
Broad-band spectroscopic studies on the spectra of XBPs with a Be star companion have until now been very limited, due to the transient nature of this class and the relatively narrow energy ranges of the previous and contemporaneous satellites. New-generation satellites having a broad-band sensitivity and a wide field of view, such as ASTRO-E, XMM, and Chandra Observatory, will greatly improve our understanding on the soft excess of XBPs.
We express our thanks to J. Hughes, who told us about the serendipitous detection of RX J0059.2$``$7138 in the GIS field and the referee for his useful comments. J.Y. is supported by JSPS Research Fellowship for Young Scientists. This research has made use of data obtained through the High Energy Astrophysics Science Archive Research Center Online Service, provided by the NASA/Goddard Space Flight Center.
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warning/0003/quant-ph0003017.html | ar5iv | text | # Non-Kolmogorov probability models and modified Bell’s inequality
## Abstract
We analyse the proof of Bell’s inequality and demonstrate that this inequality is related to one particular model of probability theory, namely Kolmogorov measure-theoretical axiomatics, 1933. We found a (numerical) statistical correction to Bell’s inequality. Such an additional term $`ϵ_\varphi `$ in the right hand side of Bell’s inequality can be considered as a probability invariant of a quantum state $`\varphi .`$ This is a measure of nonreproducibility of hidden variables in different runs of experiments. Experiments to verify Bell’s inequality can be considered as just experiments to estimate the constant $`ϵ_\varphi .`$ It seems that Bell’s inequality could not be used as a crucial reason to deny local realism. We consider deterministic as well as stochastic hidden variables models.
1. Introduction
Experimental violations of Bell’s inequality are typically (see, for example, ,) interpreted in one of two ways: (1) nonlocality: by changing the state of one particle in the EPR pair we change the state of the other particle; (2) death of reality: realism could not be used as the philosophic base of quantum mechanics (‘properties’ of quantum systems are not objective properties, i.e., the properties of an object). In particular, (2) implies that the statistical interpretation of quantum mechanics (via L. Ballentine, ) must be denied in favour of the orthodox Copenhagen interpretation. Although such a viewpoint is dominating in the quantum community, there are still some doubts that violations of Bell’s inequality must be interpreted in such a way. In particular, many scientists thought (and continue to think) that “Bell’s paradox” has purely probabilistic origin, see, for example, . Unfortunately these probabilistic considerations had merely philosophic character. In any case they did not give a new (Bell-like) inequality which has an experimental meaning.
Remark 1.1. The common opinion is that “Bell’s paradox” (experimental violations of Bell’s inequality) is just a reformulation of the EPR paradox. However, the problem is more complicated: Bell’s probabilistic reformulation of the EPR paradox contains some additional assumptions (on probability distribution of hidden variables).
In this note we follow to the general (probabilistic) attitude of . In fact, we generalize ideas of De Baere, , on connection of Bell’s inequality and “implicit reproducibility” (we arrived to such ideas independently by developing non-Kolmogorov probability formalisms, ). However, we found some statistical quantity $`ϵ_\varphi `$ which can be considered as a probability invariant of a quantum state $`\varphi .`$ It seems that the standard interpretation of violations of Bell’s inequality is a consequence of neglecting of this quantity $`ϵ_\varphi .`$
We analyse Bell’s proof and demonstrate that the possibility to derive Bell’s inequality depends crucially on the use of a particular probabilistic model, namely the model based on Kolmogorov’s axiomatics , 1933 (so called measure-theoretical approach to probability). Of course, the Kolmogorov model is the dominating mathematical model for probability theory. Therefore it is not surprising that J. Bell and many others used this approach to probability theory. However, there also exist numerous non-Kolmogorov probabilistic models (which are similar to non-Euclidean geometrical models), see, for example, and Accardi and Gudder in . In particular, I constructed probabilistic models, , which describe random phenomena in that the standard law of large numbers is violated: relative frequencies $`\nu _N=n/N`$ have no limit (stabilization), $`n\mathrm{}`$ (numerous examples of such a random behaviour can be found in ).
Remark 1.2. We want to underline that in physics the choice of the right probability model is not less important then the choice of the right geometric model. The Kolmogorov model (as well as Euclidean model) could not describe all physical phenomena.
In this note we analyse “Bell’s paradox” on the basis of the assumption that the law of large numbers can be violated for hidden variables: $`\nu _N(\lambda )`$ can fluctuate. From the physical viewpoint this means that different runs of experiments (for example, for correlated particles) can produce different “probability distributions” for hidden variables. In such a situation it would be impossible to define a Kolmogorov probability distribution $`𝐏`$ on the set of hidden variables $`\mathrm{\Lambda }.`$ Kolmogorov’s model could not be applied. We introduce a numerical measure for fluctuations $`ϵ_\varphi `$. It will be shown that “general Bell’s inequality” must contain this probability invariant of a quantum state as an additional term.
Finally, we remark that all experimental calculations are, in fact, based not on Kolmogorov model (probability as a measure), but on von Mises, 1919, model (probability as frequency). “Experimental covariation” of two observables, $`A,B,`$ is calculated as a sequence mean value: $`<A,B>_{\mathrm{fr}}=(1/N)_{i=1}^NA_iB_i,`$ where $`x=(A_1,A_2,\mathrm{},A_N)`$ and $`y=(B_1,B_2,\mathrm{},B_N)`$ are random sequences (collectives) generated by measurements of $`A`$ and $`B.`$ Hence, in fact, $`<A,B>_{\mathrm{fr}}`$ depends on $`x`$ and $`y`$: $`<A,B>_{\mathrm{fr}}=<A,B>_{xy}.`$ However, J. Bell supposed that there exists a Kolmogorov probability distribution $`𝐏`$ on the set of hidden variables $`\mathrm{\Lambda }`$ and all covariations can be written as mean values with respect to this unique measure: $`<A,B>_{\mathrm{Bell}/\mathrm{Kol}}=_\mathrm{\Lambda }A(\lambda )B(\lambda )d𝐏(\lambda ).`$ This (rather strong) assumption statistical postulate has never been verified experimentally.
2. Bell’s proof
We reproduce the proof of Bell’s inequality. Let $`𝒫=(\mathrm{\Omega },,𝐏)`$ be a Kolmogorov probability space: $`\mathrm{\Omega }`$ is a space of elementary events, $``$ is an algebra of events, $`𝐏`$ probability measure.
Theorem 1. Let $`A,B,C=\pm 1`$ be random variables on $`𝒫.`$ Then Bell’s inequality
$$|<A,B><C,B>|1<A,C>$$
(1)
holds true.
Proof. Set $`\mathrm{\Delta }=<A,B><C,B>.`$ By linearity of Lebesgue integral we obtain
$$\mathrm{\Delta }=_\mathrm{\Omega }A(\omega )B(\omega )𝑑𝐏(\omega )_\mathrm{\Omega }C(\omega )B(\omega )𝑑𝐏(\omega )=_\mathrm{\Omega }[A(\omega )C(\omega )]B(\omega )𝑑𝐏(\omega ).$$
(2)
As $`A(\omega )^2=1,`$
$$|\mathrm{\Delta }|=|_\mathrm{\Omega }[1A(\omega )C(\omega )]A(\omega )B(\omega )𝑑𝐏(\omega )|_\mathrm{\Omega }[1A(\omega )C(\omega )]𝑑𝐏(\omega ).$$
(3)
Of course, this is the rigorous mathematical proof of (1) for Kolmogorov probabilities. However, abstractness of Kolmogorov’s probability model induces serious problems, if we do not control carefully dependence of probabilities on corresponding statistical ensembles of physical systems. Bell did not control this dependence. In fact, the symbol $`𝐏`$ of probability which is used in the proof must be regarded to different statistical ensembles.
3. Fluctuating distributions of hidden variables
To simplify our considerations, we suppose that the set of hidden variables is finite: $`\mathrm{\Lambda }=\{\lambda _1,\mathrm{},\lambda _M\}.`$ For each physical observable $`U`$, the value $`\lambda `$ of hidden variables determines the value $`U=U(\lambda ).`$ Let $`U`$ and $`V`$ be physical observables, $`U,V=\pm 1.`$ We start with the consideration of the frequency (experimental) covariation $`<U,V>_{x_{UV}}`$ with respect to a random sequence $`x_{UV}=(x_1,x_2,\mathrm{},x_N,\mathrm{}),`$ where $`x_i=(u_i,v_i),`$ which is induced by measurements of the pair $`(U,V).`$ The $`x_{UV}`$ is obtained by measurements for an ensemble $`S_{UV}`$ of physical systems (for example, pairs of correlated quantum particles). Our aim is to represent experimental covariation $`<U,V>_{x_{UV}}`$ as ensemble covariation $`<U,V>_{S_{UV}}.`$ Then we shall demonstrate that in the general case it is impossible to perform for ensemble covariations Bell’s calculations which have been performed for Kolmogorov covariations. Let $`S_{UV}=\{d_1,\mathrm{},d_N\},`$ where $`i`$th measurement is performed for the system $`d_i.`$ Define a function $`i\lambda (i),`$ the value of hidden variables for $`d_i.`$ We set $`n_k(S_{UV})=|\{d_iS_{UV}:\lambda (i)=\lambda _k\}|`$ and $`𝐩_k^{UV}=𝐏_{S_{UV}}(\lambda =\lambda _k)=\frac{n_k(S_{UV})}{N}.`$ These are probabilities of hidden variables $`\lambda _k,k=1,2,\mathrm{},M,`$ in the statistical ensemble $`S_{UV}.`$ We have $`<U,V>_{x_{UV}}=\frac{1}{N}_{i=1}^NU(\lambda (i))V(\lambda (i))=_{k=1}^M𝐩_k^{UV}u_kv_k=<U,V>_{S_{UV}},`$ where $`u_k=U(\lambda _k),v_k=V(\lambda _k).`$ Thus
$$\mathrm{\Delta }=<A,B>_{x_{AB}}<C,B>_{x_{CB}}$$
$$=<A,B>_{S_{AB}}<C,B>_{S_{CB}}=\underset{k}{}(𝐩_k^{AB}a_k𝐩_k^{CB}c_k)b_k$$
and
$$<A,C>_{x_{AC}}=<A,C>_{S_{AC}}=\underset{k}{}𝐩_k^{AC}a_kc_k.$$
We now suppose that probabilities of $`\lambda _k`$ do not depend on statistical ensembles:
$$𝐩_k=𝐩_k^{AB}=𝐩_k^{CB}=𝐩_k^{AC}$$
(4)
(later we shall modify this condition to obtain statistical coincidence of probabilities, instead of the precise coincidence). Hence $`\mathrm{\Delta }=_{k=1}^M𝐩_k(a_kc_k)b_k\text{and}<A,C>_{x_{AC}}=_{k=1}^M𝐩_ka_kc_k.`$ We can now apply Theorem 1 for the discrete probability distribution $`\{𝐩_k\}_{k=1}^M`$ and obtain Bell’s inequality.
However, if condition (4) does not hold true, then equality (2) and, as a consequence, Bell’s inequality can be violated. The violation of condition (4) is the exhibition of unstable statistical structure on the level of hidden variables. Condition (4) is equivalent to a condition of implicit reproducibility which was discussed by De Baere .
Remark 3.1. ($`p`$-adic probability models, negative probabilities and Bell’s inequality). All our considerations were based on the statistical stabilization with respect to the real metric. In we considered the statistical stabilization with respect to a $`p`$-adic metric. The field of $`p`$-adic numbers $`𝐐_p,`$ where $`p>1`$ is a prime number, can be constructed (as the field of real numbers $`𝐑`$) as a completion of the field of rational numbers $`𝐐.`$ The $`p`$-adic metric differs strongly from the real one. As for finite ensembles $`S,`$ ensemble probabilities $`𝐏_S(a)=\frac{n(a)}{N}`$ are rational numbers, we can study their behaviour not only with respect to the real metric on $`𝐐,`$ but also with respect to the $`p`$-adic metric. $`p`$-adic probability theory gives numerous examples of ensemble probabilities fluctuating in the real metric and stabilizing in the $`p`$-adic metric. However, the $`p`$-adic stabilization of probabilities does not imply the possibility to repeat Bell’s proof for $`p`$-adic probabilities: these probabilities may be negative rational numbers, see (compare with Muckenheim, ).
4. Measure of statistical deviation between runs of an experiment
We introduce now a statistical analogue of the precise coincidence of ensemble probabilities for hidden variables. Let $`_1,_2`$ be two ensembles of physical systems and let $`\pi `$ be a property of elements of these ensembles. The $`\pi `$ has values $`(\alpha _1,\mathrm{},\alpha _m).`$ We define
$$\delta _\pi (_1,_2)=\underset{i=1}{\overset{M}{}}|𝐏__1(\alpha _i)𝐏__2(\alpha _i)|,$$
where $`𝐏_{}(\alpha _i)=\frac{|\{d:\pi (d)=\alpha _i\}|}{||}`$ are ensemble probabilities. We remark that the function $`\delta =\delta _\pi `$ is a pseudometric on the set of all ensembles which elements have the property $`\pi :\mathrm{\hspace{0.33em}1})\delta (_1,_2)0;2)\delta (_1,_2)=\delta (_2,_1);3)\delta (_1,_2)\delta (_1,_3)+\delta (_3,_2).`$ In our model we set $`\pi =\lambda ,`$ hidden variables. The precise reproducibility of the probability distribution of hidden variables (4) can be written as
$$\delta (S_{AB},S_{CB})=\delta (S_{AB},S_{AC})=0,$$
where $`\delta =\delta _\lambda .`$ Of course, we need not use such a precise coincidence in probabilistic considerations. Let $`\varphi `$ be a quantum state. Denote by the symbol $`T_\varphi `$ the set of all statistical ensembles $``$ which correspond to $`\varphi `$ (can be obtained with the aid of some preparation procedure corresponding to $`\varphi ).`$ Set
$$ϵ_\varphi =sup\{\delta (_1,_2):_1,_2T_\varphi \}.$$
Theorem 2. (“General Bell’s inequality”) Let $`\varphi `$ be a quantum state and let $`A,B,C`$ be physical observables such that pairs of observables $`(A,B),(C,B)`$ and $`(A,C)`$ can be measured. Then inequality
$$|<A,B><C,B>|(1+2ϵ_\varphi )<A,C>$$
(5)
holds true.
Proof. We have
$$|\mathrm{\Delta }|=|<A,B>_{x_{AB}}<C,B>_{x_{CB}}|$$
$$|\underset{k=1}{\overset{M}{}}𝐩_k^{AB}(a_kc_k)b_k|+|\underset{k=1}{\overset{M}{}}(𝐩_k^{AB}𝐩_k^{CB})c_kb_k|$$
$$ϵ_\varphi +\underset{k=1}{\overset{M}{}}𝐩_k^{AB}|a_kb_k|(1a_kc_k)(1+ϵ_\varphi )<A,C>_{S_{AC}}+\underset{k=1}{\overset{M}{}}|𝐩_k^{AC}𝐩_k^{AB}||a_kc_k|$$
$$(1+2ϵ_\varphi )<A,C>_{S_{AC}}.$$
We use the index $`N`$ to denote the cardinality of a statistical ensemble. If probabilities $`𝐏_{S_{UV}^N}(\lambda _k)`$ stabilize when $`N\mathrm{},`$
$$\underset{N\mathrm{}}{lim}𝐏_{S_{UV}^N}(\lambda _k)=𝐏(\lambda _k),$$
then $`ϵ_\varphi ^N0,N\mathrm{}.`$ This imply precise Bell’s inequality (1).
Experiments to verify Bell’s inequality can be considered as experiments to estimate the probability invariant $`ϵ_\varphi `$ for some class of quantum states. It seems that the only lesson of these experiments is that there exist quantum states $`\varphi `$ which have nonzero probability invariant $`ϵ_\varphi .`$ It may be that physical reality is nonlocal. It may be that it is even nonreal. However, it seems that Bell’s arguments did not imply neither nonlocality nor nonreality.
5. Stochastic hidden variables model
In this section it is supposed that the result of a measurement depends not only on the value $`\lambda `$ of hidden variables, but also on the state $`\omega ^U`$ of an equipment $`_U`$ which is used for measuring of $`U.`$ This the empiricists (contextualistic realists) interpretation of quantum mechanics, see, for example, W. De Muynck, W. De Baere, H. Marten, Ref. .
A measurement device $`_U`$ is a complex macroscopic system which state depends on the huge number of fluctuating parameters. Denote the ensemble of all possible states of $`_U`$ by the symbol $`\mathrm{\Sigma }_U:\mathrm{\Sigma }_U=\{\omega _1^U,\mathrm{},\omega _{L_U}^U\}.`$ The final value $`U_f`$ of an observable $`u`$ depends on both $`\lambda `$ and $`\omega :`$
$$u=U(\omega ,\lambda ).$$
We call such a model stochastic hidden variables model. Our model of stochastic hidden variables differs from the standard one, see section 6. The latter model is strongly connected with Kolmogorov’s probability model (existence of the probability distribution of hidden variables $`𝐏(\lambda )`$ and conditional probabilities $`𝐏(U,\lambda )`$ is postulated).
Let $`U`$ and $`V`$ be physical observables, $`U,V=\pm 1.`$ We start again with the consideration of the frequency covariation $`<U,V>_{x_{UV}}`$ with respect to a collective $`x_{UV}`$ induced by the measurement of the pair $`(U,V).`$ The $`x_{UV}`$ is obtained by measurements for an ensemble $`S_{UV}`$ of physical systems. Our aim is again to represent the experimental covariation $`<U,V>_{x_{UV}}`$ as ensemble covariation $`<U,V>_{S_{UV}}.`$ Then we shall demonstrate that in the general case it is impossible to perform for ensemble covariations Bell’s calculations, (2) – (3).
Let $`S_{UV}=\{d_1,\mathrm{},d_N\},`$ where $`i`$th measurement is performed for the system $`d_i.`$ Define functions $`i\lambda (i)`$ (the same function as above) and $`i\omega ^U(i),i\omega ^V(i),`$ states of apparatus $`_U`$ and $`_V,`$ respectively, at the instances, $`t_i^U`$ and $`t_i^V,`$ of measurements of $`U`$ and $`V`$ for $`i`$th system. We have
$$<U,V>_{x_{UV}}=\frac{1}{N}\underset{i=1}{\overset{N}{}}U(\omega ^U(i),\lambda (i))V(\omega ^V(i),\lambda (i)).$$
Set $`D_{ks}^U=\{i:\lambda (i)=\lambda _k,\omega ^U(i)=\omega _s^U\}`$ and $`D_{ks}^V=\{i:\lambda (i)=\lambda _k,\omega ^V(i)=\omega _s^V\},1kM,1sL_U,1qL_V.`$ Set $`l_{ksq}^{UV}=|D_{ks}^UD_{kq}^V|.`$ It is evident that
$$\underset{k=1}{\overset{M}{}}\underset{s=1}{\overset{L_U}{}}\underset{q=1}{\overset{L_V}{}}l_{ksq}^{UV}=N.$$
Hence
$$<U,V>_{x_{UV}}=\frac{1}{N}\underset{ksq}{}l_{ksq}^{UV}u_{ks}v_{kq},$$
where $`u_{ks}=U(\omega _s^U,\lambda _k),v_{kq}=V(\omega _q^V,\lambda _k).`$ We show that $`<U,V>_{x_{UV}}`$ can be represented as ensemble covariation for an appropriative ensemble of physical systems and states of measurement devices. However, a choice of such an ensemble is rather delicate problem.
First we note that $`<U,V>_{x_{UV}}<U,V>_{\mathrm{\Lambda }\times \mathrm{\Sigma }_A\times \mathrm{\Sigma }_B}`$ (compare with section 6). For the latter covariation, we have
$$<U,V>_{\mathrm{\Lambda }\times \mathrm{\Sigma }_A\times \mathrm{\Sigma }_B}=\frac{1}{ML_AL_B}\underset{k=1}{\overset{M}{}}\underset{s=1}{\overset{L_U}{}}\underset{q=1}{\overset{L_V}{}}u_{ks}v_{kq}$$
and in general $`𝐏_{\mathrm{\Lambda }\times \mathrm{\Sigma }_A\times \mathrm{\Sigma }_B}(\lambda =\lambda _k,\omega ^U=\omega _s^U,\omega ^V=\omega _q^V)=\frac{1}{ML_AL_B}\frac{l_{ksq}}{N}`$ even approximately for $`M,N,L_A,L_B\mathrm{}.`$
It is also evident that $`<U,V>_{x_{UV}}<U,V>_{S_{UV}}.`$ The latter covariation is simply not well defined, because the ‘properties’ $`\omega ^U(i)=\omega _s^U,\omega ^V(i)=\omega _q^V`$ are not objective properties of elements of the ensemble $`S_{UV}.`$ These ‘properties’ are determined by fluctuations of parameters in the apparatus $`_U`$ and $`_V.`$
To find the right ensemble, we have to introduce two new ensembles, namely, ensembles of states of the apparatus $`_U`$ and $`_V`$ (in the process of measurements for the ensemble of physical systems $`S_{UV}):`$
$$S__U=\{\alpha _1^U,\mathrm{},\alpha _N^U\},\alpha _j^U\mathrm{\Sigma }_U,S__V=\{\alpha _1^V,\mathrm{},\alpha _N^V\},,\alpha _j^V\mathrm{\Sigma }_V,$$
where $`\alpha _i^U=\omega ^U(i),\alpha _i^V=\omega ^V(i)`$ are states of $`_U`$ and $`_V`$ at the instances of $`i`$th measurements. We set
$$𝐒_{UV}=\mathrm{diag}(\mathrm{S}_{\mathrm{UV}}\times \mathrm{S}__\mathrm{U}\times \mathrm{S}__\mathrm{V})=\{\mathrm{D}_1,\mathrm{},\mathrm{D}_\mathrm{N}\},\mathrm{D}_\mathrm{j}=(\mathrm{d}_\mathrm{j},\alpha _\mathrm{j}^\mathrm{U},\alpha _\mathrm{j}^\mathrm{V}).$$
Then $`\pi (D_j)=(\lambda (j),\omega ^U(j),\omega ^V(j))`$ is an objective property of elements of the ensemble $`𝐒_{UV}`$ and
$$<U,V>_{x_{UV}}=<U,V>_{𝐒_{UV}}=\frac{1}{N}\underset{i=1}{\overset{N}{}}U(\omega ^U(i),\lambda (i))V(\omega ^V(i),\lambda (i)).$$
We set
$$𝐩_{ksq}^{UV}=𝐏_{𝐒_{UV}}(D_j:\pi (D_j)=(\lambda _k,\omega _s^U,\omega _s^V))$$
$$=\frac{|\{D_j𝐒_{UV}:\pi (D_j)=(\lambda _k,\omega _s^U,\omega _s^V)\}|}{|𝐒_{UV}|}.$$
Hence we obtained that
$$<U,V>_{x_{UV}}=<U,V>_{𝐒_{UV}}=\underset{ksq}{}𝐩_{ksq}^{UV}u_{ks}v_{kq}.$$
Thus in the general case we have
$$\begin{array}{cc}& \mathrm{\Delta }=<A,B>_{x_{AB}}<C,B>_{x_{CB}}=<A,B>_{𝐒_{AB}}<C,B>_{𝐒_{CB}}\hfill \\ & \\ & =_{ksq}𝐩_{ksq}^{AB}a_{sk}b_{kq}_{ksq}𝐩_{ksq}^{CB}c_{ks}b_{kq}\hfill \end{array}$$
and
$$<A,C>_{x_{AC}}=<A,C>_{𝐒_{AC}}=\underset{ksq}{}𝐩_{ksq}^{AC}a_{ks}c_{kq}.$$
We suppose now that probabilities $`𝐩_{ksq}^{UV}`$ do not depend on ensembles:
$$𝐩_{ksq}=𝐩_{ksq}^{AB}=𝐩_{ksq}^{CB}=𝐩_{ksq}^{AC}.$$
(6)
In particular, we suppose that all measurement devices have the same set of states (of parameters):
$$\mathrm{\Sigma }=\mathrm{\Sigma }_A=\mathrm{\Sigma }_B=\mathrm{\Sigma }_C(\text{and}L=L_A=L_B=L_C).$$
(7)
Then we obtain
$$\mathrm{\Delta }=\underset{ksq}{}𝐩_{ksq}(a_{ks}c_{ks})b_{kq}.$$
However, we could not repeat trick (3) of the proof of Bell’s inequality. The equality $`a_{ks}^2=1`$ does not give the possibility to proceed the proof. Of course, we have
$$\begin{array}{cc}& |\mathrm{\Delta }|=|_{ksq}𝐩_{ksq}(a_{ks}a_{ks}^2c_{ks})b_{kq}|_{ksq}𝐩_{ksq}|a_{ks}b_{kq}|(1a_{ks}c_{ks})\hfill \\ & \\ & 1_{ksq}𝐩_{ksq}a_{ks}c_{ks}.\hfill \end{array}$$
But in general $`_{ksq}𝐩_{ksq}a_{ks}c_{ks}`$ is not larger than $`<A,C>_{x_{AC}}=_{ksq}𝐩_{ksq}a_{ks}c_{kq}.`$
Therefore, if we keep to empiricism, then even stability condition (6) (for combined ensembles of physical systems and states of measurement apparatus) does not imply Bell’s inequality. A new source of violation of Bell’s inequality is the inconsistency of random fluctuations for two measurement devices $`_U`$ and $`_V.`$ In general $`\omega ^U(i)\omega ^V(i).`$
Suppose that it could be possible to control states of $`_U`$ and $`_V`$ and choose $`\omega `$ for $`_U`$ and $`_V`$ in the consistence way:
$$\omega =\omega ^U(i)=\omega ^V(i).$$
Then the ensemble $`𝐒_{UV}`$ would contain only triples of the form $`(\lambda _k,\omega _s,\omega _s)`$ and
$$𝐩_{ksq}^{UV}=𝐏_{𝐒_{UV}}(\lambda _k,\omega _s^U,\omega _q^V)=0,sq.$$
(8)
In such a case we obtain covariations:
$$<U,V>_{\mathrm{Ideal}}=\frac{1}{N}\underset{i=1}{\overset{N}{}}U(\omega ^U(i),\lambda (i))V(\omega ^V(i),\lambda (i))=\underset{ks}{}𝐩_{ks}^{UV}u_{ks}v_{ks},$$
where $`𝐩_{ks}^{UV}=𝐩_{kss}^{UV}.`$ If we also suppose the validity of (6), we obtain
$$|\mathrm{\Delta }_{\mathrm{Ideal}}|=|\underset{ks}{}𝐩_{ks}(a_{ks}c_{ks})b_{ks}|$$
$$1\underset{ks}{}𝐩_{ks}a_{ks}c_{ks}=1<A,C>_{\mathrm{Ideal}}.$$
However, ideal covariations have no direct connection to experimental frequency covariations.
Nevertheless, we can formulate the following mathematical theorem:
Theorem 3. Let statistical ensembles (physical systems/measurement apparatus) satisfy conditions (6) and (8). Then Bell’s inequality (1) holds true for covariations with respect to these ensembles.
Therefore, to obtain Bell’s inequality in the empiricists framework, we have to suppose: (1) statistical repeatability of ensemble distribution of hidden variables $`\lambda `$ in ensembles which are used for measurements; (3) statistical repeatability of fluctuations of states $`\omega `$ in ensembles of an equipment; (3) consistency of fluctuations of all measurement devices.
If the reader even deny the possibility of violations of (1) or (2), he must agree that condition (3) seems to be nonphysical: we could never control fluctuations of the huge number of parameters in the equipment.
Instead of precise coincidence (6), it is possible to consider (under the assumption (7)) the statistical coincidence based on the quantity:
$$\delta (𝐒_{AB},𝐒_{CB})=\underset{k=1}{\overset{M}{}}\underset{s=1}{\overset{L}{}}\underset{q=1}{\overset{L}{}}|𝐩_{ksq}^{AB}𝐩_{ksq}^{CB}|.$$
Here $`\delta =\delta _\pi `$ for the property $`\pi (i)=(\lambda (i),\omega ^U(i),\omega ^V(i)).`$ We remark that condition (6) of the precise coincidence can be written as
$$\delta (𝐒_{AB},𝐒_{CB})=0$$
for every two pairs of observable $`(A,B)`$ and $`(C,B).`$ We also introduce a new quantity which is a statistical measure of inconsistency of ensembles $`S__U`$ and $`S__V:`$
$$\sigma (𝐒_{UV})=\underset{sq}{}𝐏_{𝐒_{UV}}(\omega ^U=\omega _s,\omega ^V=\omega _q)=\underset{k}{}\underset{sq}{}𝐩_{ksq}^{UV}.$$
Condition (8) of the precise consistency for states of $`_U`$ and $`_V`$ can be written in the form:
$$\sigma (𝐒_{UV})=0.$$
Theorem 4. Let statistical ensembles (physical systems/measurement apparatus) satisfy conditions:
$$\delta (𝐒_{AB},𝐒_{CB}),\delta (𝐒_{AB},𝐒_{AC})ϵ\text{and}\sigma (𝐒_{AB}),\sigma (𝐒_{CB}),\sigma (𝐒_{AC})ϵ^{}.$$
Then inequality
$$|<A,B>_{𝐒_{AB}}<C,B>_{𝐒_{CB}}|(1+2ϵ+3ϵ^{})<A,C>_{𝐒_{AC}}$$
holds true.
Proof. We have
$$|\mathrm{\Delta }|ϵ+|\underset{ksq}{}𝐩_{ksq}^{AB}(a_{ks}c_{ks})b_{kq}|$$
$$ϵ+2ϵ^{}+\underset{ks}{}𝐩_{ks}^{AB}|(a_{ks}c_{ks})b_{ks}|ϵ+2ϵ^{}+\underset{ks}{}𝐩_{ks}^{AB}(1a_{ks}c_{ks})$$
$$ϵ+4ϵ^{}+\underset{ksq}{}𝐩_{ksq}^{AB}(1a_{ks}c_{kq})(1+2ϵ+4ϵ^{})\underset{ksq}{}𝐩_{ksq}^{AC}a_{ks}c_{kq}.$$
We remark again that experiments to test Bell’s inequality can be interpreted as just experiments to find an estimate for a constant $`C=2ϵ+4ϵ^{}.`$ From this point of view the only result of these experiments is that $`C`$ is essentially larger that zero. However, such a results could be expected: it would be rather strange if measures of statistical deviations $`\delta `$ and $`\sigma `$ would be equal to zero despite of fluctuations of parameters of measuring devices.
6. Probability distributions in stochastic hidden variables models.
Typically stochastic hidden variables models are defined as models with probabilities $`(ϵ=\pm 1)`$
$$𝐏(U=ϵ)=_\mathrm{\Lambda }𝐏(U=ϵ/\lambda )𝑑\rho (\lambda ),$$
(9)
where $`\rho (\lambda )`$ is the probability distribution of hidden variables and $`𝐏(U=ϵ/\lambda )`$ is the conditional probability to measure the value $`U=ϵ`$ for the quantum system having the hidden state $`\lambda ,`$ see, for example, .
Then (see Ref. ) the joint probability distribution can be defined (at least mathematically) as
$$𝐏(U_1=ϵ_1,U_2=ϵ_2,U_3=ϵ_3)=_\mathrm{\Lambda }𝐏(U_1=ϵ_1/\lambda )𝐏(U_2=ϵ_2/\lambda )𝐏(U_3=ϵ_3/\lambda )𝑑\rho (\lambda ).$$
(10)
In fact, to derive Bell’s inequality in the Kolmogorov framework, it is sufficient to use the existence (on the mathematical level) of the joint probability distribution (10). However, considerations in the framework of the ensemble probability theory demonstrated that ‘probabilities’ (9) has no physical meaning. These are probabilities with respect to the ensemble $`\mathrm{\Lambda }\times \mathrm{\Sigma }_U.`$ However, physical probabilities are probabilities with respect to the ensemble $`𝐒_U=\mathrm{diag}(\mathrm{S}_\mathrm{U}\times \mathrm{S}__\mathrm{U}),`$ where $`S_U=\{d_1,\mathrm{},d_N\}`$ is the ensemble of quantum system used in the measurement. We note that physical arguments against existing of representation (9) were presented by W. De Muynck, W. De Baere, H. Marten in Ref. . I think that results of this paper can strongly improve their considerations. We hope that our numerical description of nonexistence of Kolmogorov probabilities could essentially clarify the problem.
6. Other probabilistic models which do not contradict to local realism.
L. Accardi in Ref. used non-Kolmogorov model without Bayes’ formula to eliminate Bell’s inequality from considerations related to spin’s model. Recently he developed a new model which gives an explanation of violations of Bell’s inequality, see Ref. . In fact, to get “physical Bell’s inequality” we have to consider in Theorem 1 indexed random variables $`U^1`$ and $`U^2`$ corresponding to correlated particles, 1 and 2. “Physical Bell’s inequality” can be obtained only on the basis of the implicit anticorrelation: $`U^1=U^2.`$ Accardi discussed the role of this condition in Bell’s arguments.
I. Pitowsky in Ref. discussed the possibility that some nonmeasurable sets can be physical events, i.e, some physical observables may be nonmeasurable. There is no Bell’s inequality in this approach. Thus there is no problem with violations of Bell’s inequality. This model is consistent with known polarization phenomena and the existence of macroscopic magnetism. He also proposed a thought experiment which indicates a deviation from the predictions of quantum mechanics. We note that already A. N. Kolmogorov discussed ‘generalized probabilities’ on the algebra of all subsets of $`\mathrm{\Omega }.`$ Mathematicians, in particular applied mathematicians, where reluctant to take nonmeasurable sets seriously. As a result there was no mathematical theory that relates nonmeasurable distributions with relative frequencies. Such an extension of probability theory was created by I. Pitowsky and then strongly mathematically improved by S.P. Gudder in Ref. . He introduced the concept of a probability manifold $`M.`$ The global properties of $`M`$ inherited from its local structure were then considered. It was shown that a deterministic spin model due to Pitowski falls within this general framework. Finally, Gudder constructed a phase-space model for nonrelativistic quantum mechanics. These two models give the same global description as conventional quantum mechanics. However, they also give a local descriptions which is not possible in conventional quantum mechanics.
ACKNOWLEDGMENTS
This paper was completed during the visit to Moscow State University on the basis of the grant of the Royal Academy of Science (Sweden) for the collaboration with the former Soviet Union.
References
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warning/0003/math0003225.html | ar5iv | text | # Fusion rings for degenerate minimal models
## 1 Introduction
The Virasoro algebra and its minimal models are a good source of interesting vertex operator algebras. In \[W\] the rationality of the Virasoro vertex operator algebras $`L(c_{p.q},0)`$ was proved, where $`c_{p,q}=1\frac{6(pq)^2}{pq}`$ and $`(p,q)=1`$, $`p,q2`$. This result is used for the construction of the corresponding vertex tensor categories (cf. \[H1\]). A similar result is obtained for $`N=1`$ case in \[A\] and \[HM\].
In this paper we study a non–rational vertex operator algebra $`L(1,0)`$ ($`p=q`$ case) and the corresponding fusion ring for degenerate minimal models, i.e., the case “$`p=q`$”, with central charge $`c=1`$. We also consider a $`N=1`$ vertex operator superalgebra version based on $`L(\frac{3}{2},0)`$ (see below). These cases are substantially different for many reasons (let us focus on the case $`L(1,0)`$ since the same problem persists for $`L(\frac{3}{2},0)`$. The vertex operator algebra $`L(1,0)`$ is not rational (cf. \[FZ\]) but it has a distinguished family of irreducible modules (those that are not irreducible Verma modules) $`_1`$, which consists of classes of irreducible modules isomorphic to $`L(1,\frac{m^2}{4})`$ for some $`m`$. These modules have a quite simple embedding structure (\[KR\], \[FF2\]).
We show that the fusion ring for the family $`_1`$ is isomorphic to the representation ring $`ep(𝔰l(2,))`$, i.e., we “formally” have
$$L(1,\frac{n^2}{4})\times L(1,\frac{m^2}{4})=$$
$$L(1,\frac{(n+m)^2}{4})+L(1,\frac{(n+m2)^2}{4})\mathrm{}+L(1,\frac{(nm)^2}{4}),$$
where $`m,n`$ and $`nm`$.
This result seems to be known–in some form– for a while by physicists (also in \[FKRW\] is stated as a part of more general conjecture concerning fusion rings for $`W(𝔤l_N)`$ algebras–see also \[FM\]). The author of the current paper could not trace any proof in the language of vertex operator algebras. Some computations are done in \[DG\] but not complete. But instead of trying to patch missing proofs, there are two more important reasons for seeking such a proof.
* So far, not many computations of the fusion coefficients has been known for non–rational vertex operator algebras (here non–rational means non–rational in any reasonable category). In particular we offer a proof that uses universal construction (induced modules), therefore it is very general.
* As noticed by H. Li in \[L1\] and \[L2\], Frenkel–Zhu’s formula \[FZ\] does not hold for non–rational vertex operator algebras. The right formula was provided in \[L2\] but it is a non–trivial matter to use it for computational purposes in non–rational setting.
We believe that our method can be used for more complicated models–like degenerate models associated to $`𝒲`$–algebras.
We have to stress that the fusion coefficients are simply derived from the space of intertwining operators among irreducible modules. In other words it is not true that the only modules which “fuse” with $`L(1,\frac{n^2}{4})`$ and $`L(1,\frac{m^2}{4})`$ are completely reducible. This fact makes impossible to implement $`P(z)`$–tensor product construction from \[HL1\]\[HL2\]. The resolution might be to construct (a new) tensor product which takes only irreducible modules into account, but this approach will assume a good knowledge of matrix coefficients for product of intertwining operators. A different approach would be working in the larger family $`\overline{}_1`$, which consists of all quotients of Verma modules $`M(1,\frac{m^2}{4})`$. The possible constructions will be discussed elsewhere.
We also provide a different proof of the fusion formulas by constructing all intertwining operators from the lattice vertex operator algebra $`V_L`$ and its irreducible module $`V_{L+1/2}`$ (cf. \[DG\]).
A super $`N=1`$ versions of the above result stems from the $`N=1`$ Neveu–Schwarz Lie superalgebra at the level $`\frac{3}{2}`$. Again, there are essentially two approaches: one which uses the lattice construction (extended with a suitable fermionic Fock space) and the other which uses the singular vectors and projection formulas. For the future purposes we use the latter approach. We consider a set of equivalence classes of irreducible modules for the $`N=1`$ superconformal algebra (see Section 3.) with representatives $`L(\frac{3}{2},\frac{q^2}{2})`$ where $`q.`$ We proved (see Theorem 10.1 and Corollary 10.1) that the corresponding fusion ring is isomorphic to the representation ring for $`𝔬𝔰𝔭(1|2)`$, i.e., we formally have:
$$L(\frac{3}{2},\frac{r^2}{2})\times L(\frac{3}{2},\frac{q^2}{2})=$$
$$L(\frac{3}{2},\frac{(r+q)^2}{2})+L(\frac{3}{2},\frac{(r+q1)^2}{2})\mathrm{}+L(\frac{3}{2},\frac{(rq)^2}{2}),$$
for every $`r,q𝐍`$, $`rq`$, where $`\times `$ stands for the fusion product (see the last Chapter).
In particular, as in the Virasoro algebra case, these fusion coefficients are $`0`$ or $`1`$. However in \[HM\] we showed that for $`N=1`$ case has some interesting features; for some vertex operator algebras $`L(c,0)`$, fusion coefficients might be $`2`$. In Proposition 11.1 we construct a non–trivial example with $`c=\frac{15}{2}3\sqrt{5}`$.
At the very end, we construct an example of a logarithmic intertwining operator (for the definition see \[M\]) for the $`N=1`$ vertex operator superalgebra $`L(\frac{27}{2},0)`$.
n.b. These results can be extended for a more general class of vertex operator algebras $`L(c,0)`$ where $`cc_{p,q}`$; because of simplicity we treat only the case $`c=1`$ and $`c=\frac{3}{2}`$.
Acknowledgment: The author thank Prof. Haisheng Li and Prof. Yi-Zhi Huang for useful comments. Thanks go to the referee for his/her valuable remarks. This paper is a union of slightly modified preprints available at the arXive.
## 2 Representations of the Virasoro algebra at the level $`c=1`$
The representation theory for the Virasoro algebra has been studied intensively in the last two decades (\[KR\], \[FF1\]\[FF3\]). Kac’s determinant formula is the most important tool in the highest (or lowest) weight theory. From the determinant formula it follows that the lowest weight Verma module with the central charge $`c(t)=136t6t^1`$ and the weight
$$h_{p,q}(t)=\frac{1p^2}{4}t^1\frac{1pq}{2}+\frac{1q^2}{4}t,$$
has a singular vector of the weight $`h_{p.q}(t)+pq`$, $`t`$. We are interested in the case $`t=1`$, i.e $`c=1`$. It is easy to see that $`M(1,h)`$ is irreducible if and only if $`h\frac{m^2}{4}`$ for some $`m`$. In the case $`h=\frac{m^2}{4}`$ we have the following description:
###### Proposition 2.1
The Verma module $`M(1,\frac{m^2}{4})`$ has a unique singular vector of weight $`\frac{m^2}{4}+(m+1)`$. This vector generates the maximal submodule. In other words we have the following exact sequence
$$0M(1,\frac{(m+2)^2}{4})M(1,\frac{m^2}{4})L(1,\frac{m^2}{4})0.$$
(1)
Even though they do not exist in general, in the case $`h_{1,q}(t)`$, if $`p=1`$ there are explicit formulas at each level $`c(t)`$ (in particular $`t=1`$). When $`c=1`$ Benoit and S. Aubin’s formula \[BSA1\] implies that
$$P_{\mathrm{sing}}v_{1,q}=\underset{\stackrel{I=\{i_1,\mathrm{},i_n\}}{|I|=q}}{}c_q(i_1,\mathrm{},i_n)L(i_1)\mathrm{}L(i_n)v_{1.q}$$
(2)
is a singular vector for $`M(1,h_{1,q}(1))`$, where
$$c_r(i_1,\mathrm{},i_n)=\underset{\stackrel{1k<r}{k_{j=1}^si_j}}{}k(rk).$$
###### Remark 2.1
Note that every singular vector (2) has form $`L(1)^{m+1}+\mathrm{}`$, where dots represent lower degree terms (with respect to the universal enveloping algebra grading).
## 3 Vertex operator algebra $`L(1,0)`$
### 3.1 Zhu’s algebra and intertwining operators
We will use the definition of vertex operator algebra and modules as stated in \[FHL\] or \[FLM\]. Let $`L(1,0)=M(1,0)/L(1)\mathrm{𝟏}`$ be a simple vertex operator algebra associated to irreducible representation of the Virasoro algebra (cf. \[FZ\], \[W\]).
It is known that to every vertex operator algebra $`V`$, we can associate Zhu’s associative algebra $`A(V)`$ (cf. \[FZ\] and \[Z\]). In the special case $`V=L(1,0)`$, we know (see \[FZ\], \[W\]) that $`A(V)[y]`$, where $`y=[L(2)L(1)].`$ We have chosen the multiplication in $`A(V)`$ as in \[W\] (which is slightly different then the one in \[FZ\]),
$$ab=\mathrm{Res}_xY(a,x)\frac{(1x)^{\mathrm{deg}(a)}}{x}b,$$
where $`a,bA(V)`$.
By using standard techniques (see \[FZ\], \[W\]) we have the following.
###### Proposition 3.1
Every irreducible module for the vertex operator algebra $`L(1,0)`$ is isomorphic to $`L(1,h)`$ , for some $`h`$.
Proof: According to \[Z\], there is a one–to–one equivalence between equivalence classes of $``$–gradable irreducible $`L(1,0)`$–modules and irreducible $`[y]`$–modules. Every irreducible $`L(1,0)`$–module is a $`\mathrm{Vir}`$–module. Any such module is $``$–gradable and isomorphic to $`L(1,h)`$ for some $`h`$. On the other hand every finite dimensional irreducible $`[y]`$–module is one dimensional so the proof follows.
Since the notion of intertwining operator is more subtle we include here the original definition \[FHL\].
###### Definition 3.1
Let $`W_1,W_2`$ and $`W_3`$ be a triple of modules for vertex operator algebra $`V`$. A mapping
$$𝒴W_1W_2W_3\{x\},$$
is called an intertwining operator of type $`\left(\genfrac{}{}{0pt}{}{W_3}{W_1W_2}\right)`$, if it satisfies the following properties
1. The truncation property: For any $`w_iW_i`$, $`i=1,2`$,
$$(w_1)_nw_2=0,$$
for $`n`$ large enough.
2. The $`L(1)`$-derivative property: For any $`vV`$,
$$𝒴(L(1)w_1,x)=\frac{d}{dx}𝒴(w_1,x),$$
3. The Jacobi identity: In $`\mathrm{Hom}(W_1W_2,W_3)\{x_0,x_1,x_2\}`$, we have
$`x_0^1\delta \left({\displaystyle \frac{x_1x_2}{x_0}}\right)Y(u,x_1)𝒴(w_1,x_2)`$ (3)
$`x_0^1\delta \left({\displaystyle \frac{x_2x_1}{x_0}}\right)𝒴(w_1,x_2)Y(u,x_1)`$
$`=x_2^1\delta \left({\displaystyle \frac{x_1x_0}{x_2}}\right)𝒴(Y(u,x_0)w_1,x_2)`$
for $`uV`$ and $`w_1W_1`$.
We denote the space of all intertwining operators of the type $`\left(\genfrac{}{}{0pt}{}{W_3}{W_1W_2}\right)`$ by $`I\left(\genfrac{}{}{0pt}{}{W_3}{W_1W_2}\right)`$. The dimension of the space of intertwining operators (also known as “fusion rule” )of the type $`\left(\genfrac{}{}{0pt}{}{W_3}{W_1W_2}\right)`$ we denote by $`𝒩_{W_1,W_2}^{W_3}`$.
Our goal is to find the fusion rules for the degenerate minimal models, i.e.,
$$\mathrm{dim}I\left(\genfrac{}{}{0pt}{}{L(1,\frac{r^2}{4})}{L(1,\frac{p^2}{4})L(1,\frac{q^2}{4})}\right).$$
Since our modules are irreducible we want to introduce Frenkel-Zhu’s formula which gives us (roughly) a prescription for calculating fusion rules. It is not hard to see, by using the Jacobi identity, that the space $`I\left(\genfrac{}{}{0pt}{}{L(1,\frac{r^2}{4})}{L(1,\frac{p^2}{4})L(1,\frac{q^2}{4})}\right)`$ is at most one dimensional.
Now for every module $`M`$, we associate an $`A(V)`$–bimodule $`A(M):=M/O(M)`$ (cf. \[FZ\]), where $`O(M)`$ is spanned by the elements of the form
$$\mathrm{Res}_xY(u,x)\frac{(1x)^{deg(a)}}{x^2}v,$$
$`uV`$, $`vM`$. In the case $`M=M(c,h)`$,
$$O(M(c,h))=\{(L(n3)2L(n2)+L(1))v,n0,vM(c,h)\}.$$
(4)
If we let
$$y=[L(2)L(1)],x=[L(2)2L(1)+L(0)],$$
then from the formulas
$$[L(n)v]=[(nyx+\mathrm{wt}(v))v],$$
and
$$[x,y]w=0\mathrm{mod}O(M(c,h)),$$
($`[x,y]=xyyx`$) it follows that
$$A(M(c,h))[x,y],$$
as a $`[y]`$–bimodule (cf. \[L2\]) , where the lowest weight vector is identified with $`1[x,y]`$ and the actions of are
$$yp(x,y)=xp(x,y),p(x,y)y=yp(x,y),$$
for every $`p(x,y)[x,y]`$.
The Frenkel-Zhu’s formula (\[FZ\]) states that it is possible to calculate the dimension of the space $`\left(\genfrac{}{}{0pt}{}{M_3}{M_1M_2}\right)`$ by knowing $`A(V)`$, $`A(M_1)`$, $`M_2(0)`$ and $`M_3(0)`$. Instead of giving the original statement from \[FZ\], we state the following refinement obtained in \[L1\]-\[L2\]:
###### Theorem 3.1
Let $`M_1`$, $`M_2`$ and $`M_3`$ be lowest weight $`V`$–modules. Suppose that $`M_2`$ and $`M_3^{}`$ are generalized Verma $`V`$–modules (see Section 3.2). Then we have
$$𝒩_{M_1M_2}^{M_3}=\mathrm{dim}\mathrm{Hom}_{\mathrm{A}(\mathrm{V})}(\mathrm{A}(\mathrm{M}_1)_{\mathrm{A}(\mathrm{V})}\mathrm{M}_2(0),\mathrm{M}_3(0)),$$
where $`M_i(0)`$, $`i=1,2,3`$, is the “top” level of $`M_i`$, respectively, equipped with the $`A(V)`$-module structures.
This theorem is not so useful as it stands. On the other hand its proof is important. Hence it will be necessary to understand a little bit deeper assumptions on $`M_2`$ and $`M_3`$ in our situation. For warm up let us start with the “easy–half” of the Frenkel-Zhu’s formula which says:
###### Lemma 3.1
Let $`M_3`$ be an irreducible lowest weight $`V`$–module. Then
$$𝒩_{M_1M_2}^{M_3}\mathrm{dim}\mathrm{Hom}_{\mathrm{A}(\mathrm{V})}(\mathrm{A}(\mathrm{M}_1)_{\mathrm{A}(\mathrm{V})}\mathrm{M}_2(0),\mathrm{M}_3(0)).$$
Define an infinite dimensional Lie algebra $``$ spanned by
$$L(n2)2L(n1)+L(n),$$
for $`n1`$. In the case of minimal models–which is the most interesting case–the homology groups $`H_q(,L(c,h))`$ where calculated in \[FF2\]. For the Verma modules the $`0`$-th homology, $`H_0(,M(1,h))`$ with the coefficients in the Verma modules is isomorphic to $`[x,y]`$ as an $`A(L(1,0))`$–bimodule (cf. \[W\]).
The following result is an application of a more general theory \[FF1\].
###### Theorem 3.2
We have
1. $$\mathrm{H}_0(,L(1,\frac{m^2}{4})),$$
is infinite–dimensional.
2. $`\mathrm{H}_0(,L(1,\frac{m^2}{4}))`$ is finitely generated as a (left) $`A(L(1,0))`$–module.
3. $$\mathrm{Ext}_{Vir,𝒪}^1(L(1,\frac{m^2}{4}),L(1,\frac{n^2}{4}))=\{\begin{array}{c}\mathrm{if}|mn|=2\\ 0\mathrm{otherwise}\end{array}$$
where $`\mathrm{Ext}_{Vir,𝒪}^1`$ stands for the relative $`\mathrm{Ext}`$ with respect to the one–dimensional abelian subalgebra generated by $`L(0)`$.
Proof: a) Since the maximal submodule of $`M(1,\frac{m^2}{4})`$ is generated by one vector, in the projection (or homology) $`A(L(1,\frac{m^2}{4}))`$ is isomorphic to $`\frac{[x,y]}{I}`$, where $`I`$ is a cyclic submodule (with respect to the left and right actions) generated by some polynomial $`p(x,y)`$ which is a projection of $`v_{1,m}`$ in $`[x,y]`$ . It is clear that this space is infinite dimensional.
b) Note first that $`[L(1)v]=(yx+deg(v))[v]`$. By using Remark 2.1 it follows that
$$[v_{sing}]=p(x,y)=\underset{i=1}{\overset{m+1}{}}(xy+i)+q(x,y).$$
where $`\text{deg}(q)<(m+1)`$. Thus, the pure monomials in $`p(x,y)`$ with the highest powers are $`x^{m+1}`$ and $`y^{m+1}`$. Since, $`I`$ is spanned by $`p(x,y)[x]`$, here we consider only the left action, it follows that $`\frac{[x,y]}{I}`$ is finitely generated. The similar argument holds for the right action.
c) The idea is the same as in \[FF1\]. The result is however different. It is known that
$$\mathrm{Ext}_{Vir,𝒪}^{}(M,N)H^{}(Vir,𝒪,Hom(M,N)).$$
Therefore
$$H^{}(Vir,𝒪,Hom(M,N))Tor_{}^{Vir,𝒪}(N^{},M),$$
where $`N^{}`$ is the dual module. Hence we can compute our cohomology by using the tensor product of complexes
$$M(1,\frac{(m+2)^2}{4})M(1,\frac{m^2}{4}),$$
$$M(1,\frac{(n+2)^2}{4})^{opp}M(1,\frac{n^2}{4})^{opp},$$
where $`M(c,h)^{opp}`$ is the opposite Verma module (cf. \[FF1\]-\[FF2\]). The corresponding spectral sequence $`E_2^{p,q}`$ collapses at the second term Therefore
$$Tor_1^{Vir,𝒪}(L(1,\frac{n^2}{4})^{},L(1,\frac{m^2}{4})E_2^{1,0}\mathrm{or}0,$$
where non–trivial homology occurs only if the Verma module $`M(1,\frac{m^2}{4})`$ embeds inside $`M(1,\frac{n^2}{4})`$ as the maximal submodule or vice–versa. This happens if and only if $`|nm|=2`$. Therefore we have the proof <sup>1</sup><sup>1</sup>1It is crucial to notice that our cohomology is relative one, otherwise our extension are not controllable inside category $`𝒪`$. Such (non–relative) extensions are studied in \[M\]. The corresponding short–exact sequences are clearly,
$$0L(1,\frac{(m+2)^2}{4})M(1,\frac{m^2}{4})/M(1,\frac{(m+4)^2}{4})L(1,\frac{m^2}{4})0,$$
(5)
and the one obtained from 5 by applying (exact) functor $`()^{}`$ taking modules to the corresponding contragradient modules.
For every $`m,n`$ (we exclude the case $`mn=0`$), fix a multiset $`J_{m,n}=\{m+n,m+n2,\mathrm{},mn\}`$. Let $`_{\lambda ,\mu }`$ be a “density” module for the Virasoro algebra. $`_{\lambda ,\mu }`$ is spanned by $`w_r`$, $`r`$ and the action is given by
$$L_n.w_r=(\mu +r+\lambda (m+1))w_{rn}.$$
In \[FF1\] the projection formula for the singular vectors (considered as an element of the enveloping algebra) on $`_{\lambda ,\mu }`$ (more precisely $`w_0`$) was found. We want to relate the projection of the singular vectors on $`_{\lambda ,\mu }`$ with the projection inside $`A(M(1,\frac{m^2}{4}))_{C[y]}L(1,\frac{n^2}{4})`$). It is easy to see that
$`[L(j_1)\mathrm{}L(j_k)v_{m^2/4}]=`$
$`{\displaystyle \underset{r=1}{\overset{k}{}}}(j_r{\displaystyle \frac{n^2}{4}}y+\beta (r,k)).[v_{m^2/4}]=`$
$`{\displaystyle \underset{r=1}{\overset{k}{}}}(j_r{\displaystyle \frac{n^2}{4}}x+\beta (r,k))v_{m^2/4}`$ (6)
where $`v_{m^2/4}`$ is the lowest weight vector and
$$\beta (r,k)=j_{r+1}+\mathrm{}+j_k+\frac{m^2}{4}.$$
But the last factor in (3.1) is the same as the $`P(j_1,..,j_k)`$ where
$$L(j_1)\mathrm{}L(j_k).w_0=P(j_1,\mathrm{},j_k)w_{j_1+\mathrm{}+j_k},$$
and the projection is in $`_{\lambda ,\mu }`$ for $`\lambda =\frac{n^2}{4}`$ and $`\mu =\frac{n^2}{4}+\frac{m^2}{4}x`$.
In the remarkable paper \[FF2\], projection formulas for all singular vectors on the density modules were found. In the slightly different notation, for the singular vectors we consider, these formulas appeared in \[KA\]. The result is
$$v_{1,m+1}.w_0=\underset{iJ_{m,n}}{}(x\frac{i^2}{4})w_{m+1},$$
(7)
up to a multiplicative constant.
Now, by using (7) fact and the discussion above (cf. \[W\]) we obtain
###### Lemma 3.2
As a $`A(L(1,0)`$–module $`A(L(1,\frac{m^2}{4}))_{A(L(1,0))}L(1,\frac{n^2}{4})(0)`$ is isomorphic to $`\frac{[x]}{<_{iJ_{m,n}}(xi^2/4)>}.`$
If $`nm`$ notice that as an $`A(L(1,0))`$– module
$$A(L(1,\frac{m^2}{4}))_{A(L(1,0))}L(1,\frac{n^2}{4})(0)\underset{iJ_{m,n}}{}v_i,$$
(8)
where $`v_i`$ is an irreducible $`A(L(1,0))`$–module such that $`y.v_i=i^2/4v_i`$. But if $`m<n`$, then we have two-dimensional submodule in the above decomposition (and this module is not completely reducible). Thus, (8) is not symmetric if we switch $`m`$ and $`n`$.
The similar failure was already noticed in \[L1\]. Anyhow, by using Lemma 3.2 and Lemma 3.1 we obtain
###### Proposition 3.2
Let $`L(1,\frac{m^2}{4})`$, $`L(1,\frac{n^2}{4})`$ and $`M`$ an irreducible $`L(1,0)`$–modules. Then we have the following upper bounds
$`\mathrm{dim}I\left({\displaystyle \genfrac{}{}{0pt}{}{M}{L(1,\frac{m^2}{4})L(1,\frac{n^2}{4})}}\right)\{\begin{array}{c}1\mathrm{if}ML(1,\frac{r^2}{4})\mathrm{for}rJ_{m,n},\\ 0\mathrm{otherwise}\end{array}`$ (11)
where $`J_{m,n}=\{m+n,\mathrm{},mn\}`$.
Now, we shall show that the equality holds in the equation (11). We will provide two different proofs. One which uses the properties of Verma modules and the other which uses free field realization of the modules $`L(1,\frac{m^2}{4})`$.
### 3.2 Lie algebra $`g(V)`$
Let $`V`$ be a vetex operator algebra. Let $`\widehat{V}=V[t,t^1]`$, $`d=L(1)1+1\frac{d}{dt}`$ and $`g(V)=V/dV`$. It has been noticed by several authors that the space $`g(V)`$ has a Lie algebra structure if we let
$$[a(m),b(n)]=\underset{i=0}{\overset{\mathrm{}}{}}\left(\genfrac{}{}{0pt}{}{m}{i}\right)(a_ib)(m+ni).$$
If we define the grading with $`\mathrm{deg}a(m)=nm1`$, where $`aV_{(n)}`$, then we have the corresponding triangular decomposition $`g(V)=g(V)_{}g(V)_0g(V)_+`$. Let $`U`$ be any $`g(V)_0`$–module. We let (as in \[L2\])
$$F(U)=\mathrm{Ind}_{U(g(V)_+g(V)_0)}^{U(g)}U,$$
such that $`g(V)_+`$ acts as zero. We define also the quotient $`\overline{F}(U)=F(U)/J(U)`$ (the so–called generalized Verma module \[L2\]), where $`J(U)`$ is the intersection of all kernels of all $`g(V)`$–homomorphisms from $`F(U)`$ to weak modules. Now, the assumption in Theorem 3.1 on $`M_2`$ and $`M_3^{}`$ means that $`M_2\overline{F}(M_2(0))`$ and $`M_3=\overline{F}(M_3^{}(0))^{}`$.
In \[L1\]\[L2\] it was shown that every $`A(V)`$ homomorphism from $`A(W_1)_{A(V)}W_2(0)`$ to $`W_3(0)`$ does not necessary lead to an intertwining operator of the type $`\left(\genfrac{}{}{0pt}{}{W_3}{W_1W_2}\right)`$ but rather to $`\left(\genfrac{}{}{0pt}{}{F(W_3(0)^{})^{}}{W_1F(W_2(0))}\right)`$ (actually $`F(W_2(0))`$ might be replaced by $`\overline{F}(W_2(0))`$).
In the case when $`V`$ is rational $`\overline{F}(W_2(0))W_2`$ and $`\overline{F}(W_3(0)^{})^{}W_3`$ (\[L2\]). But if the vertex operator algebra $`V`$ is not rational, the main difficulty is that the generalized Verma module $`\overline{F}(W_2(0))`$ may not be isomorphic to $`W_2`$ (let alone $`F(W_2(0))`$ !) (cf. \[L2\]). Also, the spaces $`F(U)`$ and $`\overline{F}(U)`$ are extremely difficult to analyze explicitly. Still, because we are dealing with a particular example, Virasoro vertex operator algebra, we can make use of singular vectors and Verma modules to simplify the whole construction.
Let $`V=L(1,0)`$. Pick $`\omega =L(2)\mathrm{𝟏}L(1,0)`$. Then, inside $`g(L(1,0))`$, we have
$$[\omega (m+1),\omega (n+1)]=(mn)\omega (m+n+1)+\delta _{m+n,0}\frac{m^3m}{12},$$
i.e., these operators close the Virasoro algebra. From the construction of $`F(U)`$ it is clear that $`U(Vir_{})UU(g(V)_{})UF(U)`$. In particular $`M(1,h)F(M(1,h)(0))`$.
### 3.3 The fusion rules computations
Assume first that
$$mn.$$
(12)
First we replace the “big” space $`F(M(1,h))`$ with the smaller Verma module for the Virasoro algebra (we have seen already that the latter is a subspace inside $`F(M(1,h))`$).
Now, let us pick a non–trivial $`A(L(1,0))`$ homomorphism from $`A(L(1,\frac{m^2}{4}))_{A(L(1,0))}L(1,\frac{n^2}{4})(0)`$ to $`L(1,\frac{r^2}{4})(0)`$. Also let $`T=L(1,\frac{m^2}{4})[t,t^1]M(1,\frac{n^2}{4}`$ be a $`g(L(1,0))`$–module as in \[L2\]. Then the construction in \[L2\] implies that there is a bilinear pairing between $`T`$ and $`M(1,\frac{r^2}{4})F(M(1,\frac{r^2}{4})(0)^{})`$. This implies (again by applying Li’s construction in the proof of Theorem 2.11 in \[L2\]) that the corresponding intertwining operator lands in $`M(1,\frac{r^2}{4})^{}`$, i.e., it is of the type $`\left(\genfrac{}{}{0pt}{}{M(1,\frac{r^2}{4})^{}}{L(1,\frac{m^2}{4})M(1,\frac{n^2}{4})}\right)`$. Here $`M(1,\frac{r^2}{4})^{}`$ is the contragradient Verma module (cf. \[FF2\]). The contragradient module $`M(1,\frac{r^2}{4})^{}`$ is not of the lowest weight type (because $`M(1,\frac{r^2}{4})`$ is reducible). In particular, if $`v^{}`$ is the lowest weight vector
$$U(Vir)v^{}L(1,\frac{r^2}{4}),$$
i.e. we can “paste” the whole irreducible module by acting on the lowest weight subspace, but not the whole module $`M(1,\frac{m^2}{4})^{}`$. Now, the question is
How to descend from $`M(1,\frac{m^2}{4})^{}`$ to $`L(1,\frac{m^2}{4})`$ ?
Here is the proof. We have either $`nr`$ or $`r<n`$. For each of these two cases we consider
$$I\left(\genfrac{}{}{0pt}{}{M(1,\frac{r^2}{4})^{}}{L(1,\frac{m^2}{4})M(1,\frac{n^2}{4})}\right),$$
(13)
or
$$I\left(\genfrac{}{}{0pt}{}{M(1,\frac{n^2}{4})^{}}{L(1,\frac{m^2}{4})M(1,\frac{r^2}{4})}\right),$$
(14)
respectively. Notice that these two spaces are isomorphic because of $`I\left(\genfrac{}{}{0pt}{}{M_3}{M_1M_2}\right)I\left(\genfrac{}{}{0pt}{}{M_2^{}}{M_1M_3^{}}\right)`$. Suppose that $`nr`$.
Now the aim is to construct intertwining operator of the type $`\left(\genfrac{}{}{0pt}{}{M(1,\frac{r^2}{4})^{}}{L(1,\frac{m^2}{4})L(1,\frac{n^2}{4})}\right).`$ Therefore if we can check
$$w_3^{},𝒴(w_1,x)w=0,$$
(15)
for every $`wM(1,\frac{(m+2)^2}{4})M(1,\frac{m^2}{4})`$, $`w_3^{}M(1,\frac{r^2}{4})^{\prime \prime }=M(1,\frac{r^2}{4})`$ and $`w_1L(1,\frac{n^2}{4})`$, then by defining $`\overline{𝒴}(w_1,x)[w_2]:=\overline{𝒴}(w_1,x)w_2`$ where $`[w_2]M(1,\frac{m^2}{4})/M(1,\frac{(m+2)^2}{4})`$, we obtain a (well–defined) non–trivial intertwining operator of the type $`\left(\genfrac{}{}{0pt}{}{M(1,\frac{r^2}{4})^{}}{L(1,\frac{m^2}{4})L(1,\frac{n^2}{4})}\right)`$.
Let us check that (15) holds. First of all, because of the Jacobi identity and the fact that $`M(1,\frac{r^2}{4})`$ is lowest weight module, it is enough to show that
$$w_3^{},𝒴(w_1,x)v_{sing}=0$$
(16)
where $`w_3^{}M(1,\frac{r^2}{4})^{\prime \prime }(0)=M(1,\frac{r^2}{4})(0)`$ is the lowest weight vector and $`v_{sing}`$ is the singular vector that generates the maximal submodule of $`M(1,\frac{m^2}{4})`$.
$`w_3^{},𝒴(w_1,x)L(j_1)\mathrm{}L(j_k)w=`$
$`{\displaystyle \underset{i=1}{\overset{k}{}}}(x^{j_i+1}_x+(1j_i)x^{j_i}{\displaystyle \frac{n^2}{4}})w_3^{},𝒴(w_1,x)w=`$
$`{\displaystyle \underset{i=1}{\overset{k}{}}}(x^{j_i+1}_x+(1j_i)x^{j_i}{\displaystyle \frac{m^2}{4}})Cx^{\frac{r^2}{4}\frac{m^2}{4}\frac{n^2}{4}}=`$
$`(1)^{_ij_i}{\displaystyle \underset{i=1}{\overset{k}{}}}({\displaystyle \frac{r^2}{4}}{\displaystyle \frac{m^2}{4}}{\displaystyle \frac{n^2}{4}}{\displaystyle \underset{s=i+1}{\overset{k}{}}}j_s+(1j_i){\displaystyle \frac{m^2}{4}})Cx^{\frac{r^2}{4}\frac{m^2}{4}\frac{n^2}{4}}=`$
$`C{\displaystyle \underset{i=1}{\overset{k}{}}}(j_i{\displaystyle \frac{m^2}{4}}{\displaystyle \frac{r^2}{4}}+{\displaystyle \underset{s=i+1}{\overset{k}{}}}j_s+{\displaystyle \frac{n^2}{4}})Cx^{\frac{r^2}{4}\frac{m^2}{4}\frac{n^2}{4}_ij_i},`$ (17)
where $`C`$ is a constant that depends on $`𝒴`$ (we may assume that $`C`$ is equal to $`1`$). If we compare (3.3) with (3.1) we see that products appearing in both expressions are the same if we interchange $`x`$ with $`\frac{r^2}{4}`$ and $`\frac{m^2}{4}`$ with $`\frac{n^2}{4}`$. In other words the expression $`w_3^{},𝒴(w_1,x)v_{sing}=0`$ if and only if the corresponding projection inside $`A(L(1,\frac{n^2}{4})_{A(L(1,0)}A(L(1,\frac{m^2}{4}))`$ is zero (notice that now $`L(1,\frac{n^2}{4})`$ and $`L(1,\frac{m^2}{4})`$ changed positions). We know that
$$A(L(1,\frac{n^2}{4})_{A(L(1,0)}A(L(1,\frac{m^2}{4}))\frac{[x]}{_{iJ_{n,m}}x\frac{i^2}{4}}.$$
Because of (12), $`J_{n,m}J_{n,m}`$ (as multisets). Therefore
$$w_3^{},𝒴(w_1,x)v_{sing}=0$$
holds. Thus we obtain a non–trivial intertwining operator $`\overline{𝒴}`$ of the type $`\left(\genfrac{}{}{0pt}{}{M(1,\frac{r^2}{4})^{}}{L(1,\frac{m^2}{4})L(1,\frac{n^2}{4})}\right).`$ Now,
$$I\left(\genfrac{}{}{0pt}{}{M(1,\frac{r^2}{4})^{}}{L(1,\frac{m^2}{4})L(1,\frac{n^2}{4})}\right)I\left(\genfrac{}{}{0pt}{}{L(1,\frac{n^2}{4})}{L(1,\frac{m^2}{4})M(1,\frac{r^2}{4})}\right).$$
Because of our initial assumption $`nr`$, and $`mnrm+n`$ it follows that $`mrnm+r`$, therefore we can repeat the whole procedure for $`M(1,\frac{r^2}{4})`$ so we end up with a non–trivial intertwining operator of the type
$$\left(\genfrac{}{}{0pt}{}{L(1,\frac{n^2}{4})}{L(1,\frac{m^2}{4})L(1,\frac{r^2}{4})}\right).$$
If $`r<q`$ then we pick the intertwining operator (14) and the same reasoning leads to a non–trivial intertwining operator of the type
$$\left(\genfrac{}{}{0pt}{}{L(1,\frac{r^2}{4})}{L(1,\frac{m^2}{4})L(1,\frac{n^2}{4})}\right).$$
This also follows from the duality property for the intertwining operators. If we summarized everything we obtain
###### Theorem 3.3
$$\mathrm{dim}I\left(\genfrac{}{}{0pt}{}{L(1,\frac{r^2}{4})}{L(1,\frac{m^2}{4})L(1,\frac{n^2}{4})}\right)=1$$
if and only if $`r\{m+n,\mathrm{},|mn|\}`$.
###### Theorem 3.4
Let $`𝒜`$ be a free Abelian group on the set $`\{a(m):m\}`$ and
$$\times :𝒜\times 𝒜𝒜$$
a binary operation defined by the formula
$$a(m)\times a(n)=\underset{r0}{}𝒩_{L(1,\frac{m^2}{4})L(1,\frac{n^2}{4})}^{L(1,\frac{r^2}{4})}a(r).$$
Then $`𝒜`$ is a commutative associative ring with the multiplication
$$a(m)\times a(n)=a(m+n)+a(m+n2)+\mathrm{}+a(|mn|),$$
i.e. $`𝒜`$ is isomorphic to the representation ring $`ep(𝔰l(2,))`$.
###### Remark 3.1
In general if $`M`$ is any $`L(1,0)`$–module and
$$𝒴I\left(\genfrac{}{}{0pt}{}{M}{L(1,\frac{m^2}{4})L(1,\frac{n^2}{4})}\right),$$
then $`M`$ is not necessary completely reducible. Also, note that we excluded the case $`mn=0`$. If $`m`$ or $`n`$ are equal to zero then we deal with intertwining operators among two irreducible modules and vertex operator algebras, which are well known.
Another interesting fact is that in the case (12) the module $`A(L(1,\frac{m^2}{4}))_{A(L(1,0))}L(1,\frac{n^2}{4})(0)`$ is not completely reducible. This fact was exploited in \[M\] where we study logarithmic intertwining operators.
Note that in our proof we actually analyzed more carefully the failure of Frenkel-Zhu’s formula. One should not expect to apply our procedure in the more general setting, because our Virasoro vertex operator algebra has a quite simple structure. Certainly it would be interesting to study a class of vertex operator algebra for which
$$A(W_1)_{A(V)}W_2(0)A(W_2)_{A(V)}W_1(0),$$
(18)
for any choice of irreducible modules $`W_1`$ and $`W_2`$. Then we hope that for this class of vertex algebras some version of Frenkel-Zhu’s formula indeed apply. Assumption (18) turns out to be very natural since
$$I\left(\genfrac{}{}{0pt}{}{W_3}{W_1W_2}\right)I\left(\genfrac{}{}{0pt}{}{W_3}{W_2W_1}\right).$$
(19)
## 4 Construction of all intertwining operators for the family $`_1`$
### 4.1 $`V_L`$ vertex operator algebra and its irreducible modules
Let $`L`$ be a rank one even lattice with a generator $`\beta `$ normalized such that $`\beta ,\beta =1`$ and let $`\alpha =\sqrt{2}\beta `$. Thus $`\alpha ,\alpha =2`$. As in \[FLM\], \[DL\] we define $`V_L`$ as a vector space
$$V_L=M(1)[L],$$
where $`M(1)`$ is the level one irreducible module for Heisenberg algebra $`\widehat{h}_Z`$ associated to one–dimensional abelian algebra $`h=L_{}`$ and $`[L]`$ is the group algebra of $`L`$ with a generator $`e^\alpha `$. Put $`\omega =\frac{1}{2}\beta (1)^2`$. Then $`V_L`$ is a vertex operator algebra (see \[FLM\]) with the Virasoro element $`\omega `$. We have a decomposition
$$V_L=\underset{m}{}M(1)e^{m\alpha }.$$
Let $`L^o`$ be a dual lattice, $`L^o/L/2`$. Then (as in \[DL\]), for a nontrivial coset representative, we obtain an irreducible $`V_L`$–module $`V_{L+1/2}`$, which can be decomposed as
$$V_{L+1/2}=\underset{m}{}M(1)e^{m\alpha +1/2\alpha }.$$
Moreover, $`V_{L+1/2}`$, $`V_L`$ is (up to equivalence) complete list of irreducible $`V_L`$–modules. Furthermore, one can equip the space $`W=V_LV_{L+1/2}`$ (as in \[DL\]) with the structure of the generalized vertex operator algebra. We will neglect this fact in our considerations.
For every module $`W`$ for the Virasoro algebra on which $`L(0)`$ acts semisimple we define a formal character (or a $`q`$-graded dimension) by
$$ch_q(W)=\underset{n\mathrm{Spec}L(0)}{}\mathrm{dim}(W_n)q^n.$$
From the Proposition 2.1 it follows that
$$ch_q(L(1,\frac{m^2}{4}))=\frac{q^{\frac{m^2}{4}}q^{\frac{(m+2)^2}{4}}}{q^{1/24}\eta (q)}.$$
Then it is not hard to obtain
$`ch_q(V_L)={\displaystyle \underset{n0}{}}(2n+1)ch_q(L(1,n^2))`$
$`ch_q(V_{L+1/2})={\displaystyle \underset{n0}{}}(2n+2)ch_q\left(L(1,{\displaystyle \frac{(2n+1)^2}{4}})\right).`$ (20)
Consider the vectors
$$x=e^\alpha ,y=e^\alpha ,h=\alpha (1)\iota (0),$$
which span $`(V_L)_1`$. These vectors span a Lie algebra isomorphic to $`𝔰l(2,)`$. $`x_0`$, $`y_0`$ and $`h_0`$ as act derivatives on $`W`$. The following result was obtained in \[DG\].
###### Proposition 4.1
As $`(L(1,0),𝔰𝔩_2)`$–module
$$V_L\underset{m0}{}L(1,m^2)V(2m),$$
where $`V(2m)`$ is an irreducible $`2m+1`$ dimensional $`𝔰𝔩_2`$–module.
The proof uses the result from \[DLM\], \[DM1\] about the decomposition of the vertex operator algebra $`V`$ with respect to a “dual” pair $`(V^G,G)`$ where $`G=Aut(G)`$ is a compact (or finite) group and $`V^G`$ is a $`G`$–stable subvertex operator algebra. This can be modified when instead of group $`G`$ we work with the Lie algebra.
Since $`V_{L+1/2}`$ is a module for the pair $`(V_L^{𝔰𝔩_2},𝔰𝔩_2)`$ then by using (4.1) we derive
$$V_{L+1/2}\underset{m0}{}L(1,\frac{(2m+1)^2}{4})V(2m+1),$$
(21)
where $`V(2m+1)`$ is a $`2m+2`$–dimensional $`𝔰𝔩_2`$–module. It easy to see that $`V(2m+1)`$ is irreducible $`𝔰𝔩_2`$–module.
###### Remark 4.1
Note that $`V^{sl_2}`$ ($`sl_2`$–stable vertex operator algebra) is exactly $`V^G`$ where $`GSO(3)`$ is a (full) group of automorphisms of $`V_L`$. It is well known that every irreducible representation can be obtain as a representation of $`SL(2,)`$, since $`PSL(2,)SO(3)`$. In particular every such finite-dimesional representation is odd–dimensional.
Since, $`V_{L+1/2}`$ is an irreducible $`V_L`$–module we have the Jacobi identity
$`x_0^1\delta \left({\displaystyle \frac{x_1x_2}{x_0}}\right)Y(u,x_1)Y(v,x_2)w`$
$`x_0^1\delta \left({\displaystyle \frac{x_2x_1}{x_0}}\right)Y(v,x_2)Y(u,x_1)w`$
$`=x_2^1\delta \left({\displaystyle \frac{x_1x_0}{x_2}}\right)Y(Y(u,x_0)v,x_2)w,`$ (22)
for every $`uV_L`$, $`vV_{L+1/2}`$ and $`wW`$. Also, for
$$𝒴I\left(\genfrac{}{}{0pt}{}{V_L}{V_{L+1/2}V_{L+1/2}}\right),$$
we have
$`x_0^1\delta \left({\displaystyle \frac{x_1x_2}{x_0}}\right)Y(u,x_1)𝒴(v,x_2)w`$
$`x_0^1\delta \left({\displaystyle \frac{x_2x_1}{x_0}}\right)𝒴(v,x_2)Y(u,x_1)w`$
$`=x_2^1\delta \left({\displaystyle \frac{x_1x_0}{x_2}}\right)𝒴(Y(u,x_0)v,x_2)w.`$ (23)
###### Remark 4.2
Note that $`W`$ can not be equipped with a vertex operator superalgebra structure. If $`u,vV_{L+1/2}`$ then we do not get Jacobi identity in the form (4.1) or (4.1), but rather generalized identity where the delta function is suitably multiplied with the terms of the type $`\left(\frac{x_1x_0}{x_2}\right)^{1/2}`$. Studying this (generalized) Jacobi identity is useful for studying convergence and the extension properties for the intertwining operators (cf. \[H1\]).
### 4.2 Intertwining operators for the family $`_1`$.
Let $`V(i)`$, $`i`$ be an irreducible $`sl_2`$–module considered as a subspace of $`W`$ which corresponds to the decompositions in Proposition 4.1 and (21). Fix a positive integer $`j`$. We introduce a basis $`u_j(m)`$, $`m\{j,j2,\mathrm{},j\}`$ for $`V(j)`$, such that the following relations are satisfied,
$`h.u_j(m)=mu_j(m)`$
$`x.u_j(m)={\displaystyle \frac{\sqrt{(j+m+2)(jm)}}{2}}u_j(m+2)`$
$`y.u_j(m)={\displaystyle \frac{\sqrt{(j+m)(jm+2)}}{2}}u_j(m2),`$ (24)
where $`u_j(k)=0`$ for $`k\{j,\mathrm{},j\}`$. Also, we choose a dual basis $`u_j^{}(m)`$ for $`V(j)^{}`$ such that $`<u_j^{}(m),u_j(n)>=\delta _{m,n}`$. Define $`<g.u^{},v>=<u^{},g.v>`$. Then $`V(j)^{}`$ became a $`sl_2`$–module and an isomorphism from $`V(j)`$ to $`V(j)^{}`$ is given by $`\mu (u_j(m))=(1)^{jm}u_j^{}(m).`$ By using this identification, for $`j_1,j_2,j_3`$ and $`j_im_ij_i`$, $`i=1,2,3`$, we introduce real numbers (Clebsch–Gordan coefficients) $`\left(\begin{array}{ccc}j_1& j_2& j_3\\ m_1& m_2& m_3\end{array}\right),`$ such that
$$u_{j_1}(m_1)u_{j_2}(m_2)=\underset{j_3=|j_1j_2|}{\overset{j_3=j_1+j_2}{}}\left(\begin{array}{ccc}j_1& j_2& j_3\\ m_1& m_2& m_1+m_2\end{array}\right)u_{j_3}(m_1+m_2).$$
(25)
First we need an auxiliary result which is slightly modified result from \[DM1\] and \[DG\].
###### Proposition 4.2
Suppose that $`V`$ is a vertex operator algebra and $`W_1`$, $`W_2`$ and $`W_3`$ three irreducible $`V`$–modules. Let $`v_iW_1,w_iW_2`$, $`i=1,\mathrm{},k`$ be homogeneous elements such that $`v_i0`$ and $`w_i`$ are linearly independent. Then
$$\underset{i=1}{\overset{k}{}}𝒴(v_i,x)w_i0.$$
Now. let us go back to our vertex operator algebra $`V_L`$. Let $`𝒴`$ be any intertwining operator of the type
$$\left(\genfrac{}{}{0pt}{}{V_L}{V_{L+1/2}V_{L+1/2}}\right),\left(\genfrac{}{}{0pt}{}{V_{L+1/2}}{V_LV_{L+1/2}}\right)\mathrm{or}\left(\genfrac{}{}{0pt}{}{V_L}{V_LV_L}\right).$$
(26)
By using the Proposition 4.2 the map
$$𝒴(,x):V(j_1)V(j_2)W\{x\}$$
is injective. and for every $`m_1,m_2`$ and $`j_1`$, $`j_2`$ there is a $`p`$ such that
$$u_{j_1}(m_1)_pu_{j_2}(m_2)=\underset{j_3=|j_1j_2|}{\overset{j_3=j_1+j_2}{}}k(j_1,j_2,j_3,m_1,m_2,m_1+m_2)u_j(m_1+m_2),$$
where $`k(j_1,j_2,j_3,m_1,m_2,m_1+m_2)`$ is a (non–zero) multiple of
$$\left(\begin{array}{ccc}j_1& j_2& j_3\\ m_1& m_2& m_1+m_2\end{array}\right).$$
(in the special case $`𝒴=Y`$ this fact was noticed in \[DG\]).
Now it is clear that if $`\left(\begin{array}{ccc}j_1& j_2& j_3\\ m_1& m_2& m_1+m_2\end{array}\right)0`$, then the $`L(1,0)`$–module generated by $`𝒴(u_{j_1}(m_1),x)u_{j_2}(m_2)`$ contains a copy of $`L(1,\frac{j_3^2}{4})`$. Since $`L(1,0)`$ is contained in $`V_L`$ and $`L(1,\frac{m^2}{4})`$ is an $`L(1,0)`$–module then we obtain the following Jacobi identity
$`x_0^1\delta \left({\displaystyle \frac{x_1x_2}{x_0}}\right)Y(u,x_1)𝒴(v,x_2)wx_0^1\delta \left({\displaystyle \frac{x_2x_1}{x_0}}\right)𝒴(v,x_2)Y(u,x_1)w`$
$`=x_2^1\delta \left({\displaystyle \frac{x_1x_0}{x_2}}\right)𝒴(Y(u,x_0)v,x_2)w,`$ (27)
for $`uL(1,0)`$, $`vL(1,\frac{j_1^2}{4})`$ and $`wL(1,\frac{j_2^2}{4})`$ (here $`v`$ and $`w`$ lie in Vir–submodules generated by $`u_{j_1}(m_1)`$ and $`u_{j_2}(m_2)`$, respectively).
Now we can push down $`𝒴`$ to $`L(1,\frac{j_3^2}{4})`$, which is generated by the vector $`u_{j_3}(m_1+m_2)`$, since for every $`j_1,j_2`$ and $`|j_1j_2|j_3j_1+j_2`$ we can choose a pair $`m_1`$, $`m_2`$ and a $`𝒴`$ of the appropriate type (26) such that
$$\left(\begin{array}{ccc}j_1& j_2& j_3\\ m_1& m_2& m_1+m_2\end{array}\right)0.$$
We obtain an intertwining operator of the type
$$\left(\genfrac{}{}{0pt}{}{L(1,\frac{j_3^2}{4})}{L(1,\frac{j_1^2}{4})L(1,\frac{j_2^2}{4})}\right),$$
and this is the end of the construction.
## 5 Lie superalgebra $`𝔬𝔰𝔭(1|2)`$ and $`\mathrm{Rep}(𝔬𝔰𝔭(1|2))`$
The Lie superalgebra $`𝔬𝔰𝔭(1|2)`$ is a graded extension of the finite–dimensional Lie algebra $`𝔰𝔩(2,𝐂)`$. It has three even generators $`x,y`$ and $`h`$, and two odd generators $`\phi `$ and $`\chi `$, that satisfy:
$$[h,x]=2x,[h,y]=2y,[x,y]=h,$$
$$[x,\chi ]=\chi ,[x,\phi ]=\phi ,[y,\chi ]=\chi ,[y,\phi ]=\phi ,$$
$$[h,\phi ]=\phi ,[h,\chi ]=\chi ,$$
$$\{\chi ,\phi \}=2h,\{\chi ,\chi \}=2x,\{\phi ,\phi \}=2y.$$
Generators $`\{x,y,h\}`$ span a Lie algebra isomorphic to $`𝔰𝔩(2,)`$, and this fact makes the representation theory of $`𝔬𝔰𝔭(1|2)`$ quite simple. All irreducible $`𝔬𝔰𝔭(1|2)`$–modules can be constructed in the following way. Fix a positive half integer $`j`$ ($`2j𝐍`$) and a $`4j+1`$–dimensional vector space $`V(j)`$ spanned by the vectors $`\{v_j,v_{j1/2},\mathrm{},v_j\}`$, with the following actions:
$`x.v_i=\sqrt{[ji][j+i+1]}v_{i+1},`$
$`y.v_i=\sqrt{[j+i][ji+1]}v_{i1},`$
$`h.v_i=2iv_i.`$ (28)
If $`2(ij)𝐙`$ then we define
$`\phi .v_i=\sqrt{j+i}v_{i1/2},`$
$`\chi .v_i=\sqrt{ji}v_{i+1/2},`$ (29)
otherwise
$`\phi .v_i=\sqrt{ji+1/2}v_{i1/2},`$
$`\chi .v_i=\sqrt{j+i+1/2}v_{i+1/2}.`$ (30)
In all these formulas $`v_j=0`$ if $`j\{j,j\frac{1}{2},\mathrm{},j\}`$. It is easy to see that each $`V(j)`$ is an irreducible $`𝔬𝔰𝔭(1|2)`$–module and that every finite dimensional irreducible representation of $`𝔬𝔰𝔭(1|2)`$ is isomorphic to $`V(j)`$ for some $`j/2`$.
The representations with $`j`$ we call even, and the representations with $`j+\frac{1}{2}`$ we call odd. We extend this definition for an arbitrary element of $`V\mathrm{Rep}(𝔬𝔰𝔭(1|2))`$. The corresponding decomposition is $`V=V_{\mathrm{even}}+V_{\mathrm{odd}}`$.
It is a pleasant exercise to decompose the tensor product $`V(i)V(j)`$. The following result is well–known:
$$V(i)V(j)\underset{k=|ij|,k𝐍/2}{\overset{i+j}{}}V(k).$$
(31)
## 6 $`N=1`$ Neveu-Schwarz superalgebra and its minimal models
The $`N=1`$ Neveu-Schwarz superalgebra is given by
$$𝔫𝔰=\underset{n}{}L_n\underset{n}{}G_{n+1/2}C,$$
together with the following $`N=1`$ Neveu-Schwarz relations:
$`[L_m,L_n]`$ $`=`$ $`(mn)L_{m+n}+{\displaystyle \frac{C}{12}}(m^3m)\delta _{m+n,0},`$
$`[L_m,G_{n+1/2}]`$ $`=`$ $`\left({\displaystyle \frac{m}{2}}\left(n+{\displaystyle \frac{1}{2}}\right)\right)G_{m+n+1/2},`$
$`[G_{m+1/2},G_{n1/2}]`$ $`=`$ $`2L_{m+n}+{\displaystyle \frac{C}{3}}(m^2+m)\delta _{m+n,0},`$
$`[C,L_m]`$ $`=`$ $`0,`$
$`[C,G_{m+1/2}]`$ $`=`$ $`0`$
for $`m,n`$. We have the standard triangular decomposition $`𝔫𝔰=𝔫𝔰_+𝔫𝔰_0𝔫𝔰_{}`$ (cf. \[KWa\]). For every $`(h,c)𝐂^2`$, we denote by $`M(c,h)`$ Verma module for $`𝔫𝔰`$ algebra. For each $`(p,q)𝐍^2`$, $`p=q\mathrm{mod}2`$, let us introduce a family of complex ’curves’ $`(h_{p,q}(t),c(t))`$;
$$h_{p,q}(t)=\frac{1p^2}{8}t^1+\frac{1pq}{4}+\frac{1q^2}{8}t,$$
$$c(t)=\frac{15}{2}+3t^1+3t.$$
Then from the determinant formula (see \[KWa\] ) it follows that $`M(c,h)`$ is reducible if and only if there is a $`t𝐂`$ and $`p,q𝐍`$, $`p=q\mathrm{mod}2`$ such that $`c=c(t)`$ and $`h=h_{p,q}(t)`$. In this case $`M(c,h)`$ has a singular vector (i.e., a vector annihilated by $`𝔫𝔰_+`$ ) of the weight $`h+\frac{pq}{2}`$. Any such vector we denote by $`v_{\frac{pq}{2}}`$.
In this paper we are interested in the case $`t=1`$. Then $`c(1)=\frac{3}{2}`$ and $`h_{p,q}(1)=\frac{(pq)^2}{8}`$. $`h_{p,q}(1)=h_{1,pq+1}(1)`$, so we consider only the case $`h_{1,q}:=h_{1,q}(1)`$, (here $`q`$ is odd and positive). Hence, each Verma module $`M(\frac{3}{2},h_{1,q})`$ is reducible.
The following result easily follows from \[D\] (or \[AA\]) and \[KWa\]:
###### Proposition 6.1
For every odd $`q`$, $`M(\frac{3}{2},h_{1,q})`$ has the following embedding structure
$$\mathrm{}M(\frac{3}{2},h_{1,q+4})M(\frac{3}{2},h_{1,q+2})M(\frac{3}{2},h_{1,q})0.$$
(32)
Moreover, we have the following exact sequence:
$$0M(\frac{3}{2},h_{1,q+2})M(\frac{3}{2},h_{1,q})L(\frac{3}{2},h_{1,q})0,$$
(33)
where $`L(\frac{3}{2},h_{1,q})`$ is the corresponding irreducible quotient.
Benoit and Saint-Aubin (cf. \[BSA2\]) found an explicit expression for the singular vectors $`P_{\mathrm{sing}}v_{1,q}M(\frac{3}{2},h_{1,q})`$ that generates the maximal submodule:
$$\underset{N;k_1,\mathrm{},k_N}{}\underset{\sigma S_N}{}(1)^{\frac{qN}{2}}c(k_{\sigma (1)},\mathrm{},k_{\sigma (k)})G(k_1/2)\mathrm{}G(k_N/2)v_{1.q},$$
(34)
where $`S_N`$ is a symmetric group on $`N`$ letters and the first summation is over all the partitions of $`q`$ into the odd integers $`k_1,..,k_N`$ and
$$c(k_{\sigma (1)},\mathrm{},k_{\sigma (k)})=\underset{i=1}{\overset{N}{}}\left(\genfrac{}{}{0pt}{}{k_i1}{(k_i1)/2}\right)\underset{j=1}{\overset{(N1)/2}{}}\frac{4}{\sigma _{2j}\rho _{2j}},$$
where $`\sigma _j=_{l=1}^jk_l`$ and $`\rho _j=_{l=j}^Nk_l`$.
In the special case: $`q=1`$, $`h_{1,1}=0`$, $`M(\frac{3}{2},0)`$ has a singular vector $`G(1/2)v`$ which generate the maximal submodule. By quotienting we obtain a vacuum module $`L(\frac{3}{2},0)=M(\frac{3}{2},0)/G(1/2)v_{3/2,0}`$.
## 7 $`N=1`$ superconformal vertex operator superalgebra and intertwining operators
We use the definition of $`N=1`$ superconformal vertex operator superalgebra (with and without odd variables) as in \[B\] (cf. \[KV\]) and \[HM\] (see also \[KW\]).
Let $`\phi `$ be a Grassman (odd) variable such that $`\phi ^2=0`$. Every $`N=1`$ superconformal vertex operator superalgebra $`(V,Y,\mathrm{𝟏},\tau )`$ can be equipped with a structure of $`N=1`$ superconformal vertex operator algebra with an odd variable via
$`Y(,(x,\phi )):VV`$ $``$ $`V((x))[\phi ],`$
$`uv`$ $``$ $`Y(u,(x,\phi ))v,`$
where
$$Y(u,(x,\phi ))v=Y(u,x)v+\phi Y(G(1/2)u,x)v$$
for $`u,vV`$.
The same formula can be used in the case of modules for the superconformal vertex operator superalgebra $`(V,Y,\mathrm{𝟏},\tau )`$ (see \[HM\]).
It is known (see \[KW\]) that $`V(c,0):=M(c,0)/G(1/2)v_{c,0}`$ <sup>2</sup><sup>2</sup>2We write $`L(c,0)`$ if $`V(c,0)`$ is irreducible. is a $`N=1`$ superconformal vertex operator superalgebra. Also, every lowest weight $`𝔫𝔰`$–module with the central charge $`c`$, is a $`V(c,0)`$–module. If $`c=\frac{3}{2}`$ then $`V(\frac{3}{2},0)=L(\frac{3}{2},0)`$. Hence
###### Proposition 7.1
Every irreducible $`L(\frac{3}{2},0)`$–module is isomorphic to $`L(\frac{3}{2},h)`$, for some $`h𝐂`$.
Proof: The proof is essentially the same as the one in Proposition 3.1.
Among all irreducible $`L(\frac{3}{2},0)`$–modules we distinguish modules isomorphic to $`L(\frac{3}{2},h_{1,q})`$, $`q2𝐍1`$. These representations we call degenerate minimal models.
### 7.1 Intertwining operators and its matrix coefficients
The notation of an intertwining operators for $`N=1`$ superconformal vertex operator algebras is introduced in \[KW\] and \[HM\].
Let $`W_1`$, $`W_2`$ and $`W_3`$ be a triple of $`V`$–modules and $`𝒴`$ an intertwining operator of type $`\left(\genfrac{}{}{0pt}{}{W_3}{W_1W_2}\right)`$. Then we consider the corresponding intertwining operator with an odd variable (cf. \[HM\]):
$`𝒴(,(x,\phi )):W_1W_2`$ $``$ $`W_3\{x\}[\phi ]`$
$`w_{(1)}w_{(2)}`$ $``$ $`𝒴(w_{(1)},(x,\phi ))w_{(2)},`$
such that
$$𝒴(w_{(1)},(x,\phi ))w_{(2)}=𝒴(w_{(1)},x)w_{(2)}+\phi 𝒴(G(1/2)w_{(1)},x)w_{(2)}.$$
Let $`w_1`$ be a lowest weight vector for the Neveu-Schwarz algebra of the weight $`h`$. From the Jacobi identity we derive the following formulas:
$`[L(n),𝒴(w_1,x_2)]=(x_2^{n+1}{\displaystyle \frac{}{x_2}}+(1n)h)𝒴(w_1,x_2),`$
$`[G(n1/2),𝒴(w_1,x_2)]=x_2^n𝒴(G(1/2)w_1,x_2),`$
$`[L(n),𝒴(G(1/2)w_1,x_2)]=(x_2^{n+1}{\displaystyle \frac{}{x_2}}+(1n)(h+{\displaystyle \frac{1}{2}})𝒴(G(1/2)w_1,x_2),`$
$`[G(n1/2),𝒴(G(1/2)w_1,x_2)]=(x_2^n{\displaystyle \frac{}{x_2}}2nhx_2^{n1})𝒴(w_1,x_2).`$ (35)
In the odd formulation we obtain
$`[L(n),𝒴(w_1,(x_2,\phi ))]`$
$`=(x_2^{n+1}_{x_2}+(1n)x_2^n(h+1/2\phi _\phi ))𝒴(w_1,(x_2,\phi ))`$
$`[G(n1/2),𝒴(w_1,(x_2,\phi ))]`$
$`=(x_2^n(_\phi \phi _{x_2})2nx_2^{n1}(h\phi )𝒴(w_1,(x_2,\phi )),`$ (36)
where $`_\phi `$ is the odd (Grassmann) derivative.
### 7.2 Even and odd intertwining operators
In \[HM\] we proved that every intertwining operator
$$𝒴I\left(\genfrac{}{}{0pt}{}{L(c,h_3)}{L(c,h_1)L(c,h_2)}\right)$$
is uniquely determined by the operators $`𝒴(w_1,x)`$ and $`𝒴(G(1/2)w_1,x)`$, where $`w_1`$ is the lowest weight vector of $`L(c,h_1)`$. This fact will be used later in connection with the following definition.
###### Definition 7.1
Let $`||`$ denote the ($`/2`$–valued) parity operator from the union of odd and even subspaces for $`V`$–modules $`W_i`$, $`i=1,2,3`$. An intertwining operator $`𝒴I\left(\genfrac{}{}{0pt}{}{W_3}{W_1W_2}\right)`$ is:
* even, if
$$|\mathrm{Coeff}_{x^s}𝒴(w_1,x)w_2|=|w_1|+|w_2|,$$
* odd, if
$$|\mathrm{Coeff}_{x^s}𝒴(w_1,x)w_2|=|w_1|+|w_2|+1,$$
for every $`s`$ and every $`/2`$–homogeneous vectors $`w_1`$ and $`w_2`$.
The space of even (odd) intertwining operators of the type $`\left(\genfrac{}{}{0pt}{}{W_3}{W_1W_2}\right)`$ we denote by $`I\left(\genfrac{}{}{0pt}{}{W_3}{W_1W_2}\right)_{\mathrm{even}}`$ ($`I\left(\genfrac{}{}{0pt}{}{W_3}{W_1W_2}\right)_{\mathrm{odd}}`$). In general we do not have a decomposition of $`I\left(\genfrac{}{}{0pt}{}{W_3}{W_1W_2}\right)`$ into the even and odd subspaces.
### 7.3 Frenkel-Zhu’s theorem for vertex operator superalgebras
According to \[KW\] (after \[Z\]), to every vertex operator superalgebra we can associate the Zhu’s associative algebra $`A(V)`$. If $`V=L(c,0)`$, $`A(L(c,0))𝐂[y]`$. where $`y=[(L(2)L(1))\mathrm{𝟏}]=[L(2)\mathrm{𝟏}]`$ (because of the calculations that follow it is convenient to use $`y=[(L(2)L(1))\mathrm{𝟏}]`$). Also to every $`V`$–module $`W`$ we associate a $`A(V)`$–bimodule $`A(W)`$ (cf. \[KW\]). In a special case $`W=M(c,h)`$, we have
$$A(M_{𝔫𝔰}(c,h))=M_{𝔫𝔰}(c,h)/O(M_{𝔫𝔰}(c,h)),$$
where
$`O(M_{𝔫𝔰}(c,h))=\{L(n3)2L(n2)+L(1)v,`$
$`G(n1/2)G(n3/2)v,n0,vM(c,h)\}.`$ (37)
It is not hard to see that, as $`[y]`$–bimodule,
$$A(M(c,h))𝐂[x,y]𝐂[x,y]v,$$
where $`v=[G(1/2)v_h]`$ and
$$y=[L(2)L(1)],x=[L(2)2L(1)+L(0)].$$
Let $`W_1`$, $`W_2`$ and $`W_3`$ be three $`/2`$–gradable irreducible $`V`$–modules such that $`\mathrm{Spec}L(0)|_{W_i}h_i+`$, $`i=1,2,3`$ and $`𝒴I\left(\genfrac{}{}{0pt}{}{W_3}{W_1W_2}\right)`$. We define $`o(w_1):=\mathrm{Coeff}_{x^{h_3h_1h_2}}𝒴(w_1,x)`$. Because the fusion rules formula in \[FZ\] needs some modifications (cf. \[L1\]) the same modification is necessary for the main Theorem in \[KW\] (this can be done with a minor super–modifications along the lines of \[L1\]). Nevertheless (cf. \[KW\]):
###### Theorem 7.1
The mapping
$$\pi :I\left(\genfrac{}{}{0pt}{}{W_3}{W_1W_2}\right)\mathrm{Hom}_{A(V)}(A(W_1)_{A(V)}W_2(0),W_3(0)),$$
such that
$$\pi (𝒴)(w_1w_2)=o(w_1)w_2,$$
(38)
is injective.
## 8 Some Lie superalgebra homology
In this section we recall some basic definition from the homology theory of infinite dimensional Lie superalgebras which is in the scope of the monograph \[F\] (in the cohomology setting though).
Let $``$ be an any (possibly infinite dimensional) $`/2`$–graded Lie superalgebra with the $`/2`$–decomposition: $`=_0_1`$. and let $`M=M_0M_1`$ be any $`𝐙_2`$–graded $``$–module, such that the gradings are compatible. Then, we form a chain complex $`(C,d,M)`$ (for details see \[F\]),
$$0\stackrel{d_0}{}C_0(,M)\stackrel{d_1}{}C_1(,M)\stackrel{d}{}\mathrm{},$$
where
$$C_q(,M)=\underset{q_0+q_1=q}{}M\mathrm{\Lambda }^{q_0}_0S^{q_1}_1,$$
$$C_q^p(,M)=\underset{\stackrel{q_0+q_1=q}{q_1+r=p\mathrm{mod2}}}{}M_r\mathrm{\Lambda }^{q_0}_0S^{q_1}_1,$$
for $`p=0,1`$. The mappings $`d`$ are super–differentials. For $`q𝐍`$ and $`p=0,1`$, we define $`q`$–th homology with coefficients in $`M`$ as:
$$H_q^p(,M)=\mathrm{Ker}(d_q(C_q^p(,M)))_p/(d_{r+1}(C_{q+1}^p(,M)))_p.$$
(39)
In a special case $`q=0`$, we have
$$H_0^0(,M)=M_0/(_0M_0+_1M_1),$$
and
$$H_0^1(,M)=M_1/(_1M_0+_0M_1).$$
We want to calculate $`H_q(_s,L(\frac{3}{2},h_{1,q}))`$. for the Lie superalgebra
$$_s=\underset{n0}{}_s(n),$$
where $`_s(n)`$ is spanned by the vectors $`L(n3)2L(n2)+L(n1)`$ and $`G(n1/2)G(n3/2)`$ , $`n𝐍`$. From (7.3) we see (cf. \[HM\]) that $`H_0(_s,M(c,h))`$ is a $`[y]`$–bimodule such that:
$$H_0(_s,M(c,h))A(M(c,h))𝐂[x,y]𝐂[x,y]v.$$
(40)
###### Remark 8.1
It is more involved to calculate $`H_0((_s,L(c,h))`$, so we consider only the special case $`c=\frac{3}{2}`$, $`h=h_{1,q}`$, $`q`$ odd. As in the Virasoro case, it is easy to show that the space $`H_p(_s,L(\frac{3}{2},h_{1,q}))`$ is infinite dimensional for very $`p,q,s𝐍`$, and finitely generated as a $`A(L(3/2,0))`$–module. Moreover, it is not hard to see (by using the same method as in the Virasoro case) that
$$\mathrm{Ext}_{ns,𝒪}^1(L(\frac{3}{2},h_{1,q}),L(\frac{3}{2},h_{1,r}))$$
is non–trivial (and one-dimensional) if and only if $`|rq|=2`$.
In the minimal models case we expect a substantially different result (cf \[FF1\]).
###### Conjecture 8.1
Let $`c_{p,q}=\frac{3}{2}\left(12\frac{(pq)^2}{pq}\right)`$ and $`h_{p,q}^{m,n}=\frac{(npmq)^2(pq)^2}{8pq}`$. Then
$$\mathrm{dim}H_q(_s,L(c_{p,q},h_{p,q}^{m,n}))<\mathrm{},$$
for every $`q𝐍`$.
There is strong evidence that Conjecture (8.1) holds based on \[A\] and an example $`c=\frac{11}{14}`$ treated in Appendix of \[HM\].
The main difference between the minimal models and the degenerate models is the fact that the maximal submodule for a minimal model is generated by two singular vectors, compared to $`M(\frac{3}{2},h_{1,q})`$ where the maximal submodule is generated by a single singular vector.
## 9 Benoit-Saint-Aubin’s formula projection formulas
### 9.1 Odd variable formulation
We have seen before how to derive the commutation relation between generators of $`𝔫𝔰`$ superalgebra and $`𝒴(w_1,x)`$ where $`w_1`$ is a lowest weight vector for $`ns`$. We fix $`𝒴I\left(\genfrac{}{}{0pt}{}{L(\frac{3}{2},h)}{L(\frac{3}{2},h_{1,r})L(\frac{3}{2},h_{1,q})}\right)`$ and consider the following matrix coefficient,
$$w_3^{},𝒴(w_1,x,\phi )P_{\mathrm{sing}}w_2,$$
(41)
where $`P_{sing}w_2=v_{1,q}`$ (cf. (7),$`\mathrm{deg}(P_{sing})=q/2`$) and $`w_i`$, $`i=1,2,3`$ are the lowest weight vectors.
Since all modules are irreducible, by using a result from \[HM\] (Proposition 2.2), we get
$$w_3^{},𝒴(w_1,x,\phi )w_2=c_1x^{hh_{1,q}h_{1,r}}+c_2\phi x^{hh_{1,q}h_{1,r}1/2},$$
where $`c_1`$ and $`c_2`$ are constants with the property
$`c_1=c_2=0\mathrm{implies}𝒴=0.`$ (42)
From the formula (7.1)
$$w_3^{},𝒴(w_1,x,\phi )P_{\mathrm{sing}}w_2=P(_{x_2},\phi )w_3^{},𝒴(w_1,x,\phi )w_2,$$
where $`P(_{x_2},\phi )`$ is a certain super-differential operator such that
$$\mathrm{deg}(P_{\mathrm{sing}})=\mathrm{deg}P(_{x_2},\phi )=q/2.$$
Therefore
$$P(_{x_2},\phi )c_1x^{hh_{1,q}h_{1,r}}=\phi C_1(h_{1,q},h_{1,r},h)x^{hh_{1,q}h_{1,r}q/2},$$
and
$$P(_{x_2},\phi )\phi c_2x^{hh_{1,q}h_{1,r}q/2}=C_2(h_{1,q},h_{1,r},h)x^{hh_{1,q}h_{1,r}q/2}.$$
Constants $`C_1(h_{1,q},h_{1,r},h)`$ and $`C_2(h_{1,q},h_{1,r},h)`$ (in slightly different form, but in more general setting) were derived in \[BSA2\]. Considering these coefficients was motivated by deriving formulas for singular vectors from already known singular vectors. By slightly modifying result from \[BSA2\] we obtain
###### Proposition 9.1
Suppose that $`𝒴I\left(\genfrac{}{}{0pt}{}{L(\frac{3}{2},h)}{L(\frac{3}{2},h_{1,r})L(\frac{3}{2},h_{1,q})}\right)`$ and $`P(_x,\phi )`$ are as the above Then, up to a multiplicative constant,
$$C_1(h_{1,q},h_{1,r},h)=\underset{jkj}{}(hh_{1,q+4k})$$
and
$$C_2(h_{1,q},h_{1,r},h)=\underset{j+1/2kj1/2}{}(h+\frac{1}{2}h_{1,q+4k}),$$
for $`j=(r1)/4`$, $`j>0`$ (when $`j=0`$, $`C_2(h_{1,1},h_{1,r},h)=1`$).
Proof: The superdifferential operator $`P(_x,\phi )`$ is obtained by replacing generators $`L(m)`$ and $`G(n1/2)`$ by the superdifferential operators
$$L(m)(x_2^{m+1}_{x_2}+(1m)x_2^m(h_1+1/2\phi _\phi ))$$
(43)
and
$$G(n1/2)(x_2^n(_\phi \phi _{x_2})2nx_2^{n1}(h_1\phi )),$$
(44)
acting on $`w_3^{},𝒴(w_1,x,\phi )w_2`$. This action was calculated in \[BSA2\]. Their results (Formula 3.10 in \[BSA2\]) implies the statement <sup>3</sup><sup>3</sup>3In \[BSA2\] a different sign was used in the equation (43). Still, we obtain the same result if we consider an isomorphic algebra with the generators $`\stackrel{~}{L}(n):=L(n)`$. The same generators were used in \[FF2\]..
### 9.2 BSA formula without odd variables
Since Frenkel-Zhu’s formula does not involve odd variables we need a version of Proposition 9.1 without odd variables (which is of course equivalent). Again $`𝒴I\left(\genfrac{}{}{0pt}{}{L(3/2,h)}{L(3/2,h_{1,r})L(3/2,h_{1,q})}\right)`$ is the same as the above. Then
$$w_3^{},𝒴(w_1,x)P_{sing}w_2=P_2(_x)w_3^{},𝒴(G(1/2)w_1,x)w_2,$$
and
$$w_3^{},𝒴(G(1/2)w_1,x)P_{sing}w_2=P_1(_x)w_3^{},𝒴(w_1,x)w_2,$$
where $`P_1`$ and $`P_2`$ are certain differential operators. If
$$P_2(_x)c_2x^{hh_{1,q}h_{1,r}1/2}=c_2K_2(h_{1,q},h_{1,r},h)x^{hh_{1,q}h_{1,r}q/2},$$
and
$$P_1(_x)c_1x^{hh_{1,q}h_{1,r}}=c_1K_1(h_{1,q},h_{1,r},h)x^{hh_{1,q}h_{1,r}q/2},$$
then, by comparing corresponding coefficients, we obtain
$`K_1(h_{1,q},h_{1,r},h)=C_1(h_{1,q},h_{1,r},h),`$
$`K_2(h_{1,q},h_{1,r},h)=C_2(h_{1,q},h_{1,r},h).`$ (45)
Let us mention that the projection formulas from Proposition 9.1 have a simple explanation terms of super density modules for the Neveu-Schwarz superalgebra.
## 10 Fusion ring for the degenerate minimal models
In order to obtain an upper bound for the fusion coefficients (cf. Theorem 7.1) we first compute
$$A(L(\frac{3}{2},h_{1,q}))_{A(L(3/2,0)}L(\frac{3}{2},h_{1,r})(0).$$
$`/2`$–grading of the $`0`$–th homology group (39) enables us (see Theorem (10.1) to study odd and even intertwining operators (see Definition 7.1). For that purpose we introduce the following splitting:
$`A^0(L({\displaystyle \frac{3}{2}},h_{1,q})):=H_0^0(_s,L({\displaystyle \frac{3}{2}},h_{1,q})){\displaystyle \frac{𝐂[x,y]}{I_1}}`$
$`A^1(L({\displaystyle \frac{3}{2}},h_{1,q})):=H_0^1(_s,L({\displaystyle \frac{3}{2}},h_{1,q})){\displaystyle \frac{𝐂[x,y]v}{I_2}},`$ (46)
where $`I_1`$ and $`I_2`$ are cyclic submodules (the maximal submodule for $`M(\frac{3}{2},h_{1,q})`$ is cyclic !). It seems hard to obtain explicitly these polynomials. First we obtain some useful formulas Inside $`A(M(c,h))`$ (cf. \[W\]):
$`[L(n)v]=[((n1)(L(2)L(1))+L(1))v]=`$
$`[(n(L(2)L(1))(L(2)2L(1)+L(0))+L(0))v]=`$
$`(nyx+\mathrm{wt}(v))[v].`$ (47)
for every $`n𝐍`$ and every homogeneous $`vM(c,h)`$. Therefore in
$$A(M(\frac{3}{2},h_{1,q}))_{A(L(\frac{3}{2},0))}L(\frac{3}{2},h_{1,r})(0)$$
we have
$`[L(n)v]=(nh_{1,q}x+L(0))[v].`$
$`[G(n1/2)v]=[G(1/2)v].`$ (48)
Also, we have:
$`[G(n{\displaystyle \frac{1}{2}})G(m{\displaystyle \frac{1}{2}})v]=[G(1/2)G(m1/2)v]=`$
$`[(2L(m1)G(m1/2)G(1/2))v]=[(2L(m1)L(1))v]=`$
$`((2m+1)yx+\mathrm{wt}(v))[v].`$ (49)
By using (10) and (10) we obtain
$`[G(m_11/2)\mathrm{}G(m_{2r}1/2)L(n_1)\mathrm{}L(n_s)v_{1,q}]=`$
$`{\displaystyle \underset{i=1}{\overset{r}{}}}((2m_{2i}+1)h_{1,r}x+{\displaystyle \underset{p=2i+1}{\overset{2r}{}}}(m_p+1/2)+h_{1,q})`$
$`{\displaystyle \underset{j=1}{\overset{s}{}}}(n_jh_{1,r}x+{\displaystyle \underset{p=j+1}{\overset{s}{}}}n_p+h_{1,q})[v].`$ (50)
inside
$$A(M(\frac{3}{2},h_{1,q}))_{A(L(\frac{3}{2},0))}L(\frac{3}{2},h_{1,r})(0).$$
It is easy to obtain a similar formula for the vector
$$[G(m_11/2)\mathrm{}G(m_{2r+1}1/2)L(n_1)\mathrm{}L(n_s)v_{1,q}].$$
###### Lemma 10.1
Let $`[P_{\mathrm{sing}}v_{1,q}]=Q_1(x)[G(1/2)v_{1,q}]`$ and $`[G(1/2)P_{\mathrm{sing}}v_{1,q}]=Q_2(x)[v_{1,q}]`$ be the projections inside
$$A(M(\frac{3}{2},h_{1,q}))_{A(L(\frac{3}{2},0)}L(\frac{3}{2},h_{1,r})(0).$$
Then
$`Q_1(h)=K_2(h_{1,q},h_{1,r},h),`$
$`Q_2(h)=K_1(h_{1,q},h_{1,r},h),`$ (51)
for every $`h`$.
Proof: We use the notation from the section 6.2, where
$$𝒴I\left(\genfrac{}{}{0pt}{}{L(3/2,h)}{L(3/2,h_{1,r})L(3/2,h_{1,q})}\right).$$
By using (7.1), we obtain
$`w_3^{},𝒴(w_1,x)G(m_11/2)\mathrm{}G(m_{2r}1/2)L(n_1)\mathrm{}L(n_s)w_2=`$
$`{\displaystyle \underset{i=1}{\overset{r}{}}}(x^{m_{2i1}m_{2i}}{\displaystyle \frac{}{x}}2m_{2i}h_{1,r}x^{m_{2i1}m_{2i}1})`$
$`{\displaystyle \underset{j=1}{\overset{s}{}}}(x^{n_j+1}{\displaystyle \frac{}{x}}+(1n_j)h_{1,r}x^{n_j})w_3^{},𝒴(w_1,x)w_2=`$
$`c_1{\displaystyle \underset{i=1}{\overset{r}{}}}((2m_{2i}+1)h_{1,r}h+h_{1,q}+{\displaystyle \underset{p=2i+1}{\overset{2r}{}}}(m_p+1/2))`$
$`{\displaystyle \underset{j=1}{\overset{s}{}}}(n_jh_{1,r}h+{\displaystyle \underset{p=j+1}{\overset{s}{}}}n_p+h_{1,q})x^{hh_{1,q}h_{1,r}r{\scriptscriptstyle m_i}_jn_j},`$ (52)
for the constant $`c_1`$ (see Section 6.1 and 6.2) that depends only on $`𝒴`$. There is a similar expression for
$$w_3^{},𝒴(w_1,x)G(m_11/2)\mathrm{}G(m_{2r+1}1/2)L(n_1)\mathrm{}L(n_s)w_2.$$
(53)
If we compare (10) with (10) (and corresponding formulas for (53)) it follows that $`Q_1(h)`$ is, up to a non–zero multiplicative constant, equal to $`K_2(h_{1,r},h_{1,q},h)`$ (singular vector is odd!) and $`Q_2(h)`$ is, up to a multiplicative constant, equal to $`K_1(h_{1,r},h_{1,q},h)`$.
Thus, Proposition 9.1 and Theorem 10.1 gives us
###### Theorem 10.1
* As a $`A(L(3/2,0))`$–module
$`A(L(3/2,h_{1,q}))_{A(L(3/2,0))}L(3/2,h_{1,r})(0)`$ (54)
$`{\displaystyle \frac{𝐂[x]}{<_{jkj}(xh_{1,q+4k})>}}{\displaystyle \frac{𝐂[x]}{<_{j+1/2kj+1/2}(h+1/2h_{1,q+4k})>}}.`$
* The space
$$I\left(\genfrac{}{}{0pt}{}{M(3/2,h)^{}}{L(3/2,h_{1,q})M(3/2,h_{1,r})}\right),$$
is non–trivial if and only if $`h=h_{1,s}`$ for some $`s\{q+r1,q+r3,\mathrm{},qr+1\}`$.
* The space
$$I\left(\genfrac{}{}{0pt}{}{L(3/2,h)}{L(3/2,h_{1,q})L(3/2,h_{1,r})}\right),$$
is one–dimensional if and only if $`h=h_{1,s}`$, $`s\{q+r1,q+r3,\mathrm{},|qr|+1\}.`$
Proof (a): From Lemma 10.1 it follows that
$`A(L(3/2,h_{1,r}))_{A(L(3/2,0))}L(3/2,h_{1,q}){\displaystyle \frac{[x]}{Q_1(x)}}{\displaystyle \frac{[x]}{Q_2(x)}}.`$ (55)
Now we apply (9.2) and Proposition 9.1.
Proof (b): As in the Virasoro case, by examining carefully the main construction of intertwining operators in \[L1\] with a minor super–modifications, for every $`A(L(3/2,0))`$–morphism from $`A(L(3/2,h_{1,q}))_{A(L(3/2,0))}L(3/2,h_{1,r})`$ to $`L(3/2,h)(0)`$ we can construct a non–trivial intertwining operator of the form $`I\left(\genfrac{}{}{0pt}{}{M(3/2,h)^{}}{L(3/2,h_{1,q})M(3/2,h_{1,r})}\right).`$
Proof (c): The proof and all the arguments involved are the same as in the Chapter 3. so we omit the detials. We obtain a non–trivial intertwining operator of the type $`\left(\genfrac{}{}{0pt}{}{L(3/2,h)}{L(3/2,h_{1,q})L(3/2,h_{1,r})}\right)`$ if $`h=h_{1,s}`$ for $`s\{q+r1,q+r3,\mathrm{},qr+1\}\{r+q1,r+q3,\mathrm{},rq+1\}`$, i.e., $`s\{q+r1,q+r3,\mathrm{},|qr|+1\}`$.
###### Theorem 10.2
Suppose that $`qr`$ <sup>4</sup><sup>4</sup>4$`\left(\genfrac{}{}{0pt}{}{W_3}{W_1W_2}\right)\left(\genfrac{}{}{0pt}{}{W_3}{W_2W_1}\right)`$.
$$\mathrm{dim}I\left(\genfrac{}{}{0pt}{}{L(3/2,h_{1,s})}{L(3/2,h_{1,q})L(3/2,h_{1,r})}\right)_{\mathrm{even}}=1,$$
(56)
if and only if
$$s\{q+r1,q+r5,\mathrm{},qr+1\}$$
$$\mathrm{dim}I\left(\genfrac{}{}{0pt}{}{L(3/2,h_{1,s})}{L(3/2,h_{1,q})L(3/2,h_{1,r})}\right)_{\mathrm{odd}}=1$$
(57)
if and only if
$$s\{q+r3,q+r7,\mathrm{},qr+3\}.$$
Proof: By using (54) we obtain the following decomposition:
$`A^0(L({\displaystyle \frac{3}{2}},h_{1,q}))_{A(L(3/2,0))}L({\displaystyle \frac{3}{2}},h_{1,r})(0)`$
$`v_{q+r1}v_{q+r5}\mathrm{}v_{qr+1}`$
$`A^1L({\displaystyle \frac{3}{2}},h_{1,q}))_{A(L(3/2,0))}L({\displaystyle \frac{3}{2}},h_{1,r})(0)`$
$`v_{q+r3}v_{q+r7}\mathrm{}v_{qr+3},`$ (58)
where $`v_i`$ is a $`[y]`$–module such that $`y.v_i=\frac{(i1)^2}{8}v_i`$.
Claim: Let
$$\psi \mathrm{Hom}_{A(L(c,0))}(A^0(L(\frac{3}{2},h_{1,q}))_{A(L(3/2,0))}L(\frac{3}{2},h_{1,r})(0),L(\frac{3}{2},h_{1,s})(0)),$$
then the corresponding intertwining operator is even. Similarly if we start from
$$\psi \mathrm{Hom}_{A(L(c,0))}(A^1(L(\frac{3}{2},h_{1,q}))_{A(L(3/2,0))}L(\frac{3}{2},h_{1,r})(0),L(\frac{3}{2},h_{1,s})(0)),$$
the corresponding intertwining operator is odd.
Proof (of the Claim): Let us elaborate the proof when $`\psi `$ is “even”. From the construction in \[FZ\] and \[L2\] $`𝒴`$ is obtained by lifting $`\psi `$ to a mapping from $`L(3/2,h_{1,q})L(3/2,h_{1,r})(0)`$ to $`L(3/2,h_{1,s})(0)`$, such that
$$L(3/2,h_{1,q})_{\mathrm{odd}}L(3/2,h_{1,r})(0)0.$$
To extend this map to a mapping from $`L(3/2,h_{1,q})M(3/2,h_{1,r})`$ to $`M(3/2,h_{1,s})^{}`$ one uses generators and PBW so the sign is preserved. Because the isomorphism $`I\left(\genfrac{}{}{0pt}{}{W_3}{W_1W_2}\right)I\left(\genfrac{}{}{0pt}{}{W_2^{}}{W_1W_3^{}}\right)`$ preserves the sign, i.e., odd intertwining operators are mapped into odd and even into even, the result follows from the construction of intertwining operators. When $`\psi `$ is odd a similar argument works.
Let us summarize everything.
###### Corollary 10.1
Let $`𝒜_s`$ be a free abelian group with generators $`b(m),m2𝐍+1`$. Define a binary operation $`\times :𝒜_s\times 𝒜_s𝒜_s`$,
$$b(q)\times b(r)=\underset{j𝐍}{}\mathrm{dim}I\left(\genfrac{}{}{0pt}{}{L(3/2,h_{1,j})}{L(3/2,h_{1,q})L(3/2,h_{1,r})}\right)b(j).$$
Then $`𝒜_s`$ is a commutative associative ring, and the mapping $`b(m)V(\frac{m1}{4})`$ gives an isomorphism to the representation ring $`ep(𝔬sp(1|2))`$.
Proof: The proof follows from Theorem 10.1(c) and (31).
## 11 Multiplicity $`2`$ fusion rules and super logarithmic intertwiners
### 11.1 A multiplicity $`2`$ case
We have seen that in the $`c=\frac{3}{2}`$ case all fusion coefficients are $`0`$ or $`1`$. Still, we expect (according to \[HM\]) that for some vertex operator superalgebras $`L(c,0)`$, fusion coefficients are $`2`$.
Here is one example. If $`c=0`$, as in the case of the Virasoro algebra, the super vertex operator algebra $`L(0,0)=\frac{M(0,0)}{G(1/2)v_0,G(3/2)v_0}`$ is trivial. Still we can consider a vertex operator superalgebra $`V(0,0):=\frac{M(0,0)}{G(1/2)v}`$ Clearly, for every $`h`$, we have (all modules are considered to be $`V(0,0)`$–modules):
$$\mathrm{dim}I\left(\genfrac{}{}{0pt}{}{L(0,0)}{L(0,h)L(0,h)}\right)=2.$$
(59)
The previous example is little bit awkward. Here is a nice example with irrational central charge:
###### Proposition 11.1
$$\mathrm{dim}I\left(\genfrac{}{}{0pt}{}{L(\frac{15}{2}3\sqrt{5},\frac{\sqrt{5}}{2}1)}{L(\frac{15}{2}3\sqrt{5},\frac{3}{4}(\frac{\sqrt{5}}{2}1))L(\frac{15}{2}3\sqrt{5},\frac{3}{4}(\frac{\sqrt{5}}{2}1))}\right)=2.$$
(60)
Proof: It is not hard to see (by using a result form \[AA\] or \[D\]) that $`M(\frac{15}{2}3\sqrt{5},\frac{3}{4}(\frac{\sqrt{5}}{2}1))`$ has the unique submodule that is irreducible (the case $`II_+`$ in \[AA\]). If we analyze the determinant formula \[KWa\], singular vectors, and then use Theorem 9.1, we obtain (60).
### 11.2 A logarithmic intertwiner
In \[M\] we studied several examples of logarithmic intertwining operators. Roughly, logarithmic intertwiners exist if matrix coefficients yield some logarithmic solutions. Our analysis can be extended for vertex operator superalgebras.
$$\mathrm{dim}I\left(\genfrac{}{}{0pt}{}{W_2(\frac{27}{2},\frac{3}{2})}{L(\frac{27}{2},\frac{3}{2})L(\frac{27}{2},\frac{3}{2})}\right)=2,$$
(61)
where $`W_2(\frac{27}{2},\frac{3}{2})`$ is certain logarithmic module (cf. \[M\]). The proof of this result and the discussion will appear in a separate publication.
## 12 Future work and open problems
* We know that it is possible to obtain intertwining operator algebras (see \[H2\]) from the rational vertex operator algebras (satisfying some natural convergence and extension condition and an additional condition involving generalized modules). Since the notation of intertwining operator algebra can be (obviously) generalized such that fusion algebra is an infinite–dimensional associative, commutative algebra, one hopes that it is possible to construct tensor categories for degenerate minimal models. In the language of conformal field theory this involves explicit calculations of correlation functions for both products and iterates of intertwining operators (cf. Remark 4.2).
* Open problem: For rational vertex operator algebras, construct a canonical isomorphism
$$A(M_1)_{A(V)}M_2(0)A(M_2)_{A(V)}M_1(0).$$
* (N=1 case) For which triples $`L(c,h_1)`$, $`L(c,h_2)`$ and $`L(c,h_3)`$ do we have
$$\mathrm{dim}I\left(\genfrac{}{}{0pt}{}{L(c,h_3)}{L(c,h_1)L(c,h_2)}\right)=2\mathrm{?}$$
* Determine the fusion ring for degenerate minimal models for $`N=2`$ superconformal algebra by using our method (it should be related to $`\mathrm{Rep}(𝔬𝔰𝔭(2|2))`$.
* Construct an analogue of the vertex tensor categories constructed in \[HM\] (by using the main result in \[A\]), for the models studied in this paper.
Department of Mathematics, Rutgers University, 110 Frelinghuysen Rd., Piscataway, NJ 08854-8019
Current Address: Department of Mathematics, University of Arizona, Tucson, AZ 85721
E-mail address: amilas@math.rutgers.edu, milas@math.arizona.edu |
warning/0003/hep-th0003060.html | ar5iv | text | # 1 Kalb-Ramond case
## 1 Kalb-Ramond case
The Ansatz for the metric is taken to be
$$ds^2=N^2(\sigma )d\sigma ^2\eta ^2(\sigma )\overline{g}_{ij}dx^idx^jR^2(\sigma )h_{mn}dy^mdy^n,$$
(4)
where $`\overline{g}_{ij}`$ is a metric of constant curvature, i.e. $`R(\overline{g})_{ij}\lambda \overline{g}_{ij}`$, and $`h_{mn}`$ is a Ricci flat metric<sup>1</sup><sup>1</sup>1Most of the equations in what follows can easily be generalized to include a curvature for $`h_{mn}`$. on the Kähler manifold. Furthermore, the dimension of the Kähler manifold is $`d_k`$ and the dimension of the constant curvature part is $`Dd1d_k`$, where $`d`$ is the dimension of spacetime, which will be left arbitrary. Note that it is natural to take the internal space to be compact, which means that it actually is a Calabi-Yau space. For the sake of argument though, we will always refer to the internal space as a Kähler space, since this is the only property that is really needed.
Our Ansatz for the Kalb-Ramond field, in form notation, reads
$$B=f(\sigma )𝒥=\frac{1}{2}f(\sigma )𝒥_{mn}dy^mdy^n,$$
(5)
where $`𝒥`$ is the integrable almost complex structure, the Kähler form, on the Kähler manifold and as such satisfies
$$𝒥_{}^{m}{}_{p}{}^{}𝒥_{}^{p}{}_{n}{}^{}=\delta _{}^{m}{}_{n}{}^{},\overline{}_m𝒥_{np}=\mathrm{\hspace{0.33em}0},$$
(6)
where $`\overline{}`$ is the connection on the Kähler manifold. Finally, the dilaton is taken to depend on time, $`\sigma `$, only.
Calculating the equation of motion for the Kalb-Ramond field, one finds that
$$\dot{f}=\mathrm{}Ne^{2\varphi }\eta ^DR^{4d_k},$$
(7)
where $`\mathrm{}`$ is an arbitrary constant. Defining the field
$$\psi =\varphi +\frac{1}{2}\mathrm{log}\left(N\right)\frac{D}{2}\mathrm{log}\left(\eta \right)\frac{d_k}{2}\mathrm{log}\left(R\right),$$
(8)
introducing $`M=e^{2\psi }N`$ and changing variables by $`e^{2\psi }d\sigma =dt`$, one can write the equations of motion as
$`0`$ $`=`$ $`\left(\mathrm{log}R\right)^{\prime \prime }+\frac{\mathrm{}^2}{2}R^4,`$ (9)
$`0`$ $`=`$ $`\left(\mathrm{log}\eta \right)^{\prime \prime }\lambda \eta ^2M^2,`$ (10)
$`0`$ $`=`$ $`\left(\mathrm{log}M\right)^{\prime \prime }D\lambda \eta ^2M^2,`$ (11)
$`0`$ $`=`$ $`\left[\left(\mathrm{log}M\right)^{}\right]^2D\left[\left(\mathrm{log}\eta \right)^{}\right]^2d_k\left[\left(\mathrm{log}R\right)^{}\right]^2\frac{d_k\mathrm{}^2}{4}R^4,`$ (12)
where the prime indicates derivation with respect to $`t`$.
The above equations can easily be solved to give $`\lambda =0`$, $`M=e^{\alpha t}`$, $`\eta =e^{\beta t}`$ and
$$R(t)=R_0\mathrm{cosh}^{1/2}\left(\mathrm{}R_0^2t\right),$$
(13)
which after substitution in Eq. (12) leads to
$$\alpha ^2D\beta ^2=\frac{d_k\mathrm{}^2R_0^4}{4}.$$
(14)
In these coordinates the dilaton, or equivalently the string coupling squared, is $`g_s^2=e^{2\varphi }=M^1\eta ^DR^{d_k}`$ and the Kalb-Ramond field strength is $`H=\mathrm{}R^4dt𝒥`$.
Although there is a world of possibilities in the above class of solutions, perhaps the most interesting case is the one where $`\beta =0`$, since then the uncompactified part of spacetime is just Minkowski space: Taking $`\alpha =1`$ for convenience, one finds that the solution in, string, cosmological time, $`\tau `$, reads
$`ds^2`$ $`=`$ $`d\tau ^2d\stackrel{}{x}_{(D)}\mathrm{\hspace{0.17em}2}R_0^2B(\tau )^1h_{mn}dy^mdy^n,`$
$`e^{2\varphi }`$ $`=`$ $`\left(\sqrt{2}R_0\right)^{d_k}\tau ^1B(\tau )^{d_k/2},`$
$`H`$ $`=`$ $`8R_0^2d_k^{1/2}\tau ^1B(\tau )^2d\tau 𝒥,`$ (15)
$`B(\tau )`$ $`=`$ $`\tau ^{2/\sqrt{d_k}}+\tau ^{2/\sqrt{d_k}}.`$ (16)
As one can see, this is a completely regular solution, modulo the usual gravitational singularities, which smoothly interpolates between two Kasner-like regions . From the lower dimensional point of view, the Ansatz considered above corresponds to a solution of dilaton-gravity coupled to moduli , where the breathing mode, $`R`$, and $`f`$ are the scalar fields parameterizing an $`SL(2,)/U(1)`$ coset model. When $`d_k=6`$, the above solutions can be obtained from the solutions given in by applying an $`SL(2,)`$ transformation on the moduli.
## 2 RR case
In much the same way as in the foregoing subsection, we can use the RR two form in type IIA, to trigger compactification. In this case the equations of motion and the Bianchi identity imply that
$$F_{(2)}=\mathrm{}𝒥=\frac{1}{2}\mathrm{}𝒥_{mn}dy^mdy^n.$$
(17)
Applying the same steps as in the foregoing paragraph, one finds
$`0`$ $`=`$ $`\left(\mathrm{log}R\right)^{\prime \prime }+\frac{(d_k4)\mathrm{}^2}{8}M\eta ^DR^{d_k4},`$ (18)
$`0`$ $`=`$ $`\left(\mathrm{log}\eta \right)^{\prime \prime }+\frac{d_k\mathrm{}^2}{8}M\eta ^DR^{d_k4}\lambda \eta ^2M^2,`$ (19)
$`0`$ $`=`$ $`\left(\mathrm{log}M\right)^{\prime \prime }\frac{d_k\mathrm{}^2}{8}M\eta ^DR^{d_k4}D\lambda \eta ^2M^2,`$ (20)
$`0`$ $`=`$ $`\left[\left(\mathrm{log}M\right)^{}\right]^2D\left[\left(\mathrm{log}\eta \right)^{}\right]^2d_k\left[\left(\mathrm{log}R\right)^{}\right]^2\frac{d_k\mathrm{}^2}{4}M\eta ^DR^{d_k4}D\lambda M^2\eta ^2,`$ (21)
Looking at the above expressions, one sees that they simplify enormously when one considers the case $`d_k=4`$: In that case the Kähler breathing mode decouples completely and one has $`R=R_0e^{\alpha t}`$. Equating also the powers of $`M`$ and $`\eta `$ in the equations, i.e. putting $`M=\eta ^{D+2}`$, one necessarily has to impose
$$\lambda =\frac{\mathrm{}^2}{4}\left(D+3\right).$$
(22)
The two remaining equations are implied by
$$\left(\mathrm{log}\eta \right)^{}=\pm \left[\frac{\mathrm{}^2}{4}\eta ^{2(D+1)}+\frac{4\alpha ^2}{D(D+3)+4}\right]^{1/2}.$$
(23)
Now, when $`\alpha 0`$ the solution to the above equation is complex, but when $`\alpha =0`$ one finds that
$$\eta =\left(A\frac{(D+1)\mathrm{}^2}{2}t\right)^{\frac{1}{D+1}},$$
(24)
where the range of $`t`$ has to be chosen such that the function is well defined.
For a stringy cosmological observer, i.e. introducing a time coordinate $`Mdt=d\tau `$, the above solution is
$`ds^2`$ $`=`$ $`d\tau ^2\left({\displaystyle \frac{\mathrm{}^2}{2}}\tau \right)^2d\mathrm{\Omega }_\lambda ^2R_0^2h_{mn}dy^mdy^n,`$ (25)
$`e^\varphi `$ $`=`$ $`4R_0^2\mathrm{}^2\tau ^1.`$ (26)
So, we see that the solution describes a 6-dimensional open FRW with a dilaton such that the string coupling strength goes to zero when $`\tau \mathrm{}`$.
In order to find a solution for general $`d_k`$, we shall follow the same stratagem as above, i.e. we put $`R=\eta ^\alpha `$, $`M=\eta ^\beta `$ and will equate the powers on the righthand sides of Eqs. (21). It then follows that $`\beta =D+2+\alpha (d_k4)`$, using which one can calculate
$`\lambda `$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2}{16}}\left(d_k(D+3)+(d_k4)^2\right),`$ (27)
$`\alpha `$ $`=`$ $`2{\displaystyle \frac{d_k4}{d_k(D+1)+(d_k4)^2}},`$ (28)
$`\beta `$ $`=`$ $`{\displaystyle \frac{d_k(D+1)(D+2)+D(d_k4)^2}{d_k(D+1)+(d_k4)^2}},`$ (29)
The dilaton is then fixed to $`e^\varphi =\eta ^{2\alpha 1}`$.
As before, the two remaining equations are implied by
$$\left(\mathrm{log}\eta \right)^{}=\pm B\eta ^{\beta 1},$$
(30)
where
$$B^2=\frac{\mathrm{}^2}{16}\frac{\left[d_k(D+1)+(d_k4)^2\right]^2}{d_k(D+1)+(D1)(d_k4)^2}.$$
(31)
This then means that
$$\eta =\left[A(\beta 1)Bt\right]^{\frac{1}{\beta 1}},$$
(32)
where $`A`$ is some integration constant and the range of $`t`$ has to be chosen such that the above function is well-defined.
Choosing the minus-sign in the last equation in order to switch to the cosmological time, $`\tau `$, one finds that $`\eta =B\tau `$ and the solution reads
$`ds^2`$ $`=`$ $`d\tau ^2B^2\tau ^2d\mathrm{\Omega }_\lambda ^2\left(B\tau \right)^{2\alpha }h_{mn}dy^mdy^n,`$ (33)
$`e^\varphi `$ $`=`$ $`\left(B\tau \right)^{2\alpha 1}.`$ (34)
Seeing Eq. (28), one may wonder whether the breathing mode for the Kähler mode could grow faster than $`\tau `$, and thus spoil cosmological compactification. This can however happen for $`d_k=2`$ only, in which case $`\alpha =2d^1`$ which, seeing that we at least must have $`d=3`$ in order to apply the Ansatz, is always smaller than 1. A similar analysis for the dilaton shows that it can blow up at large times only when $`d_k=2`$, $`d=3`$ and is regular in the rest of the cases.
## 3 Conclusions
Although more general solutions are bound to exist, the simple examples discussed in this work show that using a Kähler form instead of a volume form is a viable option when looking for stringy cosmological solutions.
It would be interesting to see whether this kind of compactifications survive when coupling to matter is enabled and whether the dimension of the Kähler manifold will be of any importance. It would also be interesting to generalize the above work to Heterotic strings/M-theory (See e.g. ) and to investigate cosmological solutions when other type II fields are present. Work in these directions is in progress.
## Acknowledgments
The authors would like to thank E. Álvarez, T. Ortín and M.A. Vázquez-Mozo for useful discussions and the referee for very useful remarks. PM would like to thank Iberdrola and the Universidad Autónoma de Madrid for their support. |
warning/0003/quant-ph0003042.html | ar5iv | text | # Formation of ultrashort pulses with sub-Poissonian photon statistics
\[
## Abstract
A simple method for the production of ultrashort light pulses (USPs) with suppressed photon fluctuations is considered. The method is based on self-phase modulation (SPM) of an USP in a nonlinear medium (optical fibre) and subsequent transmission of pulse through a dispersive optical element.
Enlarged version of the article published in *Quantum Electronics*, 29 No. 7, 61-63 (1999)
\] Ultrashort light pulses (USPs) continue to draw the attention of investigators today. The state of the art in this area of laser physics in the late 1980s was set forth by S. A. Akhmanov *et al*. in Refs.,. In the past decade, considerable progress has been made in quantum optics in the generation of nonclassical (the so called squeezed) light fields. The formation and application of USPs in a nonclassical state makes possible to combine in experiments a high time resolution with a low level of fluctuations.
In principle, one can obtain pulsed light fields in a nonclassical state by using the same nonlinear optical interaction as those used in the case of continuous fields ,. Parametric amplification is a technique that is most extensively used for this purpose nowadays. In the case of degenerate three-frequency parametric amplification, quadrature-squeezed light is produced. However, this light is found to have super-Poissonian statistics directly at the output of the amplifier, and one needs interferometers to transform it to get sub-Poissonian statistics. One can obtain light with sub-Poissonian photon statistics with the aid of nonlinear interferometric devices in the presence of self-phase modulation (see, e.g., Ref.). Note that the self-phase modulation itself is not accompanied by a change in photon statistics.
In a recent paper, we studied SPM of a light pulse and its subsequent propagation in a dispersive linear medium (or the passage through optical compressors). In this case, pulses with sub-Poissonian photon statistics can be formed. We managed to make an accurate calculation of the process under consideration owing to the consistent quantum theory developed by us for the self-action of light pulses in a medium with inertial nonlinearity ,. The method of formation of pulses with sub-Poissonian statistics proposed here is simple for experimental realization and stable against external random effects (technical fluctuations). Note that the theory developed in Refs., for the self-effect of USPs takes into account a relaxation time of a nonlinearity that determines the region of the spectrum of quantum fluctuations which are of substantial importance in the formation of squeezed light and does not limit the amplitude of quantum fluctuation.
When analysed from the quantum point of view, SPM of USPs is described by the expression ,.
$$\widehat{A}(t,l)=e^{\widehat{O}(t)}\widehat{A}_0(t),$$
(1)
and the Hermite conjugate operator $`\widehat{A}^+(t,l)`$. Here $`\widehat{A}^+(t,l)`$ ($`\widehat{A}(t,l)`$) is the photon creation (annihilation) operator for the given section $`x`$ at a given moment of time $`t`$ (the output of a nonlinear medium is specified by the section $`x=l`$, and the input is specified by the section $`x=0`$, $`\widehat{A}(t,x=0)=\widehat{A}_0(t)`$), $`\widehat{O}(t)=i\gamma q[\widehat{n}_0(t)]`$, and $`\gamma `$ is the coefficient related to nonlinear properties of a medium and proportional to the length $`l`$ in it. The expression $`\gamma q[\widehat{n}_0(t)]`$ characterizes the nonlinear phase incursion,
$$q[\widehat{n}_0(t)]=\underset{\mathrm{}}{\overset{\mathrm{}}{}}H(|t_1|)\widehat{n}_0(tt_1)𝑑t_1,$$
(2)
and $`\widehat{n}_0(t)=\widehat{A}_0^+(t)\widehat{A}_0(t)`$ is the operator of the ”density” of the number of photons at the input of a nonlinear medium. The function $`H(t)`$ takes into account a finite response time of the nonlinearity of a medium. In view of the causality principle, $`H(t)0`$ for $`t0`$ and $`H(t)=0`$ for $`t<0`$. In expression (1) and (2), $`t=t^{^{}}x/u`$ is the time in the moving coordinate system, $`t^{^{}}`$ is the current time, and $`u`$ is the pulse velocity in a medium.
In the case of SPM (see (1)) discussed here, the quadrature-squeezed light , is formed, and the photon statistics in a medium is not changed because the operator of the number of photons $`\widehat{n}(t,l)=\widehat{n}_0(t)`$ remains unchanged. In what follows, we show that the propagation of a pulse (see (1)) through a dispersive optical device is able to transform its initial Poissonian statistics into sub- or super-Poissonian statistics.
For example, let us consider the propagation of a pulse in a dispersive linear medium. At the output of the medium, we have the following expression for the photon annihilation operator :
$$\widehat{B}(t,z)=\underset{\mathrm{}}{\overset{\mathrm{}}{}}G(tt_1,z)\widehat{A}(t_1,l)𝑑t_1.$$
(3)
It is known that the operators $`\widehat{A}(t,l)`$ and $`\widehat{B}(t,z)`$ are bound to satisfy commutation relation of the $`[\widehat{C}(t_1),\widehat{C}^+(t_2)]=\delta (t_1t_2)`$ type. In view of relation (3) between the operators under consideration, we find the condition imposed on the Green function $`G(t)`$ of the dispersive element:
$$\underset{\mathrm{}}{\overset{\mathrm{}}{}}G(t_1t,z)G^{}(t_2t,z)𝑑t=\delta (t_2t_1).$$
(4)
Let us introduce the operator of the number of photons over the measurement time $`𝒯`$,
$$\widehat{N}_𝒯(t,z)=\underset{t𝒯/2}{\overset{t+𝒯/2}{}}\widehat{N}(t_1,z)𝑑t_1,$$
(5)
and define the Mandel parameter
$$Q(t,z)=\frac{\epsilon (t,z)}{\widehat{N}_𝒯(t,z)},$$
(6)
where
$`\widehat{N}(t,z)`$ $`=`$ $`\widehat{B}^+(t,z)\widehat{B}(t,z);`$ (7)
$`\epsilon (t,z)`$ $`=`$ $`\widehat{N}_𝒯^2(t,z)\widehat{N}_𝒯(t,z)^2\widehat{N}_𝒯(t,z).`$ (8)
The parameter $`Q(t,z)`$ characterizes the difference between photon statistics and Poissonian statistics (for the latter $`Q(t,z)=0`$). The angle brackets in (8) denote averaging over the initial quantum state of the pulse.
One can calculate expression (8) by using the algebra of time dependent Bose operators ,. In this case we have the following relations:
$$\widehat{A}_0(t_1)e^{\widehat{O}(t_2)}=e^{\widehat{O}(t_2)+(t_2t_1)}\widehat{A}_0(t_1),$$
(9)
$$e^{\widehat{O}(t)}=\widehat{𝐍}\mathrm{exp}\{\underset{\mathrm{}}{\overset{\mathrm{}}{}}[e^{(t_1)}1]\widehat{n}_0(tt_1)dt_1\},$$
(10)
where $`(t)=i\gamma h(t)`$, $`h(t)=H(|t|)`$ and $`\widehat{𝐍}`$ is the operator of normal ordering.
In the case of a coherent initial pulse, a small nonlinear phase incursion per photon, and the pulse duration $`\tau _p`$ is considerably grater than the times $`𝒯`$ and $`\tau _r`$ ($`\tau _r`$ is the nonlinearity relaxation time). After some cumbersome calculations, taking into account expression (9) and (10) one gets
$$\epsilon (t,z)=2\gamma 𝒯^2Im[I_1^2(t,z)I_2^{}(t,z)],$$
(11)
where
$`I_1(t,z)`$ $`=`$ $`{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}G(tt_1,z)\alpha (t_1)e^{i\psi (t_1)}𝑑t_1;`$ (12)
$`I_2(t,z)`$ $`=`$ $`{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}G(tt_1,z)G(tt_2,z)\alpha (t_1)\alpha (t_2)`$ (13)
$`\times `$ $`h(t_1t_2)e^{i[\psi (t_1)+\psi (t_2)]}dt_1dt_2;`$ (14)
Here $`\alpha (t)`$ is the eigenvalue of the operator $`\widehat{A}_0(t)`$; $`\psi (t)=2\gamma \overline{n}_0(t)`$ is the nonlinear phase incursion, $`\overline{n}_0(t)=|\alpha (t)|^2`$ is the average density of the initial number of photons, and $`\widehat{N}_𝒯(t,z)=𝒯|I_1(t,z)|^2`$.
Let us assume that the Green function G(t) has the form
$$G(t,z)=(i2\pi k_2z)^{1/2}\mathrm{exp}\left(\frac{it^2}{2k_2z}\right),$$
(15)
where $`z`$ is the distance traveled in the dispersive linear medium. We have the coefficients $`k_2>0`$ for normal dispersion of the group velocity, and $`k_2<0`$ for anomalous dispersion. Expression (15) corresponds to the case when dispersive properties of a medium are taken into account in the second approximation of the dispersion theory .
The dispersion effect is known to have a spatial analogue in the form of diffraction spreading of a wave beam. The quantum description of diffraction in nonlinear optical processes was considered first by S. A. Akhmanov *et al*. in Ref. (see also Ref.), where it was shown that the mixing of different angular components of a parametrical amplified beam is able to cause additional nonclassical effects. One can expect similar effects in the case under consideration. The self-interaction effect of a pulse given by expression (1) is accompanied by changes in its phase (and therefore frequency) over time. In the course of its passage through dispersive linear elements, a frequency-modulated pulse of this type is compressed or stretched . It is precisely this effect that is able to change its photon statistics.
Let us determine the parameter $`Q(t,z)`$. To simplify calculations, we take the nonlinearity response function in the form $`h(t)=(1/\tau _r)\mathrm{exp}(t^2/\tau _r^2)`$. Let a pulse be initially of the Gaussian form $`\overline{n}_0(t)=\overline{n}_0\mathrm{exp}\left(t^2/\tau _p^2\right)`$. In the paraxial approximation, i.e., in the case where the phase $`\psi (t)`$ is replaced with $`\psi _0(1t^2/\tau _p^2)`$, ($`\psi _0=2\gamma \overline{n}_0`$), we obtain for $`\tau _p>\tau _r`$:
$$\widehat{N}_𝒯(t,z)=\overline{n}_0𝒯V^1(z)\mathrm{exp}\left(\frac{t^2}{V^2(z)\tau _p^2}\right),$$
(16)
$`Q(0,z)=`$ $`\psi _0`$ $`\left[{\displaystyle \frac{𝒯}{\tau _p}}\right]\mathrm{sin}\{\mathrm{arctan}\left({\displaystyle \frac{\phi (z)}{\varpi (z)}}\right)`$ (17)
$`+`$ $`{\displaystyle \frac{1}{2}}\mathrm{arctan}({\displaystyle \frac{2\phi _d(z)\varpi (z)}{2\phi (z)\phi _d(z)\varpi ^2(z)}}\left)\right\}`$ (18)
$`/`$ $`\left[\varpi ^4(z)2\phi ^2(z)\varpi ^2(z)+4\phi ^4(z)\right]^{1/4},`$ (19)
where
$`V^2(z)=\varpi ^2(z)+\phi ^2(z);`$ $`\varpi (z)=1s\psi _0\phi (z);`$ (20)
$`\phi (z)={\displaystyle \frac{z}{D}};\phi _d(z)={\displaystyle \frac{z}{d}};`$ $`D={\displaystyle \frac{\tau _p^2}{|k_2|}};d={\displaystyle \frac{\tau _r^2}{|k_2|}};`$ (21)
$`D`$ and $`d`$ are the characteristic dispersion lengths, $`s=1`$ for $`k_2<0`$, and $`s=1`$ for $`k_2>0`$. Let us restrict our consideration to the analysis of the parameter $`Q(t,z)`$ for $`t=0`$, because expression (19) for an arbitrary moment of time $`t`$ is rather cumbersome.
From (16) and (19) it follows that one can have $`Q(t,z)<0`$, i.e., obtain a pulse with sub-Poissonian photon statistics through both compression ($`s=1`$) and spreading ($`s=1`$) of a pulse. Of particular interest from the practical point of view is the case of anomalous group velocity dispersion ($`s=1`$) in which the compression of a phase-modulated pulse (see (16)) takes place. The dependence of the Mandel parameter and the average number of photons in this case are presented in Fig.1-Fig.2 and Fig.3.
One can see that the suppression of fluctuation of the number of photons becomes noticeable for the nonlinear phase $`\psi _0>1`$. For a given nonlinear phase $`\psi _0`$, there is a certain dispersive phase incursion for which the suppression of the number of photons is the greatest. According to (16), this is obtained for $`\phi _{opt}(z)=1/\psi _0`$. In this case, the Mandel parameter has the minimum value $`Q_{min}=\left[𝒯/\tau _p\right]\psi _0^2`$. Note that a pulse has the minimal duration for $`\phi (z)=\psi _0/(1+\psi _0^2)`$. For $`\psi _01`$ both extremities are almost coincident.
As follows from calculations, the photon statistics for the measurement time $`𝒯\tau _p`$ is of the Poissonian type. A situation identical in many aspects to the one considered here was observed in spectral measurements of fluctuations of femtosecond pulses transmitted through an optical fibre in the region of normal group-velocity dispersion. An increase in the spectral bandwidth of a filter caused an increase of the level of fluctuations .
Let us present a numerical example. Consider the interaction of a pulse at a wavelength of $`1`$ $`\mu m`$ in an optical fibre, and let its initial duration and maximum intensity be $`2\tau _p=10ps`$ and $`10^7Wcm^2`$. In this case, the nonlinear phase in a quartz fibre for the path length $`l=100m`$ is $`\psi _0=3`$. If the resulting USP travels through an optical compressor with anomalous dispersion $`|k_2|=10^{26}s^2cm^1`$, the dispersion phase $`\phi _{opt}(z)=1/\psi _0`$ is obtained at a distance of about $`8m`$. This example illustrates the feasibility of formation of USPs with sub-Poissonian photon statistics by the method proposed above.
The obtained result shows that it is possible in principle to form a pulse with sub-Poissonian photon statistics in the case where a high-intensity pulse undergoes SPM during its travel through a nonlinear medium and subsequently passes a dispersive optical device (an optical compressor or an optical fibre). In view of the fact that these processes occur in succession, one can choose an optimum scheme or optimum conditions at each transformation stage. In the specific case of formation of low-intensity USPs in a nonclassical state, one can attenuate high-intensity pulses at the output of a nonlinear medium. It is natural that this is accompanied by a partial loss of nonclassical properties. It should be noted that in the case of formation of optical solitons, SPM and the dispersion effect take place simultaneously .
Authors are grateful to V. A. Vysloukh for fruitful discussions. This work was supported in part by the Fundamental Metrology Program of the State Committee on Science and Technology. |
warning/0003/math0003010.html | ar5iv | text | # 1 Introduction
## 1 Introduction
For compact Lie groups, conjugacy classes are essentially eigenvalues up to the Weyl group \[Ad\]. Thus the enormous physics and mathematical literature on eigenvalues of random matrices (see \[M\] for a survey) is often a study of conjugacy classes. Although the field is rapidly evolving, perhaps the closest thing to a true probabilistic understanding of eigenvalues of such matrices are the papers of Dyson \[Dy1, Dy2\]. An equally rich theory exists for the symmetric groups. The cycles of random permutations have a probabilistic description using Poisson processes \[LShe\]. The small cycles of a randomly chosen permutation are asymptotically Poisson, the medium length cycles relate to Brownian motion \[DeP\], and the long cycles relate to stick breaking \[ScV\]. The cycle structure of random permutations has numerous applications to real world problems such as population genetics (for this and more see the reference list in \[F1\]).
Some years back Persi Diaconis observed to the author that a probabilistic understanding of conjugacy classes of finite groups of Lie type was missing and urged him to find one. The paper \[F2\] provided a useful and beautiful picture for the finite general linear and unitary groups, with connections to symmetric function theory, but gave only partial results for the symplectic and orthogonal groups. The purpose of this article is to complete the program for the finite classical groups. To this two caveats should be added. First, only odd characteristic symplectic and orthogonal groups are considered. As is clear from \[W\], the even characteristic conjugacy classes have a very complicated description; a different view is given in \[FNP\]. Second, the current paper lumps together unipotent conjugacy classes with the same underlying Jordan form. As indicated at the end of the paper, this can be remedied, but the resulting formulas seem too complicated to be useful.
Before describing the contents of this paper, it is worth remarking that the probabilistic study of Jordan forms of unipotent upper triangular matrices over a finite field has a fascinating theory behind it. From the theory of wild quivers there is a provable sense in which conjugacy classes of upper triangular matrices over a finite field have no simple description; hence the reduction to Jordan form is necessary. A lovely probabilistic description of Jordan form is given in \[K\] and is exploited in \[B\]. The survey \[F3\] links their work with symmetric function theory and potential theory.
The main motivation for the current paper is \[F4\], which gave a Markov chain description of the conjugacy classes of the finite general linear and unitary groups. It is worth recalling the general nature of that description, as a variation of it occurs here. The conjugacy classes of $`GL(n,q)`$ are parameterized by rational canonical form; for each irreducible polynomial $`\varphi z`$ over $`F_q`$, one chooses a partition $`\lambda _\varphi `$ of an integer $`|\lambda _\varphi |`$, subject to the constraint that $`_\varphi deg(\varphi )|\lambda _\varphi |=n`$. One can then define a probability measure $`M`$ on the set of all partitions of all natural numbers by taking the limit as $`n\mathrm{}`$ of $`\lambda _{z1}`$ for a uniformly chosen element of $`GL(n,q)`$. (The polynomial $`z1`$ is taken without loss of generality. For other polynomials one simply replaces $`q`$ by $`q`$ raised to the degree of the polynomial in all formulas. Furthermore, asymptotically the distributions on partitions for different polynomials are independent). Recall that partitions can be viewed geometrically. For example the diagram of the partition $`(5441)`$ is:
$$\begin{array}{ccccc}.& .& .& .& .\\ .& .& .& .& \\ .& .& .& .& \\ .& & & & \end{array}$$
The Markov chain method of sampling from $`M`$ operates by choosing the size of the first column according to a certain probability distribution; then given that column $`i`$ has size $`a`$, column $`i+1`$ will have size $`b`$ with probability $`K(a,b)`$. The remarkable fact is that the probabilities $`K(a,b)`$ are independent of $`i`$, yielding a Markov chain. An immediate consequence of this viewpoint was an elementary probabilistic proof of the Rogers-Ramanujan identities, which suggested generalizations to quivers.
The main result of this note is that a similar description occurs for the symplectic and orthogonal cases, except that now the description will require two Markov chains $`K_1`$ and $`K_2`$ defined on the natural numbers. These Markov chains have the property that they can never move up, and $`K_1`$ has the additional property that it can only move down by an even amount. For the symplectic case, steps with column number $`i`$ odd use $`K_1`$ and steps with column number $`i`$ even use $`K_2`$. For the orthogonal case, steps with column number $`i`$ odd use $`K_2`$ and steps with column number $`i`$ even use $`K_1`$. The Markov chains $`K_1,K_2`$ are the same for both cases. The only disappointing aspect of our result is that the product matrices $`K_1K_2`$ and $`K_2K_1`$ do not seem to have a simple diagonalization; this blocked us from proving Rogers-Ramanujan type identities for the symplectic and orthogonal groups.
The structure of this paper is as follows. Section 2 recalls the conjugacy classes of the symplectic and orthogonal groups and defines measures on partitions from them, giving combinatorially useful rewritings. Section 3 begins with generalizations of formulas of Rudavlis and Shinoda \[RShi\], \[Shi\] and proves the aforementioned description in terms of Markov chains. In fact it is shown that the set-up extends to a more general family of measures on partitions with a parameter $`u`$.
## 2 Conjugacy Classes and Measures on Partitions
Let $`\lambda `$ be a partition of some non-negative integer $`|\lambda |`$ into parts $`\lambda _1\lambda _2\mathrm{}`$. Let $`m_i(\lambda )`$ be the number of parts of $`\lambda `$ of size $`i`$, and let $`\lambda ^{}`$ be the partition dual to $`\lambda `$ in the sense that $`\lambda _i^{}=m_i(\lambda )+m_{i+1}(\lambda )+\mathrm{}`$. Let $`n(\lambda )`$ be the quantity $`_{i1}(i1)\lambda _i`$. It is also useful to define the diagram associated to $`\lambda `$ as the set of points $`(i,j)Z^2`$ such that $`1j\lambda _i`$. We use the convention that the row index $`i`$ increases as one goes downward and the column index $`j`$ increases as one goes to the right. So the diagram of the partition $`(5441)`$ is as in the introduction.
The following combinatorial lemma about parititions will be helpful in what follows. For a proof, one simply uses the fact that $`\lambda _i^{}=m_i(\lambda )+m_{i+1}(\lambda )+\mathrm{}`$.
###### Lemma 1
$$\underset{h<i}{}2hm_h(\lambda )m_i(\lambda )+\underset{i}{}(i1)m_i(\lambda )^2=\underset{i}{}(\lambda _i^{})^2\underset{i}{}m_i(\lambda )^2$$
We also recall the following formulas for the sizes of finite symplectic and orthogonal groups in odd characteristic.
$`|Sp(2n,q)|`$ $`=`$ $`q^{n^2}{\displaystyle \underset{i=1}{\overset{n}{}}}(q^{2i}1)`$
$`|O^\pm (2n+1,q)|`$ $`=`$ $`2q^{n^2}{\displaystyle \underset{i=1}{\overset{n}{}}}(q^{2i}1)`$
$`|O^\pm (2n,q)|`$ $`=`$ $`2q^{n^2n}(q^n1){\displaystyle \underset{i=1}{\overset{n1}{}}}(q^{2i}1)`$
### 2.1 Symplectic groups
Wall \[W\] parametrized the conjugacy classes of the finite symplectic groups and found formulas for their sizes. Let us recall his parametrization for the case of odd characteristic. Given a a polynomial $`\varphi (z)`$ with coefficients in $`F_q`$ and non vanishing constant term, define a polynomial $`\overline{\varphi }`$ by
$$\overline{\varphi }=\frac{z^{deg(\varphi )}\varphi (\frac{1}{z})}{\varphi (0)}.$$
Wall showed that a conjugacy class of $`Sp(2n,q)`$ corresponds to the following data. To each monic, non-constant, irreducible polynomial $`\varphi z\pm 1`$ associate a partition $`\lambda _\varphi `$ of some non-negative integer $`|\lambda _\varphi |`$. To $`\varphi `$ equal to $`z1`$ or $`z+1`$ associate a symplectic signed partition $`\lambda (\pm )_\varphi `$, by which is meant a partition of some natural number $`|\lambda (\pm )_\varphi |`$ such that the odd parts have even multiplicity, together with a choice of sign for the set of parts of size $`i`$ for each even $`i>0`$.
Example of a Symplectic Signed Partition
$$\begin{array}{cccccccc}& .& .& .& .& .& & \\ & .& .& .& .& .& & \\ +& .& .& .& .& & & \\ & .& .& .& & & & \\ & .& .& .& & & & \\ & .& .& & & & & \\ & .& .& & & & & \end{array}$$
Here the $`+`$ corresponds to the parts of size 4 and the $``$ corresponds to the parts of size 2. This data represents a conjugacy class of $`Sp(2n,q)`$ if and only if:
1. $`|\lambda _z|=0`$
2. $`\lambda _\varphi =\lambda _{\overline{\varphi }}`$
3. $`_{\varphi =z\pm 1}|\lambda (\pm )_\varphi |+_{\varphi z\pm 1}|\lambda _\varphi |deg(\varphi )=2n`$
Let
$$A_{Sp}(\varphi ^i)=\{\begin{array}{cc}|Sp(m_i(\lambda (\pm )_\varphi ),q)|\hfill & \text{if i odd,}\varphi =z\pm 1\hfill \\ q^{\frac{m_i(\lambda (\pm )_\varphi )}{2}}|O(m_i(\lambda (\pm )_\varphi ),q)|\hfill & \text{if i even,}\varphi =z\pm 1\hfill \\ |U(m_i(\lambda _\varphi ),q^{\frac{deg(\varphi )}{2}})|\hfill & \text{if}\varphi =\overline{\varphi }z\pm 1\hfill \\ |GL(m_i(\lambda _\varphi ),q^{deg(\varphi )})|^{\frac{1}{2}}\hfill & \text{if}\varphi \overline{\varphi }.\hfill \end{array}$$
where $`O(m_i(\lambda _\varphi ),q)`$ is the orthogonal group with the same sign as the sign associated to the parts of size $`i`$.
Theorem 1 is implicit in the discussion in \[F1\]. The three ingredients in its proof are Wall’s formulas for conjugacy class sizes \[W\], the deduction that the cycle index of the symplectic groups factors, and the fact that the formulas in the statement of Theorem 1 define probability measures (i.e. the asserted probabilities sum to one). This third fact will be deduced in the proof of Theorem 4, using only an identity of Cauchy. It is worth emphasizing that neither Steinberg’s count of unipotent elements nor the formulas of Rudvalis and Shinoda in Section 3 are needed to prove the third fact.
###### Theorem 1
Fix some value of $`u`$ with $`0<u<1`$. Then pick a non-negative even integer with the probability of getting $`2n`$ equal to $`(1u^2)u^{2n}`$ and pick uniformly in $`Sp(2n,q)`$. Let $`\mathrm{\Lambda }(\pm )_{z1},\mathrm{\Lambda }(\pm )_{z+1},\mathrm{\Lambda }_\varphi `$ be the random variables corresponding to the conjugacy class data of the chosen element of $`Sp(2n,q)`$. Then, aside from the fact that $`\mathrm{\Lambda }_\varphi =\mathrm{\Lambda }_{\overline{\varphi }}`$, any finite number of these random variables are independent, with probability laws
$`Prob(\mathrm{\Lambda }(\pm )_{z1}=\lambda (\pm )_{z1})`$
$`=`$ $`{\displaystyle \frac{_{r=1}^{\mathrm{}}(1u^2/q^{2r1})u^{|\lambda (\pm )_{z1}|}}{q^{[_{h<i}hm_h(\lambda (\pm )_{z1})m_i(\lambda (\pm )_{z1})+\frac{1}{2}_i(i1)m_i(\lambda (\pm )_{z1})^2]}_iA_{Sp}((z1)^i)}}`$
$`Prob(\mathrm{\Lambda }(\pm )_{z+1}=\lambda (\pm )_{z+1})`$
$`=`$ $`{\displaystyle \frac{_{r=1}^{\mathrm{}}(1u^2/q^{2r1})u^{|\lambda (\pm )_{z+1}|}}{q^{[_{h<i}hm_h(\lambda (\pm )_{z+1})m_i(\lambda (\pm )_{z+1})+\frac{1}{2}_i(i1)m_i(\lambda (\pm )_{z+1})^2]}_iA_{Sp}((z+1)^i)}}`$
$`Prob(\mathrm{\Lambda }_\varphi =\lambda _\varphi )`$ $`=`$ $`{\displaystyle \frac{_{r=1}^{\mathrm{}}(1+(1)^r\frac{u^{deg(\varphi )}}{q^{deg(\varphi )r/2}})u^{deg(\varphi )|\lambda _\varphi |}}{q^{deg(\varphi )[_{h<i}hm_h(\lambda _\varphi )m_i(\lambda _\varphi )+\frac{1}{2}_i(i1)m_i(\lambda _\varphi )^2]}_iA_{Sp}(\varphi ^i)}}if\varphi =\overline{\varphi }z\pm 1`$
$`Prob(\mathrm{\Lambda }_\varphi =\lambda _\varphi )`$ $`=`$ $`{\displaystyle \frac{_{r=1}^{\mathrm{}}(1\frac{u^{2deg(\varphi )}}{q^{deg(\varphi )r}})u^{2deg(\varphi )|\lambda _\varphi |}}{q^{2deg(\varphi )[_{h<i}hm_h(\lambda _\varphi )m_i(\lambda _\varphi )+\frac{1}{2}_i(i1)m_i(\lambda _\varphi )^2]}_iA_{Sp}(\varphi ^i)}}if\varphi \overline{\varphi }.`$
Furthermore, setting $`u=1`$ in these formulas yields the laws arising from the $`n\mathrm{}`$ limit of conjugacy classes of a uniformly chosen element of $`Sp(2n,q)`$, and the random variables corresponding to different polynomials are independent, up to the fact that $`\mathrm{\Lambda }_\varphi =\mathrm{\Lambda }_{\overline{\varphi }}`$.
From Theorem 1, one sees that if $`\varphi =\overline{\varphi }`$, then the corresponding measures on partitions are specializations of those for the unitary groups treated in \[F2\]. Similarly, if $`\varphi \overline{\varphi }`$, then the corresponding measures on partitions are specializations of those for the general linear groups treated in \[F2\]. As the formulas for $`z\pm 1`$ are the same, for the rest of this paper only the partition corresponding to $`z1`$ will be studied.
Combining Theorem 1 with Lemma 1 leads one to the following measure on symplectic signed partitions:
$$M_{Sp,u}^\pm (\lambda (\pm ))=\underset{r=1}{\overset{\mathrm{}}{}}(1u^2/q^{2r1})\frac{u^{|\lambda (\pm )|}}{q^{1/2[_i(\lambda (\pm )_i^{})^2_im_i(\lambda (\pm ))^2]}_iA_{Sp}((z1)^i)}.$$
Forgetting about signs (i.e. lumping together some conjugacy classes) yields a measure on underlying shapes which will be denoted by $`M_{Sp.u}`$. Using the formulas for the sizes of the finite symplectic and orthogonal groups given at the beginning of this section, one arrives at the expression:
$`M_{Sp,u}(\lambda )`$ $`=`$ $`{\displaystyle \frac{_{r=1}^{\mathrm{}}(1u^2/q^{2r1})u^{|\lambda |}}{q^{1/2[_i(\lambda _i^{})^2_im_i^2]}_{i=1mod2}(q^{\frac{m_i^2}{4}}_{l=1}^{m_i/2}(q^{2l}1))}}`$
$`{\displaystyle \frac{1}{_{\genfrac{}{}{0pt}{}{i=0mod2}{m_i=0mod2}}(q^{\frac{m_i^2}{4}\frac{m_i}{2}}_{l=1}^{m_i/2}(q^{2l}1))_{\genfrac{}{}{0pt}{}{i=0mod2}{m_i=1mod2}}(q^{\frac{m_i^2+1}{4}}_{l=1}^{(m_i1)/2}(q^{2l}1))}}`$
### 2.2 Orthogonal groups
Wall \[W\] parametrized the conjugacy classes of the finite orthogonal groups and found formulas for their sizes. Let us recall his parametrization for the case of odd characteristic. To each monic, non-constant, irreducible polynomial $`\varphi z\pm 1`$ associate a partition $`\lambda _\varphi `$ of some non-negative integer $`|\lambda _\varphi |`$. To $`\varphi `$ equal to $`z1`$ or $`z+1`$ associate an orthogonal signed partition $`\lambda (\pm )_\varphi `$, by which is meant a partition of some natural number $`|\lambda (\pm )_\varphi |`$ such that all even parts have even multiplicity, and all odd $`i>0`$ have a choice of sign. For $`\varphi =z1`$ or $`\varphi =z+1`$ and odd $`i>0`$, we denote by $`\mathrm{\Theta }_i(\lambda (\pm )_\varphi )`$ the Witt type of the orthogonal group on a vector space of dimension $`m_i(\lambda (\pm )_\varphi )`$ and sign the choice of sign for $`i`$.
Example of an Orthogonal Signed Partition
$$\begin{array}{ccccc}& .& .& .& .\\ & .& .& .& .\\ & .& .& .& \\ & .& .& & \\ & .& .& & \\ +& .& & & \\ & .& & & \end{array}$$
Here the $``$ corresponds to the part of size 3 and the $`+`$ corresponds to the parts of size 1. The data $`\lambda (\pm )_{z1},\lambda (\pm )_{z+1},\lambda _\varphi `$ represents a conjugacy class of some orthogonal group if:
1. $`|\lambda _z|=0`$
2. $`\lambda _\varphi =\lambda _{\overline{\varphi }}`$
3. $`_{\varphi =z\pm 1}|\lambda (\pm )_\varphi |+_{\varphi z\pm 1}|\lambda _\varphi |deg(\varphi )=n`$.
In this case, the data represents the conjugacy class of exactly 1 orthogonal group $`O(n,q)`$, with sign determined by the condition that the group arises as the stabilizer of a form of Witt type:
$$\underset{\varphi =z\pm 1}{}\underset{iodd}{}\mathrm{\Theta }_i(\lambda (\pm )_\varphi )+\underset{\varphi z\pm 1}{}\underset{i1}{}im_i(\lambda _\varphi )\omega ,$$
where $`\omega `$ is the Witt type of the quadratic form $`x^2\delta y^2`$ with $`\delta `$ a fixed non-square in $`F_q`$.
Let
$$A_O(\varphi ^i)=\{\begin{array}{cc}q^{m_i(\lambda (\pm )_\varphi )/2}|Sp(m_i(\lambda (\pm )_\varphi ),q)|\hfill & \text{if i even,}\varphi =z\pm 1\hfill \\ |O(m_i(\lambda (\pm )_\varphi ),q)|\hfill & \text{if i odd,}\varphi =z\pm 1\hfill \\ |U(m_i(\lambda _\varphi ),q^{\frac{deg(\varphi )}{2}})|\hfill & \text{if}\varphi =\overline{\varphi }z\pm 1\hfill \\ |GL(m_i(\lambda _\varphi ),q^{deg(\varphi )})|^{\frac{1}{2}}\hfill & \text{if}\varphi \overline{\varphi }.\hfill \end{array}$$
where $`O(m_i(\lambda _\varphi ),q)`$ is the orthogonal group with the same sign as the sign associated to the parts of size $`i`$.
Theorem 2 is implicit in \[F1\]. The three ingredients in its proof are Wall’s formulas for conjugacy class sizes \[W\], the deduction that the cycle index for the sum of $`+`$ and $``$ types of the orthogonal groups factors, and the fact that the formulas in the statement of Theorem 2 define probability measures. As in the symplectic case, this third fact will be deduced in the proof of Theorem 5, using only an identity of Cauchy.
###### Theorem 2
Fix some value of $`u`$ with $`0<u<1`$. Then pick a non-negative integer with the probability of getting $`0`$ equal to $`\frac{(1u)}{(1+u)}`$ and probability of getting $`n>0`$ equal to $`\frac{2u^n(1u)}{1+u}`$. Choose either $`O^+(n,q)`$ or $`O^{}(n,q)`$ with probability $`\frac{1}{2}`$. Finally select an element uniformly within the chosen orthogonal group and let $`\mathrm{\Lambda }(\pm )_{z1},\mathrm{\Lambda }(\pm )_{z+1},\mathrm{\Lambda }_\varphi `$ be the random variables corresponding to its conjugacy class data. Then aside from the fact that $`\mathrm{\Lambda }_\varphi =\mathrm{\Lambda }_{\overline{\varphi }}`$, any finite number of these random variables are independent, with probability laws
$`Prob(\mathrm{\Lambda }(\pm )_{z1}=\lambda (\pm )_{z1})`$
$`=`$ $`{\displaystyle \frac{_{r=1}^{\mathrm{}}(1u^2/q^{2r1})u^{|\lambda (\pm )_{z1}|}}{(1+u)q^{[_{h<i}hm_h(\lambda (\pm )_{z1})m_i(\lambda (\pm )_{z1})+\frac{1}{2}_i(i1)m_i(\lambda (\pm )_{z1})^2]}_iA_O((z1)^i)}}`$
$`Prob(\mathrm{\Lambda }(\pm )_{z+1}=\lambda (\pm )_{z+1})`$
$`=`$ $`{\displaystyle \frac{_{r=1}^{\mathrm{}}(1u^2/q^{2r1})u^{|\lambda (\pm )_{z+1}|}}{(1+u)q^{[_{h<i}hm_h(\lambda (\pm )_{z+1})m_i(\lambda (\pm )_{z+1})+\frac{1}{2}_i(i1)m_i(\lambda (\pm )_{z+1})^2]}_iA_O((z+1)^i)}}`$
$`Prob(\mathrm{\Lambda }_\varphi =\lambda _\varphi )`$ $`=`$ $`{\displaystyle \frac{_{r=1}^{\mathrm{}}(1+(1)^r\frac{u^{deg(\varphi )}}{q^{deg(\varphi )r/2}})u^{deg(\varphi )|\lambda _\varphi |}}{q^{deg(\varphi )[_{h<i}hm_h(\lambda _\varphi )m_i(\lambda _\varphi )+\frac{1}{2}_i(i1)m_i(\lambda _\varphi )^2]}_iA_O(\varphi ^i)}}if\varphi =\overline{\varphi }z\pm 1`$
$`Prob(\mathrm{\Lambda }_\varphi =\lambda _\varphi )`$ $`=`$ $`{\displaystyle \frac{_{r=1}^{\mathrm{}}(1\frac{u^{2deg(\varphi )}}{q^{deg(\varphi )r}})u^{2deg(\varphi )|\lambda _\varphi |}}{q^{2deg(\varphi )[_{h<i}hm_h(\lambda _\varphi )m_i(\lambda _\varphi )+\frac{1}{2}_i(i1)m_i(\lambda _\varphi )^2]}_iA_O(\varphi ^i)}}if\varphi \overline{\varphi }.`$
Furthermore, setting $`u=1`$ in these formulas yields the laws arising from the $`n\mathrm{}`$ limit of conjugacy classes of a uniformly chosen element of $`O(n,q)`$, where the $`+,`$ sign is chosen with probability $`1/2`$. The random variables corresponding to different polynomials are independent, up to the fact that $`\mathrm{\Lambda }_\varphi =\mathrm{\Lambda }_{\overline{\varphi }}`$.
For the same reasons as with the symplectic groups, the only case remaining to be understood is the measure of the partition corresponding to the polynomial $`z1`$. Combining Theorem 1 with Lemma 1 leads one to the following measure on orthogonal signed partitions:
$$M_{O,u}^\pm (\lambda (\pm ))=\frac{1}{1+u}\underset{r=1}{\overset{\mathrm{}}{}}(1u^2/q^{2r1})\frac{u^{|\lambda (\pm )|}}{q^{1/2[_i(\lambda (\pm )_i^{})^2_im_i(\lambda (\pm ))^2]}_iA_O((z1)^i)}.$$
Forgetting about signs (i.e. lumping together some conjugacy classes) yields a measure on underlying shapes which will be denoted by $`M_{O,u}`$. Using the formulas for the sizes of the finite symplectic and orthogonal groups given at the beginning of this section, one arrives at the expression:
$`M_{O,u}(\lambda )`$ $`=`$ $`{\displaystyle \frac{1}{1+u}}{\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}(1u^2/q^{2r1}){\displaystyle \frac{u^{|\lambda |}}{q^{1/2[_i(\lambda _i^{})^2_im_i^2]}_{i=0mod2}(q^{\frac{m_i^2}{4}\frac{m_i}{2}}_{l=1}^{m_i/2}(q^{2l}1))}}`$
$`{\displaystyle \frac{1}{_{\genfrac{}{}{0pt}{}{i=1mod2}{m_i=0mod2}}(q^{\frac{m_i^2}{4}m_i}_{l=1}^{m_i/2}(q^{2l}1))_{\genfrac{}{}{0pt}{}{i=0mod2}{m_i=1mod2}}(q^{\frac{m_i^21}{4}}_{l=1}^{(m_i1)/2}(q^{2l}1))}}`$
## 3 Markov Chain Description
This section proves the Markov chain descriptions of conjugacy classes as advertised in the introduction. The first goal is Theorem 3, which generalizes work of Rudvalis and Shinoda \[RShi\] (proved by different methods and not in the language of partitions). First, a lemma of Cauchy is needed.
###### Lemma 2
(\[An\], p.20) If $`|q|>1`$,
$$\underset{m=0}{\overset{\mathrm{}}{}}(1z/q^m)^1=1+\underset{n=1}{\overset{\mathrm{}}{}}\frac{z^n}{q^{n^2n}(11/q)(11/q^2)\mathrm{}(11/q^n)(1z)(1z/q)\mathrm{}(1z/q^{n1})}$$
Let $`G`$ denote either $`Sp`$ or $`O`$, and let $`P_{G,u}(i)`$ be the probability that a partition chosen from the measure $`M_{G,u}`$ has $`i`$ parts. Let
$`P_{Sp,u}^{}(i)`$ $`=`$ $`{\displaystyle \frac{P_{Sp,u}(i)}{_{i=1}^{\mathrm{}}(1u^2/q^{2i1})}}`$
$`P_{O,u}^{}(i)`$ $`=`$ $`{\displaystyle \frac{(1+u)P_{O,u}(i)}{_{i=1}^{\mathrm{}}(1u^2/q^{2i1})}}.`$
###### Theorem 3
$`P_{Sp,u}(2k)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}(1u^2/q^{2i1}){\displaystyle \frac{u^{2k}}{q^{2k^2+k}(1u^2/q)(11/q^2)\mathrm{}(1u^2/q^{2k1})(11/q^{2k})}}`$
$`P_{Sp,u}(2k+1)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}(1u^2/q^{2i1}){\displaystyle \frac{u^{2k+2}}{q^{2k^2+3k+1}(1u^2/q)(11/q^2)\mathrm{}(11/q^{2k})(1u^2/q^{2k+1})}}`$
$`P_{O,u}(2k)`$ $`=`$ $`{\displaystyle \frac{_{i=1}^{\mathrm{}}(1u^2/q^{2i1})}{1+u}}{\displaystyle \frac{u^{2k}}{q^{2k^2k}(1u^2/q)(11/q^2)\mathrm{}(1u^2/q^{2k1})(11/q^{2k})}}`$
$`P_{O,u}(2k+1)`$ $`=`$ $`{\displaystyle \frac{_{i=1}^{\mathrm{}}(1u^2/q^{2i1})}{1+u}}{\displaystyle \frac{u^{2k+1}}{q^{2k^2+k}(1u^2/q)(11/q^2)\mathrm{}(11/q^{2k})(1u^2/q^{2k+1})}}.`$
Proof: Using only the facts that $`M_{Sp,u}`$ and $`M_{O,u}`$ define a measure (i.e. not necessarily a probability measure), the proofs of Theorems 4 and 5 will establish the equations:
$`P_{Sp,u}^{}(a)`$ $`=`$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{ba}{abeven}}{}}{\displaystyle \frac{u^aP_{O,u}^{}(b)}{q^{(a^2b^2+2(a+1)b)/4}(q^{ab}1)\mathrm{}(q^41)(q^21)}}`$
$`P_{O,u}^{}(a)`$ $`=`$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{ba}{abeven}}{}}{\displaystyle \frac{u^aP_{Sp,u}^{}(b)q^{(ab)^2/4}}{q^{(a^2+b2a)/2}(q^{ab}1)\mathrm{}(q^41)(q^21)}}`$
$`+{\displaystyle \underset{\genfrac{}{}{0pt}{}{ba}{abodd}}{}}{\displaystyle \frac{u^aP_{Sp,u}^{}(b)q^{((ab)^21)/4}}{q^{(a^2a)/2}(q^{ab1}1)\mathrm{}(q^41)(q^21)}}.`$
To get a recurrence relation for the $`P_{Sp,u}^{}(a)`$’s, one simply plugs the second equation into the first. Similarly one obtains a recurrence relation for the $`P_{O,u}^{}(a)`$’s. These recurrences allow one to solve for $`P_{G,u}^{}(a)`$ in terms of $`P_{G,u}^{}(0)`$, implying that the formulas for $`P_{G,u}(a)`$ are proportional to the asserted values. Thus it is enough to prove that the asserted formulas for $`P_{G,u}(a)`$ satisfy the equation $`_{a0}P_{G,u}(a)=1`$. This follows readily from Lemma 2. $`\mathrm{}`$
Before continuing, we pause to indicate how the formulas of Theorem 3 can be used to deduce group theoretic results which are normally proved by techniques such as character theory and Moebius inversion. The first set of results, Corollary 1, considers only the symplectic groups. The same technique would give results for the sum of $`+,`$ type orthogonal groups. Since the cycle index for the difference of orthogonal groups also factors, one could rework all of the paper until now to give measures corresponding to the difference of $`+,`$ type orthogonal groups, apply the same technique, and then average the results to get theorems about groups over a given $`+`$ or $``$ type. This does not deserve to be done publicly.
###### Corollary 1
1. (Steinberg, pg. 156 of \[H\]) The number of unipotent elements in $`Sp(2n,q)`$ is $`q^{2n^2}`$.
2. (\[RShi\]) The probability that an randomly chosen element of $`Sp(2n,q)`$ has a $`2k`$ dimensional fixed space is
$$\frac{1}{|Sp(2k,q)|}\underset{i=0}{\overset{nk}{}}\frac{(1)^i(q^2)^{\left(\genfrac{}{}{0pt}{}{i}{2}\right)}}{|Sp(2i,q)|q^{2ik}}.$$
The probability that an randomly chosen element of $`Sp(2n,q)`$ has a $`2k+1`$ dimensional fixed space is
$$\frac{1}{|Sp(2k,q)|q^{2k+1}}\underset{i=0}{\overset{nk1}{}}\frac{(1)^i(q^2)^{\left(\genfrac{}{}{0pt}{}{i}{2}\right)}}{|Sp(2i,q)|q^{2i(k+1)}}.$$
Proof: The arguments are completely analogous to those for $`GL(n,q)`$ in \[F2\] (Corollary 1 and Theorem 6), using the cycle index of the finite symplectic groups \[F1\]. $`\mathrm{}`$
Rudvalis and Shinoda (loc. cit.) also considered the probability that the fixed space of a random element of a finite classical group has a given isometry type. For the finite unitary groups, the isometry classes are parameterized by pairs $`(s,t)`$ of natural numbers such that $`s+2tn`$. Here a subspace $`W`$ of $`V`$ has type $`(s,t)`$ if $`dim(W/rad(W))=s`$ and $`dim(rad(W))=t`$. Theorem 2 uses cycle index techniques to give new proofs of their results for the finite unitary groups. Exactly the same methods work for the finite symplectic and orthogonal groups, but we spare the reader the details.
###### Corollary 2
The probability that an element of $`U(n,q)`$ has isometry type corresponding to the pair $`(s,t)`$ is
$$\frac{_{i=0}^{n2st}\frac{(1/q)^{(t+1)i}(1/q)^{\left(\genfrac{}{}{0pt}{}{i}{2}\right)}}{(1+1/q)(11/q^2)\mathrm{}(1(1)^i/q^i)}}{q^{s^2+2st}(1+1/q)(11/q^2)\mathrm{}(1(1)^s/q^s)(1+1/q)(11/q^2)\mathrm{}(1(1)^t/q^t)}.$$
In the $`n\mathrm{}`$ limit, this converges to
$$\frac{_{r=0}^{\mathrm{}}(\frac{1}{1+1/q^{2r+1}})}{q^{s^2+2st}(1+1/q)(11/q^2)\mathrm{}(1(1)^s/q^s)(11/q^2)(11/q^4)\mathrm{}(11/q^{2t})}.$$
Proof: The most important observation (see \[FNP\] for a readable proof) is that the fixed space of an element $`\alpha `$ of $`U(n,q)`$ has isometry type $`(s,t)`$ precisely when the partition corresponding to the polynomial $`z1`$ in the rational canonical form of $`\alpha `$ satisfies $`\lambda _1^{}=s+t`$ and $`\lambda _2^{}=t`$. In other words, the partition has $`s+t`$ parts and $`s`$ $`1`$’s. Let $`[u^n]f(u)`$ denote the coefficient of $`u^n`$ in some polynomial $`f(u)`$. Then one uses the cycle index for the unitary groups as in \[F2\] to see that the sought probability for $`U(n,q)`$ is
$$[u^n]\frac{1}{1u}\underset{\lambda :\lambda _1^{}=s+t,\lambda _2^{}=t}{}M_{U,u}(\lambda ).$$
Using the fact that
$$M_{U,u}(\lambda )=\underset{r=1}{\overset{\mathrm{}}{}}(1+u/(q)^r)\frac{u^{|\lambda |}}{q^{_i(\lambda _i^{})^2}_i(1+1/q)(11/q^2)\mathrm{}(1+(1)^{m_i+1}/q^{m_i})},$$
this becomes
$`[u^n]{\displaystyle \frac{u^{s+t}}{(1u)q^{(s+t)^2}(1+1/q)(11/q^2)\mathrm{}(1(1)^s/q^s)}}{\displaystyle \underset{\lambda :\lambda _1^{}=t}{}}M_{U,u}(\lambda )`$
$`=`$ $`{\displaystyle \frac{1}{q^{(s+t)^2}(1+1/q)(11/q^2)\mathrm{}(1(1)^s/q^s)}}[u^{nst}]{\displaystyle \frac{1}{1u}}{\displaystyle \underset{\lambda :\lambda _1^{}=t}{}}M_{U,u}(\lambda )`$
$`=`$ $`{\displaystyle \frac{_{i=0}^{n2st}\frac{(1/q)^{(t+1)i}(1/q)^{\left(\genfrac{}{}{0pt}{}{i}{2}\right)}}{(1+1/q)(11/q^2)\mathrm{}(1(1)^i/q^i)}}{q^{s^2+2st}(1+1/q)(11/q^2)\mathrm{}(1(1)^s/q^s)(1+1/q)(11/q^2)\mathrm{}(1(1)^t/q^t)}},`$
where the last equality is in the proof of Theorem 6 in \[F2\].
The formula for the $`n\mathrm{}`$ limit follows from the well-known identity
$$\underset{r=1}{\overset{\mathrm{}}{}}(1v/w^r)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{(v)^n}{(w^n1)\mathrm{}(w1)}.$$
Alternatively, it follows from the principle that the limit as $`n\mathrm{}`$ of $`f(u)/(1u)`$ is $`f(1)`$ if $`f`$ has a Taylor expansion around $`0`$ converging in a circle of radius $`1`$, together with the formula for
$$\underset{\lambda :\lambda _1^{}=t}{}M_{U,1}(\lambda )$$
given in Theorem 5 of \[F2\]. $`\mathrm{}`$
Lemma 3 is crucial and motivated the combinatorial moves made in rewriting the formulas for $`M_{Sp,u}`$ and $`M_{O,u}`$ in Section 2. In all that follows $`M_{G,x}`$ will denote the probability of an event $`X`$ under the measure $`M_{G,u}`$ with $`G`$ equal to $`Sp`$ or $`O`$.
###### Lemma 3
1. If $`i`$ is odd then
$$M_{Sp,u}(\lambda _1^{}=s_1,\mathrm{},\lambda _{i1}^{}=s_{i1},\lambda _i^{}=k)=\frac{_{r=1}^{\mathrm{}}(1u^2/q^{2r1})u^{s_1+\mathrm{}+s_{i1}}P_{Sp,u}^{}(k)}{q^{(s_1^2+\mathrm{}+s_{i1}^2m_1^2\mathrm{}m_{i1}^2)/2}_{j=1}^{i1}A_{Sp}((z1)^{m_j})}.$$
2. If $`i`$ is even then
$`M_{Sp,u}(\lambda _1^{}=s_1,\mathrm{},\lambda _{i1}^{}=s_{i1},\lambda _i^{}=k)`$
$`=`$ $`{\displaystyle \frac{_{r=1}^{\mathrm{}}(1u^2/q^{2r1})u^{s_1+\mathrm{}+s_{i1}}P_{O,u}^{}(k)}{q^{(k+s_1^2+\mathrm{}+s_{i1}^2m_1^2\mathrm{}m_{i1}^2)/2}_{j=1}^{i1}A_{Sp}((z1)^{m_j})}}.`$
3. If $`i`$ is odd then
$`M_{O,u}(\lambda _1^{}=s_1,\mathrm{},\lambda _{i1}^{}=s_{i1},\lambda _i^{}=k)`$
$`=`$ $`{\displaystyle \frac{_{r=1}^{\mathrm{}}(1u^2/q^{2r1})u^{s_1+\mathrm{}+s_{i1}}P_{O,u}^{}(k)}{(1+u)q^{(s_1^2+\mathrm{}+s_{i1}^2m_1^2\mathrm{}m_{i1}^2)/2}_{j=1}^{i1}A_O((z1)^{m_j})}}.`$
4. If $`i`$ is even then
$`M_{O,u}(\lambda _1^{}=s_1,\mathrm{},\lambda _{i1}^{}=s_{i1},\lambda _i^{}=k)`$
$`=`$ $`{\displaystyle \frac{_{r=1}^{\mathrm{}}(1u^2/q^{2r1})u^{s_1+\mathrm{}+s_{i1}}q^{k/2}P_{Sp,u}^{}(k)}{(1+u)q^{(s_1^2+\mathrm{}+s_{i1}^2m_1^2\mathrm{}m_{i1}^2)/2}_{j=1}^{i1}A_O((z1)^{m_j})}}.`$
Proof: The idea for all of the proofs is the same; hence we prove part two as follows:
$`M_{Sp,u}(\lambda _1^{}=s_1,\mathrm{},\lambda _{i1}^{}=s_{i1},\lambda _i^{}=k)`$
$`=`$ $`{\displaystyle \frac{_{r=1}^{\mathrm{}}(1u^2/q^{2r1})u^{s_1+\mathrm{}+s_{i1}}}{q^{(s_1^2+\mathrm{}+s_{i1}^2m_1^2\mathrm{}m_{i1}^2)/2}_{j=1}^{i1}A_{Sp}((z1)^{m_j})}}`$
$`{\displaystyle \underset{\lambda _i^{}=k\lambda _{i+1}^{}\mathrm{}0}{}}{\displaystyle \frac{u^{_{ji}\lambda _j^{}}}{q^{_{ji}(\lambda _j^2m_j^2)/2}_{ji}A_{Sp}((z1)^{m_j})}}`$
$`=`$ $`{\displaystyle \frac{_{r=1}^{\mathrm{}}(1u^2/q^{2r1})u^{s_1+\mathrm{}+s_{i1}}}{q^{(k+s_1^2+\mathrm{}+s_{i1}^2m_1^2\mathrm{}m_{i1}^2)/2}_{j=1}^{i1}A_{Sp}((z1)^{m_j})}}`$
$`{\displaystyle \underset{\lambda _i^{}=k\lambda _{i+1}^{}\mathrm{}0}{}}{\displaystyle \frac{u^{_{ji}\lambda _j^{}}}{q^{_{ji}(\lambda _j^2m_j^2)/2}_{ji}A_O((z1)^{m_j})}}`$
$`=`$ $`{\displaystyle \frac{_{r=1}^{\mathrm{}}(1u^2/q^{2r1})u^{s_1+\mathrm{}+s_{i1}}}{q^{(k+s_1^2+\mathrm{}+s_{i1}^2m_1^2\mathrm{}m_{i1}^2)/2}_{j=1}^{i1}A_{Sp}((z1)^{m_j})}}`$
$`{\displaystyle \underset{\lambda _1^{}=k\lambda _2^{}\mathrm{}0}{}}{\displaystyle \frac{u^{_{j1}\lambda _j^{}}}{q^{_{j1}(\lambda _j^2m_j^2)/2}_{j1}A_O((z1)^{m_j})}}`$
$`=`$ $`{\displaystyle \frac{_{r=1}^{\mathrm{}}(1u^2/q^{2r1})u^{s_1+\mathrm{}+s_{i1}}}{q^{(k+s_1^2+\mathrm{}+s_{i1}^2m_1^2\mathrm{}m_{i1}^2)/2}_{j=1}^{i1}A_{Sp}((z1)^{m_j})}}P_{O,u}^{}(k),`$
as desired. Note that the meat of the lemma is the second equality, which follows from the formulas for $`A_{Sp}`$ and $`A_O`$. The third equality is simply a relabelling of subscripts. $`\mathrm{}`$
For the statements of Theorems 4 and 5, we define two Markov chains on the integers. Recall that
$`P_{Sp,u}^{}(2k)`$ $`=`$ $`{\displaystyle \frac{u^{2k}}{q^{2k^2+k}(1u^2/q)(11/q^2)\mathrm{}(1u^2/q^{2k1})(11/q^{2k})}}`$
$`P_{Sp,u}^{}(2k+1)`$ $`=`$ $`{\displaystyle \frac{u^{2k+2}}{q^{2k^2+3k+1}(1u^2/q)(11/q^2)\mathrm{}(11/q^{2k})(1u^2/q^{2k+1})}}`$
$`P_{O,u}^{}(2k)`$ $`=`$ $`{\displaystyle \frac{u^{2k}}{q^{2k^2k}(1u^2/q)(11/q^2)\mathrm{}(1u^2/q^{2k1})(11/q^{2k})}}`$
$`P_{O,u}^{}(2k+1)`$ $`=`$ $`{\displaystyle \frac{u^{2k+1}}{q^{2k^2+k}(1u^2/q)(11/q^2)\mathrm{}(11/q^{2k})(1u^2/q^{2k+1})}}.`$
The chains $`K_1,K_2`$ are defined on the natural numbers with transition probabilities
$$K_1(a,b)=\{\begin{array}{cc}\frac{u^aP_{O,u}^{}(b)}{P_{Sp,u}^{}(a)q^{\frac{a^2b^2+2(a+1)b}{4}}(q^{ab}1)\mathrm{}(q^41)(q^21)}\hfill & \text{if }ab\text{ even}\hfill \\ 0\hfill & \text{if }ab\text{ odd}\hfill \end{array}$$
$$K_2(a,b)=\{\begin{array}{cc}\frac{u^aP_{Sp,u}^{}(b)q^{(ab)^2/4}}{P_{O,u}^{}(a)q^{\frac{a^2+b}{2}a}(q^{ab}1)\mathrm{}(q^41)(q^21)}\hfill & \text{if }ab\text{ even}\hfill \\ \frac{u^aP_{Sp,u}^{}(b)q^{((ab)^21)/4}}{P_{O,u}^{}(a)q^{\frac{a^2a}{2}}(q^{ab1}1)\mathrm{}(q^41)(q^21)}\hfill & \text{if }ab\text{ odd}\hfill \end{array}$$
The fact that these transition probabilities add up to one will follow from the proof of Theorem 4.
###### Theorem 4
Starting with $`\lambda _1^{}`$ distributed as $`P_{Sp,u}`$, define in succession $`\lambda _2^{},\lambda _3^{},\mathrm{}`$ according to the rules that if $`\lambda _i^{}=a`$, then $`\lambda _{i+1}^{}=b`$ with probability $`K_1(a,b)`$ if $`i`$ is odd and probability $`K_2(a,b)`$ if $`i`$ is even. The resulting partition is distributed according to $`M_{Sp,u}`$.
Proof: The $`M_{Sp,u}`$ probability of choosing a partition with $`\lambda _i^{}=s_i`$ for all $`i`$ is
$$M_{Sp,u}(\lambda _1^{}=s_1)\underset{i=1}{\overset{\mathrm{}}{}}\frac{M_{Sp,u}(\lambda _1^{}=s_1,\mathrm{},\lambda _{i+1}^{}=s_{i+1})}{M_{Sp,u}(\lambda _1^{}=s_1,\mathrm{},\lambda _i^{}=s_i)}.$$
Since $`M_{Sp,u}(\lambda _1^{}=s_1)`$ is equal to $`P_{Sp,u}(s_1)`$ by definition, it it is enough to prove two claims: first that for every choice of $`i,a,b,s_1,\mathrm{},s_{i1}`$,
$$\frac{M_{Sp,u}(\lambda _1^{}=s_1,\mathrm{},\lambda _{i1}^{}=s_{i1},\lambda _i^{}=a,\lambda _{i+1}^{}=b)}{M_{Sp,u}(\lambda _1^{}=s_1,\mathrm{},\lambda _{i1}^{}=s_{i1},\lambda _i^{}=a)}$$
is equal to the asserted transition rule probability for moving from $`\lambda _i^{}=a`$ to $`\lambda _{i+1}^{}=b`$, and second that the transition rule probabilities sum to one.
The first claim follows from Lemma 3. For the second claim, observe that
$$\underset{ba}{}\frac{M_{Sp,u}(\lambda _1^{}=s_1,\mathrm{},\lambda _{i1}^{}=s_{i1},\lambda _i^{}=a,\lambda _{i+1}^{}=b)}{M_{Sp,u}(\lambda _1^{}=s_1,\mathrm{},\lambda _{i1}^{}=s_{i1},\lambda _i^{}=a)}=1,$$
because $`M_{Sp,u}`$ is a measure and the columns of a partition are non-increasing in size as one moves to the right. Since $`_{i0}P_{Sp,u}(i)=1`$, it follows that $`M_{Sp,u}`$ is a probability measure, as promised earlier. $`\mathrm{}`$
Theorem 5 gives the analogous result for the orthogonal groups. As the proof method is the same as for the symplectic groups, we merely record the result.
###### Theorem 5
Starting with $`\lambda _1^{}`$ distributed as $`P_{O,u}`$, define in succession $`\lambda _2^{},\lambda _3^{},\mathrm{}`$ according to the rules that if $`\lambda _i^{}=a`$, then $`\lambda _{i+1}^{}=b`$ with probability $`K_2(a,b)`$ if $`i`$ is odd and $`K_1(a,b)`$ if $`i`$ is even. The resulting partition is distributed according to $`M_{O,u}`$.
We close the paper with the following remarks.
Remarks:
1. Theorem 4 and Theorem 5 allow one to draw exact samples from the measures $`M_{Sp,u}`$ or $`M_{O,u}`$. First recall that sampling from discrete distributions $`P`$ with known formulas is straightforward; simply pick $`U`$ uniformly in $`[0,1]`$ and find the value of $`j`$ such that $`_{i=0}^jP(i)<U<_{i=0}^{j+1}P(i)`$. This allows one to sample from $`P_{Sp,u}`$ or $`P_{O,u}`$. Then move according to the appropriate Markov chains.
2. For the case of $`M_{Sp,u}`$, one can view the algorithm of Theorem 4 slightly differently. One starts with an imaginary $`0`$th column of size approaching infinity, and then gets $`\lambda _1^{}`$ by transitioning according to the chain $`K_2`$. It is straightforward to verify that the resulting distribution of the first column size agrees with $`P_{Sp,u}`$. This viewpoint was useful in the general linear and unitary cases \[F4\].
3. As noted in the introduction, it is possible that the measures $`M_{Sp,u}`$ and $`M_{O,u}`$ are related to generalizations of the Rogers-Ramanujan identities, in analogy with the corresponding measures for $`GL(n,q)`$. In this regard observe that
$`{\displaystyle \underset{\lambda :\lambda _2^{}=0}{}}M_{Sp,u}(\lambda )`$ $`=`$ $`{\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}(1u^2/q^{2r1}){\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{u^{2n}}{q^{n^2}(q^{2n}1)\mathrm{}(q^21)}}`$
$`=`$ $`{\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}(1u^2/q^{2r1}){\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{u^{2n}}{q^{2n^2+n}(11/q^2)\mathrm{}(11/q^{2n})}}.`$
For the value $`u=q^{1/2}`$, the sum has a product expansion by a Rogers Ramanujan identity.
4. Although a probabilistic understanding of $`M_{Sp,u}`$ and $`M_{O,u}`$ has been given, a viewpoint explaining the products in Theorem 3 in terms of certain random variables being independent would be desirable. This was possible for the finite general linear and unitary groups \[F2\].
5. As mentioned in the introduction, the clumping of conjugacy classes given by looking at the underlying shape was not necessary. To modify things to take signs into account, $`K_1`$ starts at a number $`a`$ but outputs an ordered pair $`(b,\pm )`$, that is a choice of sign associated with $`b`$. In the $`+`$ (resp. $``$) case, the expression $`P_{O,u}^{}(b)`$ in the numerator of the definition of $`K_1`$ is replaced by the probability that the partition under the measure $`M_{O,u}`$ has $`b`$ parts and a $`+`$ (resp. $``$) choice for the parts of size 1. Unfortunately we do not know of simple formulas for these probabilities analogous to Theorem 3. Formulas can be inferred from the paper \[RShi\], but the result involves unpleasant sums and does not seem useful. The transition probabilities for $`K_2`$ are also affected: the $`P_{O,u}^{}(a)`$ are replaced the same way as for $`K_1`$, and letting $`ϵ`$ be the sign associated to $`a`$, the transition probabilities are multiplied by an additional factor of
$$\frac{1}{|O^ϵ(ab,q)|}\frac{1}{\frac{1}{|O^+(ab,q)|}+\frac{1}{O^{}(ab,q)|}}.$$
It is not necessarily surprising that the theory is nicer when conjugacy classes are lumped; the cycle index only factors for sums or differences of orthogonal groups.
## Acknowledgements
This work is a natural follow-up to the author’s Ph.D. thesis; the author expresses deep gratitude to Persi Diaconis for having introduced him to this lovely part of mathematics. A summer of conversations with Peter M. Neumann and Cheryl E. Praeger broadened the author’s understanding of conjugacy classes. This work received the financial support of an NSF Postdoctoral Fellowship. |
warning/0003/hep-ph0003233.html | ar5iv | text | # Exclusive production of pion pairs in 𝛾^∗𝛾 collisions at large 𝑄² Work supported by Department of Energy contract DE–AC03–76SF00515 and by TMR contracts FMRX–CT96–0008 and FMRX–CT98–0194
## I Introduction
Exclusive hadron production in two-photon collisions provides a tool to study a variety of fundamental aspects of QCD and has long been a subject of great interest (cf., e.g., and references therein). Recently a new facet of this has been pointed out, namely the physics of the process $`\gamma ^{}\gamma \pi \pi `$ in the region where $`Q^2`$ is large but $`W^2`$ small . This process factorizes into a perturbatively calculable, short-distance dominated scattering $`\gamma ^{}\gamma q\overline{q}`$ or $`\gamma ^{}\gamma gg`$, and non-perturbative matrix elements measuring the transitions $`q\overline{q}\pi \pi `$ and $`gg\pi \pi `$. We have called these matrix elements generalized distribution amplitudes (GDAs) to emphasize their close connection to the distribution amplitudes introduced many years ago in the QCD description of exclusive hard processes .
Indeed it is instructive to consider $`\gamma ^{}\gamma \pi \pi `$ as a generalization of the process $`\gamma ^{}\gamma \pi ^0`$, where the distribution amplitude of a single pion appears. The $`\gamma `$$`\pi `$ transition form factor has been the subject of detailed theoretical studies . The experimental data are well reproduced by a description based on QCD factorization and provide one of the best constraints so far on the form of the single-pion distribution amplitude.
From a different point of view $`\gamma ^{}\gamma \pi \pi `$ is the crossed channel of virtual Compton scattering on a pion. The kinematical region we consider here is closely related to deeply virtual Compton scattering (DVCS), which has attracted considerable attention in the context of skewed parton distributions .
Our reaction can also be seen as the exclusive limit of a hadronization process. The hadronization of a $`q\overline{q}`$-pair originating from a hard, short-distance process such as a $`\gamma ^{}\gamma `$ collision is usually formulated in terms of fragmentation functions which describe in a universal way semi-inclusive reactions, specifically the transition from a quark or antiquark to a final-state hadron when one integrates over all final states containing this hadron. We specialize here to the case where the final state consists of two mesons with specified four-momenta, and nothing else.
Like other hadronic matrix elements the GDAs are process independent. It has recently been pointed out that they occur in the hard exclusive process $`\gamma ^{}p\pi \pi p`$, where the pion pair is or is not the decay product of a $`\rho `$ meson, and that the analysis of that reaction would benefit from the measurement of the two-pion GDA in $`\gamma ^{}\gamma \pi \pi `$.
All these aspects lead us to consider GDAs as a promising new tool for hadronic physics, which may be used to unveil some of the mysteries of hadronization and the confining regime of QCD. The process $`\gamma ^{}\gamma \pi \pi `$ is well suited to access these quantities experimentally. In the present paper, we develop in detail the phenomenology of this reaction and emphasize the feasibility of its investigation at existing $`e^+e^{}`$ colliders.
In Sect. II we discuss the kinematics of our process, recall its main properties in the factorization regime we are interested in, and sketch the crossing relation between $`\gamma ^{}\gamma \pi \pi `$ and deep virtual Compton scattering. In Sect. III we list the general properties of generalized distribution amplitudes and in particular review their QCD evolution equations. These properties lead us to construct a simple model of the two-pion GDA, which is described in Sect. IV. Section V gives a comparison between one-pion and two-pion production in $`\gamma ^{}\gamma `$ collisions. Relations with the inclusive production of hadrons, commonly described by the photon structure function, are discussed in Sect. VI. The phenomenology of our process in $`e\gamma `$ collisions is described in detail in Sect. VII, with special emphasis on the information contained in angular distributions and in the interference with the bremsstrahlung mechanism. In Sect. VIII we give estimates for the cross section for various experimental setups at existing $`e^+e^{}`$ colliders. Section IX contains our conclusions. In Appendix A we specify our sign conventions for pion states, and in Appendix B we discuss what additional information can be obtained with polarized beams.
## II The process $`\gamma ^{}\gamma \pi \pi `$
### A Kinematics in the $`\gamma ^{}\gamma `$ center of mass.
The reaction we are interested in is
$$e(k)+\gamma (q^{})e(k^{})+\pi ^i(p)+\pi ^j(p^{}),$$
(1)
where four-momenta are indicated in parentheses. We further use
$$q=kk^{},Q^2=q^2,P=p+p^{},W^2=P^2.$$
(2)
The pions may be charged ($`i=+`$, $`j=`$) or neutral ($`i=j=0`$), and the lepton $`e`$ may be an electron or a positron. Scattered with large momentum transfer this lepton radiates a virtual photon $`\gamma ^{}(q)`$, and for the $`\gamma ^{}\gamma `$ subprocess we introduce the Bjorken variable
$$x=\frac{Q^2}{2qq^{}}=\frac{Q^2}{Q^2+W^2}.$$
(3)
In $`e^+e^{}`$ collisions the photon $`\gamma (q^{})`$ can be obtained by bremsstrahlung from the other beam lepton, so that the overall process is
$$e(k)+e(l)e(k^{})+e(l^{})+\pi ^i(p)+\pi ^j(p^{})$$
(4)
with $`q^{}=ll^{}`$. In the spirit of the equivalent photon approximation we approximate $`q^2`$ as zero and the momenta $`q^{}`$ and $`l`$ as collinear. We write $`E_1=k^0`$, $`E_2=l^0`$ and $`q^0=x_2l^0`$ for the energies in the laboratory frame.<sup>*</sup><sup>*</sup>* We neglect the small finite crossing angle between the beams at BELLE, so that in our parlance the lepton beams are collinear in the “laboratory frame”. For the c.m. energies of the $`ee`$ and $`e\gamma `$ collisions we have
$$s_{ee}=(k+l)^2,s_{e\gamma }=(k+q^{})^2=x_2s_{ee}.$$
(5)
Let us now discuss the kinematics in the $`\gamma ^{}\gamma `$ center of mass frame. We use a coordinate system with the $`z`$ axis along $`𝐪`$, and with $`x`$ and $`y`$ axes such that $`𝐩`$ lies in the $`x`$-$`z`$ plane and has a positive $`x`$ component, i.e.,
$$q=(q^0,0,0,|𝐪|),p=(p^0,|𝐩|\mathrm{sin}\theta ,0,|𝐩|\mathrm{cos}\theta ),$$
(6)
where we have introduced the polar angle $`\theta `$ of $`𝐩`$. Another natural variable for our process in this frame is the azimuth $`\phi `$ of $`𝐤^{}`$, which is the angle between the leptonic and the hadronic planes, cf. Fig. 1. In terms of Lorentz invariants these angles can be obtained from
$`\mathrm{cos}\theta `$ $`=`$ $`{\displaystyle \frac{2q(p^{}p)}{\beta (Q^2+W^2)}},`$ (7)
$`\mathrm{cos}\phi `$ $`=`$ $`{\displaystyle \frac{2k(p^{}p)(Q^2+W^2)+\beta \mathrm{cos}\theta [Q^2(s_{e\gamma }Q^2W^2)s_{e\gamma }W^2]}{2\beta \mathrm{sin}\theta \sqrt{s_{e\gamma }Q^2W^2(s_{e\gamma }Q^2W^2)\text{}}}},`$ (8)
$`\mathrm{sin}\phi `$ $`=`$ $`{\displaystyle \frac{4ϵ_{\alpha \beta \gamma \delta }(p+p^{})^\alpha p^\beta k^\gamma q^\delta }{\beta \mathrm{sin}\theta \sqrt{s_{e\gamma }Q^2W^2(s_{e\gamma }Q^2W^2)\text{}}}}`$ (9)
with $`ϵ_{0123}=+1`$ and the velocity
$$\beta =\sqrt{1\frac{4m_\pi ^2}{W^2}}$$
(10)
of the pions. A further quantity we will use is the usual $`y`$-variable for the $`e\gamma `$ collision,
$$y=\frac{qq^{}}{kq^{}}=\frac{Q^2+W^2}{s_{e\gamma }},$$
(11)
which can be traded for
$$ϵ=\frac{1y}{1y+y^2/2},$$
(12)
the ratio of longitudinal to transverse polarization of the virtual photon $`\gamma ^{}(q)`$.
We finally define light cone components $`a^\pm =(a^0\pm a^3)/\sqrt{2}`$ for any four-vector $`a`$ and introduce the fraction
$$\zeta =\frac{p^+}{P^+}=\frac{1+\beta \mathrm{cos}\theta }{2}$$
(13)
of light cone momentum carried by $`\pi ^i(p)`$ with respect to the pion pair.
### B Factorization at large $`Q^2`$ and small $`W^2`$
Let us briefly review how $`\gamma ^{}\gamma \pi \pi `$ factorizes in the kinematical regime we are interested in. Firstly, we require $`Q^2`$ to be large compared with the scale $`\mathrm{\Lambda }^21`$ GeV<sup>2</sup> of soft interactions, thus providing a hard scale for the process. Secondly, we ask $`W^2`$ to be small compared with this large scale $`Q^2`$. In this regime the dynamics of the process is conveniently represented in the Breit frame, obtained by boosting from the $`\gamma ^{}\gamma `$ center of mass along the $`z`$ axis. The spacetime cartoon of the process one can derive from power counting and factorization arguments is shown in Fig. 2.
In the Breit frame the real photon moves fast in the negative $`z`$ direction and is scattered into an energetic hadronic system moving in the positive $`z`$ direction. The hard part of this process takes place at the level of elementary constituents, and the minimal number of quarks and gluons compatible with conservation laws (color etc.) are produced. At Born level one simply has $`\gamma ^{}\gamma q\overline{q}`$, but through a quark box the photons can also couple to two gluons. Each quark or gluon carries a fraction $`z`$ or $`1z`$ of the large light-cone momentum component $`P^+`$. Subsequently the soft part of the reaction, i.e., hadronization into a pion pair, takes place.
At leading order in $`\alpha _S`$ the amplitude is given by the diagram of Fig. 3 (a) and the one where the two photon vertices are interchanged. One calculates for the hadronic tensor
$$T^{\mu \nu }=id^4xe^{iqx}\pi (p)\pi (p^{})|TJ_{\mathrm{em}}^\mu (x)J_{\mathrm{em}}^\nu (0)|0=g_T^{\mu \nu }\underset{q}{}\frac{e_q^2}{2}_0^1𝑑z\frac{2z1}{z(1z)}\mathrm{\Phi }_q^{\pi \pi }(z,\zeta ,W^2),$$
(14)
where $`g_T^{\mu \nu }`$ denotes the metric tensor in transverse space ($`g^{11}=1`$). The sum on the r.h.s. runs over all quarks flavors, $`e_q`$ is the charge of quark $`q`$ in units of the positron charge $`e`$, and $`eJ_{\mathrm{em}}^\mu (x)`$ is the electromagnetic current. While the expression of the hard subprocess $`\gamma ^{}\gamma q\overline{q}`$ is explicit in Eq. (14), the soft part of $`\gamma ^{}\gamma \pi \pi `$ is parameterized by the generalized distribution amplitude
$$\mathrm{\Phi }_q^{\pi \pi }(z,\zeta ,W^2)=\frac{dx^{}}{2\pi }e^{iz(P^+x^{})}\pi (p)\pi (p^{})|\overline{q}(x^{})\gamma ^+q(0)|0$$
(15)
for each quark flavor $`q`$. We work in light cone gauge $`A^+=0`$, otherwise the usual path ordered exponential of gluon potentials appears between the quark fields. $`\mathrm{\Phi }_q`$ depends on the light-cone fraction $`z`$ of the quark with respect to the pion pair, on the kinematical variables $`\zeta `$ and $`W^2`$ of the pions, and on a factorization scale. The latter dependence, not displayed in Eq. (15), will be discussed in Sect. III B.
In Eq. (14) a scaling behavior for our process is manifest: at fixed $`\zeta `$ and $`W^2`$ the $`\gamma ^{}\gamma `$ amplitude is independent of $`Q^2`$, up to logarithmic scaling violations from radiative corrections to the hard scattering and from the evolution of the two-pion distribution amplitude. This scaling property is central to all processes where a factorization theorem holds, and it is the basic signature one looks for when testing whether the asymptotic analysis developed here applies to an experimental situation at finite $`Q^2`$. There will of course be power corrections in $`\mathrm{\Lambda }/Q`$ and $`W/Q`$ to this leading mechanism. Examples are the hadronic component of the real photon, and the effect in the hard scattering of the transverse momentum of the produced parton pair. We note that the crossed channel, i.e., virtual Compton scattering has been analyzed in detail within the operator product expansion , which may provide a framework for a systematic study of higher twist effects.
Contracting the hadronic tensor (14) with the photon polarization vectors we see that in order to give a nonzero $`\gamma ^{}\gamma \pi \pi `$ amplitude the virtual photon must have the same helicity as the real one. As in the case of deep virtual Compton scattering this is a direct consequence of chiral invariance in the collinear hard-scattering process and is valid at all orders in $`\alpha _S`$. In the case of the $`\gamma ^{}\gamma gg`$ subprocess the photon helicities can also be opposite . In any case the virtual photon must be transverse. As a consequence nonleading twist effects can be studied in the amplitude for longitudinal $`\gamma ^{}`$ polarization, without any “background” from leading twist pieces. We will develop in Sect. VII how the different $`\gamma ^{}\gamma `$ helicity amplitudes are experimentally accessible.
As we already mentioned, there is a close analogy of two-pion production in the region $`Q^2W^2,\mathrm{\Lambda }^2`$ with the one-pion channel, commonly described in terms of the $`\gamma `$$`\pi `$ transition form factor. There again a factorization theorem holds, which allows the hadronic tensor $`T^{\mu \nu }`$ to be expressed in terms of the single-pion distribution amplitude $`\varphi ^\pi `$ as
$$T^{\mu \nu }=id^4xe^{iqx}\pi ^0|TJ_{\mathrm{em}}^\mu (x)J_{\mathrm{em}}^\nu (0)|0=ϵ_T^{\mu \nu }\underset{q}{}\frac{e_q^2}{2}_0^1𝑑z\frac{1}{z(1z)}\varphi _q^\pi (z)$$
(16)
to leading order in $`\alpha _S`$, where $`ϵ_T^{\mu \nu }`$ is the antisymmetric tensor in transverse space ($`ϵ_T^{12}=1`$) and
$$\varphi _q^\pi (z)=i\frac{dx^{}}{2\pi }e^{iz(P^+x^{})}\pi ^0(P)|\overline{q}(x^{})\gamma ^+\gamma _5q(0)|0.$$
(17)
Notice the different Dirac structures in the matrix elements (15) and (17), due to the different parity transformation properties of one- and two-pion states .
The theoretical analysis of this process has been highly developed . Its generalization to the production of $`\eta `$ and $`\eta ^{}`$ is also important, in particular with respect to the $`SU(3)`$ flavor structure of the QCD evolution equations and the mixing of the quark singlet and gluon channels . In Sect. V we will further compare the production of a single pion with that of a pion pair.
### C Relation with deep virtual Compton scattering and parton distributions in the pion
The process $`\gamma ^{}\gamma \pi \pi `$ at large $`Q^2`$ and $`sQ^2`$ is related by $`s`$$`t`$ crossing to deep virtual Compton scattering on a pion, i.e., to $`\gamma ^{}\pi \gamma \pi `$ at large $`Q^2`$ and $`tQ^2`$. It turns out that factorization works in completely analogous ways for both cases, as is shown in Fig. 3. The non-perturbative matrix elements occurring in the Compton process are skewed parton distributions , defined in the pion case as
$$H_q(x,\xi ,t)=\frac{1}{2}\frac{dz^{}}{2\pi }e^{ix(P^+z^{})}\pi (p^{})|\overline{q}(z^{}/2)\gamma ^+q(z^{}/2)|\pi (p)$$
(18)
with $`P=(p+p^{})/2`$. They have been recognized as objects of considerable interest and have triggered intensive theoretical and experimental work. The processes $`\gamma ^{}\gamma \pi \pi `$ and $`\gamma ^{}\pi \gamma \pi `$ share many common features, from their scaling behavior and the details of their helicity selection rules to the possibilities of phenomenological analysis, which we will develop in Sect. VII.
The imaginary part of the forward virtual Compton amplitude, $`\gamma ^{}\pi \gamma ^{}\pi `$, obtained from Fig. 3 (b) by replacing the $`\gamma `$ with a second $`\gamma ^{}`$, gives the cross section for inclusive deep inelastic scattering, $`\gamma ^{}\pi X`$, where the ordinary parton distributions in a pion occur.
As observed in it is useful to implement crossing at the level of moments in momentum fractions ($`z`$ and $`\zeta `$ for GDAs, $`x`$ and $`\xi `$ for SPDs), which depend only on a factorization scale and a Lorentz invariant ($`s`$ for GDAs, $`t`$ for SPDs). The moments of GDAs and of SPDs are connected by analytic continuation in that invariant. In particular, analytic continuation to the point $`t=0`$ leads to moments of the ordinary parton distributions in the pion, which we will use as an input for our model of GDAs in Sect. IV.
## III General properties of GDAs
### A Charge conjugation and isospin properties
Let us start by compiling some symmetry properties which will be useful in the following. For the quark GDAs (15) the invariance of strong interactions under charge conjugation $`C`$ implies
$$\mathrm{\Phi }_q^{\pi \pi }(z,\zeta ,W^2)=\mathrm{\Phi }_q^{\pi \pi }(1z,1\zeta ,W^2).$$
(19)
It is useful to project GDAs for charged pions on eigenstates of $`C`$ parity,
$$\mathrm{\Phi }_q^\pm (z,\zeta ,W^2)=\frac{1}{2}\left(\mathrm{\Phi }_q^{\pi ^+\pi ^{}}(z,\zeta ,W^2)\pm \mathrm{\Phi }_q^{\pi ^+\pi ^{}}(z,1\zeta ,W^2)\right),$$
(20)
so that
$$\mathrm{\Phi }_q^{\pi ^+\pi ^{}}=\mathrm{\Phi }_q^+(z,\zeta ,W^2)+\mathrm{\Phi }_q^{}(z,\zeta ,W^2).$$
(21)
In the $`C`$ even sector Eq. (19) reduces to
$$\mathrm{\Phi }_q^+(z,\zeta ,W^2)=\mathrm{\Phi }_q^+(1z,\zeta ,W^2).$$
(22)
Our process is only sensitive to the $`C`$ even part of $`\mathrm{\Phi }_q^{\pi ^+\pi ^{}}`$ since the initial state two-photon state has positive $`C`$ parity. Of course a $`\pi ^0\pi ^0`$ pair has positive $`C`$ parity as well, so that $`\mathrm{\Phi }_q^{\pi ^0\pi ^0}`$ has no $`C`$-odd part at all.
Let us now turn to isospin symmetry. The $`C`$ odd component of a two-pion state has total isospin $`I=1`$, whereas its $`C`$ even component contains both $`I=0`$ and $`I=2`$ pieces. The quark operator in $`\mathrm{\Phi }_q^{\pi \pi }`$ has only components with isospin $`I=0`$ or $`I=1`$. Hence it is a consequence of the leading twist production mechanism and of isospin invariance that in our process the pion pair is in a state of zero isospin, i.e., that no component with $`I=2`$ is produced. Another consequence of isospin invariance is that
$$\mathrm{\Phi }_q^{\pi ^0\pi ^0}(z,\zeta ,W^2)=\mathrm{\Phi }_q^+(z,\zeta ,W^2),$$
(23)
so that the production amplitudes for neutral and charged pion pairs are equal. Deviations from isospin symmetry in the present reaction would be interesting, but since one can expect them to be small we will assume isospin invariance to hold throughout the rest of our study. Isospin invariance also implies that
$$\mathrm{\Phi }_u^+=\mathrm{\Phi }_d^+,\mathrm{\Phi }_u^{}=\mathrm{\Phi }_d^{},$$
(24)
so that in the $`C`$ even sector we only need to know the $`SU(2)`$ flavor singlet combination $`\mathrm{\Phi }_u^++\mathrm{\Phi }_d^+`$.
The connection between the notation $`\mathrm{\Phi }_{||}^{I=0,1}`$ of Polyakov and ours is
$$\mathrm{\Phi }_{||}^{I=0}=\mathrm{\Phi }_u^+,\mathrm{\Phi }_{||}^{I=1}=\mathrm{\Phi }_u^{}.$$
(25)
We remark that the second term in Eq. (2.6) of Ref. should come with a minus sign . Our relation $`\mathrm{\Phi }_{||}^{I=1}=\mathrm{\Phi }_u^{}`$ takes this correction into account.
Notice that the signs in Eqs. (23) and (25) depend on the choice of relative phases in the definition of charged pion states. We specify our convention in Appendix A.
### B Evolution equation
In the process of factorization generalized distribution amplitudes acquire a scale dependence in the same way as usual distributions do. This scale dependence can be computed within perturbative QCD, and there is nothing special with multiparticle states since the scale dependence is a property of the nonlocal product of fields under consideration, rather than one of a particular hadronic matrix element (see for an approach exploiting this feature). The scale dependence of GDAs can be cast in the form of an ERBL evolution equation , and the only complication in the channel we are concerned with here is the mixing of quark and gluon distribution amplitudes. The leading-logarithmic form of the evolution equations has been studied in detail for the parity-odd sector , where the relevant quark operator is $`\overline{q}\gamma ^+\gamma _5q`$. Our application to pion pairs leads us to consider the parity-even sector, where the quark operator is $`\overline{q}\gamma ^+q`$ instead, see our remark after Eq. (17). For completeness we give here the basic steps for deriving and solving the evolution equation in this channel, following the procedure outlined in . Taking into account the different normalization conventions we find agreement with the results of Baier and Grozin , who reported a sign discrepancy with Chase for the gluon evolution kernel.
We are then concerned with the generalized quark and gluon distribution amplitudes in $`A^+=0`$ gauge:
$`\mathrm{\Phi }_q(z,\zeta ,W^2)`$ $`=`$ $`{\displaystyle \frac{dx^{}}{2\pi }e^{iz(P^+x^{})}\pi (p)\pi (p^{})|\overline{q}(x^{})\gamma ^+q(0)|0},`$ (26)
$`\mathrm{\Phi }_g(z,\zeta ,W^2)`$ $`=`$ $`{\displaystyle \frac{1}{P^+}}{\displaystyle \frac{dx^{}}{2\pi }e^{iz(P^+x^{})}\pi (p)\pi (p^{})|F^{+\mu }(x^{})F_{\mu }^{}{}_{}{}^{+}(0)|0},`$ (27)
$`=`$ $`z(1z)P^+{\displaystyle \frac{dx^{}}{2\pi }e^{iz(P^+x^{})}\pi (p)\pi (p^{})|A^\mu (x^{})A_\mu (0)|0}.`$ (28)
Our gluon distribution amplitude $`\mathrm{\Phi }_g(z,\zeta ,W^2)`$ coincides with $`\mathrm{\Phi }^G(z,\zeta ,W^2)`$ introduced in . From the definition (26) one readily obtains
$$\mathrm{\Phi }_g(z,\zeta ,W^2)=\mathrm{\Phi }_g(1z,\zeta ,W^2),$$
(29)
and from $`C`$ invariance one has
$$\mathrm{\Phi }_g(z,\zeta ,W^2)=\mathrm{\Phi }_g(1z,1\zeta ,W^2).$$
(30)
Here we have given definitions for a two-pion state, but as stated above the evolution equation for DAs and GDAs is not specific to the details of the hadronic system. The considerations of this and the following subsection thus apply to any state in the parity even sector which has four-momentum $`P`$ and total angular momentum $`J_z=0`$ along the axis defining the light cone variables.
We now study the evolution of the distributions for gluons and of quarks in the singlet combination of $`n_f`$ flavors. For convenience we introduce
$`z\overline{z}f_Q(z)`$ $`=`$ $`{\displaystyle \underset{q=1}{\overset{n_f}{}}}\mathrm{\Phi }_q(z),`$ (31)
$`z^2\overline{z}^2f_G(z)`$ $`=`$ $`\mathrm{\Phi }_g(z),`$ (32)
where we use the notation $`\overline{z}=1z`$. In the end we will return to the amplitudes $`\mathrm{\Phi }_q`$ and $`\mathrm{\Phi }_g`$.
The scale dependence is controlled by the parameter
$$\xi (\mu ^2,\mu _0^2)=\frac{2}{\beta _1}\mathrm{ln}\left(\frac{\alpha _S(\mu _0^2)}{\alpha _S(\mu ^2)}\right),$$
(33)
where $`\alpha _S`$ is the one-loop running coupling and $`\beta _1=112n_f/3`$. This parameter describes how the distribution amplitude evolves when one changes the factorization point from $`\mu _0`$ to $`\mu `$. The evolution equation takes the form
$$\frac{}{\xi }f(z,\xi )=Vf=_0^1𝑑uV(z,u)f(u,\xi ).$$
(34)
where $`f`$ is a two-component vector
$$f=\left(\begin{array}{c}f_Q\\ f_G\end{array}\right),$$
(35)
and $`V`$ is the $`2\times 2`$ matrix kernel
$$V=\left(\begin{array}{cc}V_{QQ}& V_{QG}\\ V_{GQ}& V_{GG}\end{array}\right).$$
(36)
To obtain the leading logarithmic evolution equation it is sufficient to consider one-loop corrections to the scattering amplitude. The latter is depicted in Fig. 4 and has the form $`Hf`$, where $`H=(H_Q,H_G)`$ denotes the hard-scattering kernels. It turns out that in light cone gauge $`A^+=0`$ the relevant one-loop diagrams consist of an insertion between $`H`$ and $`f`$ of the graphs shown in Fig. 5 (a) to (e), supplemented by (renormalized) self-energy insertions on each line connecting $`H`$ to $`f`$ in Fig. 4. Calling the sum of these insertions $`\xi V`$ the one-loop diagrams have the structure $`H\xi Vf`$.
The evolution from zeroth to first order of the generalized distribution amplitude may thus be written as
$$f^{(1)}(z)=f^{(0)}(z)+\xi _0^1𝑑uV(z,u)f^{(0)}(u).$$
(37)
In the computation of the diagrams, the $`\kappa ^{}`$ integral is performed by the Cauchy method of contour integration in the complex plane, and $`\xi `$ is the result of the integral over transverse momentum from $`\kappa _T=\mu _0`$ to $`\kappa _T=\mu `$:
$$\xi (\mu ^2,\mu _0^2)=_{\mu _0^2}^{\mu ^2}\frac{d\kappa _T^2}{\kappa _T^2}\frac{\alpha _S(\kappa _T^2)}{2\pi }.$$
(38)
Despite the presence of $`\alpha _S`$ in Eq. (38), $`\xi `$ is not small if $`\mu ^2\mu _0^2`$, and this signals the necessity of an all-order analysis. This analysis leads to the evolution equation, with the feature that $`V`$ is the same matrix in Eqs. (34) and (37). We refer the reader to the literature for a general discussion .
The integration over $`\kappa ^+`$ may be reexpressed as an integral over the incoming light cone fraction $`u`$. The evolution kernels contain the remaining part of the dynamics, in particular they describe the change of light cone fractions from $`u`$ to $`z`$. We get
$`V_{QQ}(z,u)`$ $`=`$ $`C_F[\theta (zu){\displaystyle \frac{u}{z}}(1+{\displaystyle \frac{1}{zu}})+\{u\overline{u},z\overline{z}\}]_+,`$ (39)
$`V_{QG}(z,u)`$ $`=`$ $`2n_fT_F[\theta (zu){\displaystyle \frac{u}{z}}(2zu)\{u\overline{u},z\overline{z}\}],`$ (40)
$`V_{GQ}(z,u)`$ $`=`$ $`{\displaystyle \frac{C_F}{z\overline{z}}}[\theta (zu){\displaystyle \frac{u}{z}}(\overline{z}2\overline{u})\{u\overline{u},z\overline{z}\}],`$ (41)
$`V_{GG}(z,u)`$ $`=`$ $`{\displaystyle \frac{C_A}{z\overline{z}}}[\theta (zu)({\displaystyle \frac{u\overline{u}}{zu}}u\overline{u}{\displaystyle \frac{u}{2z}}[(2z1)^2+(2u1)^2])+\{u\overline{u},z\overline{z}\}]_+{\displaystyle \frac{2}{3}}n_fT_F\delta (uz),`$ (42)
where the color factors are $`C_F=4/3`$, $`T_F=1/2`$ and $`C_A=3`$. The subscript $`+`$ stands for the $`+`$ distributions, whose action on a function $`f`$ may be expressed symbolically as
$$[\mathrm{}]_+f(u)=[\mathrm{}](f(u)f(z)).$$
(43)
The kernels (42) give the finite parts that remain after the cancellation of infrared divergences between graph (a), resp. (d), and quark self-energy, resp. gluon self-energy insertions. A simple way to obtain self-energy corrections is to notice their relation to parton splitting , that is
$`f_Q^{(1)}(z)|_{SE}`$ $`=`$ $`\left[1\xi {\displaystyle 𝑑xP_{QQ}(x)}\right]f_Q^{(0)}(z)=\left[1\xi {\displaystyle 𝑑xP_{GQ}(x)}\right]f_Q^{(0)}(z)`$ (44)
$`f_G^{(1)}(z)|_{SE}`$ $`=`$ $`\left[1\xi {\displaystyle 𝑑x\left(\frac{1}{2}P_{GG}(x)+n_fP_{QG}(x)\right)}\right]f_G^{(0)}(z),`$ (45)
with the unregularized DGLAP splitting functions
$`P_{QQ}(x)`$ $`=`$ $`C_F{\displaystyle \frac{1+x^2}{1x}},`$ (46)
$`P_{QG}(x)`$ $`=`$ $`T_F\left[x^2+(1x)^2\right],`$ (47)
$`P_{GQ}(x)`$ $`=`$ $`C_F{\displaystyle \frac{1+(1x)^2}{x}},`$ (48)
$`P_{GG}(x)`$ $`=`$ $`2C_A\left[{\displaystyle \frac{x}{1x}}+{\displaystyle \frac{1x}{x}}+x(1x)\right].`$ (49)
The integrals (45) are not defined in the limit $`x0,1`$, which is a manifestation of the infrared divergence of self-energy graphs.
### C Solution
We will now solve the evolution equation (34). Given our application we restrict ourselves to the $`C`$ even parts $`\mathrm{\Phi }_q^+`$ of the quark distributions, the gluon distribution being of course even under $`C`$ from the start.
We look for solutions of the form
$$f(z,\xi )=f(z)e^{\gamma \xi }.$$
(50)
To this end it is convenient to change variables, introducing $`y=2u1`$ and $`x=2z1`$, and to study the convolution of the matrix kernel $`V`$ with
$$\left(\begin{array}{c}x^n\\ 0\end{array}\right),\left(\begin{array}{c}0\\ x^{n1}\end{array}\right),$$
(51)
where $`n`$ is an odd integer to accommodate the symmetry properties (22) and (29). One finds
$`V_{QQ}y^n=\gamma _{QQ}(n)x^n+O(x^{n2}),V_{QG}y^{n1}=\gamma _{QG}(n)x^n+O(x^{n2}),`$ (52)
$`V_{GQ}y^n=\gamma _{GQ}(n)x^{n1}+O(x^{n3}),V_{GG}y^{n1}=\gamma _{GG}(n)x^{n1}+O(x^{n3}),`$ (53)
with anomalous dimensions
$`\gamma _{QQ}(n)`$ $`=`$ $`C_F\left({\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{(n+1)(n+2)}}+2{\displaystyle \underset{k=2}{\overset{n+1}{}}}{\displaystyle \frac{1}{k}}\right),`$ (54)
$`\gamma _{QG}(n)`$ $`=`$ $`n_fT_F{\displaystyle \frac{n^2+3n+4}{n(n+1)(n+2)}},`$ (55)
$`\gamma _{GQ}(n)`$ $`=`$ $`2C_F{\displaystyle \frac{n^2+3n+4}{(n+1)(n+2)(n+3)}},`$ (56)
$`\gamma _{GG}(n)`$ $`=`$ $`C_A\left({\displaystyle \frac{1}{6}}{\displaystyle \frac{2}{n(n+1)}}{\displaystyle \frac{2}{(n+2)(n+3)}}+2{\displaystyle \underset{k=2}{\overset{n+1}{}}}{\displaystyle \frac{1}{k}}\right)+{\displaystyle \frac{2}{3}}n_fT_F.`$ (57)
Since for a given $`n_0`$ the space of solutions with $`nn_0`$ is stable under the application of the kernel one can find polynomials $`p_n(x)`$ and $`q_{n1}(x)`$ satisfying
$`V_{QQ}p_n=\gamma _{QQ}(n)p_n,V_{QG}q_{n1}=\gamma _{QG}(n)p_n,`$ (58)
$`V_{GQ}p_n=\gamma _{GQ}(n)q_{n1},V_{GG}q_{n1}=\gamma _{GG}(n)q_{n1}.`$ (59)
The symmetry properties of the kernels
$`(1x^2)V_{QQ}(x,y)`$ $`=`$ $`(1y^2)V_{QQ}(y,x),`$ (60)
$`2C_F(1x^2)V_{QG}(x,y)`$ $`=`$ $`n_fT_F(1y^2)^2V_{GQ}(y,x),`$ (61)
$`(1x^2)^2V_{GG}(x,y)`$ $`=`$ $`(1y^2)^2V_{GG}(y,x),`$ (62)
then imply that the $`(p_n)`$ are orthogonal polynomials on the interval $`[1,\mathrm{\hspace{0.17em}1}]`$ with weight $`1x^2`$, i.e., they are proportional to the Gegenbauer polynomials $`C_n^{(3/2)}(x)`$, whereas the $`(q_{n1})`$ are orthogonal on $`[1,\mathrm{\hspace{0.17em}1}]`$ with weight $`(1x^2)^2`$, that is, proportional to the Gegenbauer polynomials $`C_{n1}^{(5/2)}(x)`$. To complete the identification it is necessary to take into account the standard normalization of Gegenbauer polynomials. One finds that $`p_n=C_n^{(3/2)}`$ and $`q_{n1}=C_{n1}^{(5/2)}`$ fulfill Eq. (59), provided one makes the replacements
$$\gamma _{QG}^{}(n)\gamma _{QG}^{}(n)=\frac{n}{3}\gamma _{QG}^{}(n),\gamma _{GQ}^{}(n)\gamma _{GQ}^{}(n)=\frac{3}{n}\gamma _{GQ}^{}(n).$$
(63)
The final step is to diagonalize the $`2\times 2`$ anomalous dimension matrices for each value of $`n`$. The eigenvalues are
$$\mathrm{\Gamma }_n^{(\pm )}=\frac{1}{2}\left[\gamma _{QQ}^{}(n)+\gamma _{GG}^{}(n)\pm \sqrt{[\gamma _{QQ}^{}(n)\gamma _{GG}^{}(n)]^2+4\gamma _{QG}^{}(n)\gamma _{GQ}^{}(n)}\right],$$
(64)
and the eigenvectors of the kernel matrix read
$$v_n^{(\pm )}(x)=\left(\begin{array}{c}C_n^{(3/2)}(x)\\ g_n^{(\pm )}C_{n1}^{(5/2)}(x)\end{array}\right),$$
(65)
where
$$g_n^{(\pm )}=\frac{\mathrm{\Gamma }_n^{(\pm )}\gamma _{QQ}(n)}{\gamma _{QG}^{}(n)}.$$
(66)
The general $`C`$ even solution of Eq. (34) may then be written as
$$f(x,\xi )=\underset{\mathrm{odd}n}{}\left\{A_n^{(+)}v_n^{(+)}(x)e^{\mathrm{\Gamma }_n^{(+)}\xi }+A_n^{()}v_n^{()}(x)e^{\mathrm{\Gamma }_n^{()}\xi }\right\}$$
(67)
with integration constants $`A_n^{(\pm )}`$.
We now return to the amplitudes $`\mathrm{\Phi }_q`$, $`\mathrm{\Phi }_g`$ and explicitly express $`\xi `$ in terms of $`\mu `$ and $`\mu _0`$. The key result of this section then reads
$`{\displaystyle \underset{q=1}{\overset{n_f}{}}}\mathrm{\Phi }_q^+(z,\mu ^2)`$ $`=`$ $`z(1z){\displaystyle \underset{\mathrm{odd}n}{}}A_n(\mu ^2)C_n^{(3/2)}(2z1),`$ (68)
$`\mathrm{\Phi }_g(z,\mu ^2)`$ $`=`$ $`z^2(1z)^2{\displaystyle \underset{\mathrm{odd}n}{}}A_n^{}(\mu ^2)C_{n1}^{(5/2)}(2z1),`$ (69)
with
$`A_n(\mu ^2)`$ $`=`$ $`A_n^{(+)}\left({\displaystyle \frac{\alpha _S(\mu ^2)}{\alpha _S(\mu _0^2)}}\right)^{K_n^{(+)}}+A_n^{()}\left({\displaystyle \frac{\alpha _S(\mu ^2)}{\alpha _S(\mu _0^2)}}\right)^{K_n^{()}},`$ (70)
$`A_n^{}(\mu ^2)`$ $`=`$ $`g_n^{(+)}A_n^{(+)}\left({\displaystyle \frac{\alpha _S(\mu ^2)}{\alpha _S(\mu _0^2)}}\right)^{K_n^{(+)}}+g_n^{()}A_n^{()}\left({\displaystyle \frac{\alpha _S(\mu ^2)}{\alpha _S(\mu _0^2)}}\right)^{K_n^{()}},`$ (71)
and exponents $`K_n^{(\pm )}=2\mathrm{\Gamma }_n^{(\pm )}/\beta _1`$, which are positive except for $`K_1^{()}=0`$. For $`n_f=2,3,4`$, one explicitly finds
$$K_1^{(+)}=\frac{32+6n_f}{996n_f}=0.51,0.62,0.75,K_3^{()}=0.71,0.76,0.82,K_3^{(+)}=1.45,1.64,1.85.$$
(72)
¿From Eq. (71) we easily see that the integration constants $`A_n^{(\pm )}`$ depend on the starting scale $`\mu _0`$ of the evolution through a factor $`\alpha _S(\mu _0^2)^{K_n^{(\pm )}}`$.
### D Expansion in $`\zeta `$
For a two-meson state, the coefficients $`A_n^{}`$ and $`A_n^{}`$ are functions of the factorization scale $`\mu ^2`$ and of the remaining kinematical variables $`\zeta `$ and $`W^2`$. From the definition of GDAs in term of fields given in Eq. (26) one obtains moments
$`{\displaystyle _0^1}𝑑zz^n\mathrm{\Phi }_q(z)`$ $`=`$ $`{\displaystyle \frac{1}{(P^+)^{n+1}}}\left[(i^+)^n\pi (p)\pi (p^{})|\overline{q}(x)\gamma ^+q(0)|0\right]_{x=0},`$ (73)
$`{\displaystyle _0^1}𝑑zz^{n1}\mathrm{\Phi }_g(z)`$ $`=`$ $`{\displaystyle \frac{1}{(P^+)^{n+1}}}\left[(i^+)^{n1}\pi (p)\pi (p^{})|F^{+\mu }(x)F_{\mu }^{}{}_{}{}^{+}(0)|0\right]_{x=0}.`$ (74)
These local matrix elements are the plus-components of tensors that can be decomposed on a basis built up with the metric $`g^{\mu \nu }`$ and the vectors $`(p+p^{})^\mu `$ and $`(pp^{})^\mu `$. Since $`(p+p^{})^+=P^+`$ and $`(pp^{})^+=(2\zeta 1)P^+`$ the moments (74) are then polynomials in $`2\zeta 1`$ with degree at most $`n+1`$. The $`A_n^{}`$ and $`A_n^{}`$ are Gegenbauer moments of $`_q\mathrm{\Phi }_q`$ and $`\mathrm{\Phi }_g`$, respectively, and therefore have the same polynomiality properties in $`\zeta `$. Following we expand them on the Legendre polynomials, writing
$$A_n(\zeta ,W^2)=6n_f\underset{\mathrm{even}l}{\overset{n+1}{}}B_{nl}(W^2)P_l(2\zeta 1)$$
(75)
and the analogous expression for $`A_n^{}`$ with coefficients $`B_{nl}^{}`$. The $`C`$ invariance properties (19) and (30) restrict $`l`$ to even integers in the $`C`$ even sector. The expansion coefficients $`B_{nl}`$ are linear combinations of the local operator matrix elements in Eq. (74) and are therefore analytic functions in $`W^2`$. As we mentioned in Sect. II C their continuation to zero or spacelike $`W^2`$ leads to the moments of parton distributions in the pion.
From Eq. (71) the factorization scale dependence of the $`B_{nl}`$ may be written as
$$B_{nl}(W^2,\mu ^2)=B_{nl}^{(+)}(W^2)\left(\frac{\alpha _S(\mu ^2)}{\alpha _S(\mu _0^2)}\right)^{K_n^{(+)}}+B_{nl}^{()}(W^2)\left(\frac{\alpha _S(\mu ^2)}{\alpha _S(\mu _0^2)}\right)^{K_n^{()}},$$
(76)
with an analogous equation for $`B_{nl}^{}`$ involving the factors $`g_n^{(\pm )}`$.
In the limit $`\mu \mathrm{}`$ only the terms with the smallest exponent $`K_1^{()}=0`$ in the coefficients (71) survive. The asymptotic form of the distribution amplitudes thus has only $`n=1`$ in the Gegenbauer expansion (69) and reads
$`{\displaystyle \underset{q=1}{\overset{n_f}{}}}\mathrm{\Phi }_q^+(z,\zeta ,W^2)`$ $`=`$ $`18n_fz(1z)(2z1)\left[B_{10}^{()}(W^2)+B_{12}^{()}(W^2)P_2(2\zeta 1)\right],`$ (77)
$`\mathrm{\Phi }_g(z,\zeta ,W)`$ $`=`$ $`48z^2(1z)^2\left[B_{10}^{()}(W^2)+B_{12}^{()}(W^2)P_2(2\zeta 1)\right],`$ (78)
where $`P_2(2\zeta 1)=16\zeta (1\zeta )`$. Note that $`B_{10}^{()}`$ and $`B_{12}^{()}`$ do not depend on a starting scale $`\mu _0`$ because $`K_1^{()}=0`$. For reasons that will become clear we will also keep the terms with the first nonzero exponent $`K_1^{(+)}`$ in our model for the GDAs to be developed in Sect. IV. For the quark distribution amplitudes this simply amounts to replacing $`B_{10}^{()}`$ and $`B_{12}^{()}`$ in the first line of Eq. (78) with the $`\mu `$-dependent coefficients $`B_{10}`$ and $`B_{12}`$.
Let us finally remark that, as discussed in , there is another generalized gluon distribution amplitude, with an operator different from the one in Eq. (26). It corresponds to pion pairs with angular momentum $`J_z=\pm 2`$ and gives the leading-twist part of the amplitudes $`\gamma ^{}\gamma \pi \pi `$ where the photon helicities are opposite. The evolution of this helicity-two distribution amplitude does not mix with any quark distribution. Its smallest anomalous dimension is positive, so that this distribution amplitude tends logarithmically to zero as $`\mu \mathrm{}`$. The study of this distribution would be very interesting. Nothing is, however, known about its size at present, and in our phenomenological analysis we will neglect its contribution.
### E Partial wave expansion
The decomposition of generalized distribution amplitudes on Legendre polynomials performed in the previous section translates into a partial waves decomposition if one transforms from polynomials $`P_l(2\zeta 1)`$ to $`P_l(\mathrm{cos}\theta )`$ using that $`2\zeta 1=\beta \mathrm{cos}\theta `$. The rearranged series reads
$$\underset{q=1}{\overset{n_f}{}}\mathrm{\Phi }_q^+=6n_fz(1z)\underset{\genfrac{}{}{0pt}{}{n=1}{\mathrm{odd}}}{\overset{\mathrm{}}{}}\underset{\genfrac{}{}{0pt}{}{l=0}{\mathrm{even}}}{\overset{n+1}{}}\stackrel{~}{B}_{nl}(W^2)C_n^{(3/2)}(2z1)P_l(\mathrm{cos}\theta )$$
(79)
for quarks, where the coefficients $`\stackrel{~}{B}_{nl}(W^2)`$ are linear combinations of the form
$$\stackrel{~}{B}_{nl}=\beta ^l\left[B_{nl}+c_{l,l+2}B_{n,l+2}+\mathrm{}+c_{l,n+1}B_{n,n+1}\right]$$
(80)
with polynomials $`c_{l,l^{}}`$ in $`\beta ^2`$. Keeping only $`n=1`$ in the Gegenbauer expansion one is restricted to an $`S`$\- and a $`D`$-wave:
$`{\displaystyle \underset{q=1}{\overset{n_f}{}}}\mathrm{\Phi }_q^+`$ $`=`$ $`18n_fz(1z)(2z1)\left[B_{10}(W^2)+B_{12}(W^2)P_2(2\zeta 1)\right]`$ (81)
$`=`$ $`18n_fz(1z)(2z1)\left[\stackrel{~}{B}_{10}(W^2)+\stackrel{~}{B}_{12}(W^2)P_2(\mathrm{cos}\theta )\right]`$ (82)
with
$`\stackrel{~}{B}_{10}(W^2)`$ $`=`$ $`B_{10}(W^2){\displaystyle \frac{1\beta ^2}{2}}B_{12}(W^2),`$ (83)
$`\stackrel{~}{B}_{12}(W^2)`$ $`=`$ $`\beta ^2B_{12}(W^2).`$ (84)
It is a remarkable consequence of the condition $`ln+1`$ that the presence of high partial waves implies a departure of the two-pion distribution amplitude from its asymptotic form. The $`\theta `$-distribution of the produced pion pair thus contains information about the dependence of the GDAs on $`z`$, which as a loop variable is integrated over in the amplitude of the process, cf. Eq. (14).
One-meson distribution amplitudes are real valued functions due to time reversal invariance. This is not true for generalized distribution amplitudes: the two-pion “out” state in the definition (15) of $`\mathrm{\Phi }^{\pi \pi }`$ is transformed into an “in” state under time reversal, and these states are different because hadrons interact with each other. Below the inelastic threshold, however, two-pion “in” and “out” states with definite angular momentum are related in a simple way via the phase shifts of elastic $`\pi \pi `$-scattering. With the aid of Watson’s theorem one then obtains the relation $`\stackrel{~}{B}_{nl}^{}=\stackrel{~}{B}_{nl}^{}\mathrm{exp}(2i\delta _l)`$ . This fixes the phase of the expansion coefficient $`\stackrel{~}{B}_{nl}`$ up to its overall sign:
$$\stackrel{~}{B}_{nl}=\eta _{nl}|\stackrel{~}{B}_{nl}|\mathrm{exp}(i\delta _l),\eta _{nl}=\pm 1,$$
(85)
where $`\delta _l`$ is the $`\pi \pi `$ phase shift for the $`l`$-th partial wave in the $`I=0`$ channel.
### F Momentum sum rule
Of particular interest are the moments
$`{\displaystyle _0^1}𝑑z(2z1)\mathrm{\Phi }_q^+(z,\zeta ,W^2)`$ $`=`$ $`{\displaystyle \frac{2}{(P^+)^2}}\pi ^+(p)\pi ^{}(p^{})|T_q^{++}(0)|0,`$ (86)
$`{\displaystyle _0^1}𝑑z\mathrm{\Phi }_g(z,\zeta ,W^2)`$ $`=`$ $`{\displaystyle \frac{1}{(P^+)^2}}\pi ^+(p)\pi ^{}(p^{})|T_g^{++}(0)|0,`$ (87)
where $`T_q^{\mu \nu }(x)`$ and $`T_g^{\mu \nu }(x)`$ respectively denote the Belinfante improved energy-momentum tensors for quarks of flavor $`q`$ and for gluons. After summing (86) over all flavors these moments project out the coefficients $`B_{10}(W^2)`$, $`B_{12}(W^2)`$ and $`B_{10}^{}(W^2)`$, $`B_{12}^{}(W^2)`$.
To proceed one decomposes $`\pi ^+(p)\pi ^{}(p^{})|T_q^{\mu \nu }(0)|0`$ on form factors. Their analytical continuation to zero or negative $`W^2`$ leads to the form factors of the matrix elements $`\pi ^+(p)|T_q^{\mu \nu }(0)|\pi ^+(p^{})`$ between one-pion states, with $`W^2=0`$ corresponding to $`p=p^{}`$. At that point we get from Eq. (86)
$$B_{12}(0)=\frac{10}{9n_f}R_\pi ,$$
(88)
where $`R_\pi `$ is the fraction of light-cone momentum carried by quarks and antiquarks in the pion. No constraint on $`B_{10}(0)`$ is obtained this way, since the corresponding form factor in the decomposition of $`\pi ^+(p)|T_q^{\mu \nu }(0)|\pi ^+(p^{})`$ is multiplied by a tensor that vanishes for $`p=p^{}`$. In an analogous fashion one obtains an expression for $`B_{12}^{}(0)`$ from the sum rule (87).
We emphasize that both sides of Eq. (88) depend on the renormalization scale $`\mu `$. Only the total energy-momentum tensor, i.e., the sum $`T^{\mu \nu }=_qT_q^{\mu \nu }+T_g^{\mu \nu }`$ over quarks and gluons is conserved, so that its matrix elements are renormalization scale independent. The appropriate sum of the moments (86) and (87) leads to a linear combination of $`B_{12}`$ and $`B_{12}^{}`$ where the scale dependent term with $`B_{12}^{(+)}`$ indeed drops out and only $`B_{12}^{()}`$ is left. The normalization of $`\pi ^+(p)|T^{\mu \nu }(0)|\pi ^+(p)`$ thus fixes the expansion coefficient
$$B_{12}^{()}(0)=\frac{10}{9n_f+48},$$
(89)
which through the relation (88) gives the asymptotic value
$$R_\pi \stackrel{\mu \mathrm{}}{}\frac{3n_f}{3n_f+16},$$
(90)
in agreement with the well-known result from the evolution of singlet parton distributions .
## IV A simple model of the GDA
So far no experimental information exists on the two-pion GDA. In the numerical studies to follow we will therefore use a simple ansatz for $`\mathrm{\Phi }_q^+(z,\zeta ,W^2)`$, which is based on the general properties we have discussed in the previous section.
We only consider the contributions from $`u`$\- and $`d`$-quarks, i.e., we take $`n_f=2`$. As already mentioned we will use the isospin relations (23) and (24), and take the asymptotic form of the $`z`$ dependence given in Eq. (82). It thus remains to make an ansatz for the coefficients $`B_{10}(W^2)`$ and $`B_{12}(W^2)`$, or equivalently for $`\stackrel{~}{B}_{10}(W^2)`$ and $`\stackrel{~}{B}_{12}(W^2)`$ introduced in Eq. (84).
For their phases, given by Eq. (85), we use simple parameterizations of the isosinglet $`S`$\- and $`D`$-wave phase shifts $`\delta _0`$ and $`\delta _2`$ obtained in . They are shown in Fig. 6, where for later use the phase shift $`\delta _1`$ of the $`P`$-wave is also displayed. The result (85) only holds below the inelastic threshold in $`\pi \pi `$ scattering, therefore we restrict all our studies to invariant masses $`W`$ below 1 GeV. Around that mass, corresponding to the $`K\overline{K}`$ threshold, the phase shift $`\delta _0`$ of the $`S`$-wave drastically increases. While the analysis of stops at $`W=0.97`$ GeV and does not exhibit this abrupt change, the investigations in Ref. find values of order 200 at $`W=1`$ GeV. Our parameterization of $`\delta _0`$ in that region is meant to be indicative rather than a precise description of this quantity. Through interference effects, the rapid variation of a phase shift leads to a characteristic behavior in the $`W`$-spectrum of appropriate observables in our process, as we shall see in Sect. VII.
The analyticity properties of the $`\stackrel{~}{B}_{nl}`$ and the phase information from Watson’s theorem (85) may be used to obtain the $`W^2`$-dependence of $`\stackrel{~}{B}_{nl}`$ via dispersion relations, which has been exploited in . Note, however, that while the complex phases are simple for the $`\stackrel{~}{B}_{nl}`$, it is the $`B_{nl}`$ that have simple analytic properties in the $`W^2`$-plane, given their definition through operator matrix elements. The transformation from $`B_{nl}`$ to $`\stackrel{~}{B}_{nl}`$ introduces extra poles at $`W^2=0`$, cf., e.g., the factors $`\beta ^2=(W^24m_\pi ^2)/W^2`$ in Eq. (84). Furthermore, the evaluation of the integrals that solve the dispersion relations requires knowledge of the phases at energies above the value of $`W`$ where $`\stackrel{~}{B}_{nl}`$ is evaluated. This further restricts the range of $`W`$ where $`\stackrel{~}{B}_{nl}`$ can be obtained using the $`\pi \pi `$ phase shifts as input.
To keep our model simple we will make a less sophisticated ansatz. We keep the energy dependent phases $`\delta _0`$ and $`\delta _2`$ from Watson’s theorem (85). To determine $`|\stackrel{~}{B}_{10}|`$, $`|\stackrel{~}{B}_{12}|`$, and the overall signs $`\eta _{10}`$, $`\eta _{12}`$ in Eq. (85), we retain only the kinematical factors $`\beta ^2`$ in the relation (84) and replace $`B_{10}(W^2)`$ and $`B_{12}(W^2)`$ with their values at $`W=0`$. Close to $`W=1`$ GeV one will not expect this to be a good approximation for the $`S`$-wave, given the presence of the $`f_0(980)`$. Below this there is however no prominent $`\pi \pi `$ resonance in the $`I=0`$ channel, and the phase shifts show a smooth behavior. It seems therefore reasonable to assume that the isosinglet form factors $`\stackrel{~}{B}_{10}`$ and $`\stackrel{~}{B}_{12}`$ do not have a strong energy dependence in that region, certainly not as strong as the electromagnetic pion form factor $`F_\pi `$ with its large variations in modulus and phase due to the $`\rho (770)`$. We do however not claim our simple model to be better than, say, a factor of 2.
For the input value of $`B_{12}(0)`$ we use the constraint (88) with $`R_\pi `$ evaluated from the parton distributions in the pion. Taking the LO parameterization of GRS we find $`R_\pi `$ ranging from 0.5 to 0.6 at a factorization scale $`\mu ^2`$ between 1 GeV<sup>2</sup> and 20 GeV<sup>2</sup>. In our numerical studies we use $`R_\pi =0.5`$. Note that this is very far from the asymptotic value (90), which for $`n_f=2,3,4`$ is $`R_\pi =0.27`$, $`0.36`$ and $`0.43`$, respectively. While using the asymptotic form of the $`z`$-dependence of the GDA for simplicity (and lack of experimental information) we thus retain a clear non-asymptotic effect in the coefficient $`B_{12}(0)`$. We also remark that in the GRS LO parameterization the contribution of strange quarks and antiquarks to $`R_\pi `$ is at the level of 5% to 10% in a wide range of the factorization scale. This corroborates our restriction to $`u`$\- and $`d`$-quarks in the GDA, although with the caveat that the sea quark distribution in the pion is not constrained from experimental data .
For the coefficient $`B_{10}(0)`$ we make use of the relation
$$B_{10}(0)=B_{12}(0),$$
(91)
which has been obtained in using chiral symmetry in the form of a soft-pion theorem. Notice that our ansatz then has the property that for $`\beta 1`$ the $`S`$\- and $`D`$-wave components of the GDA have equal size and opposite sign, as is easily seen from Eq. (84).
Putting everything together, we will take the following model GDAs in our numerical studies:
$$\mathrm{\Phi }_u^+=\mathrm{\Phi }_d^+=10z(1z)(2z1)R_\pi \left[\frac{3\beta ^2}{2}e^{i\delta _0(W^2)}+\beta ^2e^{i\delta _2(W^2)}P_2(\mathrm{cos}\theta )\right]$$
(92)
with $`R_\pi =0.5`$.
With this we can easily calculate the scattering amplitude for $`\gamma ^{}\gamma \pi \pi `$ to leading order in $`\alpha _S`$. We shall neglect here the radiative corrections to the hard scattering, which have been worked out to one loop in . Taking the asymptotic form (78) of the quark and gluon GDAs, including the asymptotic value (90) of the ratio $`R_\pi `$, they were found to reduce the leading-order amplitude for equal photon helicities by 30% if $`\alpha _S=0.3`$, with most of the correction being due to the contribution from $`\mathrm{\Phi }_g`$. Finally, we recall from the end of Sect. III D that we will neglect the contribution of the helicity-two gluon GDA to the photon double helicity-flip amplitude, which is also a one-loop effect.
## V Comparison with $`\gamma ^{}\gamma \pi ^0`$
Given the close analogy between the production of one and of two pions it is natural to compare the production rates of these two processes. Since our estimations for $`\pi \pi `$ production are at lowest order in $`\alpha _S`$ we will compare with the corresponding expression for the one-pion case for consistency, although experimental data and more refined theory analyses are available there. From the leading-order expression (16) we obtain the cross section for the process $`e\gamma e\pi ^0`$ as
$$\frac{d\sigma _{e\gamma e\pi ^0}}{dQ^2}=\frac{\alpha ^3}{s_{e\gamma }^2}\frac{1}{Q^2(1ϵ)}\mathrm{\hspace{0.17em}2}\pi ^2f_\pi ^2$$
(93)
where we have used the asymptotic distribution amplitude $`\varphi _u^\pi =\varphi _d^\pi =3\sqrt{2}f_\pi z(1z)`$ with $`f_\pi 131`$ MeV. For a lowest-order approximation, the cross section (93) is in fair agreement with the data .
To compare with two-pion production, we integrate the cross section for $`e\gamma e\pi ^0\pi ^0`$ from threshold up to $`W_{\mathrm{𝑚𝑎𝑥}}`$. With our model GDA (92) we find
$`{\displaystyle \frac{d\sigma _{e\gamma e\pi ^0\pi ^0}}{dQ^2}}`$ $`=`$ $`{\displaystyle \frac{25\alpha ^3}{72s_{e\gamma }^2}}{\displaystyle \frac{1}{Q^2(1ϵ)}}{\displaystyle _{4m_\pi ^2}^{W_{\mathrm{𝑚𝑎𝑥}}^2}}𝑑W^2\sqrt{1{\displaystyle \frac{4m_\pi ^2}{W^2}}}\left(|\stackrel{~}{B}_{10}|^2+{\displaystyle \frac{1}{5}}|\stackrel{~}{B}_{12}|^2\right)`$ (94)
$`=`$ $`{\displaystyle \frac{125\alpha ^3}{243s_{e\gamma }^2}}{\displaystyle \frac{1}{Q^2(1ϵ)}}R_\pi ^2m_\pi ^2\sqrt{1{\displaystyle \frac{4m_\pi ^2}{W_{\mathrm{𝑚𝑎𝑥}}^2}}}\left({\displaystyle \frac{W_{\mathrm{𝑚𝑎𝑥}}^2}{4m_\pi ^2}}{\displaystyle \frac{3}{4}}{\displaystyle \frac{m_\pi ^2}{W_{\mathrm{𝑚𝑎𝑥}}^2}}\right).`$ (95)
A consequence of the identical scaling behavior of the two processes is that the ratio of the cross sections (95) and (93) is independent of $`Q^2`$ in the Born approximation.
Fig. 7 shows the ratio of the cross sections (95) and (93) as a function of the upper integration limit $`W_{\mathrm{𝑚𝑎𝑥}}`$. We see that, even when integrating up to $`W=1`$ GeV, the single-pion production comes out as clearly dominant. We remark that the measured production rates for a single $`\eta `$ or $`\eta ^{}`$ are comparable to that of a $`\pi ^0`$. With our isospin relation (23) the cross section for $`\gamma ^{}\gamma \pi ^+\pi ^{}`$ is twice that of $`\gamma ^{}\gamma \pi ^0\pi ^0`$, the relative factor $`1/2`$ for $`\pi ^0\pi ^0`$ being due to the phase space of identical particles. Due to phase space one does not expect the production of more than two pions to be important for $`W`$ below 1 GeV, except for the decays $`\eta 3\pi `$ and $`\eta ^{}5\pi `$. The picture thus emerges that with our estimation for $`\gamma ^{}\gamma \pi \pi `$ the production of hadrons in $`\gamma ^{}\gamma `$ collisions up to 1 GeV is dominated by the pseudoscalar channel, in other words by the parity-odd sector as opposed to the parity-even one. This is reminiscent of the special role played by the axial current in low-energy QCD.
At this point we wish to comment on the end-point regions of the integrals over $`z`$ in the factorized expressions (14) and (16) for two-pion and one-pion production. For $`z0`$ and $`z1`$ the hard-scattering kernels are divergent, corresponding to the quark exchanged between the $`\gamma `$ and $`\gamma ^{}`$ going on-shell. These poles are canceled by the end-point zeroes of the two-pion and one-pion distribution amplitudes, so that the end-point regions give a finite contribution to the scattering amplitude in both cases. Quantitatively, the quark virtualities in the hard-scattering diagrams are $`zQ^2`$ and $`(1z)Q^2`$, and it is clear that for a given finite $`Q^2`$ there is a region in $`z`$ where our leading-order expressions should receive important corrections. At small virtualities the strong coupling becomes large, increasing the size of $`\alpha _S`$ corrections, and when $`zQ^2`$ or $`(1z)Q^2`$ becomes comparable to the square of typical transverse quark momenta in a pion, then power corrections due to the effect of the transverse momentum of the produced $`q\overline{q}`$\- pair will be important. We recall in this context that various theoretical attempts to evaluate such corrections lead to fair agreement with the data for the $`\gamma `$$`\pi `$ transition form factor down to rather low values of $`Q^2`$.
For pion pair production both the hard-scattering kernel and the distribution amplitude are zero at $`z=1/2`$ due to the constraints from charge conjugation invariance, so that compared to the single-pion case the integral in $`z`$ is more sensitive to the end-point regions. We thus expect that for intermediate values of $`Q^2`$ corrections to the lowest-order results will be more important in $`\gamma ^{}\gamma \pi \pi `$ than they are in $`\gamma ^{}\gamma \pi ^0`$. The experimental comparison of the $`Q^2`$-dependence of these two processes will therefore be interesting and may help us to better understand the origin of these corrections, which are a subject of considerable importance in the physics of exclusive processes.
Taking the asymptotic $`z`$-dependence of the distribution amplitudes as an example, we can explicitly see how important the end-point contributions are in the leading-order expressions (14) and (16). For single pion production the integrand in Eq. (16) is a constant then, so that 50% of the $`z`$-integral comes from the regions where $`z`$ or $`1z`$ is smaller than 0.25. For two-pion production the integrand is proportional to $`(2z1)^2`$, and 50% of the integrand comes from the regions with $`z`$ or $`1z`$ smaller than $`(12^{1/3})/20.1`$. Given these numbers, one can expect that corrections to our leading-order calculation will not be negligible for $`Q^2`$ around 4 GeV<sup>2</sup>, which is the lowest value considered in our numerical estimates in Sect. VIII.
## VI Relations with the photon structure function
The exclusive process we consider here contributes of course to the inclusive reaction $`\gamma ^{}\gamma X`$. As we mentioned in the previous section, the inclusive process is built up from a limited number of exclusive channels in the mass region of $`W`$ below 1 GeV. Let us examine the connection between our discussion of one- and two-pion production with the familiar description of inclusive $`\gamma ^{}\gamma `$ scattering in the kinematical limit we are taking here.
The unpolarized cross section for inclusive deep inelastic scattering on a photon, $`e\gamma eX`$, can be parameterized by two photon structure functions $`F_T`$ and $`F_L`$ as
$$\frac{d\sigma _{e\gamma eX}}{dQ^2dW^2}=\frac{2\pi \alpha ^2}{s_{e\gamma }^2}\frac{1}{xQ^2(1ϵ)}\left(2xF_T(x,Q^2)+ϵF_L(x,Q^2)\right),$$
(96)
where $`F_T`$ and $`F_L`$ respectively give the contribution from transverse and longitudinal polarization of the exchanged $`\gamma ^{}`$. The transverse structure function $`F_T`$ is often traded for $`F_2=2xF_T+F_L`$.
At the level of partons inclusive hadron production is described by $`\gamma ^{}\gamma q\overline{q}`$ to leading order in $`\alpha _S`$, which gives the well-known expressions
$`F_T^{q\overline{q}}`$ $`=`$ $`{\displaystyle \frac{3\alpha }{2\pi }}{\displaystyle \underset{q}{}}e_q^4\left\{\mathrm{ln}{\displaystyle \frac{1+\beta _q}{1\beta _q}}\left[x^2+(1x)^2+4x(1x){\displaystyle \frac{m_q^2}{Q^2}}8x^2{\displaystyle \frac{m_q^4}{Q^4}}\right]\beta _q\left[(12x)^2+4x(1x){\displaystyle \frac{m_q^2}{Q^2}}\right]\right\},`$ (97)
$`F_L^{q\overline{q}}`$ $`=`$ $`{\displaystyle \frac{12\alpha }{\pi }}{\displaystyle \underset{q}{}}e_q^4x^2(1x)\left[\beta _q{\displaystyle \frac{2m_q^2}{W^2}}\mathrm{ln}{\displaystyle \frac{1+\beta _q}{1\beta _q}}\right],`$ (98)
where $`\beta _q=(14m_q^2/W^2)^{1/2}`$. Note that $`m_q`$ is to be understood here as a cutoff parameter, which regulates the collinear divergence in the box diagram with massless quarks.
The limit of large $`Q^2`$ at fixed small $`W^2`$ we are taking here implies $`x1`$ and is different from the Bjorken limit, where $`W^2`$ is scaled up with $`Q^2`$ so that $`x`$ remains constant. Neglecting terms of order $`1xW^2/Q^2`$ and $`m_q^2/Q^2`$ the expressions (97) become
$$F_T^{q\overline{q}}=\frac{3\alpha }{2\pi }\underset{q}{}e_q^4\left\{\mathrm{ln}\frac{1+\beta _q}{1\beta _q}\beta _q\right\},F_L^{q\overline{q}}=O\left(\frac{W^2}{Q^2}\right).$$
(99)
We observe that in our limit the leading-order expression for $`F_T`$ becomes independent of $`Q^2`$, i.e., it has the same scaling behavior as the exclusive channels $`\gamma ^{}\gamma \pi `$ and $`\gamma ^{}\gamma \pi \pi `$. This is to be contrasted with the Bjorken limit, where $`\mathrm{ln}[(1+\beta _q)/(1\beta _q)]\mathrm{ln}[Q^2/m_q^2]+\mathrm{ln}[(1x)/x]`$ gives rise to the well-known logarithmic scaling violation of $`F_T`$ at zeroth order in $`\alpha _S`$.
Just as in the case of $`\gamma ^{}\gamma \pi \pi `$, the contribution $`F_L`$ from longitudinal photons is power suppressed in our limit. Let us add that in the Bjorken limit the hadronic part of $`F_T`$, often parameterized using vector meson dominance, is only suppressed by a factor $`\mathrm{ln}Q^2`$ with respect to the pointlike part (97), but does not survive our limiting procedure here: since hadronic structure functions typically decrease like a power of $`1x`$ for $`x1`$, it becomes a correction in $`W^2/Q^2`$.
The contribution of our process to the structure functions is, with our ansatz (92) for $`\mathrm{\Phi }_q^+`$,
$$F_T^{\pi ^+\pi ^{}+\pi ^0\pi ^0}=\frac{25\alpha }{96\pi }\beta \left(|\stackrel{~}{B}_{10}|^2+\frac{1}{5}|\stackrel{~}{B}_{12}|^2\right)=\frac{625\alpha }{3456\pi }R_\pi ^2\beta \left(1\frac{2}{3}\beta ^2+\frac{1}{5}\beta ^4\right).$$
(100)
As a function of $`W`$ this quickly rises from the threshold at $`2m_\pi `$, levels off for $`W`$ around 400 to 500 MeV, and then remains flat with a value $`F_T^{\pi ^+\pi ^{}+\pi ^0\pi ^0}/\alpha 0.0077`$. Let us compare this with the result (99) of the $`q\overline{q}`$ calculation for $`u`$\- and $`d`$-quarks (including strange quarks would only lead to a minute change due to the charge factor $`e_q^4`$). With the quark masses $`m_u=m_d=290`$ MeV from the parameterization of the photon structure function by Gordon and Storrow we get a value of $`F_T^{q\overline{q}}/\alpha 0.15`$ at $`W=1`$ GeV, much larger than the one we obtain for pion pairs.
It is worth remembering that $`\gamma ^{}\gamma q\overline{q}`$ also is the hard-scattering subprocess in our factorized expression for $`\gamma ^{}\gamma \pi \pi `$. As we discussed at the end of the previous section, the collinear divergence of this process shows up as singularities at the end-points of the $`z`$-integration in Eq. (14) and is canceled by the end-point zeroes of the GDA, i.e., by the hadronization process. In the calculation of open $`q\overline{q}`$ production no such cancellation takes place and the divergence of the diagram has to be regulated. This reflects the fact that even in the limit $`Q^2\mathrm{}`$ inclusive hadron production from $`\gamma ^{}\gamma `$ cannot be calculated in perturbation theory alone (unlike for instance inclusive hadron production from a single timelike photon) and that the separation of $`F_T`$ into a perturbative pointlike and a non-perturbative hadronic part is not unambiguous. While more sophisticated procedures have been elaborated in the literature, we consider it sufficient for our purpose to use the quark mass regulator in Eq. (99). One might also take massless quarks and a lower cutoff $`\kappa _{}`$ on the transverse quark momentum, obtaining the same result (99) with $`m_q`$ replaced by $`\kappa _{}`$ in the expression of $`\beta _q`$. While keeping us away from the region where perturbation theory breaks down, such phenomenological regulators become of course inadequate as one approaches the “threshold” where $`\beta _q=0`$. For our numbers this is at $`W=580`$ MeV. One should bear this in mind when using the expression (99) for invariant masses $`W`$ around 1 GeV.
On the other hand we saw in Sect. V that with our estimate of two-pion production the hadronic mass spectrum below 1 GeV is dominated by the single-meson states $`\pi ^0`$, $`\eta `$, $`\eta ^{}`$. It is clear that in such a region the parton-level result can only hold in the sense of parton-hadron duality, averaged over a sufficiently large interval in $`W`$. We therefore integrate the cross section for $`e\gamma eX`$ over $`W`$ from threshold up to 1 GeV. The parton-level result, obtained from Eq. (99) with $`m_u=m_d=290`$ MeV, amounts to 2.42 times the cross section (93) for one-pion production. This factor should be compared with the factor $`1+0.26+0.97+2.64`$ for the individual contributions of the exclusive channels $`\pi +\pi \pi +\eta +\eta ^{}`$. Here we used Eq. (100) for two-pion production, whereas for $`\eta `$ and $`\eta ^{}`$ we replaced $`f_\pi =131`$ MeV in Eq. (93) with the respective values 129 MeV and 213 MeV taken from the analysis of . Given the caveats of parton-hadron duality (below 1 GeV there are very few resonances in the two-photon channel, and $`W=1`$ GeV is just above the $`\eta ^{}`$ threshold) and those of the parton-level calculation itself (discussed above), we find the agreement remarkably fair.
## VII Phenomenology
We will now discuss the phenomenology of our process in $`e\gamma `$ and in $`e^+e^{}`$ collisions. The production of neutral and charged pion pairs is rather different in this respect, since $`\pi ^0\pi ^0`$ is only produced by the $`\gamma ^{}\gamma `$ subprocess we have discussed so far, whereas for $`\pi ^+\pi ^{}`$ production this process interferes with bremsstrahlung, i.e., the production of the pion pair from a timelike photon radiated off the beam lepton. We start with the simpler case of neutral pions, and then discuss charged pairs. In the following we will restrict ourselves to unpolarized photon and lepton beams. A brief discussion of beam polarization will be given in Appendix B.
### A Helicity amplitudes
The building blocks of our investigation are the helicity amplitudes for $`\gamma ^{}\gamma \pi \pi `$, which describe the dynamics of this process in a model independent way. They are obtained from the hadronic tensor $`T^{\mu \nu }`$ by multiplying the reduced amplitudes
$$A_{ij}(Q^2,W^2,\theta )=ϵ_i^\mu T_{\mu \nu }ϵ^{}{}_{j}{}^{\nu },i=+,0,,j=+,$$
(101)
with the squared elementary charge $`e^2`$. In the $`\gamma ^{}\gamma `$ c.m. our photon polarization vectors read
$$ϵ_0=\frac{1}{Q}(|𝐪|,0,0,q^0),ϵ_\pm =\frac{1}{\sqrt{2}}(0,1,i,0)$$
(102)
for the virtual and
$$ϵ^{}{}_{\pm }{}^{}=\frac{1}{\sqrt{2}}(0,1,+i,0)$$
(103)
for the real photon, where we have used the coordinate system described in Sect. II A. By parity invariance, there are only three independent helicity amplitudes, which we choose to be $`A_{++}`$, $`A_+`$ and $`A_{0+}`$.
Each of these three amplitudes plays a distinctive dynamical role in the kinematical region $`Q^2W^2,\mathrm{\Lambda }^2`$. It is $`A_{++}`$ that receives the leading twist contribution we have discussed in detail, and which in the scaling regime gives access to the generalized quark distribution amplitudes $`\mathrm{\Phi }_q^{\pi \pi }`$,
$$A_{++}=\underset{q}{}\frac{e_q^2}{2}_0^1𝑑z\frac{2z1}{z(1z)}\mathrm{\Phi }_q^{\pi \pi }(z,\zeta ,W^2)$$
(104)
to zeroth order in $`\alpha _S`$. The amplitude $`A_+`$ has a leading-twist part at order $`\alpha _S`$, due to the helicity-two gluon GDA. We briefly discussed this at the end of Sect. III D; for more detail we refer to . Finally, the contribution $`A_{0+}`$ from a longitudinal $`\gamma ^{}`$ is nonleading twist. The predicted power behavior in $`Q^2`$ at fixed $`W^2`$ and $`\zeta `$ is therefore that $`A_{++}`$ becomes independent of $`Q^2`$, whereas $`A_{0+}`$ decreases at least like $`1/Q`$. The amplitude $`A_+`$ should become $`Q^2`$-independent. If the helicity-two gluon GDA is however not sufficiently large, $`A_+`$ may be dominated by higher-twist contributions at accessible values of $`Q^2`$ and should decrease like a power of $`1/Q`$ in the corresponding $`Q^2`$-range. Of course all these predictions are to be understood as up to corrections in $`\mathrm{log}Q^2`$. At sufficiently large $`Q^2`$, the longitudinal amplitude $`A_{0+}`$ is thus predicted to be small compared with $`A_{++}`$. One can also expect that $`A_+`$ will be smaller than $`A_{++}`$, since its leading-twist part is suppressed by $`\alpha _S`$.
To discuss the different partial waves in which the pion pair can be produced, we expand each of the amplitudes $`A_{++}`$, $`A_{0+}`$, $`A_+`$ as
$$A_{ij}=\underset{\genfrac{}{}{0pt}{}{l=ji}{\mathrm{even}}}{\overset{\mathrm{}}{}}A_{ijl}(Q^2,W^2)P_l^{ji}(\mathrm{cos}\theta ),i=+,0,,j=+,$$
(105)
where $`P_l^m`$ denotes the associated Legendre polynomial corresponding to the value of $`J_z`$ of the $`\pi \pi `$ system in its c.m.
### B The $`\gamma ^{}\gamma `$ subprocess and $`\pi ^0\pi ^0`$ production
The differential $`e\gamma `$ cross section for neutral pion pair production reads
$`{\displaystyle \frac{d\sigma _{e\gamma e\pi \pi }}{dQ^2dW^2d(\mathrm{cos}\theta )d\phi }}|_G={\displaystyle \frac{\alpha ^3}{16\pi }}{\displaystyle \frac{\beta }{s_{e\gamma }^2}}{\displaystyle \frac{1}{Q^2(1ϵ)}}`$ $`(`$ $`|A_{++}|^2+|A_+|^2+2ϵ|A_{0+}|^2`$ (106)
$``$ $`\mathrm{cos}\phi \mathrm{Re}\left\{A_{++}^{}A_{0+}^{}A_+^{}A_{0+}^{}\right\}2\sqrt{ϵ(1+ϵ)}`$ (107)
$``$ $`\mathrm{cos}2\phi \mathrm{Re}\left\{A_{++}^{}A_+^{}\right\}2ϵ),`$ (108)
where the subscript $`G`$ indicates that the pions are produced in a $`\gamma ^{}\gamma `$ subprocess. For $`\pi ^0\pi ^0`$ production the phase space in Eq. (108) is understood as restricted to $`\mathrm{cos}\theta (0,1)`$, $`\phi (0,2\pi )`$ because there are two identical particles in the final state. We notice the close similarity of the expression (108) with the cross section of the crossed channel process of virtual Compton scattering, and much of what we discuss in the following has its counterpart there .
To obtain the $`e^+e^{}`$ cross section we use the equivalent photon approximation ,
$$\frac{d\sigma _{eeee\pi \pi }}{dQ^2dW^2d(\mathrm{cos}\theta )d\phi dx_2}=\frac{\alpha }{\pi }\frac{1}{x_2}\left(\frac{1+(1x_2)^2}{2}\mathrm{ln}\left[\frac{Q_{\mathrm{𝑚𝑎𝑥}}^2(x_2)}{Q_{\mathrm{𝑚𝑖𝑛}}^2(x_2)}\right](1x_2)\right)\frac{d\sigma _{e\gamma e\pi \pi }}{dQ^2dW^2d(\mathrm{cos}\theta )d\phi },$$
(109)
where $`Q_{\mathrm{𝑚𝑖𝑛}}^2`$ and $`Q_{\mathrm{𝑚𝑎𝑥}}^2`$ are the minimal and maximal virtuality of the photon $`q^{}`$, respectively. We have a lower kinematical limit $`Q_{\mathrm{𝑚𝑖𝑛}}^2=x_2^2m_e^2/(1x_2)`$ determined by the electron mass $`m_e`$, whereas $`Q_{\mathrm{𝑚𝑎𝑥}}^2`$ depends on experimental cuts and will be discussed in more detail in Sect. VIII A. We remark that for a given $`ee`$ collider energy the variables $`x_2`$ and $`y`$ are not independent at fixed $`Q^2`$ and $`W^2`$, since
$$yx_2=\frac{Q^2+W^2}{s_{ee}},$$
(110)
and that in Eq. (109) one can easily trade $`dx_2`$ for $`dy`$.
Since the helicity amplitudes $`A_{ij}`$ are independent of $`\phi `$ they can be partially disentangled from the $`\phi `$-dependence of the cross section, which is completely explicit in Eq. (108). In particular, the relative size and the $`Q^2`$-behavior of the $`\phi `$-independent term and of the terms with $`\mathrm{cos}\phi `$ and $`\mathrm{cos}2\phi `$ allow detailed tests of the scaling predictions. This provides indicators on how close one is to the asymptotic regime at finite values of $`Q^2`$. The $`\phi `$-independent term in the large parentheses of Eq. (108) receives contributions from leading-twist amplitudes and should thus display scaling behavior. The coefficient of $`\mathrm{cos}\phi `$ is the interference of leading-twist and non-leading twist amplitudes and should be suppressed by at least one power of $`1/Q`$. Finally, the $`\mathrm{cos}2\phi `$ term should scale or be power suppressed depending on the size of the helicity-two gluon GDA.
Apart from standard fitting techniques a way to separate terms with different angular dependence is the use of weighted cross sections. Weighting each event with a function $`w(\phi ,\theta )`$ we define
$$S_{e\gamma }(w)=𝑑Q^2𝑑W^2𝑑\mathrm{\Omega }\frac{d\sigma _{e\gamma }}{dQ^2dW^2d\mathrm{\Omega }}w(\phi ,\theta ),$$
(111)
where $`d\mathrm{\Omega }=d(\mathrm{cos}\theta )d\phi `$. Notice that since it is not normalized, $`S_{e\gamma }(w)`$ is not just the average value of the function $`w(\phi ,\theta )`$, and includes information about the size of the cross section itself. Interpreting $`S_{e\gamma }(w)`$ as a statistical variable one can calculate its standard deviation and finds for its relative statistical error (cf., e.g., )
$$\delta (w)=\frac{1}{\sqrt{N}}\frac{\sqrt{{\displaystyle 𝑑Q^2𝑑W^2𝑑\mathrm{\Omega }\frac{d\sigma _{e\gamma }}{dQ^2dW^2d\mathrm{\Omega }}w^2(\phi ,\theta )}}\sqrt{{\displaystyle 𝑑Q^2𝑑W^2𝑑\mathrm{\Omega }\frac{d\sigma _{e\gamma }}{dQ^2dW^2d\mathrm{\Omega }}}}}{\left|{\displaystyle 𝑑Q^2𝑑W^2𝑑\mathrm{\Omega }\frac{d\sigma _{e\gamma }}{dQ^2dW^2d\mathrm{\Omega }}w(\phi ,\theta )}\right|},$$
(112)
where
$$N=𝑑Q^2𝑑W^2𝑑\mathrm{\Omega }\frac{d\sigma _{e\gamma }}{dQ^2dW^2d\mathrm{\Omega }}$$
(113)
is the expected number of events for a given integrated luminosity $``$. Eq. (112) generalizes the well-known result that the relative statistical error of the cross section, i.e., of $`S_{e\gamma }(1)`$, is $`1/\sqrt{N}`$. We emphasize that the method of weighted cross sections is very flexible, and that the choice of weights $`w(\phi ,\theta )`$ can for instance be adapted to experimental conditions such as limited angular acceptance or cuts. One can of course take weighting functions that depend on other variables than only $`\theta `$ and $`\phi `$. In the following we will also use weighted differential cross sections, where only some of the kinematical variables have been integrated out while others are held fixed. In a data analysis, one may thus use the weighting technique for some variables and fitting for others.
The weighting technique is convenient to project out different terms in the cross section. As an immediate example we note that the terms constant in $`\phi `$, with $`\mathrm{cos}\phi `$, and with $`\mathrm{cos}2\phi `$ in the $`e\gamma `$ cross section are obtained from
$$\frac{dS_{e\gamma }(\mathrm{cos}m\phi )}{dQ^2dW^2d(\mathrm{cos}\theta )}$$
(114)
with $`m=0`$, 1 and 2, respectively. If the moments with $`m=1`$ and 2 are measured to be small compared with the moment $`m=0`$, this can be because any two of the amplitudes $`A_{++}`$, $`A_{0+}`$, $`A_+`$ are much smaller than the third, or it may be due to their relative phases. From the theoretical considerations in Sect. VII A the most natural hypothesis in this case is however that $`A_{0+}`$ and $`A_+`$ are small compared with $`A_{++}`$.
While the $`\phi `$-dependence of the cross section (108) gives access to the various helicity combinations of the real and virtual photon, its dependence on $`\theta `$ contains information on the angular momentum states in which the pion pair is produced. A priori there can be arbitrarily high partial waves, but to analyze the $`\theta `$-distribution in practice one will assume that at a given $`W`$ only a finite number of them is important, if only for reasons of phase space. It is easy to see from Eqs. (105) and (108) that for a superposition of partial waves $`l=0`$, 2, …$`L`$ the moment of $`\mathrm{cos}m\phi `$ in (114) is a linear combination of polynomials $`P_{2l}^m(\mathrm{cos}\theta )`$ with highest degree $`2L`$. Weighting the cross section with $`\mathrm{cos}(m\phi )P_{2L+2}^m(\mathrm{cos}\theta )`$ and integrating over $`\phi `$ and $`\theta `$ then gives a zero result. Using these weights thus provides one way to estimate from experimental data how many partial waves are relevant. Let us recall the physical relevance of this information: in the scaling regime the highest partial wave relevant in $`A_{++}`$ provides a constraint on how far the two-pion distribution amplitude is from its asymptotic form, as we discussed in Sect. III E.
Let us assume that only partial waves with $`lL`$ effectively contribute in the cross section (108). The $`\theta `$-dependence of the moments in (114) is then determined by $`L+1`$ coefficients for $`m=0`$, $`L`$ coefficients for $`m=1`$ and $`L`$ coefficients for $`m=2`$, corresponding to the number of polynomials $`P_{2l}^m(\mathrm{cos}\theta )`$ with $`lL`$. On the other hand, there are $`3L/2+1`$ complex amplitudes $`A_{ijl}`$ with $`lL`$ in the expansion (105), so that there are $`3L+2`$ real quantities one would like to determine. A global phase is however unobservable in the cross section (108), and one may for instance refer all phases to the phase of $`A_{++0}`$. The $`3L+1`$ coefficients one can extract from the dependence of the cross section on $`\phi `$ and $`\theta `$ thus allow one to reconstruct the $`|A_{ijl}|`$ and their relative phases. Since the relation between the angular coefficients and the amplitudes is quadratic, there will however be multiple solutions in general. More information can be obtained with polarized beams, which we briefly discuss in Appendix B.
The situation is simplest if the $`\theta `$-dependence of the cross section is compatible with the $`\pi ^0\pi ^0`$ system being produced only in an $`S`$\- and a $`D`$-wave, and if in addition the $`\phi `$-dependence is flat. Assuming that $`A_{0+}`$ and $`A_+`$ are negligible compared to $`A_{++}`$, one can then decompose
$$\frac{d\sigma _{e\gamma e\pi ^0\pi ^0}}{dQ^2dW^2d(\mathrm{cos}\theta )}=C_{00}+C_{02}P_2(\mathrm{cos}\theta )+C_{22}[P_2(\mathrm{cos}\theta )]^2$$
(115)
and project out the coefficients, using that $`C_{ll^{}}=dS_{e\gamma }(w_{ll^{}})/(dQ^2dW^2)`$ with weights
$`w_{00}`$ $`=`$ $`{\displaystyle \frac{5}{16}}(142\mathrm{cos}^2\theta +49\mathrm{cos}^4\theta ),`$ (116)
$`w_{02}`$ $`=`$ $`{\displaystyle \frac{35}{8}}(16\mathrm{cos}^2\theta +5\mathrm{cos}^4\theta ),`$ (117)
$`w_{22}`$ $`=`$ $`{\displaystyle \frac{35}{16}}(330\mathrm{cos}^2\theta +35\mathrm{cos}^4\theta ).`$ (118)
¿From these coefficients one can readily extract the amplitudes $`|A_{++0}|`$, $`|A_{++2}|`$, and the cosine of their relative phase.
In Fig. 8 (a) we show the coefficients $`C_{00}`$, $`C_{02}`$ and $`C_{22}`$ obtained with our model GDA (92). The interference term between the $`S`$\- and $`D`$-waves contains a factor $`\mathrm{cos}(\delta _0\delta _2)`$ and thus is sensitive to the phase shifts. Characteristic features in the $`W`$-dependence of $`C_{02}`$ are the point where $`\delta _0\delta _2=90^{}`$, and the sudden change just below $`W=1`$ GeV due to the behavior of the $`S`$-wave. To explore the dependence of these observables on our input GDA we have made an ad hoc modification, changing the sign in the prediction (91) from chiral dynamics and taking instead $`B_{10}(0)=B_{12}(0)`$ with $`B_{12}(0)`$ fixed by the constraint (88) as before. Notice that this flips the overall sign $`\eta _{10}`$ of the $`S`$-wave in our model. The result is shown in Fig. 8 (b) and illustrates the sensitivity, especially of the $`S`$-$`D`$ interference, to the detailed dynamics of the $`\gamma ^{}\gamma `$ process.
### C Production of $`\pi ^+\pi ^{}`$ and interference with bremsstrahlung
For the production of $`\pi ^+\pi ^{}`$ pairs in $`e\gamma `$ collisions, the $`\gamma ^{}\gamma `$ reaction we want to study competes with bremsstrahlung, where the pion pair originates from a virtual photon radiated off the lepton , cf. Fig. 9. This process produces the pion pair in the $`C`$-odd channel and hence does not contribute for $`\pi ^0\pi ^0`$. Its amplitude can be fully computed for values of $`W`$ where the timelike electromagnetic pion form factor $`F_\pi (W^2)`$ is known. The modulus of $`F_\pi `$ has been well measured in $`e^+e^{}\pi ^+\pi ^{}`$. By Watson’s theorem its phase is equal to the $`P`$-wave phase shift $`\delta _1`$, provided that $`W`$ is in the range where $`\pi \pi `$ scattering is elastic. This is rather well satisfied for $`W`$ up to 1 GeV. In our numerical studies we use for $`F_\pi `$ the parameterization $`N=1`$ of , which is in good agreement with the data for $`|F_\pi |^2`$ shown in Fig. 10. It also gives a fair description of the phase of $`F_\pi `$ in the $`W`$-range where we use it, as we see from the comparison with the phase shift $`\delta _1`$ in Fig. 6.
The contribution of the $`\gamma ^{}\gamma `$ subprocess to the cross section of $`e\gamma e\pi ^+\pi ^{}`$ has the same form (108) as for $`e\gamma e\pi ^0\pi ^0`$. We recall that with the isospin relation (23) the leading-twist helicity amplitude $`A_{++}`$ in (104) is the same for neutral and for charged pion pairs. The bremsstrahlung contribution reads
$`{\displaystyle \frac{d\sigma _{e\gamma e\pi \pi }}{dQ^2dW^2d(\mathrm{cos}\theta )d\phi }}|_B={\displaystyle \frac{\alpha ^3}{16\pi }}{\displaystyle \frac{\beta }{s_{e\gamma }^2}}{\displaystyle \frac{2\beta ^2}{W^2ϵ}}|F_\pi (W^2)|^2`$ $`(`$ $`[12x(1x)]\mathrm{sin}^2\theta +4x(1x)ϵ\mathrm{cos}^2\theta `$ (119)
$`+`$ $`\mathrm{cos}\phi \sqrt{2x(1x)}(2x1)\sqrt{ϵ(1+ϵ)}\mathrm{\hspace{0.17em}2}\mathrm{sin}\theta \mathrm{cos}\theta `$ (120)
$``$ $`\mathrm{cos}2\phi x(1x)\mathrm{\hspace{0.17em}2}ϵ\mathrm{sin}^2\theta ).`$ (121)
Finally, the interference term of the two subprocesses can be written asA C program containing the expressions (108), (119), (122), (130), as well as the amplitude $`A_{++}`$ calculated with our model GDA (92), can be obtained from the authors.
$$\frac{d\sigma _{e\gamma e\pi \pi }}{dQ^2dW^2d(\mathrm{cos}\theta )d\phi }|_I=2e_l\frac{\alpha ^3}{16\pi }\frac{\beta }{s_{e\gamma }^2}\frac{\sqrt{2}\beta }{\sqrt{W^2Q^2ϵ(1ϵ)}}\left(C_0+C_1\mathrm{cos}\phi +C_2\mathrm{cos}2\phi +C_3\mathrm{cos}3\phi \right)$$
(122)
with $`e_l=1`$ for positrons and $`1`$ for electrons, and coefficients
$`C_0`$ $`=`$ $`\mathrm{Re}\left\{F_\pi ^{}A_{++}\right\}\sqrt{2x(1x)}\sqrt{ϵ(1+ϵ)}\mathrm{cos}\theta `$ (124)
$`+\mathrm{Re}\left\{F_\pi ^{}A_{0+}\right\}(1x)\sqrt{ϵ(1+ϵ)}\mathrm{sin}\theta ,`$
$`C_1`$ $`=`$ $`\mathrm{Re}\left\{F_\pi ^{}A_{++}\right\}[1(1x)(1+ϵ)]\mathrm{sin}\theta `$ (127)
$`\mathrm{Re}\left\{F_\pi ^{}A_{0+}\right\}\sqrt{2x(1x)}\mathrm{\hspace{0.17em}2}ϵ\mathrm{cos}\theta `$
$`+\mathrm{Re}\left\{F_\pi ^{}A_+\right\}(1x)\mathrm{sin}\theta ,`$
$`C_2`$ $`=`$ $`\mathrm{Re}\left\{F_\pi ^{}A_{0+}\right\}x\sqrt{ϵ(1+ϵ)}\mathrm{sin}\theta `$ (129)
$`\mathrm{Re}\left\{F_\pi ^{}A_+\right\}\sqrt{2x(1x)}\sqrt{ϵ(1+ϵ)}\mathrm{cos}\theta ,`$
$`C_3`$ $`=`$ $`\mathrm{Re}\left\{F_\pi ^{}A_+\right\}xϵ\mathrm{sin}\theta .`$ (130)
Remember that in our kinematical limit $`1xW^2/Q^2`$ is small. The structure of the bremsstrahlung contribution (119) then becomes rather simple, since at $`Q^2W^2`$ the terms in large parentheses reduce to $`\mathrm{sin}^2\theta `$. With the scaling predictions for the $`\gamma ^{}\gamma `$ amplitudes discussed in Sect. VII A we also obtain the $`Q^2`$-behavior for each of the coefficients $`C_n`$ in the interference term (122).
The relative dependence on $`Q^2`$, $`W^2`$ and on $`ϵ`$ of the three contributions to the cross section is controlled by the prefactors
$$\frac{1}{Q^2(1ϵ)},\frac{2\beta ^2}{W^2ϵ},\frac{\sqrt{2}\beta }{\sqrt{W^2Q^2ϵ(1ϵ)}}$$
(131)
for $`\gamma ^{}\gamma `$, bremsstrahlung and their interference, respectively, and by the pion form factor $`F_\pi (W^2)`$, which appears linearly in the interference and squared in the pure bremsstrahlung term. The factors $`Q^2`$ and $`W^2`$ in (131) can be traced back to the propagator of the virtual photon in each subprocess, and the extra factor of $`\beta `$ in the bremsstrahlung amplitude reflects the fact that the pion pair is produced in the $`P`$-wave there.
From the factors (131) it follows that the $`\gamma ^{}\gamma `$ contribution decreases faster with $`Q^2`$ than bremsstrahlung. On the other hand the $`\gamma ^{}\gamma `$ process is enhanced at large $`ϵ`$, whereas bremsstrahlung profits from small $`ϵ`$. To study the amplitudes $`A_{ij}`$ either in the $`\gamma ^{}\gamma `$ contribution or in the interference term, one will therefore go to larger values of $`ϵ`$, corresponding to small or intermediate values of $`y`$ (notice that $`ϵ=0.8`$ corresponds to $`y=0.5`$). The behavior in $`Q^2`$ and $`y`$ of the different contributions to the $`ee`$ cross section can be seen in Figs. 11 and 12, respectively, which have again been obtained with our model GDA (92). Notice that apart from the factors (131) just discussed, there is a global dependence on $`y`$ and $`Q^2`$ through the factor $`1/s_{e\gamma }^2`$ in the $`e\gamma `$ cross section and through the variable $`x_2`$ in the real photon flux, cf. Eqs. (11) and (110).
A very strong effect on the relative weight of the different contributions comes from the pion form factor $`F_\pi (W^2)`$. As one can anticipate from Fig. 10 it leads to a considerable enhancement of the bremsstrahlung term in a broad $`W`$ interval around the $`\rho `$ mass, thereby also enhancing the interference. The $`W^2`$-dependence of the different terms, obtained with of our ansatz (92) for the GDA, are shown in Fig. 13. As we discussed in Sect. IV this ansatz most likely oversimplifies the $`W^2`$-dependence of the coefficients $`B_{10}`$ and $`B_{12}`$ in $`\mathrm{\Phi }_q^+`$, but the corresponding error in estimating the $`W^2`$-behavior of $`A_{++}`$ should not change the qualitative picture of Fig. 13.
In the limit of large $`Q^2`$ the different contributions to the cross section have distinctive dependences on $`\phi `$. The $`\gamma ^{}\gamma `$ contribution is predicted to be constant in $`\phi `$ with a $`\mathrm{cos}2\phi `$ modulation due to the product $`A_{++}^{}A_+^{}`$. The bremsstrahlung term should be flat, and the interference between them should be dominated by $`\mathrm{cos}\phi `$ and $`\mathrm{cos}3\phi `$, going with $`A_{++}`$ and $`A_+`$, respectively. We show examples of the $`\phi `$-behavior in Fig. 14, remembering that in our model $`A_+`$ is zero because we have neglected the contribution of the helicity-two gluon GDA. We notice that the $`\mathrm{cos}2\phi `$ term in bremsstrahlung, which is kinematically suppressed by $`1xW^2/Q^2`$, is clearly visible at the larger energy $`W=800`$ MeV. The $`\theta `$-dependence, shown in Fig. 15, is also quite different for the three components of the cross section. For the $`\gamma ^{}\gamma `$ term and the interference it depends in detail on the coefficients of the different partial waves contributing to the amplitudes $`A_{ij}`$.
### D Studying the $`\gamma ^{}\gamma `$ subprocess through the interference term
The interference between the $`\gamma ^{}\gamma `$ and bremsstrahlung subprocesses provides an opportunity to study the $`\gamma ^{}\gamma `$ contribution at *amplitude* level. On one hand this means that one can completely separate the contributions $`A_{++}`$, $`A_+`$ and $`A_{0+}`$ from different photon polarizations. On the other hand it gives access to the phases of these amplitudes relative to the phase of the pion form factor $`F_\pi `$, which is equal to the $`\pi \pi `$ phase shift $`\delta _1`$ in the range of $`W`$ we are considering. In kinematical regions where the bremsstrahlung amplitude is large, especially for $`W`$ around the $`\rho `$ mass peak, the interference can also be used to “amplify” the $`\gamma ^{}\gamma `$ signal.
For this to be useful it is essential that one can cleanly separate the interference term (122) from the pure $`\gamma ^{}\gamma `$ and bremsstrahlung contributions in the cross section. This is possible since the $`\gamma ^{}\gamma `$ collision produces the pion pair in the $`C`$-even channel, whereas in bremsstrahlung $`\pi \pi `$ occurs in the $`C`$-odd projection. The interference term can therefore be separated by reversing the charge of the lepton in the $`e\gamma `$ collision, a possibility that is automatically provided at $`e^+e^{}`$ colliders. Alternatively, any observable that is odd under exchange of the $`\pi ^+`$ and $`\pi ^{}`$ momenta is only sensitive to the interference term, which in turn drops out in any observable even under this exchange. In terms of the variables we are using, this exchange corresponds to the substitution $`(\theta ,\phi )(\pi \theta ,\pi +\phi )`$. This means that we have direct access to the interference through the angular distribution of the pion pair in its rest frame. We emphasize that on the experimental level this does not require a perfect angular measurement, but only that the detection and reconstruction does not introduce a bias between positive and negative pions.
From the $`\phi `$-dependence of the cross section one can extract the four coefficients $`C_n`$ in Eq. (122), which determine the three quantities $`\mathrm{Re}\{F_\pi ^{}A_{++}\}`$, $`\mathrm{Re}\{F_\pi ^{}A_{0+}\}`$ and $`\mathrm{Re}\{F_\pi ^{}A_+\}`$. In fact, they over-determine them, and one can for instance use only $`C_1`$, $`C_2`$, $`C_3`$, and keep the information from $`C_0`$ for a cross check. We remark in passing that this is owed to the fact that pions have zero spin, otherwise there would be more helicity amplitudes for the $`\gamma ^{}\gamma `$ reaction than independent observables one can extract from the $`\phi `$-dependence. Using the $`\phi `$-moments (114) with $`m=1`$, 2, 3 and inverting the relation between $`C_1`$, $`C_2`$, $`C_3`$ and the helicity amplitudes we obtain
$`{\displaystyle \frac{K}{1(1x)(1+ϵ)}}{\displaystyle \frac{dS_{e\gamma }(w_+)}{dQ^2dW^2d(\mathrm{cos}\theta )}}+\{\theta \pi \theta \}`$ $`=`$ $`2\mathrm{Re}\left\{F_\pi ^{}A_{++}\right\}\mathrm{sin}^3\theta ,`$ (132)
$`{\displaystyle \frac{K}{x\sqrt{ϵ(1+ϵ)}}}{\displaystyle \frac{dS_{e\gamma }(w_0)}{dQ^2dW^2d(\mathrm{cos}\theta )}}+\{\theta \pi \theta \}`$ $`=`$ $`2\mathrm{Re}\left\{F_\pi ^{}A_{0+}\right\}\mathrm{sin}^2\theta \mathrm{cos}\theta ,`$ (133)
$`{\displaystyle \frac{K}{xϵ}}{\displaystyle \frac{dS_{e\gamma }(w_{})}{dQ^2dW^2d(\mathrm{cos}\theta )}}+\{\theta \pi \theta \}`$ $`=`$ $`2\mathrm{Re}\left\{F_\pi ^{}A_+\right\}\mathrm{sin}\theta `$ (134)
with a global factor
$$K(Q^2,W^2,ϵ)=e_l\left(\frac{\alpha ^3}{8}\frac{(\beta xy)^2}{Q^4}\frac{\sqrt{2}}{\sqrt{W^2Q^2ϵ(1ϵ)}}\right)^1$$
(135)
and weights
$`w_+`$ $`=`$ $`\mathrm{sin}^2\theta \mathrm{cos}\phi \sqrt{{\displaystyle \frac{2(1x)}{x}}}\sqrt{{\displaystyle \frac{ϵ}{1+ϵ}}}\mathrm{\hspace{0.17em}2}\mathrm{cos}\theta \mathrm{sin}\theta \mathrm{cos}2\phi +{\displaystyle \frac{1x}{xϵ}}(\mathrm{sin}^2\theta +4ϵ\mathrm{cos}^2\theta )\mathrm{cos}3\phi ,`$ (136)
$`w_0`$ $`=`$ $`\mathrm{sin}\theta \mathrm{cos}\theta \mathrm{cos}2\phi +\sqrt{{\displaystyle \frac{2(1x)}{x}}}\sqrt{{\displaystyle \frac{1+ϵ}{ϵ}}}\mathrm{cos}^2\theta \mathrm{cos}3\phi ,`$ (137)
$`w_{}`$ $`=`$ $`\mathrm{cos}3\phi .`$ (138)
By taking weights that are odd under the exchange of the $`\pi ^+`$ and $`\pi ^{}`$ momenta and summing over configurations with $`\theta `$ and $`\pi \theta `$ we have canceled the contributions from the pure $`\gamma ^{}\gamma `$ and bremsstrahlung terms in the cross section. We remark that our method can easily be adapted to the case where one does not have full acceptance in $`\phi `$, since the moments of $`\mathrm{cos}\phi `$, $`\mathrm{cos}2\phi `$ and $`\mathrm{cos}3\phi `$ are always linear combinations of $`\mathrm{Re}\{F_\pi ^{}A_{ij}\}`$.
The functions $`w_i`$ have been chosen such that they are finite, because the use of unbounded weighting functions is problematic. As a consequence, the terms $`\mathrm{Re}\{F_\pi ^{}A_{ij}\}`$ on the r.h.s. of Eq. (132) are still multiplied with functions of $`\theta `$. One can avoid the rather strong suppression of angles $`\theta `$ near 0 or $`\pi `$ in $`\mathrm{Re}\{F_\pi ^{}A_{++}\}\mathrm{sin}^3\theta `$ if the measurement of the moments (132) indicates that $`A_+`$ and $`A_{0+}`$ are small compared with $`A_{++}`$. In this case one may replace the weight $`w_+`$ with $`\mathrm{cos}\phi `$, whose moment is dominated by $`A_{++}\mathrm{sin}\theta `$ with corrections of order $`\sqrt{1x}A_{0+}`$ and $`(1x)A_+`$. Alternatively, the moment of
$$w^{}{}_{+}{}^{}=\mathrm{sin}\theta \mathrm{cos}\phi \sqrt{\frac{2(1x)}{x}}\sqrt{\frac{ϵ}{1+ϵ}}\mathrm{\hspace{0.17em}2}\mathrm{cos}\theta \mathrm{cos}2\phi ,$$
(139)
projects on $`A_{++}\mathrm{sin}^2\theta `$ with corrections only of order $`(1x)A_+`$. In a similar way the moment of $`\mathrm{cos}\theta \mathrm{cos}2\phi `$ approximately projects on $`A_{0+}\mathrm{sin}\theta \mathrm{cos}\theta `$ if $`A_+`$ is sufficiently small.
The $`\theta `$-dependence of the moments (132) contains information on the partial wave decomposition of the pion pair. One way to extract the partial waves is of course to fit the $`\theta `$-dependence of the weighted differential cross sections (132). Alternatively, one can use weighted cross sections integrated over both $`\phi `$ and $`\theta `$. The weight $`\mathrm{cos}3\phi P_l^2(\mathrm{cos}\theta )/\mathrm{sin}\theta `$ readily projects out the $`l`$th partial wave in $`A_+`$ as we easily see from Eq. (132). Note that, since $`P_l^2(\mathrm{cos}\theta )\mathrm{sin}^2\theta `$, this weighting function is a trigonometric polynomial. Similarly, $`\mathrm{cos}2\phi P_l^1(\mathrm{cos}\theta )/\mathrm{sin}\theta `$ can be used to obtain the $`l`$th partial wave in $`A_{0+}`$ if the contribution from $`A_+`$ is small enough.
For $`A_{++}`$ the situation is more complicated, because the functions $`w_+P_l(\mathrm{cos}\theta )/\mathrm{sin}^3\theta `$, $`w^{}{}_{+}{}^{}P_{l}^{}(\mathrm{cos}\theta )/\mathrm{sin}^2\theta `$ and $`\mathrm{cos}\phi P_l(\mathrm{cos}\theta )/\mathrm{sin}\theta `$ are all unbounded. The same problem occurs for the function $`w_0P_l^1(\mathrm{cos}\theta )/(\mathrm{sin}^2\theta \mathrm{cos}\theta )`$. In practice one may proceed as we discussed in Sect. VII B and restrict the analysis to a finite number of partial waves, which has to be determined from the data. Decomposing the coefficient $`C_n`$ in (130) on polynomials $`P_{l+1}^n(\mathrm{cos}\theta )`$ one can see that if only partial waves with $`lL`$ are relevant in the amplitudes $`A_{ij}`$, then weighting the cross section with $`\mathrm{cos}n\phi P_{L+3}^n(\mathrm{cos}\theta )`$ and integrating over $`\phi `$ and $`\theta `$ must give zero. For a restricted number of partial waves one can then find weights to project out the corresponding amplitudes. In the case where $`A_+`$ and $`A_{0+}`$ are negligible and only the partial waves $`l=0`$ and $`l=2`$ are important in $`A_{++}`$, we have for instance
$$\frac{K}{1(1x)(1+ϵ)}\frac{dS_{e\gamma }(w_{+l})}{dQ^2dW^2}=\mathrm{Re}\left\{F_\pi ^{}A_{++l}\right\},l=0,2$$
(140)
with
$$w_{+0}=\frac{4}{3\pi }\mathrm{cos}\phi (1+2\mathrm{cos}^2\theta ),w_{+2}=\frac{16}{3\pi }\mathrm{cos}\phi (14\mathrm{cos}^2\theta ).$$
(141)
In Fig. 16 we show the moments of $`\mathrm{cos}\phi `$, $`w_{+0}`$ and $`w_{+2}`$ as a function of $`W`$ for our model GDA (92) and also for the alternative ansatz described at the end of Sect. VII B. We clearly see the sensitivity of our observables to the detailed phase structure of the $`\gamma ^{}\gamma `$ amplitude.
### E Comparison with lepton pair production
In this section we compare our process $`e\gamma e\pi ^+\pi ^{}`$ with the production of a muon pair, $`e\gamma e\mu ^+\mu ^{}`$, in the same kinematics. This is interesting in itself because $`\mu ^+\mu ^{}`$ production is the QED analogue of the reaction we are studying, but also because it constitutes an experimental background to the extent that a muon pair can be misidentified as a pair of charged pions.
The helicities of the muons can couple to 0 or $`\pm 1`$ along the direction of the $`\mu ^+`$ momentum in the $`\gamma ^{}\gamma `$ c.m. From angular momentum conservation in the subprocesses $`\gamma ^{}\gamma \mu ^+\mu ^{}`$ and $`\gamma ^{}\mu ^+\mu ^{}`$ (the latter occurring in bremsstrahlung) it is clear that the dependence on $`\theta `$ and $`\phi `$ must be different in the cross sections for pion and for muon pair production. We therefore restrict ourselves here to the cross sections integrated over these angles. For the bremsstrahlung contribution we have
$`{\displaystyle \frac{d\sigma _{e\gamma eX}}{dQ^2dW^2}}|_B={\displaystyle \frac{\alpha ^3}{3s_{e\gamma }^2}}{\displaystyle \frac{12x(1x)(1ϵ)}{ϵ}}f_B^X(W^2),`$ (142)
where
$$f_B^{\pi ^+\pi ^{}}=\frac{\beta ^3|F_\pi (W^2)|^2}{W^2},f_B^{\mu ^+\mu ^{}}=\frac{2\beta _\mu (3\beta _\mu ^2)}{W^2}$$
(143)
with the muon velocity $`\beta _\mu =(14m_\mu ^2/W^2)^{1/2}`$ in the $`\gamma ^{}\gamma `$ c.m. For the $`\gamma ^{}\gamma `$ process we can easily adapt the result (99) for open $`q\overline{q}`$-production to the $`\mu ^+\mu ^{}`$ case and find
$$\frac{d\sigma _{e\gamma eX}}{dQ^2dW^2}|_G=\frac{\alpha ^3}{4s_{e\gamma }^2}\frac{1}{Q^2(1ϵ)}f_G^X(W^2),$$
(144)
where
$$f_G^{\pi ^+\pi ^{}}=\left(\frac{25R_\pi }{18}\right)^2\beta \left(1\frac{2}{3}\beta ^2+\frac{1}{5}\beta ^4\right),f_G^{\mu ^+\mu ^{}}=8\left(\mathrm{ln}\frac{1+\beta _\mu }{1\beta _\mu }\beta _\mu \right),$$
(145)
up to corrections of order $`W^2/Q^2`$. Notice that both for bremsstrahlung and for $`\gamma ^{}\gamma `$, the $`Q^2`$-dependence is the same in the pion and the muon case.
The functions $`f_B^X`$ and $`f_G^X`$ are compared in Fig. 17. We see that for the bremsstrahlung contribution pion production is enhanced by the strong resonance effect around the $`\rho `$ mass, as manifested in $`F_\pi (W^2)`$. In the $`\gamma ^{}\gamma `$ subprocess, on the other hand, we find that with our estimate of the GDA, pion production is suppressed compared to muons by a factor 50 to 100. This is mostly due to the numerical constants in the expressions (145). In part it also comes from the logarithm $`\mathrm{log}(1\beta _\mu )`$ in $`f_G^{\mu ^+\mu ^{}}`$, which is generated by the collinear regions around $`\theta =0`$ and $`\pi `$ as discussed in Sect. VI. Notice that for this reason the $`\mu ^+\mu ^{}`$ cross section will be relatively sensitive to cuts that affect $`\theta `$. The same will apply to the interference between bremsstrahlung and $`\gamma ^{}\gamma `$, which drops of course out after angular integration. From the results on $`f_B^X`$ and $`f_G^X`$ we expect that the ratio of muon to pion pair production will be appreciable in the interference term.
Another experimental background, again due to particle misidentification, is $`e^\pm \gamma e^\pm e^+e^{}`$. Compared with $`\mu ^+\mu ^{}`$ production there are further Feynman diagrams, which can be obtained from the muon case by interchanging the lines with momenta $`k^{}`$ and either $`p`$ or $`p^{}`$, now corresponding to identical particles. We shall not analyze these diagrams here, but will at least assess the contributions from those diagrams that are also present in muon production. Replacing $`\beta _\mu `$ with $`\beta _e`$ we obtain velocities extremely close to 1. Nothing dramatic happens in the bremsstrahlung part (143), but the logarithm in the $`\gamma ^{}\gamma `$ subprocess (145) is now much larger than for muons. This large logarithm is however generated by transverse momenta $`p_{}`$ of order $`m_e`$ in the $`\gamma ^{}\gamma `$ c.m., which correspond to extremely small angles $`\theta `$ of order $`m_e/W`$. For any cut that effectively leads to a minimum angle $`\theta _{\mathrm{𝑚𝑖𝑛}}`$ much larger than that, one has to replace $`\beta _\mu `$ with $`\mathrm{cos}\theta _{\mathrm{𝑚𝑖𝑛}}`$ in Eq. (145), which can significantly reduce the size of the logarithm.
We finally note that the differential cross sections for $`e^+e^{}e^+e^{}e^+e^{}`$ and $`e^+e^{}e^+e^{}\mu ^+\mu ^{}`$ have been fully calculated to first order in QED and are available in the form of Monte Carlo generators .
## VIII Cross section estimates
### A Laboratory kinematics and experimental cuts
Before giving our estimates for the cross section of our process at various $`e^+e^{}`$ colliders, we give a brief discussion of the kinematics in the laboratory frame and the effects of some experimental cuts. Starting with the kinematics of the scattered lepton $`k^{}`$, we remark that there is a simple transformation between the variables $`(Q^2,y)`$ and $`(E_1{}_{}{}^{^{}},\alpha _{1L})`$, where $`E_1^{^{}}`$ and $`\alpha _{1L}`$ respectively are the energy and scattering angle of $`k^{}`$ in the laboratory frame. Imposing minimum values on both quantities we have
$$y=1+\frac{Q^2}{4E_1^2}\frac{E_1^{}}{E_1}1+\frac{Q^2}{4E_1^2}\frac{E_1^{min}}{E_1}$$
(146)
and
$$y=1\frac{Q^2}{4E_1^2}\frac{1+\mathrm{cos}\alpha _{1L}}{1\mathrm{cos}\alpha _{1L}}1\frac{Q^2}{4E_1^2}\frac{1+\mathrm{cos}\alpha _{1L}^{min}}{1\mathrm{cos}\alpha _{1L}^{min}}.$$
(147)
The condition (146) cuts on large values of $`y`$ and is generally not very serious, because most information on the $`\gamma ^{}\gamma `$ process is obtained from low or intermediate $`y`$ as we discussed after Eq. (131). The lower cut (147), on the other hand, severely restricts the interesting $`y`$-range in some experimental setups if $`Q^2`$ is not large enough. We will encounter an example of this in Sect. VIII B.
The transformation of the pion momenta into the laboratory system leads to rather lengthy expressions, which we will not give here. Notice that the lepton $`k^{}`$ has a large transverse momentum $`k_L^{}=Q\sqrt{1y}`$ in the laboratory, which must be compensated by the two pions. Even though the $`\pi \pi `$ system has a rather low invariant mass, the pions thus carry large transverse momentum which helps to detect them. An exception are configurations with the the c.m. angle $`\theta `$ close to 0 or $`\pi `$, which in the laboratory correspond to an asymmetric sharing of momentum between the two pions. This is illustrated in Fig. 18.
It is instructive to consider the point where there the $`\pi \pi `$ system has zero longitudinal momentum $`P_L^3`$ in the laboratory. With the approximation $`W^2Q^2`$ we find
$$P_L^3=yE_1\frac{1y}{y}\frac{Q^2}{4E_1},$$
(148)
so that $`P_L^3=0`$ when $`y`$ equals
$$y_0=\frac{Q}{2E_1}\left(\sqrt{1+\frac{Q^2}{16E_1^2}}\frac{Q}{4E_1}\right).$$
(149)
For $`QE_1`$ this simplifies to $`y_0=Q/(2E_1)`$. If $`y`$ is very different from $`y_0`$ the $`\pi \pi `$ system is strongly boosted along the beam axis, and if this boost is too large then one or both pions will go out of the detector acceptance.
We finally have to discuss the kinematics of the scattered lepton $`l^{}`$ in the laboratory. In terms of its scattering angle $`\alpha _{2L}`$ we have, up to electron mass corrections,
$$l_L^{}=(1x_2)E_2\mathrm{sin}\alpha _{2L}$$
(150)
for the transverse component of $`l^{}`$, and
$$Q^2=q^2=(1x_2)E_2^2\left(2\mathrm{sin}\frac{\alpha _{2L}}{2}\right)^2$$
(151)
for the photon virtuality. For small $`\alpha _{2L}`$ we obtain the simple relation
$$Q^2=\frac{l_L^2}{1x_2}.$$
(152)
It turns out that an antitagging condition on the lepton $`l^{}`$, i.e., $`\alpha _{2L}^{}\alpha _{2L}^{max}`$ with $`\alpha _{2L}^{max}`$ determined by the acceptance of a lepton in the detector, is not enough to keep $`Q^2`$ small. With the parameters $`E_2`$ and $`\alpha _{2L}^{max}`$ in Tables II and III we find that, except in the region of $`x_2`$ very close to 1, the maximum values of $`Q^2`$ and $`l_L^2`$ are a few GeV<sup>2</sup>. Under such circumstances it is clearly inappropriate to approximate $`q^2`$ as zero and the momenta $`l`$, $`l^{}`$, $`q^{}`$ as collinear, which we have done throughout this work. Both the kinematical transformation from the $`e\gamma `$ frame to the laboratory and the calculation of the cross section have to be modified then. One must not only recalculate the two-photon and bremsstrahlung processes of Fig. 9 but also include further diagrams contributing to the reaction $`e^+e^{}e^+e^{}\pi \pi `$. Although this is possible in principle, we wish to retain here the simpler expressions for the cross section with one real photon. We therefore require that $`Q^2`$ be small compared with the other kinematical invariants in our problem.
A way to achieve this, suggested by Eq. (152), is to impose an upper cut on $`l_L^{}`$, i.e., in practical terms on the sum $`|𝐤_L^{}+𝐩_L^{}+𝐩_L^{}|`$ of the reconstructed transverse momenta, possibly supplemented by a lower cut on $`1x_2`$. In our numerical studies we determine the maximum virtuality $`Q_{\mathrm{𝑚𝑎𝑥}}^2`$ in the photon flux of Eq. (109) through Eqs. (150) and (151) by requiring both $`\alpha _{2L}^{}\alpha _{2L}^{max}`$ and $`l_L^{}l_L^{\mathrm{𝑚𝑎𝑥}}=100`$ MeV. This leads to considerably smaller virtualities than the antitagging condition alone, although for $`x_2`$ very close to 1 the resulting $`Q_{\mathrm{𝑚𝑎𝑥}}^2`$ is still not very much smaller than $`W^2`$. In practice one may therefore consider an additional cut on $`x_2`$, but we have refrained from this in our estimates. Notice that the $`Q^2`$-spectrum of the photon flux is logarithmic so that a substantial part of the cross section comes from $`Q^2`$ much smaller than $`Q_{\mathrm{𝑚𝑎𝑥}}^2`$.
### B $`B`$-factories
We have now all elements to give cross section estimates for existing $`e^+e^{}`$ facilities. We start with the $`B`$-factories, BABAR, BELLE and CLEO, running at a c.m. energy $`\sqrt{s_{ee}}`$ around 10 GeV. Using our model GDA (92) we calculate the integrated cross section $`\sigma `$ and the individual contributions $`\sigma _G`$ and $`\sigma _B`$ from the $`\gamma ^{}\gamma `$ and bremsstrahlung subprocesses. To project out their interference term we take simple examples of weighted $`e^+e^{}`$ cross sections, $`S_{ee}(\mathrm{sgn}(\mathrm{cos}\phi ))`$ and $`S_{ee}(\mathrm{cos}\phi )`$, defined in complete analogy with the weighted $`e\gamma `$ cross sections (111). We remark that $`S_{ee}(\mathrm{sgn}(\mathrm{cos}\phi ))`$ is simply the left-right asymmetry of the pions in their c.m. We integrate over $`y`$ from its lower kinematical limit
$$y\frac{Q^2+W^2}{4E_1E_2}$$
(153)
up to $`y=0.5`$. Choosing a larger value increases the cross section, but the gain is mainly due to bremsstrahlung. Up to which values of $`y`$ one can extract useful information on the $`\gamma ^{}\gamma `$ process depends of course on the detailed kinematics and must be studied in each particular case. The same is true for the upper limit of the $`Q^2`$-integration. For its lower limit we take 4 GeV<sup>2</sup> as a minimum value where one might expect a lowest-order calculation to be reliable, cf. our discussion in Sect. V. To determine the value of $`Q_{\mathrm{𝑚𝑎𝑥}}^2`$ in the equivalent photon flux we impose the cuts discussed at the end of Sect. VIII A. Our results for $`e^+e^{}e^+e^{}\pi ^+\pi ^{}`$ are given in Table II, where apart from the quantities just discussed we also give the coefficients in the relative statistical errors $`\delta (w)\mathrm{const}/\sqrt{N}`$ of the weighted cross sections $`S_{ee}(w)`$. We see that the results for the different kinematical situations are practically identical. This indicates that it is the cut $`l_L^{}100`$ MeV which determines the real photon flux in most of the relevant parameter space, and not the cut on $`\alpha _{2L}`$, which is different in each of the five cases. We also find that $`S_{ee}(\mathrm{cos}\phi )`$ has a slightly smaller relative statistical error than $`S_{ee}(\mathrm{sgn}(\mathrm{cos}\phi ))`$ and thus greater sensitivity to the interference term.
To estimate the effects of experimental acceptance for the detected particles we impose
* a cut $`\alpha _{1L}^{min}\alpha _{1L}^{}\alpha _{1L}^{max}`$ on the scattering angle $`\alpha _{1L}`$ of the tagged lepton $`k^{}`$,
* a cut $`\theta _L^{min}(\theta _L^{},\theta _L^{})\theta _L^{max}`$ on the polar angles $`\theta _L^{}`$ and $`\theta _L^{}`$ of the pion momenta $`p`$ and $`p^{}`$, measured with respect to the direction of the initial beam lepton $`k`$,
* a minimum transverse momentum of 100 MeV for the tagged lepton and for each of the pions.
All quantities refer of course to the laboratory frame. The results are shown in Table II.
Comparing with Table II we see that the effects of these cuts are generally quite moderate. The strongest effect is observed for BABAR kinematics in the case where the $`e^{}`$ is tagged. This can be traced back to the constraint $`\alpha _{1L}^{min}\alpha _{1L}^{}`$. The minimum value of $`y`$ implied by Eq. (147) for $`Q^2=4`$ GeV<sup>2</sup> is 0.46 in this case, which effectively cuts away all phase space where the $`\gamma ^{}\gamma `$ process is relevant. The situation improves rapidly as $`Q^2`$ goes up, and for $`Q^2=6`$ GeV<sup>2</sup> our cut implies $`y0.19`$. For the other experimental configurations the same cut is much less restrictive: for BABAR kinematics with a tagged $`e^+`$ our cut on $`\alpha _{1L}`$ implies $`y0.18`$ at $`Q^2=4`$ GeV<sup>2</sup>, whereas in the cases of BELLE and CLEO there is not restriction on $`y`$ from the inequality (147) at all, not even at $`Q^2=4`$ GeV<sup>2</sup>.
We find that in the kinematics of $`B`$-factories the interference term is clearly larger than the contribution from $`\gamma ^{}\gamma `$ alone. With several 10 fb<sup>-1</sup> integrated luminosity our estimated cross sections give event rates of order 10,000. As we see from the tables, the relative statistical error on the interference term, extracted through the moments $`S_{ee}(\mathrm{sgn}(\mathrm{cos}\phi ))`$ or $`S_{ee}(\mathrm{cos}\phi )`$ is about 8 to 10 times larger than for integrated cross sections (where it is $`1/\sqrt{N}`$), so that the interference could be measured with statistical errors in the 10% range.
For the production of neutral pion pairs we easily obtain the cross section without cuts by multiplying $`\sigma _G`$ in Table II with a factor 1/2, due to the restricted phase space of identical particles. We refrain from a discussion of the experimental reconstruction of the four-photon state coming from two pion decays, but for an order-of-magnitude indication of event rates one may take half of the cross sections $`\sigma _G`$ in Table II. We then estimate hundreds of events with several 10 fb<sup>-1</sup>, corresponding again to a statistical error around 10%. Thus studies of both charged and neutral pair production seem promising to us.
### C LEP
Let us now investigate the situation at high-energy colliders, taking as examples LEP1 at $`E_1=E_2=45`$ GeV and LEP2 at $`E_1=E_2=95`$ GeV.
In the columns labeled “no cuts” in Table III we list our predicted cross sections, with cuts only on $`l_L^{}`$ and $`\alpha _{2L}`$ so that the real photon flux is defined. For the kinematics we have chosen, the cross sections come out about a factor 2 to 3 larger than at the $`B`$-factories. Luminosities at LEP are however much smaller, so that unfortunately we estimate rather low achievable event rates, and it is not clear to what extent studies of our process in this kinematical regime will be feasible.
To see the effect of cuts on the detected particles we require
* $`\alpha _{1L}^{min}\alpha _{1L}\pi \alpha _{1L}^{min}`$ with $`\alpha _{1L}^{min}=30`$ mrad and $`E_1{}_{}{}^{^{}}0.7E_1`$ for the tagged lepton,
* $`\theta _L^{min}(\theta _L^{},\theta _L^{})\pi \theta _L^{min}`$ with $`\theta _L^{min}=262`$ mrad, corresponding to pseudorapidities $`|\eta |2`$, and a minimum transverse momentum of 100 MeV for each of the pions.
The results are given in the columns “with cuts” of Table III. The most serious restriction here is the cut on the pion angles $`\theta _L^{}`$ and $`\theta _L^{}`$. This can be understood from our considerations after Eq. (149). The value of $`y`$ where the $`\pi \pi `$ system has zero longitudinal momentum in the laboratory is $`Q/(2E_1)`$ and thus of order 0.01 to 0.05 here. Over most of the $`y`$-range the pions are therefore so strongly boosted in the lab that they appear under extremely small angles and cannot be detected. We observe in fact in Table III that the effect of cuts is stronger at LEP2 with its higher beam energy, and that it is more pronounced for bremsstrahlung than for the $`\gamma ^{}\gamma `$ process, the latter being less affected by a loss of events at larger $`y`$.
At LEP1 the cut on $`\alpha _{1L}`$ puts no restriction on $`y`$, but for LEP2 we find that for $`Q^2=4`$ GeV<sup>2</sup> it implies $`y>0.5`$, so that one must go to larger $`Q^2`$. For $`Q^2`$ of about 8 GeV<sup>2</sup> there is no restriction on $`y`$ from the constraint (147) any more.
We finally note that at the very large values of $`Q^2`$ accessible at high-energy colliders one can afford invariant masses $`W`$ well above 1 GeV, while still fulfilling the basic condition $`W^2Q^2`$ of our study. We have not explored this mass region, since our model for the pion GDA is not applicable there. It is however clear that there will be a strong enhancement of the GDAs at $`W`$ around the masses of $`C`$-even resonances, such as the $`f_2`$(1270).
## IX Summary and outlook
In this paper we have analyzed in detail the process $`\gamma ^{}\gamma \pi \pi `$ in the domain where the virtuality $`Q`$ of the $`\gamma ^{}`$ is much larger than the invariant mass $`W`$ of the two-pion system. It factorizes into a parton-level subprocess, which is under perturbative control, and non-perturbative matrix elements called generalized distribution amplitudes. This makes the reaction a laboratory to study the non-perturbative dynamics of a two-pion system forming from a well-defined partonic state, namely from a quark-antiquark or a two-gluon pair produced at small distance. The perturbative stage of the overall process is completely analogous to the one in single-meson production, well studied in the case of a $`\pi ^0`$, $`\eta `$ and $`\eta ^{}`$. It results in a scaling behavior of the amplitude as $`Q^2`$ increases at fixed $`W^2`$, selects characteristic helicity combinations of the two photons, and predicts that the two pions are produced with total isospin zero. The dynamical content of the non-perturbative matrix elements, on the other hand, is more complex than for a single particle. Even the lowest Fock state of $`|\pi |\pi `$, that is, $`q\overline{q}q\overline{q}`$, contains more partons that the initial $`q\overline{q}`$ or $`gg`$ system from which the two pions are formed. In this sense a GDA describes the transition between different parton configurations in the non-perturbative regime. The two-pion distribution amplitude contains the full strong interactions between the two pions, leading to dynamical phases which, by Watson’s theorem, are identical to the phase shifts in elastic $`\pi \pi `$ scattering as long as $`W`$ is below the inelastic threshold. We use this relation as an input for our model GDA, and therefore restrict our study to the $`W`$-region up to 1 GeV.
The evolution equation giving the factorization scale dependence of the GDAs is more complex than for a single pion due to the mixing of $`q\overline{q}`$ or $`gg`$ amplitudes, and we have given the relevant splitting functions and anomalous dimensions for the quantum numbers of relevance here. A simultaneous expansion of $`\mathrm{\Phi }(z,\zeta ,W^2)`$ in the parton momentum fraction $`z`$ and partial waves of the pion system leads to local matrix elements between the vacuum and a two-pion state. By analytic continuation they are related to the moments of the parton distribution functions of the pion. We have used the quark momentum fraction $`R_\pi `$ in the pion, determined from a global fit of these distributions, as an input for our model of $`\mathrm{\Phi }(z,\zeta ,W^2)`$. The corresponding value of $`R_\pi `$ is well below its asymptotic value under perturbative evolution, which may be an indication that the lowest non-asymptotic terms in the crossed-channel quantity $`\mathrm{\Phi }(z,\zeta ,W^2)`$ are not small at factorization scales in the GeV range. We emphasize that the question of how close one is to the asymptotic result of evolution is particularly interesting, because in the case of light pseudoscalars the single-meson distribution amplitudes may be surprisingly close to their asymptotic form even at low scales .
From a theory point of view it is also interesting to consider $`\mathrm{\Phi }_q(z,\zeta ,W^2)`$, defined by the matrix element in Eq. (15), for values of $`W`$ much larger than the scale of non-perturbative interactions. While the dynamics in $`\mathrm{\Phi }_q(z,\zeta ,W^2)`$ is entirely soft for small $`W`$, part of it becomes hard when $`W`$ increases. In the limit $`W1`$ GeV and to leading order in $`\alpha _S`$ one can explicitly write $`\mathrm{\Phi }_q(z,\zeta ,W^2)`$ in terms of a perturbative subprocess and the $`q\overline{q}`$ distribution amplitudes for each separate pion . The resulting $`\mathrm{\Phi }_q(z,\zeta ,W^2)`$ is very far from the asymptotic form in $`z`$. It receives substantial contributions from high partial waves of the $`\pi \pi `$ system, has a power-law falloff like $`1/W^2`$, and its imaginary part is small compared to its real part.
We have constructed a model for the GDA at $`W`$ below 1 GeV, using simple structure as a guide, and $`R_\pi `$ and the $`\pi \pi `$ phase shifts as phenomenological inputs. Comparing the rates for the production of $`\pi \pi `$ and of a single pseudoscalar meson, we found that the hadron spectrum in $`\gamma ^{}\gamma `$ collisions below 1 GeV is strongly dominated by the single resonances $`\pi ^0`$, $`\eta `$, and $`\eta ^{}`$.
We have further compared our process with open $`q\overline{q}`$ production, which at higher invariant masses $`W`$ is commonly used to describe the part of the total hadronic $`\gamma ^{}\gamma `$ cross section due to the pointlike part of the real photon. Interestingly, we find that in our particular kinematical limit, the corresponding scattering amplitude has the same scaling behavior and helicity structure as the one for the exclusive processes $`\gamma ^{}\gamma \pi `$ and $`\gamma ^{}\gamma \pi \pi `$. The main difference is that in the $`\pi `$ and $`\pi \pi `$ cases the collinear divergence of the lowest-order hard scattering diagrams is regulated by the hadronization process. This is encapsulated in the distribution amplitudes, which vanish at the end points $`z=0`$ and 1. In the open $`q\overline{q}`$ calculation, on the other hand, the divergence has to be regulated explicitly. We also note that the sensitivity to the soft end-point region may be larger for pion-pair production than for a single pion, because for two pions the hard scattering and the distribution amplitudes vanish at $`z=1/2`$ for symmetry reasons. Thus one may expect the onset of the scaling behavior to occur at different $`Q^2`$ in the two cases, an issue that will be interesting to study in experiment.
An investigation of the structure of the cross section shows that in $`e\gamma `$ and $`e^+e^{}`$ collisions information on the $`\gamma ^{}\gamma `$ process can be obtained either through the square of the $`\gamma ^{}\gamma `$ amplitude, or from its interference with the bremsstrahlung process if the pions are charged. This interference can readily be projected out by appropriate $`C`$-odd observables, and it offers the opportunity to separate the different $`\gamma ^{}\gamma `$ helicity amplitudes. It further provides direct access to their dynamical phases, although a full phase reconstruction requires polarized beams (cf. Appendix B).
The angular distribution of the pion pair in its c.m. contains detailed information about the dynamics of the $`\gamma ^{}\gamma `$ process. The dependence on the azimuth $`\phi `$ separates the different helicity combinations of the real and virtual photon, each of which plays a distinct role in the scaling limit. In particular it permits one to study leading-twist and non-leading twist amplitudes at the same time, which should provide additional insight into how far one is from the asymptotic regime. The $`\theta `$-dependence, on the other hand, gives access to the partial waves in which the two pions are produced. It is sensitive to the phases, which reflect the dynamics of the $`\pi \pi `$ system and its resonances. Even though one will probably not be able to perform a full extraction of the $`\pi \pi `$ phase shifts in this way, our process provides constraints on these quantities that are independent of the analyses of elastic $`\pi \pi `$ scattering. The presence of higher partial waves would in itself be very interesting, since it gives indirect information on the deviation of $`\mathrm{\Phi }(z,\zeta ,W^2)`$ from its asymptotic form in $`z`$.
We have restricted ourselves to the production of pion pairs in this work, but it is clear that many of our results are also valid for other exclusive systems. The most obvious generalization is to charged or neutral $`K\overline{K}`$ pairs, whose comparison with $`\pi \pi `$ would allow one to study aspects of flavor $`SU(3)`$ breaking in the context of the quark-hadron transition. At even higher values of $`W^2`$ there is the production of $`p\overline{p}`$, where extra spin degrees of freedom come in, as in the well-studied case of the parton distributions of the nucleon.
Another very similar process is the production of $`\mu ^+\mu ^{}`$ pairs, i.e., the QED analogue of our reaction. Comparing the rates of $`e\gamma e\mu ^+\mu ^{}`$ with our estimate for $`e\gamma e\pi ^+\pi ^{}`$ we find that the bremsstrahlung mechanism prefers pions if $`W`$ is in the vicinity of the $`\rho `$ mass, reflecting the strong resonance effect in the $`\pi \pi `$ system. For the production from $`\gamma ^{}\gamma `$, on the other hand, the cross section is considerably larger in the case of muon pairs. We remark that this could not be anticipated from a dimensional analysis. The amplitudes for $`\gamma ^{}\gamma \mu \mu `$ and for $`\gamma ^{}\gamma \pi \pi `$ have the same $`Q^2`$-dependence in our kinematical limit, and the two-pion distribution amplitude, which describes that pions are not pointlike but have internal structure, is a dimensionless quantity.
Using our model GDA to calculate the cross section for $`e^+e^{}e^+e^{}\pi \pi `$, we find encouraging rates for the kinematics and luminosity of $`B`$-factories. Thus there should be enough statistics for detailed studies at these facilities. Our estimates of the effect of cuts also indicate that in the kinematical region interesting in our context, the pions and the tagged lepton are well within the experimental acceptance. For high-energy colliders such as LEP, our predictions are less optimistic, at least in the range of $`W`$ below 1 GeV which we have studied here, due both to the lesser luminosity and the strong longitudinal boost of the pion system.
In conclusion, we find that the process $`\gamma ^{}\gamma \pi \pi `$ can offer valuable insight into the interactions between quarks, gluons and hadrons, and that it should well be measurable at existing $`e^+e^{}`$ facilities.
###### Acknowledgements.
It is a pleasure to thank P. Aurenche, S.J. Brodsky, T. Feldmann, M. Fontannaz, P. Hoyer, L. Mankiewicz, O. Nachtmann, M. Polyakov and O.V. Teryaev for discussions, and H. Marsiske, C. Munger, V. Savinov, S. Söldner-Rembold, S. Uehara and M. Wang for their interest and valuable information about experimental aspects. M.D. thanks CPhT and LPNHE of École Polytechnique for kind invitations.
## A Pion isospin states
We specify in this appendix our sign convention for the definition of pion states. The relative sign for $`\pi ^+`$ and $`\pi ^{}`$ is relevant because it determines the relative sign of the GDAs for charged and neutral pion pairs.
In terms of eigenstates $`|\pi ^i`$ of the isospin operators $`I^i`$ ($`i=1,2,3`$) we define
$$|\pi ^+=\frac{1}{\sqrt{2}}\left(|\pi ^1+i|\pi ^2\right),|\pi ^{}=\frac{1}{\sqrt{2}}\left(|\pi ^1i|\pi ^2\right),|\pi ^0=|\pi ^3.$$
(A1)
Notice that the sign for $`|\pi ^+`$ is opposite to the usual convention for eigenstates of $`SU(2)`$. This has to be remembered when writing down two-pion states with definite isospin using the Clebsch-Gordan coefficients.
The convention (A1) is in line with the customs of field theory, see for instance Sect. 12.5 of . If, starting from the real scalar fields associated with $`|\pi ^1`$ and $`|\pi ^2`$, one constructs the complex scalar field $`\phi `$ which creates $`|\pi ^{}`$ out of the vacuum, then $`|\pi ^+`$ is created by the conjugated field $`\phi ^{}`$. If one used the opposite sign in defining $`|\pi ^+`$, which is more natural in the context of isospin, then there would be an extra minus sign between the fields creating $`|\pi ^{}`$ and $`|\pi ^+`$. Through the LSZ reduction formula this sign would show up in crossing relations. With our definition (A1) this does not happen, and we have for instance that the spacelike pion form factor
$$\pi ^+(p)|J_{\mathrm{em}}^\mu (0)|\pi ^+(p^{})=(p+p^{})^\mu F_\pi (t)$$
(A2)
with $`t=(pp^{})^2`$ becomes
$$\pi ^+(p)\pi ^{}(p^{})|J_{\mathrm{em}}^\mu (0)|0=(pp^{})^\mu F_\pi (s)$$
(A3)
with $`s=(p+p^{})^2`$ in the timelike region. We remark in passing that if one uses the isospin relation (24) and neglects the contributions from strange and heavy quarks, one has the sum rule
$$𝑑z\mathrm{\Phi }_u^{}(z,\zeta ,W^2)=(2\zeta 1)F_\pi (W^2).$$
(A4)
The choice (A1) also leads to a convenient relation for the action of the charge conjugation operator $`C`$, namely
$$C|\pi ^+=|\pi ^{},C|\pi ^0=|\pi ^0.$$
(A5)
The impossibility to find a sign convention that is natural for both charge conjugation and the isospin algebra is discussed at length in Chapt. 5, §7 of (where the other sign in defining $`|\pi ^+`$ was chosen). We also remark that the definition (A1) implies
$$\pi ^+|\overline{u}_\alpha (x)d_\beta (0)|0=\pi ^{}|\overline{d}_\alpha (x)u_\beta (0)|0,$$
(A6)
and therefore a relative plus sign between the distribution amplitudes for $`\pi ^+`$ and $`\pi ^{}`$.
Our definition is the same as the one chosen by Polyakov et al., cf. , and it was also adopted in . We finally mention that the definition leading to Eq. (15) of has the opposite sign for $`|\pi ^+`$.
## B Beam polarization
As we have shown in Sects. VII B and VII D, the unpolarized $`e\gamma `$ cross section contains detailed information on the $`\gamma ^{}\gamma `$ helicity amplitudes $`A_{ij}`$. From Eqs. (108) and (130) it is however clear that this information is not sufficient to fully reconstruct the three independent complex amplitudes $`A_{++}`$, $`A_{0+}`$ and $`A_+`$. For completeness we give in this appendix the expressions of the cross section with longitudinally polarized lepton and photon beams, and discuss what additional information can be obtained from single and double polarization asymmetries.
Starting with the $`\gamma ^{}\gamma `$ contribution, we have
$`{\displaystyle \frac{d\sigma _{e\gamma e\pi \pi }}{dQ^2dW^2d(\mathrm{cos}\theta )d\phi }}|_G=\mathrm{eq}.(\text{108})+{\displaystyle \frac{\alpha ^3}{16\pi }}{\displaystyle \frac{\beta }{s_{e\gamma }^2}}{\displaystyle \frac{1}{Q^2(1ϵ)}}`$ $`(`$ $`P_l\mathrm{sin}\phi \mathrm{Im}\left\{A_{++}^{}A_{0+}^{}A_+^{}A_{0+}^{}\right\}2\sqrt{ϵ(1ϵ)}`$ (B1)
$`+`$ $`P_\gamma \mathrm{sin}\phi \mathrm{Im}\left\{A_{++}^{}A_{0+}^{}+A_+^{}A_{0+}^{}\right\}2\sqrt{ϵ(1+ϵ)}`$ (B2)
$`+`$ $`P_\gamma \mathrm{sin}2\phi \mathrm{Im}\left\{A_{++}^{}A_+^{}\right\}2ϵ`$ (B3)
$`+`$ $`P_lP_\gamma \left\{|A_{++}|^2|A_+|^2\right\}\sqrt{1ϵ^2}`$ (B4)
$``$ $`P_lP_\gamma \mathrm{cos}\phi \mathrm{Re}\{A_{++}^{}A_{0+}^{}+A_+^{}A_{0+}^{}\}2\sqrt{ϵ(1ϵ)}),`$ (B5)
where $`P_l`$ and $`P_\gamma `$ respectively denote the longitudinal polarization of the lepton and photon beam, ranging from $`1`$ to 1. Together with Eq. (108) we see that if both lepton and photon are polarized, one has enough independent terms to reconstruct the real and imaginary parts of the interferences $`A_{++}^{}A_+^{}`$, $`A_{++}^{}A_{0+}^{}`$ and $`A_+^{}A_{0+}^{}`$. Furthermore, the squared terms $`|A_{++}|^2`$ and $`|A_+|^2`$ come with a different relative sign in the unpolarized cross section and the double polarization asymmetry.
The bremsstrahlung contribution to the cross section reads
$`{\displaystyle \frac{d\sigma _{e\gamma e\pi \pi }}{dQ^2dW^2d(\mathrm{cos}\theta )d\phi }}|_B=\mathrm{eq}.(\text{119})+{\displaystyle \frac{\alpha ^3}{16\pi }}{\displaystyle \frac{\beta }{s_{e\gamma }^2}}{\displaystyle \frac{2\beta ^2}{W^2ϵ}}|F_\pi (W^2)|^2P_lP_\gamma `$ $`(`$ $`(2x1)\sqrt{1ϵ^2}\mathrm{sin}^2\theta `$ (B6)
$`+`$ $`\mathrm{cos}\phi \sqrt{2x(1x)}\sqrt{ϵ(1ϵ)}\mathrm{\hspace{0.17em}2}\mathrm{sin}\theta \mathrm{cos}\theta ).`$ (B7)
Notice that it only contributes to the unpolarized cross section and the double polarization asymmetry, but not to single polarization asymmetries. Finally, the interference term can be written as
$`{\displaystyle \frac{d\sigma _{e\gamma e\pi \pi }}{dQ^2dW^2d(\mathrm{cos}\theta )d\phi }}|_I=\mathrm{eq}.(\text{122})2e_l{\displaystyle \frac{\alpha ^3}{16\pi }}{\displaystyle \frac{\beta }{s_{e\gamma }^2}}{\displaystyle \frac{\sqrt{2}\beta }{\sqrt{W^2Q^2ϵ(1ϵ)}}}`$ $`[`$ $`P_l\left(C_1^l\mathrm{sin}\phi +C_2^l\mathrm{sin}2\phi \right)`$ (B8)
$`+`$ $`P_\gamma \left(C_1^\gamma \mathrm{sin}\phi +C_2^\gamma \mathrm{sin}2\phi +C_3^\gamma \mathrm{sin}3\phi \right)`$ (B9)
$`+`$ $`P_lP_\gamma (C_0^{l\gamma }+C_1^{l\gamma }\mathrm{cos}\phi +C_2^{l\gamma }\mathrm{cos}2\phi )]`$ (B10)
with coefficients
$`C_1^l`$ $`=`$ $`\mathrm{Im}\left\{F_\pi ^{}A_{++}\right\}x\sqrt{1ϵ^2}\mathrm{sin}\theta `$ (B12)
$`+\mathrm{Im}\left\{F_\pi ^{}A_+\right\}(1x)\sqrt{1ϵ^2}\mathrm{sin}\theta ,`$
$`C_2^l`$ $`=`$ $`\mathrm{Im}\left\{F_\pi ^{}A_{0+}\right\}x\sqrt{ϵ(1ϵ)}\mathrm{sin}\theta `$ (B14)
$`\mathrm{Im}\left\{F_\pi ^{}A_+\right\}\sqrt{2x(1x)}\sqrt{ϵ(1ϵ)}\mathrm{cos}\theta `$
for lepton polarization,
$`C_1^\gamma `$ $`=`$ $`\mathrm{Im}\left\{F_\pi ^{}A_{++}\right\}[1(1x)(1ϵ)]\mathrm{sin}\theta `$ (B17)
$`+\mathrm{Im}\left\{F_\pi ^{}A_{0+}\right\}\sqrt{2x(1x)}\mathrm{\hspace{0.17em}2}ϵ\mathrm{cos}\theta `$
$`\mathrm{Im}\left\{F_\pi ^{}A_+\right\}(1x)\mathrm{sin}\theta `$
$`C_2^\gamma `$ $`=`$ $`\mathrm{Im}\left\{F_\pi ^{}A_{0+}\right\}x\sqrt{ϵ(1+ϵ)}\mathrm{sin}\theta `$ (B19)
$`+\mathrm{Im}\left\{F_\pi ^{}A_+\right\}\sqrt{2x(1x)}\sqrt{ϵ(1+ϵ)}\mathrm{cos}\theta ,`$
$`C_3^\gamma `$ $`=`$ $`\mathrm{Im}\left\{F_\pi ^{}A_+\right\}xϵ\mathrm{sin}\theta `$ (B20)
for photon polarization, and
$`C_0^{l\gamma }`$ $`=`$ $`\mathrm{Re}\left\{F_\pi ^{}A_{++}\right\}\sqrt{2x(1x)}\sqrt{ϵ(1ϵ)}\mathrm{cos}\theta `$ (B22)
$`\mathrm{Re}\left\{F_\pi ^{}A_{0+}\right\}(1x)\sqrt{ϵ(1ϵ)}\mathrm{sin}\theta ,`$
$`C_1^{l\gamma }`$ $`=`$ $`\mathrm{Re}\left\{F_\pi ^{}A_{++}\right\}x\sqrt{1ϵ^2}\mathrm{sin}\theta `$ (B24)
$`\mathrm{Re}\left\{F_\pi ^{}A_+\right\}(1x)\sqrt{1ϵ^2}\mathrm{sin}\theta ,`$
$`C_2^{l\gamma }`$ $`=`$ $`\mathrm{Re}\left\{F_\pi ^{}A_{0+}\right\}x\sqrt{ϵ(1ϵ)}\mathrm{sin}\theta `$ (B26)
$`+\mathrm{Re}\left\{F_\pi ^{}A_+\right\}\sqrt{2x(1x)}\sqrt{ϵ(1ϵ)}\mathrm{cos}\theta `$
if both lepton and photon are polarized. We see that with polarized photons one can extract $`\mathrm{Im}\{F_\pi ^{}A_{++}\}`$, $`\mathrm{Im}\{F_\pi ^{}A_{0+}\}`$ and $`\mathrm{Im}\{F_\pi ^{}A_+\}`$, which together with the unpolarized interference term makes it possible to reconstruct all three complex $`\gamma ^{}\gamma `$ amplitudes for values of $`W`$ where the pion form factor $`F_\pi `$ is known. One cannot achieve the same with a polarized lepton beam alone, since there are only two terms in the $`\phi `$-dependence. In this case one can still use the suppression by $`1x`$ of the second term in $`C_1^l`$ in order to approximately extract $`\mathrm{Im}\{F_\pi ^{}A_{++}\}`$. Finally, the double polarization asymmetry gives access to the same quantities one can already obtain in the unpolarized case. |
warning/0003/astro-ph0003215.html | ar5iv | text | # Persistent fluctuations and scaling properties in galaxy number counts
## 1 Introduction
In the counts of galaxies, large fluctuations from field to field and from author to author, both in faint and bright counts, and in different spectral bands, have been reported (e.g. Shanks et al 1989, Tyson 1988, Cowie et al. 1990, Maddox et al. 1990, Metcalfe et al., 1991, Picard 1991, Weir et al. 1995, Bertin & Dennefeld 1997, Arnout et al 1997). These fluctuations can be as large as a factor of two. There have been controversy as to whether these fluctuations are due to real clustering or to differences in the magnitude zero point of the various surveys. Hence, in order to avoid possible systematic errors, it is very important to understand the nature of fluctuations in a given field of a single survey, once the magnitude system and zero point have been carefully calibrated. It is, in fact, possible that discrepancies among these surveys are not due mostly to differences in photometric systems or in data reduction effects, but rather to real effects, i.e. large scale structures. In this Letter we propose a method to verify this latter possibility in the actual data. The slope and the amplitude of the counts are shown to be compatible with a fractal distribution of galaxies, and we point out that fundamental information about clustering can obtained by studying the fluctuations of counts as a function of apparent magnitude.
## 2 Average number counts in a fractal distribution
As suggested in Baryshev (1981), and proposed in a series of papers (Sylos Labini et al. 1996, Montuori et al. 1997, Sylos Labini Montuori & Pietronero 1998) number counts versus apparent magnitude can be used to test whether the large scale distribution of galaxies can be compatible with a fractal or with an homogeneous behavior. In this context, we discuss the case in which the joint space-luminosity distribution $`\nu (\stackrel{}{r},L)`$ can be factorized as the product of the number spatial density $`n(\stackrel{}{r})`$ and the luminosity function $`\varphi (L)`$ <sup>1</sup><sup>1</sup>1which we take to be Schechter like (Schechter 1976) but its actual function does not change the final results (Binggeli, Tammann & Sandage 1988):
$$\nu (\stackrel{}{r},L)d^3rdL=n(\stackrel{}{r})\varphi (L)d^3rdL.$$
(1)
This is known to be a good approximation in the case of small redshift ($`z1`$). All the eventual corrections to Eq. 1 (space geometry, K-corrections, evolution, etc.) are in fact proportional to $`z`$ (Yoshii & Takahara 1988, Sandage 1995).
In the case of a fractal distribution, the average density seen from a galaxy (averaged over enough many observing galaxies) can be written as $`n(\stackrel{}{r})\mathrm{\Gamma }(r)=Br^{D3}`$ (Pietronero 1987, Sylos Labini et al. 1998) where $`D`$ is the fractal dimension; then:
$$\nu (\stackrel{}{r},L)=\mathrm{\Gamma }(r)\varphi (L)=Br^{D3}AL^\delta e^{\frac{L}{L_{}}}.$$
(2)
In this case one ends up with a very simple relation for the integrated counts as a function of apparent flux ($`f=L/(4\pi r^2`$)), for unit of steradian:
$$N(>f)=N_0f^{\frac{D}{2}}.$$
(3)
By using the transformation between apparent flux and magnitude (Peebles 1993)
$$f=\frac{L_{}}{4\pi (10\text{pc})^2}10^{0.4(M_{}m)},$$
(4)
where $`M_{}`$ is the cut-off of the luminosity function $`L_{}`$ in terms of magnitude, one obtains
$$N(<m)=\stackrel{~}{N}_010^{\frac{D}{5}m},$$
(5)
and hereafter we denote $`\alpha D/5`$.
Note that, from what concerns the average behavior, the case of a homogeneous distribution is included in the fractal case with $`D=3`$. Eq.2 has been tested (Sylos Labini, Montuori & Pietronero 1998; Joyce, Montuori & Sylos Labini 1999) to be a rather good approximation in local redshift surveys. Thus, the exponent of the average counts is simply related to the fractal dimension of galaxies in the three dimensional space (see also Sandage, Tammann & Hardy 1972, Peebles 1993). In Eq. 2 $`A`$ is a normalizing constant such that
$$A=\frac{1}{_{L_{min}}^{\mathrm{}}L^\delta e^{L/L_{}}𝑑L},$$
(6)
where $`L_{min}`$ is the faintest object observed in current surveys. Such a lower cut-off, larger than zero is necessary to avoid divergences for $`\delta 1`$. Therefore, Eq. 2 depends on a combination of five different parameters which can be independently measured. Three parameters are related with the luminosity function: the exponent $`\delta `$, the luminosity cut-off $`L_{}`$ and the lower cut-off $`L_{min}`$. These three quantities have been measured with good precision in different redshift surveys (Binggeli, Sandage & Tammann 1988, Efsthatiou, Ellis & Peterson 1988). The fourth parameter is the fractal dimension $`D`$ and the last one, $`B`$, is the absolute normalization of the fractal distribution. This latter can be, for example, defined as the average number of galaxies of any luminosity as seen by an average observer in a ball of radius $`1h^1Mpc`$ and can be measured in redshift surveys (see Sylos Labini, Montuori & Pietronero 1998, Joyce, Montuori & Sylos Labini 1999 for a more detailed discussion of the subject).
The amplitude $`N_0`$ in Eq.3 is given by
$$N_0=\frac{AB}{2(4\pi )^{\frac{D2}{2}}}L_{}^{\delta +\frac{D+2}{2}}\mathrm{\Gamma }_e\left(\delta +\frac{D}{2}\right),$$
(7)
where $`\mathrm{\Gamma }_e`$ is the Euler’s Gamma function.
In view of Eq.2 and Eq.4, one can compute the average redshift of a galaxy with apparent magnitude $`m`$. We obtain
$$z=\frac{h}{310^8}\frac{\mathrm{\Gamma }_e\left(\frac{D+3}{2}+\delta \right)}{\mathrm{\Gamma }_e\left(\frac{D+2}{2}+\delta \right)}10^{0.2(mM_{})},$$
(8)
where $`h`$ is the normalized Hubble’s constant.
From current data both the amplitude and the slope of counts can be estimated. In general in the standard $`B_J`$ photometric system<sup>2</sup><sup>2</sup>2We adopt hereafter the standard Johnson-Cousins system following the choice of Arount et al (1997) and of Bertin & Dennefeld (1997)., and in the range of magnitude from $`11^m`$ to $`19^m`$ (see Tab.1), one has $`\alpha =0.50\pm 0.04`$ corresponding to $`D=2.5\pm 0.2`$. The corresponding range of average redshift (Eq. 8) is $`10^3\mathrm{}<z\mathrm{}<10^1`$. Note that, in the faint end part of counts, where the cosmological corrections are known to be relevant, the slopes are consistent with the bright end (see Tab.2).
Such a value of $`\alpha `$ (and hence of $`D`$) is slightly larger than the value of $`D`$ found in nearby redshift surveys, which is $`D=2.1\pm 0.1`$ up to $`30÷50h^1Mpc`$. Whether such a difference is due to an increase of fractal dimension with scale or it is related to some systematic effects in the counts will be discussed in forthcoming papers (e.g. Gabrielli & Sylos Labini 2000). It is worth to note that Teerikorpi et al. (1988) have found a dimension $`D=2.35\pm 0.05`$ up to $`100h^1Mpc`$ by counting galaxies in real space and in volume limited samples.
Note that in the range of $`z1`$ one expects eventual cosmological and luminosity evolution corrections to be negligible. However, we propose a further test to discriminate the importance of these effects.
More specifically, we propose to study in detail the fluctuations around the average behavior of number counts as a function of apparent magnitude. In fact, as shown below, this test can discriminate between the fractal or smooth cosmological nature of the deviation of the $`\alpha `$ exponent from Euclidean behavior ($`\alpha =0.6`$). This study is motivated by the fact that through number counts we can analyze much larger space volumes than in redshift surveys. In fact, the deepest actual red-shift surveys where the fractal dimension has been estimated (Joyce et al., 1999) contains some thousand galaxies, whereas magnitude limited surveys can have as many as some millions of galaxies up to very faint magnitudes and deep scales (e.g. POSS-II).
## 3 Fluctuations
A very illustrative and simple case is a poissonian homogeneous distribution of galaxies. In this case the difference between the number of points in two equal non overlapping volumes is of the order of the square root of the average number. The variance of counts can be easily computed from the probabilistic definition of Poisson distribution $`n(\stackrel{}{r})`$ and by using again Eq.1, obtaining
$$\sigma _m^2=\frac{(N(<m)N(<m))^2}{N(<m)^2}10^{0.6m},$$
(9)
where $`N(<m)`$ is the number of galaxies with apparent magnitude brighter than $`m`$. The average $`N(<m)`$ is given by Eq.5 with $`D=3`$. A more rigorous derivation, considering three point correlation function, can be found in Gabrielli & Sylos Labini (2000). Thus, in the poissonian case, relative fluctuations decrease exponentially at faint magnitudes. The pre-factor in Eq.9 is simply related to a combination of the parameters in Eq.2.
In a fractal distribution the typical fluctuation of the number of points $`N(r)`$ in a sphere of radius $`r`$, with respect to the average value over different observers $`N(r)`$, is always of the same order of the average number (e.g. Mandelbrot 1977):
$$\delta N(r)=\sqrt{(N(r)N(r))^2}N(r).$$
(10)
This property is very important for counts, which are not averaged over different observers (Sylos Labini Montuori & Pietronero 1998). Eq. 10 means that, in a fractal, at any scale, one expects to find a void or a structure, the extension of which is of the same order of the scale itself: this is the source of geometrical self-similarity. This property implies that fluctuations in the number of points (differential or integral) should be, in absolute value, always proportional to the average number itself and never decreases with distance.
From Eq. 10 and Eq. 1 one obtains that the relative fluctuation in the counts as a function of apparent magnitude has a constant amplitude:
$$\sigma _mconst.>0.$$
(11)
Eq. 11 describes the “persistent” character of fluctuations in number counts induced by the fractal nature of the spatial distribution. The numerical value of $`\sigma _m`$ depends now on the same parameters in Eq.1, and on some other morphological characteristics of the specific studied fractal. In fact, fluctuations are characterized by higher order correlation functions (Blumenfeld & Ball 1993; Gabrielli, Sylos Labini & Pellegrini 1999) and the fractal dimension does not determine them univocally. Note that $`\sigma _m`$ can be also very small: its striking feature being in fact that it is constant as a function of $`m`$, and not its absolute amplitude. By using simple approximations, it is possible to relate the constant $`\sigma _m`$ to three point correlation function of the distribution (Gabrielli & Sylos Labini 2000). In a deterministic fractal, fluctuations have a nearly constant amplitude with a log-periodical modulation (Sornette, 1998) as a function of scale, because the algorithm generating such a structure is a deterministic one. In the more realistic case of stochastic fractals, the oscillations are in general a superposition of waves, which are periodic in log-space, but which have different frequencies and amplitudes.
The poissonian case describes also the situation in which one has a spatial distribution of galaxies with a small crossover scale $`\lambda _0`$ to homogeneity and a finite correlation length $`r_c`$ (Gaite et al. 1999, Gabrielli, Sylos Labini & Durrer 2000). A different situation occurs in the case of a spatial distribution with a finite homogeneity scale, but an infinite correlation length. This case can be thought as obtained by a superposition of a fractal distribution to a dominating flat constant density. In this case $`N(<m)`$ is again given by Eq. 7 with $`D=3`$ (i.e. it is dominated by the flat constant distribution). On the other hand the absolute fluctuation $`(N(<m)N(<m))^2`$ is dominated by the fractal scale invariant correlations. Consequently, the normalized varaince $`\sigma _m^2`$ is again an exponentially decreasing function of $`m`$, even if with a slower behavior than a poissonian distribution
$$\sigma _m^210^{(0.2(3D)m)},$$
(12)
where $`D<3`$ is the dimension of the fractal superimposed to the constant density.
Note that, in general one can have more complex situations, but the case described by Eq. 11 is an unambiguous indication of persistent and scale invariant real space fluctuations typical of statistically self similar irregular distributions.
## 4 Discussion and Conclusions
Up to $`B_J18^m`$ (i.e. $`z<0.1`$) cosmological models, in the framework of the Friedmann solutions, predict an exponent $`\alpha =0.6`$. This is because, one assumes an homogeneous distribution starting at very small scale, i.e. $`5÷20h^1Mpc`$ (Peebles 1993, Davis 1997, Wu, Lahav & Rees 1999). Evolutionary effects, as other cosmological corrections, become efficient at $`z1`$ (Yoshii & Takahara 1989, Sandage 1995). Such a situation is nearly independent on the value of $`q_0`$, the amount of K-corrections, the possible evolution of galaxies with redshift, and the photometric band chosen. It is important to note that $`N(m,q_0)`$ is degenerate to $`z`$ in first order: i.e. it is independent on $`q_0`$ for small redshift. For instance, at $`z=0.1`$ (i.e. $`B_J18^m`$) the deviation from the Euclidean $`\alpha =0.6`$ slope is less than $`10\%`$ for any value of $`q_0`$. Moreover, the slope at fainter magnitudes should be a rapidly varying function of the magnitude itself. Clearly, this is not the case for the data shown in Tab.1 and Tab.2.
From an experimental point of view, we propose to study the fluctuations with respect to the average in the integrated number counts $`N(<m)`$ instead of in the differential one $`N(m)`$, in order to avoid problems with shot noise in magnitude bins. In such a way it is clear that at bright magnitudes $`\sigma _m`$ shows an initial decay due to the paucity of bright galaxies in small solid angle fields. Then, after the integrated number of points has reached a large enough value, one should be able to detect only the effect of eventual intrinsic fluctuations.
However, calibration errors or other systematic field-to-field possible biases can affect the measurement. For this reason we suggest first to focus the study to a single sky field at time, instead of considering fluctuations in different sky fields. That is, after having determined the best fit as in Eq. 3 (or by adding eventual cosmological corrections) one can study fluctuations in a single well calibrated sky field. The instrisic nature of fluctuations reveals as soon as the shot noise contribution becomes negligible. Clearly, there is a transient between the shot noise regime and the instrinsical fluctuations, where the there is a combination of these two effects. The range of magnitudes of such a combination of shot noise and intrinsic fractal fluctuations, depends not only on the number of points, but also on the solid angle of the survey (see Sylos Labini, Montuori & Pietronero 1998 for a more detailed discussion about the shot noise effect on the galaxy counts).
It is important to note that the presence of eventual persistent and scale-invariant fluctuations, in the $`\mathrm{log}N(<m)`$ vs. $`m`$ plot, cannot be due to any smooth correction to the data as cosmological and evolution effects, but they can be the outcome exclusively of strongly correlated fractal fluctuations. The reason being that smooth linear corrections are not able to produce persistent scale-invariant fluctuations on $`N(<m)`$ of the same order of $`N(<m`$) itself.
It is important to note that the exponent of the counts can be very sensible to the photometric band chosen, due to the different K-correction and kind of objects selected. However, it is important to stress that the nature of fluctuation structures must be the same in all different photometric bands. In other words, since these fluctuations are intrinsic, then they should not depend on the photometric band used. On the contrary other possible intervening effects, like galactic extinction fluctuations, strongly depend on the photometric band chosen.
From the present discussion an important challenge for the new generation of experiments is represented by the following questions: (i) Why the slope and the amplitude of the counts remain nearly constant beyond $`B_J18^m`$? (ii) Why there is no clear sign of change of slope due to galaxy evolution, space-time geometry effects and K-correction, even when the average redshift becomes to be of order unity ? (iii) The last question concerns the detection of persistent and scale-invariant fluctuations in the counts. The new generation of redshift surveys like SSDS and 2dF, together with the new POSS-II (Djorgovski et al., 2000) photometric survey, will be able to answer to these fundamental questions.
## Acknowledgements
F.S.L. warmly acknowledge continuous and detailed comments of Y.V. Baryshev. We are grateful to E. Bertin for illuminating discussions about data. We also thank H. Di Nella, G. Djorgovski, R. Durrer, J.-P. Eckmann, M. Joyce, G. Mamon, M. Montuori, G. Paturel, L. Pietronero and P. Teerikorpi for very useful comments and suggestions. This work is partially supported by the EC TMR Network ”Fractal structures and self-organization” ERBFMRXCT980183 and by the Swiss NSF. |
warning/0003/astro-ph0003221.html | ar5iv | text | # The reality of old moving groups - the case of HR 1614 Based on observations with the ESA Hipparcos satellite
## 1 Introduction
Globular and open clusters provide useful probes of the longterm chemical and dynamical evolution of the Milky Way. The globular clusters probe the formation and early evolution of the spheroidal components of the Milky Way while the open clusters provide a useful tool to study the evolution of the galactic disk. However, the paucity of very old (old is here taken to mean $`10^8`$ years, the time scale on which an open cluster will be dissolved, Spitzer 1958) open clusters in the disk forces us to consider the considerably more loosely arranged moving groups to probe the earlier evolution of the disk. It was Olin Eggen who first introduced the concept of moving or stellar kinematic groups, of which the Hyades is a well known example. The basic idea behind the moving groups is that stars form in clusters and thus with similar space motion, on top of which the random motions of single stars are added, resulting in a modest velocity dispersion within the group. Through the orbital motion within the galactic potential the group will be stretched out into a tube-like structure and finally, after several galactic orbits, dissolve. The result of the stretching is that the stars will appear, if the Sun happens to be inside the tube, all over the sky but may be identified as a group through their common space velocity. Thus the moving groups may provide the essential, and so far largely un-utilized, link between cluster and field stars. These are the assumptions, but are moving groups observable realities? A large stumbling block for assessing the reality of moving groups has been the lack of large numbers of reliable parallaxes. This has now been largely overcome by the observations from the Hipparcos satellite (ESA 1997). This has, in fact, resulted in a small burst of recent papers studying, mainly young, moving groups, e.g. Asiain et al. (1999), Barrado y Navascués (1998), Odenkirchen et al. (1998), Skuljan et al. (1997), and Dehnen (1998).
Eggen defined moving groups as stars that all share the same velocity in the direction of galactic rotation, i.e. V-velocity. Specifically the velocities required to be constant were corrected for the stars differing radial distance from the Sun in order to make the circular orbits iso-periodic (Eggen 1998b). However, using the Hipparcos parallaxes it is noted that firstly the groups get more compact and secondly that stars identified as group members do not form flat bars or ellipses with small $`\sigma _V`$ but are in fact structures tilted in the $`UV`$-plane (Skuljan et al. 1997 Fig. 1). Part of this shape can be attributed to the errors in the parallaxes themselves and their transformation into errors in the $`UV`$-plane. However, through dynamical simulations Skuljan et al. (1997) show that all of the tilt cannot be attributed to the errors in parallaxes but also has a physical basis.
In view of these new possibilities it is now appropriate to re-asses the reality and membership criteria for the HR 1614 moving group. Eggen (1998b) has compared Hipparcos and cluster parallaxes for stars in his sample (Eggen 1992) of HR 1614 moving group member stars. However, he disregards the Hipparcos parallaxes in favour of cluster parallaxes, also when the discrepancies are large, without further discussion.
The article is organized as follows: Sects. 2, 3, and 4 describe the search for the HR 1614 moving group in the Hipparcos catalogue, as well as dynamical simulations of old moving groups and their characteristics today. Sect. 5 reviews previous work on the HR 1614 moving group. A new selection criterion for HR 1614 moving group is developed in Sect. 6 and used to derive its age. In Sect. 7. we derive, from data in the literature, a metallicity for the moving group. The Eggen (1992, 1998b) sample is revisited in Sect. 8 and discussed in detail. Sect. 9 contains a discussion primarily of possible sources of contamination in our sample. Sect. 10 provides a brief summary of the main results of this paper.
## 2 Finding the HR 1614 moving group – Searching the Hipparcos catalogue
We perform an unbiased search in the Hipparcos catalogue complemented with radial velocities from the Hipparcos Input Catalogue (ESA 1992), Grenier et al. (1999), and Barbier-Brossat et al. (1994) over a wide area in $`UV`$-space (Fig. 1 shows the distribution of stars) around the probable values for the HR 1614 moving group to see if we can find any signature of what appears to be stars with similar, high, metallicities and with correlation in the $`UV`$-plane. First we divide the $`UV`$-plane into seven boxes and construct the corresponding HR-diagrams, Figs. 1 and 2.
The four boxes 1, 3, 4, and 6, all show stellar population with metallicities around solar and below. Box 3 has what appears to be a younger population as well, while the other boxes show exclusively old populations. None of these boxes show any trace of a large population significantly more metal-rich than the Sun and we will disregard them from further discussions and concentrate on the remaining three boxes.
The three boxes 2, 5, and 7 all show well populated HR-diagrams. As we move from box 2 to 7 over box 5 the turn-off age of stars with solar-like metallicities increases. The canonical view of the stellar populations in the solar neighbourhood implies that we would expect subsequently more and more metal-poor stars as we move further away from the Sun in velocity space. Especially we would expect box 5 to be older and more metal-poor than box 2 and box 7 even older and more metal-poor as we sample more and more halo and thick disk stars and less of the disk stars.
The visual impression from Fig. 2 is the presence of a a richly populated metal-rich isochrone present in box 2 and 5 but not in any other box. However the relative number of metal-rich stars appears much larger in box 5 than in box 2. Is this significant? To find out we combine our Hipparcos catalogue with Strömgren photometry (Hauk & Mermilliod 1998) and calculate \[Me/H\] using the calibration in Schuster & Nissen (1989). All stars flagged as possible binaries in the Hipparcos catalogue were excluded. This includes both systems detected by Hipparcos itself and system previously known (e.g. from radial velocity variation) and included in the CCDM catalogue of multiple stars (ESA 1997). These stars have to be excluded since it is not possible to derive metallicities from their Strömgren photometry. The results are shown in Figs. 3 and 4 where we also distinguish between three major metallicity ranges, \[Me/H\] $`<0.1`$, $`0.1`$\[Me/H\]$`<+0.1`$, and $`+0.1`$ \[Me/H\]. Again box 5 shows a clear high-metallicity population not at all present in the other field. The resulting normalized metallicity-distributions for box 2, 5, and 7 are shown in Fig. 5. Clearly box 7 has a large, metal-poor tail, as expected since this box should contain many halo stars as well as thick disk stars. Box 2 on the other hand, which should be the most solar-like box, has almost no metal-poor tail and a distinct peak at $`0.2`$ dex, just what is expected for the solar neighbourhood, Wyse & Gilmore (1995). Box 5 on the other hand is more metal-rich than the solar neighbourhood. This does not fit into the canonical picture of the general stellar populations in the Galaxy. It should be more metal-poor (on average) than box 2 and more metal-rich than box 7.
The background distribution of stars in Fig. 1 shows that the maximum density occurs in several clumps rather close to the local standard of rest (LSR) but also that a secondary maximum exists centered on $`U20`$ km s<sup>-1</sup>, $`V40`$ km s<sup>-1</sup> (see also Dehnen 1999, Fig. 1 for a similar smoothed density distribution plot). This stellar population with a mostly negative U-velocity, was already tentatively identified in pre-Hipparcos data and named the U-anomaly. In the Hipparcos data it is much more well defined. Raboud et al. (1998) who obtained Geneva photometry for stars in this area of the $`UV`$-space, found a metallicity distribution similar to that of the bulge. Supported by dynamical simulations they concluded that the stars formed in the inner disk were scattered by the bar into the solar neighborhood. These stars make up a large part of the old stars found in box 2. Figure 4 shows the velocity distribution of stars with Strömgren metallicities in box 2 and 5. The U-anomaly is clearly seen in box 2 and found to be dominated by stars of solar metallicity and below.
From this and our previous considerations it is clear that the super-solar metallicity stars are predominantly found in box 5, the area of the $`UV`$-plane in which the HR 1614 moving group is supposed to be found. We thus conclude that there exists a distinct stellar population more metal-rich than the average background centered at $`U10`$ and $`V60`$ km s<sup>-1</sup> and tilted in the $`UV`$-plane. This population is not found in the adjacent areas in the $`UV`$-plane. Can this stellar population be identified with the moving group HR 1614 found by Eggen?
## 3 Finding the HR 1614 moving group – Dynamic simulations
To answer this question we now turn to a dynamical simulation of the evolution of old moving groups to find out whether or not the structure we observe in box 5 can indeed be identified with an old moving group.
A galactic potential model consisting of a stellar and gaseous disk, thick disk, bulge and dark halo (Dehnen & Binney 1998) is used to integrate the orbits of the individual stars in a dissolving moving group. In the model the solar distance from the galactic centre is R$`{}_{}{}^{}=8.0`$ kpc, the height above the plane is z$`{}_{}{}^{}=8`$ pc, the solar peculiar velocity components are $`U=10`$ km s<sup>-1</sup>, $`V=5`$ km s<sup>-1</sup>, $`W=7`$ km s<sup>-1</sup> and the circular velocity $`v_\mathrm{c}(R_{})=219`$ km s<sup>-1</sup>. The starting point for the simulation is attained by taking a velocity vector and position from a star assumed to belong to the moving group e.g. HR 1614 itself and integrate the orbit backwards for the assumed age of the moving group. At this position an ensemble of stars is placed and followed forward into the present time. The stars are treated as test particles in the static potential and the dispersion processes is modeled in two different ways. Either the moving group is born unbound like an OB association with a certain velocity dispersion and no further dispersion processes are active or the group is born as a bound system like an open cluster with identical velocities and the dispersion is gradually built up by stars becoming unbound to the cluster and starting to experience orbital diffusion (Wielen 1977).
Both methods are rather crude approximations to the real processes that affect the evolution of a dissolving moving group, but they can be seen as limiting cases to the real processes at work. Our main use of the simulations is also to test the assumption that the sample of stars identified as probable members due to metallicity and kinematics have a common origin in time and space. The two examples in Fig. 6 shows that the structure in the $`UV`$-plane outlined by the metal-rich stars in Fig. 4 with a width in U-velocity of 60 km s<sup>-1</sup> and in $`V`$ of 20 km s<sup>-1</sup> and not confined to a single V-velocity is a natural result of the dynamically simulated dispersion processes. The classical configuration, used by Eggen, with all stars belonging to the group lined up along a single V-velocity only appears when the Sun is located very close to the centre of the group, which is not likely. Otherwise, either if the Sun is located to the outside or the inside of the tube-orbit defined by the moving group, we get the tilted structure in the $`UV`$-plane shown in Fig. 6. In the simulation in Fig. 6a with an original velocity dispersion of 6 km s<sup>-1</sup>, the center of the group had a velocity of $`U=15`$ km s<sup>-1</sup>, $`V=58`$ km s<sup>-1</sup> and $`W=7`$ km s<sup>-1</sup> when it passed the present Solar position 40 Myrs ago. In the simulation with Wielen diffusion, we used a time and velocity independent diffusion coefficient $`D=2.010^7`$ km s<sup>-1</sup> yr<sup>-1</sup> and the center of the group had a velocity of $`U=0`$ km s<sup>-1</sup>, $`V=58`$ km s<sup>-1</sup> and $`W=7`$ km s<sup>-1</sup> when it passed the present Solar position only 5 Myrs ago. These two examples have their peak density in the $`UV`$-plane at approximately the same position as the observed distribution found in Sect. 2.
## 4 Finding the HR 1614 moving group – Conclusions
We close the first part of our paper with a few remarks on the reality of old moving groups. Before the Hipparcos mission the numbers of reliable parallaxes were too small to address the reality of most proposed old moving groups successfully. Recent studies (e.g. Barrado y Navascués 1998, and Skuljan et al. 1997) have shown that well known moving groups such as the Pleiades and the Hyades but also several other young moving groups, e.g. Castor ($`200\pm 100`$ Myr, Barrado y Navascués 1998), are well identified as physical entities using the new data from the Hipparcos mission.
Are the proposed old moving groups a reality? We conclude that at least one old moving group exists and that it’s possible to find other ones using velocity information in combination with metallicities based on Strömgren photometry. The study by Dehnen (1998) further supports our findings. He recovered many maxima, using a maximum likelihood solution, in the velocity distribution of nearby stars using the Hipparcos catalogue. Several of these are identifiable with known moving groups. In particular he found several maxima that exclusively contained red stars, indicating an old age. He identifies one of these maxima with the HR 1614 moving group.
However, HR 1614 might be a rather special case. It is particularly metal-rich compared to the majority of stars in the part of the $`UV`$-plane it resides in. In other parts of the $`UV`$-plane (e.g. close to the local standard of rest) the group would have been completely obscured by other metal-rich stars.
In summary, we conclude that at least one old moving group exists and it’s possible to find others using our simple method if they stand out in terms of metallicity and/or age from the ambient background of stars in their space of the $`UV`$-plane.
## 5 Review of previous work on the HR 1614 moving group
The HR 1614 moving group stands out among stellar moving groups in terms of age and metallicity. The age and metallicity have been estimated to be roughly similar to that of the old open cluster NGC 6791, Hufnagel & Smith (1994) and Eggen (1998a). However, this is challenged by the recent determination by Chaboyer at al. (1999) who found NGC 6791 to have a metallicity of $`+0.4`$ dex and an age of 8 Gyr.
It was Eggen (1978) who, following leads in earlier studies by Eggen (1971) and Hearnshaw (1974), first identified the presence of a moving group associated with the K dwarf star HR 1614 (HD 32147) by studying a sample of stars which were selected as being within $`\pm 10`$ km s<sup>-1</sup> of the V-velocity of the star HR 1614 (then estimated to $`58`$ km s<sup>-1</sup>). He found that in a sample of 44 stars with $`(U,V)`$<sup>1</sup><sup>1</sup>1These are the $`U`$ and $`V`$ velocities in a right handed system. Eggen use a left handed system, however, we conform with the now common practice and use a right handed system where U points in the direction towards the galactic centre and $`V`$ in the direction of galactic rotation. $`(0,60)`$ km s<sup>-1</sup> 60 % were over abundant with respect to the Sun (based on a few spectroscopic studies, mainly Oinas (1974) and the $`m_1`$ index). The stars appear as metal-rich as the Hyades in the $`M_V,RI`$ diagram. They also showed strong blanketing effects in $`by`$ (an excess of $`0.03`$ compared to the Hyades group, which was utilized as a selection criterion) and made up a colour magnitude diagram resembling that of an old stellar cluster (Eggen 1978 Fig. 1a).
Smith (1983) subsequently obtained DDO photometry of 19 suggested member stars and found many of them to have enhanced cyanogen bands similar to those found in so called super-metal-rich stars (SMR, see Taylor 1996 for a discussion of SMR stars). The group also showed anomalously many CN-rich stars when compared with a random sample of stars. Derived \[Fe/H\] generally confirmed the high metallicity of the member stars, however, two giant stars were found to be significantly lower in metallicity casting some doubt on Eggen’s (1978) $`by`$ criterion for membership. A stricter application of membership criteria, i.e. that the stars should show the same behaviour in both UBV, DDO and $`by`$, showed that many of Eggen’s original candidates did not belong to the group. Smith (1983) further suggests that the abundances and kinematics of the HR 1614 moving group stars are consistent with the abundance gradient observed in old open clusters (if it formed at perigalacticon) and thus the lack of observations of similar clusters would be due to that either they no longer exist or that observations up to 1983 had not sampled the inner regions of the Galaxy well enough. Utilizing the CN-enhancement Eggen (1992) found a total of 39 main-sequence members of the HR 1614 moving group by isolating them through their V-velocity. All these dwarfs are within 40 pc of the Sun. A further 19 red giants within 200 pc of the Sun were also identified.
No high resolution spectroscopic study has targeted the probable member stars of the HR 1614 moving group.
## 6 A new selection of probable HR 1614 members and its age
We now proceed with the second part of this paper and provide a new sample of stars with high probability of being member stars. From this new sample we are able to derive an age as well as a metallicity for the moving group.
### 6.1 Defining a new selection criterion
Is it possible to define a selection criterion in the $`UV`$-plane for the moving group HR 1614 that is stricter than the boundaries of box 5? From Fig. 4 we find that metal-rich stars in box 5 predominantly fall along a diagonal band going from the lower left to the upper right-hand corner. In order to find out if this tilted structure is significant we select only stars with small relative errors in the parallaxes, $`\sigma _\pi /\pi <0.05`$, and with \[Me/H\]$`>0.1`$ dex. The latter selection is based on the assumption that, if existing (as we assume in this section), the HR 1614 group should have a metallicity around $`+0.2`$ dex. For example the star HR 1614 itself has \[Fe/H\] $`=+0.28`$ dex from high resolution spectroscopy (Feltzing & Gustafsson 1998). Fig. 7a show the $`UV`$-plane for these stars with error-bars on the $`U`$ and $`V`$ velocities arising from errors in the parallaxes and radial velocities. Apart from one star (HIP 87116, which we consider as belonging to the general background and is excluded from the following discussion) the stars tend to fall along a diagonal line. The star at the top ($`U24`$ km s<sup>-1</sup>) could be contamination from the U-anomaly (see text above), but we leave it in the sample. A simple least square fit to these data give $`\mathrm{\Delta }V/\mathrm{\Delta }U=+0.19`$, Fig. 7 a. We also show lines representing $`\pm `$ one rms. In Fig. 7 b we show all stars in box 5 with \[Fe/H\] $`>0.1`$ dex and $`\sigma _\pi /\pi <0.05`$ as well as all stars with $`0.0<`$ \[Fe/H\] $`<0.1`$ dex. As can be seen all of the stars with \[Fe/H\] $`>0.1`$ dex could be regarded as falling inside this tilted band if the errors on the velocities are taken into account. This finding combined with our simulations suggests that we have identified a metal-rich co-moving stellar sample with a common origin in space and time. The simulations indicate that the structures may not only be tilted but also curved in the $`UV`$-plane, however, our data set is too small to address that question and we have to be satisfied with a “straight line”.
This selection procedure now helps us to define a new sample of probable HR 1614 moving group member stars from our Strömgren catalogue, Table 1. Our new sample is limited by the fact that the calibration from Schuster & Nissen (1989) used to calculate \[Me/H\] is only valid for dwarf stars with $`0.32BV1.0`$. Because of these limitations K dwarf stars are excluded and for example HR 1614 itself is not included although it perfectly obeys our selection criteria in all other respects. However, see discussions in Sect. 7.2 and 8.
We thus propose that future searches for member stars of the HR 1614 moving groups should be directed to the area in the $`UV`$-plane defined by our tilted band.
### 6.2 A new estimate of the age of the moving group HR 1614
We are now in a position to determine a new age estimate of the HR 1614 moving group from the HR-diagram of stars with Strömgren metallicities and based on the selection described in the previous section. Thus from Fig. 8 we derive an age estimate of 2 Gyr using the Bertelli et al. (1994) isochrones.
The iron content \[Fe/H\] of the Z = 0.05 isochrone from Bertelli et al. (1994) which is found to be \[Fe/H\] = +0.40 dex from their relation (11), is higher then the observed one for our sample of HR 1614 member stars which is $`+0.19\pm 0.06`$. This discrepancy could be due to an error in the colour transformation from the $`\mathrm{log}T_\mathrm{e},M_{\mathrm{bol}}`$ to the $`BV,M_\mathrm{V}`$ plane where a metal-rich model atmosphere is used compared to the solar metallicity case which is much better constrained. Other possible error sources are the assumed $`\mathrm{\Delta }Y/\mathrm{\Delta }Z`$, where a lower value would give a lower metallicity for the isochrone and the assumed scaled solar abundances. If other elements than iron is over-abundant this would lower the iron content for the same assumed metallicity.
In an attempt to estimate an error on this age estimate we have compared the the Bertelli et al. (1994) isochrones with those from Chaboyer et al. (1999). Their metal-rich isochrones present excellent fits to observations from solar metallicity (M67) to the very high metallicity of \[Fe/H\] $`=+0.4`$ for NGC 6791, both of which are older than the HR 1614 moving group. For solar metallicity the Bertelli et al. (1994) and Chaboyer et al. (1999) isochrones are very similar up to the sub-giant region. The two sets of isochrones agree well for ages applicable to NGC 6791. Chaboyer at al. (1999) find this open cluster to be 8 Gyr. Furthermore, and most importantly here, they also investigate the effect of changing $`\mathrm{\Delta }Y/\mathrm{\Delta }Z`$. Using values between 1 and 3 they find the impact on the derived age to be small. Bertelli et al. (1994) use $`\mathrm{\Delta }Y/\mathrm{\Delta }Z=2.5`$. It is quite likely that the HR 1614 moving group has a $`\mathrm{\Delta }Y/\mathrm{\Delta }Z2.5`$. From our comparison of the isochrones we conclude that errors on the derived age, from the possibility that our stars have $`\mathrm{\Delta }Y/\mathrm{\Delta }Z2.5`$, is less than 1 Gyr.
The conclusion by Hufnagel & Smith (1994) that the age of the HR 1614 group is $``$3 Gyr based on the MgII chromospheric emission index is consistent with our isochrone determined age of 2 Gyr in the light of the work by Rocha-Pinto & Maciel (1998). They show that the strength of the chromospheric emission has a strong dependence of the metallicity of the stars, with metal-rich stars having higher chromospheric ages than isochrone ages. Using their relation between the error in the chromospheric age and the iron abundance, for a assumed \[Fe/H\] of +0.20 for HR 1614, a chromospheric age of 3 Gyr corresponds to a isochrone age of 1 Gyr.
## 7 The metallicity of the HR 1614 moving group
We now proceed to determine the metallicity of the HR 1614 moving group from \[Fe/H\] measurements found in the literature.
### 7.1 The spectroscopic catalogue
The basis for our spectroscopic catalogue is the kinematic catalogue described in Sect. 2 supplemented by iron abundances as well as other elemental abundances from the catalogues by Fuhrmann (1998), Favata et al. (1997), Edvardsson et al. (1993), and Feltzing & Gustafsson (1998).
These studies were chosen because they all have large samples of stars, high-resolution has been used, they all have small internal errors, and a comparison of the several stars in common between the studies reveals a high internal consistency between the studies, $`\sigma `$(\[Fe/H\])$`=\pm 0.06`$ dex. All kinematical data have been recalculated using the new parallaxes and proper motions from Hipparcos.
The Edvardsson et al. (1993) sample was selected in order to study F and early G main sequence stars, which were evenly distributed in metallicity between $`1.0`$ dex and $`+0.3`$ dex, and slightly evolved off the ZAMS so that ages could be estimated. The selection of stars was done from Olsen (1988) which includes almost all F and early G stars brighter than $`V8.3`$. The stars in Feltzing & Gustafsson (1998) were selected from the same photometric catalogue with updates in order to study the metal-rich population of disk dwarf stars on solar orbits but also of stars with fairly eccentric orbits and a subsample was selected to have $`V_{\mathrm{LSR}}<50`$ km s<sup>-1</sup> and/or total space velocity $`>60`$ km s<sup>-1</sup>. We may therefore expect to find most probable HR 1614 moving group members in this study. Note that in the selection of stars for these two studies no attempt was made to cover the probable member stars of the HR 1614 moving group.
Favata et al. (1997) provide a volume limited sample of G and K dwarf stars drawn from the Gliese catalogue. Fuhrmann (1998) provide iron abundances for 50 nearby F and G type stars on the main-sequence, turn-off, and sub-giant branch.
Comparing \[Fe/H\] from Edvardsson et al. (1993) and Feltzing & Gustafsson (1998) with \[Me/H\] derived using the calibration by Schuster & Nissen (1989) we find that $`[\mathrm{Fe}/\mathrm{H}][\mathrm{Me}/\mathrm{H}]=+0.015\pm 0.076`$.
### 7.2 The metallicity
We now apply the $`UV`$-selection criterion to our spectroscopic catalogue. The resulting HR diagram and also the $`UV`$-plot for the full catalogue are shown in Fig. 9.
As noted in Sect. 6.1 our final sample of probable member stars presented in Table 1 only contained stars which fell inside the calibration by Schuster & Nissen (1989), however, for example HR 1614 itself would fall outside this calibration but does indeed fulfill the selection criteria. We thus include two K dwarf stars in our final sample of probable member stars on the basis of their spectroscopically measured metallicities.
For the 8 stars with \[Fe/H\]$`>0.05`$ dex $`<`$\[Fe/H\]$`>`$ $`=0.19\pm 0.06`$ using the \[Fe/H\] for HR 1614 from Feltzing & Gustafsson (1998). For the 11 rejected stars (i.e. both stars inside the strip with \[Fe/H\] $`<0.05`$ dex and stars outside the strip) we find $`<`$\[Fe/H\]$`>`$ $`=0.19\pm 0.38`$. The three that fall just outside our selection criterion in the $`UV`$-plane (HIP 22336, HIP 26437, HIP 70470) have \[Fe/H\] = 0.14, 0.15, 0.06 dex. Including them into the sample we arrive at $`0.17\pm 0.06`$ dex.
Obviously the spectroscopic catalogue is subject to selection biases. However, none of the investigations included in the catalogue have been biased against metal-rich stars in this part of the $`UV`$-plane. In fact Feltzing & Gustafsson (1998) did actively include metal-rich stars with large negative $`V_{\mathrm{LSR}}`$. Note though that no attempt was made in their investigation to preferentially include HR 1614 members. We may thus expect to sample the metal-rich stars well, while the metal-poor stars are less well sampled.
In conclusion we find that the HR 1614 moving group has a metallicity, determined from high resolution spectroscopy, of $`0.170.19\pm 0.06`$ dex depending on the exact selection in the $`UV`$-plane.
## 8 The Eggen sample revisited
Having defined a new selection criterion for membership in the $`UV`$-plane we now revisit the Eggen (1992, 1998b) sample of HR 1614 moving group stars. Fig. 10 show the sample of stars that Eggen (1998b) indicate as probable members of the HR 1614 moving group. In Fig. 10a we use his cluster parallaxes, in Fig. 10b the Hipparcos parallaxes. The use of the cluster parallaxes forces the stars onto more or less one isochrone. Eggen (1998b) does not discuss the apparent metallicity spread in Fig. 10a. We note especially that the turn-off region, crucial for age determination, shows a larger scatter in Fig. 10b as compared to Fig. 10a.
When we apply our $`UV`$-selection criteria to Eggen’s sample (with re-calculated velocities and using the Hipparcos parallaxes) a total of 21 stars with $`M_V>2`$ falls inside the band in Fig. 7. 8 of these are in common with our selection from the Strömgren catalogue. These stars are marked in the first part of Table 1. Of the remaining 13 stars one is HR 1614 and is already included in our sample on the basis of its \[Fe/H\] measured from high resolution spectroscopy. The two stars HIP 11575 and HIP 35872 appears in our Strömgren catalogue and have \[Me/H\]$`=0.12`$ and $`0.09`$, respectively, disqualifying them as probable member stars of the HR 1614 moving group. Four more stars were excluded due to binarity, HIP 7143 is a spectroscopic binary, HIP 94570 and HIP 96037 were resolved as binary systems by Hipparcos and HIP 112426 have an accelerated solution, probably due to orbital motion (ESA 1997).
This leaves, after having considered binarity and metallicity when information is available, an additional 6 stars from Eggen’s (1992, 1998b) sample to be included in our final sample of probable member stars of the HR 1614 moving group. The stars are detailed in the final part of Table 1.
The HR-diagram of the 15 stars from Eggen (1992, 1998b) thus full-filling our criteria are plotted in Fig. 10c.
## 9 Discussion
As already touched upon, the HR-diagram of the HR 1614 group is “contaminated” by stars from the general field and possibly also by stars from other moving groups and open clusters.
The general background of stars in the $`UV`$-plane using the Hipparcos results have been studied by Dehnen (1999) and Raboud et al. (1998). Dehnen (1999) shows that there exists an over-density of stars in the $`UV`$-plane centered on $`U=20`$ and $`V=45`$ km s<sup>-1</sup>. He associates this stellar over-density with stars thrown out from the inner disk by the galactic bar. In particular stars close to the outer Lindblad resonance are susceptible to this and the phenomenon a clear indication of the non-axisymmetry of the galactic potential. Thus we should expect to see a generally more metal-rich stellar population present in the upper left half of our boxes in Fig. 2. We do indeed do so, however, we also note that the contamination is mainly in box 1 and 2 and does not effect box 5. This is further born out by the $`UV`$-plots using the Strömgren sample, Sect. 2.
The moving group Wolf 630 is present in the upper right hand corner of box 5 ($`50<U<+25`$ and $`75<V<45`$ km s<sup>-1</sup>, see also Fig. 1 in Skuljan et al. 1997 for the position of Wolf 630). However, since the Wolf 630 moving group has a mean \[Fe/H\] of $`0.14`$ dex (Boesgaard & Friel 1990) almost none of the stars would remain in our Strömgren sample. Schuster & Nissen (1989) estimated the scatter of their calibration to 0.16 dex. When we compare the estimated metallicities with spectroscopic abundances the differences are much smaller, around 0.06 dex. Thus we may expect these stars to be removed also in the Strömgren sample. In the spectroscopic sample all would be removed. We conclude that the contamination of our sample by stars from the Wolf 630 moving group is negligible.
Eggen (1998a) suggested that the open cluster NGC 6791 might be part of the HR 1614 moving group. This is ruled out by two recent studies by Tripico et al. (1995) and Chaboyer et al. (1999) which both agree that the open cluster has an age around $`810`$ Gyr. This means that even if we allow for the uncertainties in our age determination due to uncertainties in the modeling of stellar evolutionary tracks the age of the HR 1614 moving group is $`>`$ 5 Gyr younger than the old open cluster NGC 6791. Also, the metallicity of NGC 6791 is almost 0.2 dex higher than that of the HR 1614 moving group. Combining these facts the association of the HR 1614 moving group with NGC 6791 must be incorrect.
Our sample provides a sample for future studies of the abundance profiles in old moving groups to address the question whether or not todays field stars originate in stellar clusters subsequently dissolved. It is at the higher metallicities, i.e. $`[\mathrm{Fe}/\mathrm{H}]0.1`$ dex, that differences in star formation histories will manifest themselves in the abundance ratios (e.g. Pagel 1997, fig 8.6, and Matteucci 1991). If metal-rich stars in the field originate in clusters then abundance ratios for stars in a metal-rich moving group should be identical to those of the field stars in the solar neighbourhood. On the other hand if the moving groups are not the source of the field stars then it is most likely that their star formation rate was different and thus the resulting elemental abundances will be different. If the star formation rate was more rapid in the cluster than what is typical for the places where the metal-rich stars, now in the solar neighborhood, were born then the \[$`\alpha `$/Fe\] will be larger at say \[Fe/H\]$`0.2`$ dex for the group stars than for the field stars. If on the other hand the star formation rate was lower \[$`\alpha `$/Fe\] will be smaller, Pagel (1997 Fig. 8.6).
## 10 Summary
Utilizing the new possibilities given by the Hipparcos mission we perform, for the first time, an unbiased search in the $`UV`$-plane of all stars with measured radial velocities and find an over-density of stars more metal-rich than the Sun close to the $`U`$ and $`V`$ values associated with the moving group HR 1614 (Eggen 1998b). Supported by dynamical simulations we find that old moving groups in fact does exist, at least the HR 1614 moving group. The selection criterion for the HR 1614 moving group in the $`UV`$-space is further refined using metallicities derived from Strömgren photometry. This new criterion is applied to the large catalogue of all Hipparcos stars with measured radial velocities and Strömgren photometry resulting in a new determination of the probable age of the moving group.
We derive an age of 2 Gyr using the Bertelli et al. (1994) isochrones and a iron abundance of $`+0.19\pm 0.06`$ dex using available data, from high resolution spectroscopy, in the literature.
We also provide a comparison of \[Me/H\] derived from Strömgren photometry of the Hauk & Mermilliond catalogue (1998) using the Schuster & Nissen (1989) calibration. The comparison shows that the Strömgren \[Me/H\] and \[Fe/H\] for F and G type dwarf and sub-giant stars are in extremely good agreement with $`[\mathrm{Fe}/\mathrm{H}][\mathrm{Me}/\mathrm{H}]=+0.015\pm 0.076`$.
The sample of bright main-sequence and turn-off stars that are probable members of the moving group HR 1614 is presented in Table 1. They make up a sample for further investigations into the abundance profile of this old, metal-rich moving group.
###### Acknowledgements.
This work has made use of the SIMBAD on-line database search facility maintained by the CDS, Strasbourg. JH acknowledges the financial support of the Swedish National Space Board. The Royal Physiographic Society in Lund is thanked for providing funding for computer facilities. The anonymous referee is thanked for providing contructive remarks that lead to a revision and considerable improvement of the paper. |
warning/0003/astro-ph0003171.html | ar5iv | text | # Ekman layers and the damping of inertial r-modes in a spherical shell: application to neutron stars
## 1 Introduction
Recently much work has been devoted to the study of the rotational instability of neutron stars resulting from a coupling between gravitational radiation and the so-called “r-modes” of a rotating star Andersson (1998); Friedman and Morsink (1998); Lindblom et al. (1998); Kokkotas and Stergioulas (1999). Such an instability may indeed play a key role in the distribution of rotation periods of neutron stars as well as it may be an important source of gravitational radiation.
In this paper, we shall first clarify a point of history concerning “r-modes” which are in fact a special class of inertial modes; we shall then review their singular properties which have been clarified only very recently in Rieutord and Valdettaro (1997), Rieutord et al. (2000a) and Rieutord et al. (2000b). The last section will present the analytical derivation of the damping rate of inertial r-modes in a neutron star with a crust and/or a core through the boundary layer analysis within the framework of newtonian theory. We conclude on the stability of crusted neutron stars when modeled by an incompressible viscous fluid in a rotating sphere.
## 2 A short point of history
The very first work on rotating fluids oscillations which are presently known as inertial modes dates back to Thomson<sup>1</sup><sup>1</sup>1later Lord Kelvin (1880) who analysed the case of a fluid contained in a cylinder. However, another impetus to the study of these oscillations was given soon after by Poincaré’s (1885) work on the stability of rotating self-gravitating masses, a work applied to MacLaurin spheroids by Bryan (1889)<sup>2</sup><sup>2</sup>2 but see the recent rederivation by Lindblom and Ipser (1999). and later continued by Cartan (1922) who christened the equation of inertial modes as “Poincaré equation”. In these studies, however, the effect of rotation is combined to the one of gravity through (for an incompressible fluid) surface gravity waves. In fact, except for the work of Thomson, investigations on the oscillations specific of rotating fluids seem to have started with the work of Bjerknes et al. (1933) where they are called “elastoid-inertial oscillations” since conservation of angular momentum makes axis-centered rings of fluid behave elastically; but see Fultz (1959) or Aldridge (1967) for an account on this part of history. In the sixties, much work has been devoted to these oscillations, mainly by Greenspan who introduced the terminology of “inertial oscillations”. The presently used denomination “inertial modes” has been “officially” given by Greenspan’s book (Greenspan, 1969).
However, inertial modes are somewhat too general for applications in some specific domains like atmospheric sciences. In this field indeed motions are essentially two dimensional and inertial modes may be simplified into the well-known Rossby (or planetary) waves.
The introduction of r-modes by Papaloizou and Pringle (1978) was quite unfortunate since they associated eigenmodes of rotating fluids with a very special class of inertial modes, namely purely toroidal inertial modes. This lead following authors to introduce weird names such as “hybrid modes” or “generalized r-modes” (Lockitch and Friedman, 1999) for describing the general class of inertial modes. We therefore encourage authors to use, as fluid dynamicists, inertial modes unless they discuss the very specific r-modes.
## 3 The present theory of inertial modes
Inertial modes are a class of modes of oscillation of rotating fluids which owe their existence to the Coriolis force. This force of inertia has indeed a restoring action on perturbations of rotating fluids since it insures the global conservation of angular momentum. These modes have many properties similar to those of gravity modes of stably stratified fluids (Rieutord and Noui, 1999).
The dynamics of inertial modes may be appreciated when all other effects are suppressed: no compressibility, no magnetic fields, no gravity, etc… only an incompressible inviscid rotating (like a solid body) fluid. In this case, perturbations of velocity $`\delta 𝐯`$ and pressure $`\delta P`$ obey
$$\frac{\delta 𝐯}{t}+2𝛀\times \delta 𝐯=\mathrm{NewA}\delta P,\mathrm{NewA}\delta 𝐯=0$$
(1)
where $`𝛀`$ is the angular velocity of the fluid. Concentrating on time-periodic oscillations and choosing $`(2\mathrm{\Omega })^1`$ as the time scale, (1) can be written
$$i\omega 𝐮+𝐞_z\times 𝐮=\mathrm{NewA}p,\mathrm{NewA}𝐮=0$$
(2)
with non-dimensional variables; $`\omega `$ is the non-dimensional (real) frequency. When the velocity $`𝐮`$ is eliminated in favor of the pressure perturbation $`p`$, one is left with
$$\mathrm{\Delta }p\frac{1}{\omega ^2}\frac{^2p}{z^2}=0$$
(3)
which is known as Poincaré equation since Cartan (1922). This equation is remarkable in the fact that it is hyperbolic spatially since $`|\omega |1`$ (Greenspan, 1969). As the solution of (3) must meet boundary conditions, namely $`𝐮𝐧=0`$, we see that inertial modes are solutions of an ill-posed boundary value problem<sup>3</sup><sup>3</sup>3Such boundary conditions eliminate any distorsion of the surface due to the fluid motion; for a free surface, these distorsion are surface gravity waves (see Rieutord 1997) but their inclusion (in order to be more realistic) would not modify the ill-posed nature of the problem.. This property means that, in general, inertial modes are singular; in other words they cannot exist physically if the fluid is strictly inviscid. These properties are detailed in Rieutord and Valdettaro (1997) and Rieutord et al. (2000a); to make a long story short, one may summarize the situation as follows.
Let us first recall that in hyperbolic systems, energy propagates along the characteristics of the equation. For the Poincaré problem, these are straight lines in a meridional plane. A way to approach the solutions of this difficult problem is to examine the propagation of characteristics as they reflect on the boundaries. They define trajectories which depend strongly on the container. Let us therefore concentrate on the case of a spherical shell as a container; this configuration is relevant for neutron stars with a central core due to some phase transition of the nuclear matter (see Haensel, 1996). In this case, it may be shown that characteristic trajectories generically converge towards attractors which are periodic orbits. It may be shown (Rieutord et al., 2000a) that in this case, the associated solutions are singular, namely the velocity field is not square-integrable. However, inertial r-modes are still solutions of the problem since they meet the boundary conditions ($`u_r=0`$); in fact they are the only regular (square-integrable) solutions of the Poincaré problem in a spherical shell. In a more mathematical way, we may say that the spectrum of eigenvalues of the Poincaré problem in a spherical shell is empty except for the inertial r-modes. In this sense, these modes are quite exceptional. This situation occurs because there exists no system of coordinate in which the dependent variables of the Poincaré equation can be separated. This is a consequence of the conflict between the symmetry of the Coriolis force (cylindrical) and the geometry of the boundaries. Thus, when this constraint is relaxed, like in the case of a cylindrical container, regular solutions exist and a dense spectrum of eigenvalues appears in the allowed frequency range, namely $`[0,2\mathrm{\Omega }]`$. In the case the container is a full sphere, attractors also disappear and eigenmodes exist; they are also related to a dense spectrum of eigenfrequencies. In this case, Poincaré equation is exactly solvable (Greenspan, 1969).
However, real fluids have viscosity ($`\nu `$) and equation (2) should be transformed into
$$\lambda 𝐮+𝐞_z\times 𝐮=\mathrm{NewA}p+E\mathrm{\Delta }𝐮,\mathrm{NewA}𝐮=0$$
(4)
where $`\lambda `$ is the complex eigenvalue and $`E=\nu /2\mathrm{\Omega }R^2`$ is the Ekman number ($`R`$ is the outer radius of the shell).
Using no-slip ($`𝐮=\mathrm{𝟎}`$) or stress-free boundary conditions, (4) yields a well-posed problem. Yet, the singularities of the associated inviscid solutions show up through the existence of shear layers. As shown by fig. 1, the shape of inertial modes is deeply influenced by the underlying singularity of the inviscid solution. We have shown (Rieutord et al., 2000a) that these shear layers are in fact nested layers with different scales for their inner part scales as $`E^{1/3}`$ and their outer part seems to scale with $`E^{1/4}`$. Because of these internal shear layers, these modes are strongly damped.
We therefore see that according to whether a neutron star has a central core or not, the damping of inertial modes will be extremely different. If there is a central core, the only regular modes are the inertial r-modes which will be by far the least damped; if there isn’t any core then a dense spectrum exists (Greenspan, 1969; Lockitch and Friedman, 1999) but inertial r-modes remain the most unstable because of their simple structure.
## 4 Damping of toroidal inertial modes in a sphere or a shell
We shall now give the expression of the viscous damping of inertial r-modes when one of the boundary is solid therefore when the dissipation is due to Ekman boundary layers; we shall thus complete the works of Bildsten and Ushomirsky (2000) and Andersson et al. (2000) by giving the rigorous estimate of the damping rate; the method which we use here has been outlined in Greenspan (1969).
The damping rate is given by:
$$\gamma =\mathrm{}e(\lambda )=E\frac{(\mathrm{NewA}:𝐮)^2dV}{𝐮^2𝑑V}$$
(5)
where $`(\mathrm{NewA}:𝐮)^2`$ stands for the squared rate-of-strain tensor (see below). The velocity field of r-modes is
$$\overline{u}_\theta =Ar^m(\mathrm{sin}\theta )^{m1}\mathrm{sin}(m\varphi +\omega _mt)$$
$$\overline{u}_\phi =Ar^m(\mathrm{sin}\theta )^{m1}\mathrm{cos}\theta \mathrm{cos}(m\varphi +\omega _mt)$$
The kinetic energy integral may be evaluated explicitly
$$𝐮^2𝑑V=\pi A^2\left(1\eta ^{2m+3}\right)\frac{2^{m+1}(m+1)!}{m(2m+3)!!}$$
(6)
where $`\eta `$ is the ratio of the radius of the inner boundary to the one of the outer boundary.
The dissipation integral need more work if one of the boundary is no-slip. In this case, dissipation is essentially coming from the Ekman layers and thus we need to derive the flow in these layers. The method has been given by Greenspan from whom we know that the boundary layer correction $`\stackrel{~}{u}`$, is related to the interior solution $`\overline{u}`$ by
$$\stackrel{~}{u}_\theta +i\stackrel{~}{u}_\phi =(\overline{u}_\theta +i\overline{u}_\phi )_{r=r_b}e^{\zeta \sqrt{i\mathrm{cos}\theta \pm i\omega }}$$
(7)
where $`\zeta `$ is the radial scaled variable $`(rr_b)/\sqrt{E}`$ with $`r_b`$ as the radius of the boundary (1 or $`\eta `$). The complete solution is then $`𝐮=\stackrel{~}{u}+\overline{u}`$; setting $`\beta =\omega t+m\phi `$ we have
$`\overline{u}_\theta +i\overline{u}_\phi ={\displaystyle \frac{Ar^m(\mathrm{sin}\theta )^{m1}}{2i}}`$
$`\times \left\{(1\mathrm{cos}\theta )e^{i\beta }(1+\mathrm{cos}\theta )e^{i\beta }\right\}`$
from which it follows that
$`\stackrel{~}{u}_\theta +i\stackrel{~}{u}_\phi ={\displaystyle \frac{Ar_b^m(\mathrm{sin}\theta )^{m1}}{2i}}\{(1+\mathrm{cos}\theta )e^{i\beta \zeta \sqrt{i\mathrm{cos}\theta i\omega }}`$
$`(1\mathrm{cos}\theta )e^{i\beta \zeta \sqrt{i\mathrm{cos}\theta +i\omega }}\}`$
Now we need the expression of the square of the rate-of-strain tensor $`s_{ij}=_iv_j+_jv_i`$ in spherical coordinates, viz
$$(\mathrm{NewA}:𝐮)^2=s_{rr}^2+s_{\theta \theta }^2+s_{\varphi \varphi }^2+2(s_{r\theta }^2+s_{r\varphi }^2+s_{\theta \varphi }^2)$$
Since the radial derivatives dominate, this expression reduces to the contribution of the tangential stresses. Using the scaled coordinate, $`\zeta =|rr_b|/\sqrt{E}`$, we have
$$(\mathrm{NewA}:𝐮)^2=\frac{2}{E}\{\left(\frac{u_\theta }{\zeta }\right)^2+\left(\frac{u_\phi }{\zeta }\right)^2\}_{r=r_b}$$
We now set $`p=\mathrm{cos}\theta \omega `$ and $`q=\mathrm{cos}\theta +\omega `$. We thus get
$`(\mathrm{NewA}:𝐮)^2={\displaystyle \frac{A^2r_b^{2m}(\mathrm{sin}\theta )^{2m2}}{2E}}\{(1+\mathrm{cos}\theta )^2|p|e^{\zeta \sqrt{2|p|}}`$
$`+(1\mathrm{cos}\theta )^2|q|e^{\zeta \sqrt{2|q|}}`$
$`2\mathrm{sin}^2\theta \sqrt{pq}\mathrm{}e\left(e^{2i\beta \zeta (\sqrt{iq}+\sqrt{ip})}\right)\}`$
Integrating over the $`\phi `$-variable yields
$`{\displaystyle _0^{2\pi }}(\mathrm{NewA}:𝐮)^2d\phi ={\displaystyle \frac{\pi A^2r_b^{2m}(\mathrm{sin}\theta )^{2m2}}{E}}`$
$`\times \left\{(1+\mathrm{cos}\theta )^2\right|p|e^{\zeta \sqrt{2|p|}}+(1\mathrm{cos}\theta )^2|q|e^{\zeta \sqrt{2|q|}}\}`$ (8)
We now integrate over the radial variable:
$`{\displaystyle _\eta ^1}{\displaystyle _0^{2\pi }}(\mathrm{NewA}:𝐮)^2d\phi r^2dr=`$
$`{\displaystyle \frac{\pi A^2(\mathrm{sin}\theta )^{2m2}}{E}}{\displaystyle _\eta ^1}r^{2m+2}f(\zeta )𝑑r`$
with $`f(\zeta )=(1+\mathrm{cos}\theta )^2|p|e^{\zeta \sqrt{2|p|}}+(1\mathrm{cos}\theta )^2|q|e^{\zeta \sqrt{2|q|}}`$; since $`r=\eta +\sqrt{E}\zeta `$ or $`r=1\sqrt{E}\zeta `$ according to which side of the integral is chosen, it turns out that
$`{\displaystyle _\eta ^1}{\displaystyle _0^{2\pi }}(\mathrm{NewA}:𝐮)^2d\phi r^2dr`$
$`=`$ $`K{\displaystyle _0^{\mathrm{}}}[(1+\mathrm{cos}\theta )^2|p|e^{\zeta \sqrt{2|p|}}`$
$`+(1\mathrm{cos}\theta )^2|q|e^{\zeta \sqrt{2|q|}}]d\zeta `$
$`=`$ $`{\displaystyle \frac{K}{\sqrt{2}}}[(1+\mathrm{cos}\theta )^2\sqrt{|\mathrm{cos}\theta \omega |}`$
$`+(1\mathrm{cos}\theta )^2\sqrt{|\mathrm{cos}\theta +\omega |}]`$
with $`K=\pi A^2(\mathrm{sin}\theta )^{2m2}P(\eta )/\sqrt{E}`$, where $`P(\eta )`$ is a function depending on the boundary conditions (see table 1). Finally integrating over $`\theta `$, we find
$$(\mathrm{NewA}:𝐮)^2dV=\frac{2\pi A^2P(\eta )}{\sqrt{2E}}_m$$
(9)
with
$$_m=_0^\pi (1+\mathrm{cos}\theta )^2\sqrt{|\mathrm{cos}\theta \omega |}\mathrm{sin}^{2m1}\theta d\theta $$
(10)
Finally grouping (6) and (9), we find the damping rate
$$\gamma =\frac{m(2m+3)!!}{2^{m+3/2}(m+1)!}Q(\eta )_m\sqrt{E}$$
(11)
where $`Q(\eta )=P(\eta )/\left(1\eta ^{2m+3}\right)`$.
For the cases $`m=1`$ and $`m=2`$ we evaluated the expression of $`_m`$, viz
$$_1=\frac{\sqrt{2}}{35}\left(3^{5/2}+19\right)$$
$$_2=4\left(\frac{2}{3}\right)^{11/2}\frac{3401+2176\sqrt{2}}{5\times 7\times 9\times 11}$$
Other values are computed numerically and given in table 2.
The values given by (11) may be compared to other derivations, in particular that of Greenspan (1969) for $`m=1`$ who finds $`\gamma /\sqrt{E}=2.62/\sqrt{2}=1.8526`$ ! For $`m=2`$, a direct numerical calculation, similar to that of Rieutord and Valdettaro (1997), gives $`2.482\sqrt{E}`$ at $`E=10^8`$ which is in good agreement with the analytical formula.
## 5 Application to neutron stars and conclusions
Let us now apply these results to the case of rapidly rotating neutron stars. We take the viscosity from Bildsten and Ushomirsky (2000), $`\nu =1.8f/T_8^2`$ m<sup>2</sup>/s where $`f`$ is a dimensionless parameter taking into account the different transport mechanisms in the fluid (superfluid phases for instance) and $`T_8`$ is the temperature in 10<sup>8</sup>K unit. Using a radius of 12.53 km and an angular frequency of $`2\pi \times `$1kHz, we find an Ekman number $`10^{12}`$ which is indeed very small and thus boundary layer theory applies.
We may now estimate the charateristic time scale for the damping of the $`m=2`$-mode. We find
$$T_d=26.7\mathrm{s}\frac{T_8}{\sqrt{f}}\left(\frac{R}{10\mathrm{km}}\right)\left(\frac{1\mathrm{kHz}}{\nu _s}\right)^{1/2}$$
(12)
which is a somewhat smaller value than the previous estimate of Bildsten and Ushomirsky (2000) and Andersson et al. (2000) who find a characteristic time of 100 s and 200 s respectively. Our disagreement with these authors comes from their approximate evaluation of the boundary layer dissipation and from the resulting functional dependence with respect to mass and density. Let us first evaluate the damping rate according to Landau and Lifchitz (1989); it turns out that
$$2\gamma =\left(\frac{\omega E}{2}\right)^{1/2}\frac{\left(_{4\pi }𝐮^2\mathrm{sin}\theta d\theta d\phi \right)(r=1)}{_{(V)}𝐮^2𝑑V}$$
(13)
where we used our non-dimensional units. Since the radial dependence of the modes is in $`r^m`$ and $`\omega =1/(m+1)`$, we easily find that
$$\gamma =\frac{2m+3}{2\sqrt{2m+2}}\sqrt{E}$$
When this expression is applied to the $`m=2`$-mode, we find that $`\gamma =1.429\sqrt{E}`$ which is a factor $`1.74`$ weaker than the correct result.
If we use, as previous authors, a step function for describing the density difference between that of the crust and the mean density, we find that the damping rate reads
$$\gamma _{Ek}=2.4876\sqrt{E}\frac{\rho _b}{\overline{\rho }}\mathrm{\hspace{0.17em}2}\mathrm{\Omega }=0.0346\frac{\rho _b}{\overline{\rho }}\frac{\sqrt{f\mathrm{\Omega }_{}}}{T_8}\mathrm{s}^1$$
(14)
where $`\rho _b`$ is the density of the fluid just below the crust and $`\mathrm{\Omega }_{}=\mathrm{\Omega }/\sqrt{\pi G\overline{\rho }}`$.
Our calculation therefore shows that the window of instability in the $`\mathrm{\Omega },T`$ plane is smaller than previously estimated for crusted neutron stars.
Considering a 1.4 M neutron star with a radius of 12.53 km as a test case, the growth rate of the mode due to gravitational radiation is $`\gamma _{gw}=0.658\mathrm{s}^1\mathrm{\Omega }_{}^6`$ (we use the expression given in Lindblom et al., 1998); although, it is not relevant for an incompressible fluid, we take into account the damping rate due to bulk viscosity in order to ease comparison with previous work; from Lindblom et al. (1999), we find $`\gamma _{bulk}=\mathrm{2.2\hspace{0.17em}10}^{12}\mathrm{s}^1T_9^6\mathrm{\Omega }_{}^2`$. From (14), we have $`\gamma _{Ek}=\mathrm{1.53\hspace{0.17em}10}^3\mathrm{s}^1\mathrm{\Omega }_{}^{1/2}/T_9`$ where we took $`\rho _b=\mathrm{1.5\hspace{0.17em}10}^{17}`$kg/m<sup>3</sup>; solving the equation
$$\gamma _{gw}+\gamma _{Ek}+\gamma _{Bulk}=0$$
for different values of the temperature yields the curves displayed in figure 2.
As expected, we see that the window of instability narrows compared to Andersson et al. (2000): for a given temperature, the critical angular velocity raises by $``$10% typically.
Another interesting conclusion of this work is that the presence of a solid inner core does not change the damping rates very much unless its radius is close to unity. The reason for that is to be found in the shape of the inertial r-modes whose amplitudes are concentrated near the outer boundary. Therefore, the rotating instability of rapidly rotating stars is quite insensitive to the presence of a solid core and more generally to any phase transition which does not occur close to the surface.
I am very grateful to S. Bonazzola and E. Gourgoulhon for drawing my attention on these questions and for helpful discussions. I am also very grateful to Ian Jones for his note about the density profile used in models of neutron stars in previous work. Part of the calculations have been carried out on the Nec SX5 of IDRIS at Orsay and on the CalMip machine of CICT in Toulouse which are gratefully acknowledged. |
warning/0003/hep-ph0003026.html | ar5iv | text | # The use of neutrino beams from muon storage rings 1footnote 11footnote 1Invited plenary talks given at 23rd Johns Hopkins Workshop on Current Problems in Particle Theory: Neutrinos in the Next Millennium, Baltimore, MD, 10-12 June 1999, and The International Conference on Physics Potential and Development of 𝜇⁺𝜇⁻ Colliders, San Francisco, CA, 15-17 December 1999.
## I Motivation
The purpose of this talk is to overview physics goals for the short baseline experiments utilizing high energy and high intensity neutrino beams. These include standard model electroweak physics, novel tests of QCD, and rare processes sensitive to physics beyond the standard model (SM).
The standard model electroweak parameters that are conventionally measured in neutrino experiments are $`\mathrm{sin}^2\theta _W`$ and the Cabibbo-Kobayashi-Maskawa (CKM) quark mixing matrix elements $`V_{cd}`$ and $`V_{cs}`$. Given the intense high energy beam of neutrinos these measurements will certainly yield new precise values for these quantities. In addition, completely new measurements, like precision studies of $`V_{ub}`$ and $`V_{cb}`$ CKM matrix elements, will be possible. On the QCD side, neutrino-nucleon interactions are potentially the best probes of various valence parton distribution functions, both unpolarized and polarized, as well as the strong coupling constant $`\alpha _s`$. One can also study various non-perturbative parameters, such as fragmentation functions. Finally, there are interesting new physics scenarios that can be tested in neutrino interactions. These include supersymmetric extensions of the standard model with broken $``$ parity (or any models with leptoquarks), new heavy neutral leptons or gauge bosons, etc.
A proposed muon storage ring should provide a highly collimated, high-intensity $`\nu `$ beam from muons decaying in the accelerator tunnel. The neutrino spectra can be easily calculated Bigietal ; for instance, the $`\nu _\mu `$ energy spectrum is given by $`dN_{\nu _\mu }/dx6x^24x^3`$ with $`x=2E_\nu ^{}/m_\mu `$ being the normalized neutrino energy in the $`\mu `$ rest frame. $`E_\nu ^{}`$ is easily related to the neutrino energy in the lab frame, $`E_\nu =xE_\mu (1+\mathrm{cos}\theta )/2`$, where $`\theta `$ is a neutrino angle in the muon rest frame. Another advantage of this facility is that for sufficiently high muon energies the neutrinos are produced in thin pencil-like beams with an opening half-angle $`\theta _\nu m_\mu /E_\mu `$.
## II QCD and Electroweak studies
### II.1 Measurements of $`\mathrm{sin}^2\theta _W`$
One of the most important parameters of the standard model is the weak mixing angle $`\theta _W`$ which represents the angle of rotation from the “gauge” basis to the “physical” basis where the mass matrix of the gauge $`Z`$ boson and the photon is diagonal. One of the many possible definitions is the on-shell definition of $`\mathrm{sin}^2\theta _W`$
$$\mathrm{sin}^2\theta _W^{os}1\frac{M_W^2}{M_Z^2},$$
(1)
In neutrino-nucleon interactions $`\mathrm{sin}^2\theta _W`$ can be extracted using the Llewellyn Smith LS83 or Pascos-Wolfenstein relations PW73 . These methods involve measuring three total cross sections and forming three ratios,
$$R_\nu =\frac{\sigma (\nu N\nu X)}{\sigma (\nu N\mu ^{}X)},R_{\overline{\nu }}=\frac{\sigma (\overline{\nu }N\overline{\nu }X)}{\sigma (\overline{\nu }N\mu ^+X)},r=\frac{\sigma (\nu N\mu ^+X)}{\sigma (\nu N\mu ^{}X)}.$$
(2)
In the approach of Llewellyn Smith, these can be combined to obtain $`\mathrm{sin}^2\theta _W`$:
$`R_\nu `$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{sin}^2\theta _W+{\displaystyle \frac{5}{9}}(1+r)\mathrm{sin}^4\theta _W+C_\nu `$
$`R_{\overline{\nu }}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{sin}^2\theta _W+{\displaystyle \frac{5}{9}}(1+r^1)\mathrm{sin}^4\theta _W+C_{\overline{\nu }},`$ (3)
where $`C_\nu `$ and $`C_{\overline{\nu }}`$ represent known QCD and electroweak corrections. Alternatively, a Pascos-Wolfenstein construction can be used to extract $`\mathrm{sin}^2\theta _W`$:
$`R^\pm `$ $`=`$ $`{\displaystyle \frac{R_\nu \pm rR_{\overline{\nu }}}{1\pm r}},R^{}={\displaystyle \frac{1}{2}}\mathrm{sin}^2\theta _W+{\displaystyle \frac{C_\nu rC_{\overline{\nu }}}{1\pm r}}`$
$`R^+`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{sin}^2\theta _W+{\displaystyle \frac{10}{9}}\mathrm{sin}^2\theta _W+{\displaystyle \frac{C_\nu +rC_{\overline{\nu }}}{1\pm r}}.`$ (4)
This method is actually “cleaner” as the QCD and electroweak corrections partially cancel out in Eq. (II.1). These relations are now used to extract $`\mathrm{sin}^2\theta _W`$ by CCFR/NuTeV collaboration and will be used again at $`\nu `$FMSR.
In addition to the methods described above, intense neutrino beams from the muon storage ring ($`\nu `$FMSR) should allow for another measurement of $`\mathrm{sin}^2\theta _W`$, which involves neutrino-electron scattering. This method is theoretically “cleaner”, as it involves scattering of two leptons. This measurement involve investigation of four neutrino-electron elastic cross sections $`\nu _i(\overline{\nu }_i)e^{}\nu _i(\overline{\nu }_i)e^{}`$ for $`i=e,\mu `$. The involved cross section are much smaller then the corresponding DIS cross sections described above, but theoretical clearness of this process and much improved neutrino beam intensity makes this measurement a realistic possibility. Of course, future determinations of $`\mathrm{sin}^2\theta _W`$ from $`\nu `$FMSR should be comparable or better than the projected result of the SLAC E158 Moller scattering experiment, i.e. should measure $`\mathrm{sin}\theta _W`$ with relative accuracy of better then $`afew\times 10^4`$. Preliminary studies Bigietal show that it is quite realistic.
### II.2 CKM and quark densities
Extraction of the matrix elements of the Cabibbo-Kobayashi-Maskawa quark mixing matrix is one of the outstanding challenges in phenomenology of the standard model. It is most likely that Nature has chosen only three generations of quarks, so
$`V_{CKM}=S_L^uS_L^d=\left(\begin{array}{ccc}V_{ud}& V_{us}& V_{ub}\\ V_{cd}& V_{cs}& V_{cb}\\ V_{td}& V_{ts}& V_{tb}\end{array}\right).`$ (5)
Currently, the charm quark sector of CKM is best studied in $`\nu N`$ charged current interactions, where neutrino interacts with valence and sea quarks of the nucleon. For example, the CCFR collaboration has provided a direct measurement of $`|V_{cd}|=0.232_{0.019}^{+0.017}`$. Independent knowledge of the strange sea quark density should also provide an independent measurement of $`|V_{cs}|`$ as well. In the framework of a “naive” parton model,
$`{\displaystyle \frac{d^2\sigma (\nu N\mu ^+\mu ^{}X)}{d\xi dy}}`$ $`=`$ $`{\displaystyle \frac{G_F^2ME_\nu }{\pi }}\{[\xi u(\xi )+\xi d(\xi )]|V_{cd}|^2`$ (6)
$`+`$ $`2\xi s(\xi )|V_{cs}|^2\}[1{\displaystyle \frac{m_c^2}{2ME_\nu \xi }}]D(z)B_c,`$
where $`\xi =x(1+m_c^2/Q^2)`$ in a “slow rescaling” model of Georgi and Politzer GePo76 and $`D(z)`$ is a charm fragmentation function. Thus, measuring $`d^2\sigma `$ at different values of $`\xi `$ provides an independent measurement of $`|V_{cd}|`$ and $`|V_{cs}|`$, if quark densities are known well enough Conrad:1997ne . Otherwise, a multiparameter fit can be performed to determine both $`q(x)`$ and CKM matrix elements.
Even though in the above discussion a parton model was used, a problem of semiinclusive single particle production can be addressed model-independently using the formalism of perturbative QCD factorization theorems. In this framework, the essential problem of having a heavy quark in the final state is the fact that its mass brings an additional scale to the problem at hand. The presence of this scale might affect theoretical predictions by inducing large logarithms involving $`m_Q`$, which have to be resummed in order for perturbative expansion to make sense. A practical recipe for such resummation is provided by the prescription of Aivazis et. al. (ACOT) ACOT .
A future measurement utilizing high intensity neutrino beams would provide accurate determinations of various $`q(x)`$. It is interesting to note that future neutrino factories would be able to study production of the $`b`$-flavored mesons, which would allow for an accurate determination of the intrinsic charm content of the nucleon. If an $`x`$-dependent high statistics measurement of $`b`$-quark production becomes available, an independent determination of $`|V_{ub}|`$ and $`|V_{cb}|`$ CKM matrix elements will be possible as well Bigietal .
Of course, heavy quark production is not the only way of studying the nucleon structure. Quark densities are usually measured in DIS-type experiments. These measurements are naturally performed in neutrino-nucleon interactions. Here, $`\nu `$FMSR offers a tantalizing possibility to measure parity-violating polarized nucleon structure functions. These measurements were considered hopeless in $`\nu N`$ experiments due to the enormous technical difficulties in polarizing heavy targets, the only possible targets for neutrino accelerator experiments if sufficient statistics is expected. At $`\nu `$FMSR , light targets (like $`H_2`$ or $`D_2`$), which are relatively easy polarized, can be used due to the large density of neutrino beam. A number of unique measurements (such as the measurement of parity-violating polarized structure function of neutron) is possible Bigietal .
### II.3 Charmonium production and gluon density
High intensity neutrino beams would also allow studies of charmonium production, a sensitive probe of the gluon distribution function in the nucleon. Contrary to the open-flavor meson production, the production of charmonium states can be described in a model-independent fashion using the factorization theorems of Non-Relativistic QCD (NRQCD),
$$\sigma (A+BH_{c\overline{c}}+X)=\underset{n}{}\frac{F_n}{m_c^{d_n4}}0|𝒪_n^H|0,$$
(7)
which separates short-distance physics, represented by the coefficients $`F_n`$ (which might be sensitive to various parton distribution functions) from the long-distance physics, represented by the NRQCD matrix elements
$$0|𝒪_n^H|0=\underset{X}{}\underset{m_J}{}0|𝒦_n|H_{m_J}+XH_{m_J}+X|𝒦_n^{}|0,$$
(8)
and determine the probabilities of charm quarks produced in the various angular momentum and color (singlet and octet) states by action of NRQCD operators $`𝒦_n^{()}`$ to evolve into a physical charmonium state, like a $`J/\psi `$. At the moment, these matrix elements cannot be computed model-independently. However, they are universal (i.e. process-independent), so they can be extracted from other experiments. Clearly, $`J/\psi `$ produced in sufficiently high numbers can be used to study gluon distribution function in the wide range of $`x`$ Petrov:1999fm .
A major advantage of using the neutrino beam is that, at leading order in $`\alpha _s`$, the spin structure of the $`\nu Z`$ coupling selects a certain combination of octet operators. The largest contribution is from the one with the quantum numbers $`{}_{}{}^{3}S_{1}^{(8)}`$. The differential cross section was calculated in Petrov:1999fm :
$`{\displaystyle \frac{d\sigma (s,Q^2)}{dQ^2}}`$ $`=`$ $`{\displaystyle \frac{\pi ^2\alpha ^2\alpha _s}{3\mathrm{sin}^42\theta _W}}{\displaystyle \frac{1}{\left(Q^2+m_{Z}^{}{}_{}{}^{2}\right)^2}}`$ (9)
$`\times `$ $`{\displaystyle \underset{n}{}}{\displaystyle \frac{0|𝒪_n|0}{m_c^3}}{\displaystyle _{\frac{Q^2+4m_c^2}{s}}^1}𝑑xf_{g/N}(x,Q^2)h_n(y,Q^2),`$
where $`s`$ is the total invariant mass of the $`\nu N`$ system, $`x`$ is the momentum fraction of the incoming gluon, $`Q^2`$ is the momentum-squared transferred from the leptonic system, $`y=\frac{Q^2+4m^2}{sx}`$, and $`f_{g/N}(x,Q^2)`$ is a gluon distribution function in the nucleon. The charmonium structure functions are given by
$`h_{{}_{}{}^{1}S_{0}^{(8)}}(y,Q^2)`$ $`=`$ $`\left(g_V^c\right)^2\times 6{\displaystyle \frac{Q^2m_c^2}{\left(Q^2+4m_c^2\right)^2}}(y^22y+2)`$
$`h_{{}_{}{}^{3}S_{1}^{(8)}}(y,Q^2)`$ $`=`$ $`\left(g_A^c\right)^2\times 2m_c^2{\displaystyle \frac{Q^2(y^22y+2)+16(1y)m_c^2}{\left(Q^2+4m_c^2\right)^2}}`$
$`h_{{}_{}{}^{3}P_{0}^{(8)}}(y,Q^2)`$ $`=`$ $`\left(g_V^c\right)^2\times 2Q^2{\displaystyle \frac{\left(Q^2+12m_c^2\right)^2}{\left(Q^2+4m_c^2\right)^4}}(y^22y+2)`$ (10)
$`h_{{}_{}{}^{3}P_{1}^{(8)}}(y,Q^2)`$ $`=`$ $`\left(g_V^c\right)^2\times 4Q^4{\displaystyle \frac{Q^2(y^22y+2)+16(1y)m_c^2}{\left(Q^2+4m_c^2\right)^4}}`$
$`h_{{}_{}{}^{3}P_{2}^{(8)}}(y,Q^2)`$ $`=`$ $`\left(g_V^c\right)^2\times {\displaystyle \frac{4}{5}}Q^2[{\displaystyle \frac{(y^22y+2)Q^4}{\left(Q^2+4m_c^2\right)^2}}`$
$`+{\displaystyle \frac{48(1y)Q^2m_c^2+96(y^22y+2)m_{c}^{}{}_{}{}^{4}}{\left(Q^2+4m_c^2\right)^2}}],`$
where $`g_A^c=\frac{1}{2}`$ and $`g_V^c=\frac{1}{2}\left(1\frac{8}{3}\mathrm{sin}^2\mathrm{\Theta }_W\right)`$ are the vector and axial couplings of the $`c`$-quark. Clearly, the coupling constants favor the $`{}_{}{}^{3}S_{1}^{(8)}`$ contribution, which is due to the large axial coupling (a similar contribution is, of course, absent in the case of $`J/\psi `$ lepto- and photoproduction). Indeed a numerical estimate Petrov:1999fm shows that this matrix element dominates the total cross section, and also the differential cross section unless $`Q^2m_c^2`$. At large $`Q^2`$, the relative $`Q^4`$ enhancement of the $`P`$-wave structure functions makes them dominant. These structure functions should be incorporated in the specific Monte Carlo generators built for the particular detector design.
An important question to address is the expected event rate of $`J/\psi `$ production. Computing the total cross sections for the $`J/\psi `$ production (Table (1)), a simple calculation shows that currently running neutrino experiments NOMAD and NuTeV could collect a few $`J/\psi `$ events (due to either low energy of the neutrino beam or particular detector configuration) and “confirm” the color octet mechanism. On the contrary, a neutrino experiment at the future Muon Collider would collect about $`3\times 10^3`$ events/year and provide precise measurement of various NRQCD matrix elements and/or the gluon distribution function.
### II.4 Neutrino factory $`=`$ charm factory?
It is clear from the preceding discussion that charm production plays an important role in the studies of nucleon structure and electroweak parameters. It is also important that with the estimated $`10^8`$ well-reconstructed charm events Bigietal $`\nu `$FMSR is also an impressive charm factory. Charm physics is an important complement to the $`B`$-physics program at $`B`$-factories (see, e.g. BaBar ). Besides testing our understanding of QCD effects in charmed particle decays, it also offers an opportunity to look for the effects of new physics in rare decays of charmed mesons, CP-violating asymmetries and $`D\overline{D}`$ mixing studies, as the standard model background to this processes is tiny AP .
It is interesting to see if $`\nu `$FMSR has any advantages over the existing charm experiments. One important advantage of $`D\overline{D}`$ mixing analysis performed at $`\nu `$FMSR that is not available elsewhere involves initial $`D`$ flavor tagging. In particular, $`D^0`$ mesons produced in charged current interactions receive an automatic initial flavor tag in the form of the final state lepton charge. Correlation studies of the charges of the “tag” lepton and, say, lepton from the semileptonic charm decay would offer experimentally clean signatures of $`D\overline{D}`$ mixing.
## III Rare processes
Neutrino-nucleon processes at low momentum transfer are sensitive to generic four-fermion contact terms produced by the high energy neutral current interactions. These four-fermion interactions can be associated with supersymmetric theories with $``$-parity nonconservation, new vector bosons, quark compositness or even loop effects associated with the new flavor-changing neutral current interactions Bigietal . Consider, for instance, the low energy remnant of a generic high energy electron-quark neutral current interaction. It can be represented by
$`_{NC}`$ $`=`$ $`{\displaystyle \underset{q}{}}[\eta _{LL}^{eq}\left(\overline{e_L}\gamma _\mu e_L\right)\left(\overline{q_L}\gamma ^\mu q_L\right)+\eta _{RR}^{eq}\left(\overline{e_R}\gamma _\mu e_R\right)\left(\overline{q_R}\gamma ^\mu q_R\right)`$ (11)
$`+\eta _{LR}^{eq}\left(\overline{e_L}\gamma _\mu e_L\right)\left(\overline{q_R}\gamma ^\mu q_R\right)+\eta _{RL}^{eq}\left(\overline{e_R}\gamma _\mu e_R\right)\left(\overline{q_L}\gamma ^\mu q_L\right)].`$
A similar equation can be written for a direct neutrino-quark interactions. One can use $`SU(2)`$ symmetry to relate $`\nu `$ and $`e`$ couplings
$`\eta _{LL}^{\nu u}`$ $`=`$ $`\eta _{LL}^{ed},\eta _{LL}^{\nu d}=\eta _{LL}^{eu},`$
$`\eta _{LR}^{\nu u}`$ $`=`$ $`\eta _{LR}^{eu},\eta _{LR}^{\nu d}=\eta _{LR}^{ed},`$
so that $`\nu N`$ interactions can be used to constrain $`\eta `$’s of the Lagrangian of Eq. (11). A particular example of a high-energy model that leads to the low-energy Lagrangian of this type is provided by $``$-parity violating SUSY, where at low values of transferred momenta one can integrate out heavy $`\stackrel{~}{e}_{_R^L}^i`$ and $`\stackrel{~}{d}_{_R^L}^i`$ to rewrite the Lagrangian in terms of local four-fermion interactions. Assuming that the squarks of first two generations are degenerate and imposing $`SU(2)`$ symmetry constraints,
$$\eta _{LR}^{ed}=\frac{(\lambda _{1j1}^{})^2}{2m_{\stackrel{~}{u}_L^j}^2}=\frac{(\lambda _{1j1}^{})^2}{2m_{\stackrel{~}{d}_L^j}^2}=\eta _{LR}^{\nu d}.$$
(12)
Here $`\lambda _{ijk}^{}`$ is a parameter of the original $`\text{/}`$ SUSY Lagrangian. Indeed, the best constraint on this coupling, $`\eta _{LR}^{ed}<0.07_{0.24}^{+0.24}`$ comes from the analysis of neutrino nucleon scattering experiments Zeppenfeld:1998un . Other new physics scenarios involve new heavy neutral leptons (models with $`H_L^0\nu _\mu `$ mixing) Gronau:1984ct or new neutral gauge bosons like $`Z^{}`$ which appears in many superstring-motivated models.
## IV Conclusions
New experiments utilizing high energy and intensity neutrino beams would offer a unique opportunity to perform new precision studies of QCD and electroweak interactions. New types of measurements, like charmonium production and extractions of $`V_{cb}`$ and $`V_{ub}`$ CKM matrix elements, will become possible. Interesting new physics scenarios can also be explored. As a result, a high intensity neutrino facility could prove to be a very useful addition to the Muon Collider physics program. |
warning/0003/cond-mat0003153.html | ar5iv | text | # Statistics of the one-dimensional Riemann walk
## 1 Introduction
The Lévy flight is a random walk in continuous space whose step size distribution has a power law tail and is therefore sometimes called a ”Lévy distribution”. The ubiquity of such distributions has been emphasized by many authors, and is a consequence of the power law tail being invariant under convolution. Many interesting instances of the occurrence of Lévy distributions are given by Tsallis and Tsallis et al. . These range from applications in physics (superdiffusion, chaotic fluid flow) and engineering (leaking taps) through studies of the physiology of heart activity, all the way to descriptions of fluctuations of financial markets.
A one-dimensional lattice version of the Lévy flight may be constructed as follows. Let a random walk consist of independent steps, and let the probability $`p(\mathrm{})`$ for a displacement of $`\mathrm{}`$ lattice units in a single step be given by $`p(0)=0`$ and
$$p(\mathrm{})=A|\mathrm{}|^{1\alpha }(\mathrm{}=\pm 1,\pm 2,\mathrm{})$$
(1.1)
Here $`\alpha >0`$ is the Lévy exponent. Normalization of $`p`$ implies that $`A^1=2\zeta (1+\alpha )`$ where $`\zeta `$ is the Riemann zeta function. This random walk was first studied by Gillis and Weiss in 1970. It is called the Riemann walk by Hughes (Ref. , p. 154) and we will conform to that terminology. More generally we call of Riemann type any one-dimensional lattice walk whose $`p(\mathrm{})`$ is asymptotically proportional to $`|\mathrm{}|^{1\alpha }`$ when $`|\mathrm{}|\mathrm{}`$.
Riemann type walks were reviewed in detail by Hughes . Of particular interest is the exponent regime $`0<\alpha 2`$, where these walks have a mean square displacement per step, $`\mathrm{}^2`$, which is infinite. There then exists, at least for certain global walk features, a correspondence between simple random walk on a $`d`$-dimensional lattice and one-dimensional Riemann type walks of exponent $`\alpha =2/d`$. In some ways the fraction $`\frac{2}{\alpha }`$ acts as the walk’s effective dimensionality. But whereas analytical results for noninteger dimension $`d`$ cannot be checked by computer simulations, the full continuum of $`\alpha `$ values is accessible to Monte Carlo studies.
Much interest has centered around the following question. Let there be a $`t`$ step Riemann walk. Then what are the statistical properties of its support $`𝒮(t)`$, i.e., of the set of sites that the walk has visited? There appears immediately an important difference between the exponent regimes $`0<\alpha <1`$ and $`1<\alpha <2`$. In the former regime the Riemann walk is transient and it is easy to show (see Sec. 2.5) that $`𝒮(\mathrm{})`$ is a set of fractal dimension $`d_𝒮=\alpha `$. In the latter case the Riemann walk is recurrent , $`𝒮(\mathrm{})`$ coincides with the full one-dimensional lattice, and $`d_𝒮=1`$. The existing literature deals with the different question of finding the properties of $`𝒮(t)`$ for asymptotically $`t`$; the results reflect, nevertheless, the same distinction between $`0<\alpha <1`$ and $`1<\alpha <2`$. The borderline case $`\alpha =1`$ is more subtle.
Gillis and Weiss study the number $`S(t)`$ of distinct sites in the support. They find, among other results, that for $`t\mathrm{}`$ the average of this random variable behaves as $`\overline{S(t)}t`$ for $`0<\alpha <1`$ and as $`\overline{S(t)}t^{1/\alpha }`$ for $`1<\alpha <2`$, where $``$ indicates asymptotic proportionality. For $`\alpha =1`$ and $`\alpha =2`$ power laws with logarithmic correction factors appear . For $`\alpha >2`$ the result $`\overline{S(t)}t^{1/2}`$ is identical to that for the simple random walk in $`d=1`$.
A recent extension of this work is due to Berkolaiko et al. . Pursuing a question initially asked for the case of the simple random walk by Larralde et al. , these authors investigate the number $`S_N(t)`$ of distinct sites visited by $`N`$ independent $`t`$ step Riemann type walks all starting on the same lattice site. Again power laws appear, both for $`t\mathrm{}`$ at fixed $`N`$ and for $`N\mathrm{}`$ at fixed $`t`$.
The present work extends the investigations of Gillis and Weiss into a different direction. We limit ourselves to the Riemann walk defined by Eq. (1.1), with $`\alpha `$ in the regime of greatest interest, that is, $`0<\alpha <2`$. Our results may be summarized under three headings.
1. Variance $`\overline{\mathrm{\Delta }S^2(t)}`$. For any quantity $`X(t)`$ we will denote its instantaneous deviation from average by $`\mathrm{\Delta }X(t)X(t)\overline{X(t)}`$. Traditionally in this field the calculation of the average number $`\overline{S(t)}`$ of distinct sites visited has been followed by a calculation of the variance $`\overline{\mathrm{\Delta }S^2(t)}`$ of that number. Thus, for the simple random walk $`\overline{S(t)}`$ was first calculated by Dvoretzky and Erdös in 1951, and $`\overline{\mathrm{\Delta }S^2(t)}`$ by Jain and Pruitt in 1970. For the one-dimensional Riemann walk the present work supplements the 1970 results due to Gillis and Weiss for $`\overline{S(t)}`$ by the corresponding ones for the variances $`\overline{\mathrm{\Delta }S^2(t)}`$ in the regime $`0<\alpha <2`$.
2. Variables other than $`S(t)`$. We use the powerful generating function method (GFM), which was introduced into the field of random walks by Montroll and Montroll and Weiss . Overviews of this method are given by Weiss and by Hughes . The first calculation of a variance by the GFM, viz. that of $`S(t)`$ for the simple random walk, is due to Torney in 1986.
In 1994 Coutinho et al. performed Monte Carlo simulations of, among other things, the number of unvisited islands enclosed by the support of the $`t`$ step simple random walk in two dimensions. This led Caser and Hilhorst to analytically determine the asymptotic behavior of the average number of islands. Subsequently Van Wijland et al. developed a compact GFM based analytical scheme for calculating simultaneously the averages, variances, and correlations of a large class of observables characteristic of the support, generically denoted by the symbol $`M(t)`$. In $`d=2`$ this class includes also the total boundary length of the support, and in $`d=3`$ its surface area and Euler index.
Here we bring this scheme to bear on the one-dimensional Riemann walk. The support of this walk consists of alternating sequences of visited and unvisited sites. Among the most prominent members of the class of observables $`M(t)`$ is, next to $`S(t)`$, the number of visited sequences, that we will denote by $`I(t)`$. Islands in $`d=1`$ just are unvisited sequences enclosed by the support, of which there are $`I(t)1`$; the support furthermore has $`2I(t)`$ boundary sites (= visited sites adjacent to an unvisited one). Table I summarizes our results for the asymptotic laws of the averages, variances, and correlations involving $`S(t)`$ and $`I(t)`$. Beyond their intrinsic interest these laws may serve in heuristic arguments in reaction–diffusion processes, e.g., to estimate the trapping probability of an atom that diffuses in a random absorbing environment, or the effective reaction rate between two diffusing species. We defer further comments to Sec. 7.
3. Universality of fluctuations. The deviations from average $`\mathrm{\Delta }S(t)`$,$`\mathrm{\Delta }I(t)`$, $`\mathrm{}`$,$`\mathrm{\Delta }M(t)`$,$`\mathrm{}`$ are randomly time-dependent variables that one would a priori expect to exhibit some degree of correlation. One calls these fluctuations universal – by lack of a better name – when in the limit $`t\mathrm{}`$ all $`\mathrm{\Delta }M(t)`$ are asymptotically equal (up to a proportionality constant) to a single $`M`$ independent stochastic process. For the simple random walk universality was shown to hold in dimensions $`d=2`$ and $`d=3`$ , and not to hold in $`d=4,5,\mathrm{}`$. For the $`d=1`$ Riemann walk we find that universality holds in the exponent regime $`\frac{2}{3}\alpha <2`$, but not for $`0<\alpha <\frac{2}{3}`$. A novelty with respect to the case of the simple random walk is that for $`1<\alpha <2`$ not a single, but two $`M`$ independent processes are needed to describe the universal fluctuations: one applies to bulk and the other to surface observables. The precise statements are given in Sec. 6.
This article is set up as follows. Sec. 2 describes those elements of our analysis that are common to the full exponent interval $`0<\alpha <2`$. Secs. 3 and 4 deal more in particular with the exponent regimes $`0<\alpha <1`$ and $`1<\alpha <2`$, respectively, and derive the asymptotic behavior of averages, variances, and correlations. In Sec. 5 we do the same for the exceptional values $`\alpha =\frac{2}{3}`$ and $`\alpha =1`$. In Sec. 6 we discuss the universality properties. In Sec. 7 we provide some additional interpretation of our results and conclude.
## 2 Observables, averages, and correlations
### 2.1 Observables $`M(t)`$
Our analysis is based on first writing quantities of interest in terms of the field $`m(x,t)`$ of ”complementary occupation numbers” defined by $`m(x,t)=1`$ if site $`x`$ has not yet been visited at time $`t`$, and $`m(x,t)=0`$ otherwise. The expressions of $`S`$ and $`I`$ in terms of $`m`$ are
$`S(t)`$ $`=`$ $`{\displaystyle \underset{x=\mathrm{}}{\overset{\mathrm{}}{}}}[1m(x,t)]`$ (2.1)
$`I(t)`$ $`=`$ $`{\displaystyle \underset{x=\mathrm{}}{\overset{\mathrm{}}{}}}m(x,t)[1m(x+1,t)]`$ (2.2)
$`S`$ and $`I`$ are representatives of a general class of ”observables” $`M`$ that are sums on $`x`$ of a summand to which each lattice site contributes a factor $`m`$$`1m`$, or $`1`$, i.e., the summand tests for the presence of a specific pattern of visited (”black”) and unvisited (”white”) sites. The following slightly more abstract characterization of the $`M`$ will be needed. Let $`A=\{a\}`$ be a finite set of distinct nonnegative integers $`a`$, such that either $`A=\mathrm{}`$ or, if not, $`A`$ includes the element $`a=0`$. The general observable $`M(t)`$ that we will consider is
$$M(t)=\underset{x=\mathrm{}}{\overset{\mathrm{}}{}}\underset{A}{}\mu _A\underset{aA}{}m(x+a,t)$$
(2.3)
where for $`A=\mathrm{}`$ the product is equal to unity and where $`\{\mu _A\}`$ is a set of numerical coefficients characteristic of $`M`$. When their $`M`$ dependence needs to be indicated we will write $`\mu _A[M]`$. Eqs. (2.2) and (2.1) show that $`S(t)`$ and $`I(t)`$ are of the form of Eq. (2.3) with only two nonzero coefficients, as shown in Table II.
Two further examples of observables of type (2.3) are the total number $`S_1(t)`$ of visited sequences consisting of only a single site, and the total number $`I_1(t)`$ of single-site unvisited sequences. Their coefficients $`\mu _A`$ involve sets $`A`$ of up to three elements; they are easily determined and have also been listed in Table II.
The following remarks, important for later, are verified without much effort. The coefficient $`A_{\mathrm{}}`$ is nonzero if and only if $`M`$ is built up exclusively out of factors $`m`$. Since these correspond to visited sites, that make up the ”bulk” of the support, we will call an $`M`$ of this type a bulk observable. Observables built up exclusively out of factors $`m`$ do not occur, since their expectation value on an infinite lattice is infinite. Hence the remaining observables refer to patterns consisting of both visited and unvisited sites, and we will therefore call them surface observables. \[In the terminology of Refs. these are ”black” and ”black-and-white” observables. They might also be called ”S-like” and ”I-like”, respectively.\] The distinction between these two subclasses will play a role only in the exponent regime $`1\alpha <2`$.
### 2.2 Basic formulas for averages and correlations
In this work we will first evaluate the $`t\mathrm{}`$ behavior of the averages $`\overline{M(t)}`$. Then we turn to the covariance matrix $`\overline{\mathrm{\Delta }M(t)\mathrm{\Delta }M^{}(t)}`$, where $`M^{}(t)`$ is a second observable with coefficients $`\mu _A^{}`$. Although the authors of Refs. and deal with the simple random walk, the larger part of their formal developments also holds for the Riemann walk.
The averages $`\overline{M(t)}`$ and $`\overline{M(t)M^{}(t)}`$ can be obtained as follows . Let $`G(x,t)`$ be the Green function of the one-dimensional Riemann walk, that is, the probability for a walker starting at the origin to occupy site $`x`$ after $`t`$ steps. Let $`\widehat{G}(x,z)=_{t=0}^{\mathrm{}}z^tG(x,t)`$ denote its generating function and let G$`{}_{A}{}^{}(z)`$ be the $`|A|\times |A|`$ matrix of elements $`\widehat{G}(aa^{},z)`$ with $`a,a^{}A`$. From this matrix one constructs the ”inverse sum” $`𝖦_A(z)`$ defined by
$$𝖦_A^1(z)=\underset{a,a^{}A}{}[\text{G}_A^1(z)]_{aa^{}}$$
(2.4)
These scalars satisfy certain elementary relations stated in Appendix A as Properties 1–3. Two functions $`C_M(z)`$ and $`C_{MM^{}}(z)`$ are defined in terms of the $`𝖦_A(z)`$ according to
$`C_M(z)`$ $`=`$ $`{\displaystyle \underset{A\mathrm{}}{}}\mu _A{\displaystyle \frac{1}{𝖦_A(z)}}`$ (2.5)
$`C_{MM^{}}(z)`$ $`=`$ $`{\displaystyle \underset{A\mathrm{}}{}}{\displaystyle \underset{B\mathrm{}}{}}\mu _A\mu _B^{}{\displaystyle \underset{r=\mathrm{}}{\overset{\mathrm{}}{}}}\left[{\displaystyle \frac{1}{𝖦_{A(r+B)}(z)}}{\displaystyle \frac{1}{𝖦_A(z)}}{\displaystyle \frac{1}{𝖦_B(z)}}\right]`$ (2.6)
Here $`A(r+B)`$ denotes the union of the set $`B`$, translated by $`r`$, and $`A`$. The averages $`\overline{M(t)}`$ and $`\overline{M(t)M^{}(t)}`$ are then obtained as
$`\overline{M(t)}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi \text{i}}}{\displaystyle \frac{\text{d}z}{z^{t+1}}\frac{1}{(1z)^2}C_M(z)}`$ (2.7)
$`\overline{M(t)M^{}(t)}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi \text{i}}}{\displaystyle \frac{\text{d}z}{z^{t+1}}\frac{1}{(1z)^2}C_{MM^{}}(z)}`$ (2.8)
where the integrations are counterclockwise around the origin.
The identities (2.4)–(2.8) are fundamental to random walk theory; they hold for any translationally invariant random walk, whether on a finite lattice with periodic boundary conditions or on an infinite lattice. They allow for the calculation, in a very compact way, of many known and new results.
Special cases. When $`M(t)=S(t)`$, the following simplifications occur. The sums on $`A`$ and on $`B`$ in Eqs. (2.5) and (2.6) then have only the single term with $`A=\{0\}`$ and $`B=\{0\}`$, respectively. Furthermore G$`{}_{\{0\}}{}^{}(z)=\widehat{G}(0,z)`$, the matrix G$`{}_{\{0\}(r+\{0\})}{}^{}(z)`$ is two by two, and an easy calculation leads to G$`{}_{\{0\}(r+\{0\})}{}^{}(z)=\frac{1}{2}(\widehat{G}(0,z)+\widehat{G}(r,z))`$. When $`M(t)=I(t)`$, the sum in Eq. (2.5) involves G$`{}_{\{0\}}{}^{}(z)`$ and G$`{}_{\{0,1\}}{}^{}(z)=\frac{1}{2}(\widehat{G}(0,z)+\widehat{G}(1,z))`$. The sums on $`A`$ and $`B`$ in Eq. (2.6) then lead to four terms, which may be evaluated with a little more effort.
### 2.3 Limit $`t\mathrm{}`$ and scaling limit
Explicit evaluation of the general expressions (2.4) -(2.8) is limited in practice by the calculation of the inverse sums $`𝖦_A`$, which require the inversion of a matrix of dimension $`|A|`$. Similarly, evaluation of $`𝖦_{A(r+B)}`$ is an inversion problem of dimension $`|A|+|B|`$ (when $`A`$ and $`r+B`$ have an empty intersection). It turns out that the sum on $`r`$ in Eq. (2.6) can be performed only in the scaling limit
$$z1,|r|\mathrm{}\text{with}\xi =r(1z)^{\frac{1}{\alpha }}\text{fixed}$$
(2.9)
Finally, it will be possible to evaluate the integrals in Eqs. (2.7) and (2.8) only asymptotically for $`t\mathrm{}`$, a limit already implied by Eq. (2.9).
In order to prepare for these limits we rewrite the preceding expressions as follows. Using the simplified notation $`G_0(z)=\widehat{G}(0,z)`$ we split the generating function $`\widehat{G}(x,z)`$ up according to
$$\widehat{G}(x,z)=G_0(z)g(x,z)$$
(2.10)
In full analogy to G$`{}_{A}{}^{}(z)`$ we define g$`{}_{A}{}^{}(z)`$ as the matrix of elements $`g(aa^{},z)`$ with $`a,a^{}A`$, and $`g_A^1`$ as the sum of all elements of g$`{}_{}{}^{1}{}_{A}{}^{}`$. Let now J be the square matrix of elements $`J_{aa^{}}=1`$. Then
$$\text{G}_A(z)=G_0(z)\text{J}\text{g}_A(z)$$
(2.11)
and, by Property 1 of Appendix A,
$$𝖦_A(z)=G_0(z)g_A(z)$$
(2.12)
This splitup will be useful for studying the $`z1`$ behavior of $`𝖦_A(z)`$. Although $`G_0(z)`$ may $`(1\alpha <2)`$ or may not $`(0<\alpha <1)`$ diverge as $`z1`$, the functions $`g(x,z)`$ and $`g_A(z)`$ remain finite in that limit.
We now turn to the inverse sum $`𝖦_{A(r+B)}`$ constructed from the matrix G<sub>A∪(r+B)</sub>. The dimension of this matrix is typically $`|A|+|B|`$. Let J<sup>AB</sup> be the $`|A|\times |B|`$ matrix with all $`J_{ab}^{AB}=1`$. We then have (for $`AB=\mathrm{}`$)
$$\text{G}_{A(r+B)}(z)=\left(\begin{array}{cc}\text{G}_A(z)& \widehat{G}(r,z)\text{J}^{AB}\\ \widehat{G}(r,z)\text{J}^{BA}& \text{G}_B(z)\end{array}\right)+\left(\begin{array}{cc}0& \text{V}\\ \text{V}^T& 0\end{array}\right)$$
(2.13)
where V is the matrix of elements
$$V_{a,r+b}=\widehat{G}(r+ba,z)\widehat{G}(r,z)aA,bB$$
(2.14)
and V<sup>T</sup> is its transpose. The first matrix on the RHS of (2.13) has the form (A.4) of Appendix A. Applying Property 3 to that matrix we conclude that
$$\frac{1}{𝖦_{A(r+B)}(z)}=\frac{𝖦_A(z)+𝖦_B(z)2\widehat{G}(r,z)}{𝖦_A(z)𝖦_B(z)\widehat{G}^2(r,z)}+𝒪(\text{V}^2)$$
(2.15)
where we anticipate, and will have to show later, that V is small, that the correction terms are of order $`𝒪(𝐕^2)`$, and that they are negligible for our purpose.
Further analysis depends on the exponent $`\alpha `$. We consider the two main regimes $`0<\alpha <1`$ and $`1<\alpha <2`$ in Secs. 3 and 4, respectively. The exceptional values $`\alpha =\frac{2}{3}`$ and $`\alpha =1`$ are discussed in Sec. 5.
### 2.4 Riemann walk Green function
All quantities of interest have been expressed above in terms of the Riemann walk Green function $`\widehat{G}(x,z)`$. An elementary calculation yields
$`\widehat{G}(x,z)`$ $`=`$ $`{\displaystyle _\pi ^\pi }{\displaystyle \frac{\text{d}q}{2\pi }}{\displaystyle \frac{\text{e}^{\text{i}qx}}{1z\lambda (q)}}`$ (2.16)
$`\lambda (q)`$ $`=`$ $`{\displaystyle \frac{1}{\zeta (1+\alpha )}}{\displaystyle \underset{\mathrm{}=1}{\overset{\mathrm{}}{}}}\mathrm{}^{1\alpha }\mathrm{cos}\mathrm{}q`$ (2.17)
The $`q0`$ behavior of $`\lambda (q)`$ is crucial for the large scale features of the Riemann walk. It is known that
$$\lambda (q)=1C_\alpha |q|^\alpha +𝒪(q^2)(q0)$$
(2.18)
where for completeness we state the explicit expression
$$C_\alpha ^1=2\zeta (1+\alpha )\mathrm{\Gamma }(1+\alpha )/[\pi \mathrm{sin}(\alpha \pi /2)]$$
(2.19)
One finds by standard methods (see e.g. ) that in the limit $`z1`$ the Green function in the origin $`G_0(z)`$ has the asymptotic expansion
$`G_0(z)`$ $`=`$ $`G_0(1)B_\alpha (1z)^{\frac{1}{\alpha }1}+𝒪(1z)(0<\alpha <1;\alpha \frac{1}{2})`$ (2.20)
$`G_0(z)`$ $`=`$ $`\frac{1}{3}\mathrm{log}[c(1z)^1]+𝒪(1z)(\alpha =1)`$ (2.21)
$`G_0(z)`$ $`=`$ $`A_\alpha (1z)^{1+\frac{1}{\alpha }}+𝒪(1)(1<\alpha <2)`$ (2.22)
where $`B_\alpha `$ and $`A_\alpha `$ are the constants
$`B_\alpha `$ $`=`$ $`C_\alpha ^{1/\alpha }/[2\mathrm{sin}(\pi /\alpha )](\frac{1}{2}<\alpha <1)`$ (2.23)
$`A_\alpha `$ $`=`$ $`1/[2\alpha C_\alpha ^{1/\alpha }\mathrm{sin}(\pi /\alpha )](1<\alpha <2)`$ (2.24)
and $`c`$ is a constant such that there is no $`𝒪(1)`$ term in Eq. (2.21). In Eqs. (2.20)–(2.22) and elsewhere we use the following convention. The symbol $`𝒪(X)`$ indicates terms that are of order X in the applicable limit ($`X0`$ or $`X\mathrm{}`$); this however is not to say that all preceding terms are larger. Thus, the nonanalytic term in Eq. (2.20) is larger than the $`𝒪(1z)`$ terms only for $`\frac{1}{2}<\alpha <1`$. For $`0<\alpha <\frac{1}{2}`$ it is present only as a correction to the $`𝒪(1z)`$ terms; the expression for its coefficient $`B_\alpha `$ in that regime is different from Eq. (2.23) but will not be needed. For the borderline case $`\alpha =\frac{1}{2}`$, excluded from Eq. (2.20), we have $`G_0(z)=G_0(1)\frac{2}{\pi }(1z)\mathrm{log}(1z)^1+𝒪(1z)`$; but this special nonanalytic behavior will stay subdominant everywhere in the remainder.
From Eqs. (2.16) and (2.18) one deduces that in the scaling limit (2.9)
$$\widehat{G}(r,z)(1z)^{\frac{1}{\alpha }1}F(\xi )(0<\alpha <2)$$
(2.25)
where $`\xi =r(1z)^{\frac{1}{\alpha }}`$ and $`F(\xi )`$ is the scaling function
$$F(\xi )=_{\mathrm{}}^{\mathrm{}}\frac{\text{d}k}{2\pi }\frac{\text{e}^{\text{i}k\xi }}{1+C_\alpha |k|^\alpha }$$
(2.26)
For $`\xi 0`$ it behaves as
$`F(\xi )`$ $``$ $`2\alpha \zeta (1+\alpha )/[\pi \mathrm{sin}(\alpha \pi )]\xi ^{1+\alpha }(0<\alpha <1)`$ (2.27)
$`F(\xi )`$ $``$ $`\frac{1}{3}\mathrm{log}\xi ^1(\alpha =1)`$ (2.28)
$`F(\xi )`$ $`=`$ $`A_\alpha +𝒪(\xi ^{\alpha 1})(1\alpha <2)`$ (2.29)
In Secs. 4 and 5 we will also use the function $`f(r,z)`$ defined by
$$\widehat{G}(r,z)=G_0(z)f(r,z)(1\alpha <2)$$
(2.30)
In the scaling limit one has $`f(r,z)f(\xi )=F(\xi )/F(0)`$ when $`1<\alpha <2`$.
### 2.5 Support at $`t=\mathrm{}`$
Whereas the remainder of this paper deals with the large $`t`$ behavior, we briefly comment here on the structure of the support $`𝒮(t)`$ at $`t=\mathrm{}`$.
As is well-known ), random walks are recurrent (are transient) if $`G_0(1)=\mathrm{}`$ (if $`G_0(1)<\mathrm{}`$). The Riemann walk of this work is recurrent for $`1\alpha <2`$, which means that all sites are visited with probability 1, and that at $`t=\mathrm{}`$ the support $`𝒮(\mathrm{})`$ coincides with the full one-dimensional lattice.
For $`0<\alpha <1`$, however, the Riemann walk is transient, so that at $`t=\mathrm{}`$ the support $`𝒮(\mathrm{})`$ will still be only a subset of the full lattice. We may estimate the average number of visited sites $`\mathrm{\Sigma }_L`$ between $`x=L`$ and $`x=L`$ in $`𝒮(\mathrm{})`$. According to standard random walk theory
$$\mathrm{\Sigma }_L=\underset{x=L}{\overset{L}{}}\frac{\widehat{G}(x,1)}{G_0(1)}$$
(2.31)
Upon substituting (2.16) in (2.31) one easily evaluates $`\mathrm{\Sigma }_L`$ for asymptotically large $`L`$, with the result that $`\mathrm{\Sigma }_LL^\alpha `$. It follows that the support $`𝒮(\mathrm{})`$ has fractal dimension $`d_𝒮=\alpha `$.
## 3 Riemann walk of exponent $`\mathrm{\hspace{0.17em}0}<\alpha <1`$
### 3.1 Averages
The large time behavior of $`\overline{M(t)}`$ comes from the behavior of $`C_M(z)`$, defined in Eq. (2.5), in the limit $`z1`$. From Eq. (2.10) and the explicit expressions (2.16) and (2.17) it may be shown that $`g(x,z)=g(x,1)+𝒪(1z)`$ for all $`0<\alpha <1`$. Hence
$$g_A(z)=g_A(1)+𝒪(1z)$$
(3.1)
after which it follows from Eqs. (2.12), (2.22), and (3.1) that
$$𝖦_A(z)=𝖦_A(1)B_\alpha (1z)^{\frac{1}{\alpha }1}+𝒪(1z)$$
(3.2)
Inverting this relation and substituting in Eq. (2.5) gives
$$C_M(z)=m_1B_\alpha m_2(1z)^{\frac{1}{\alpha }1}B_\alpha ^2m_3(1z)^{\frac{2}{\alpha }2}+\mathrm{}+𝒪(1z)$$
(3.3)
where the $`m_n`$ are determined by the coefficients $`\mu _A`$ of the observable $`M`$ according to
$$m_n[M]=\underset{A\mathrm{}}{}\frac{\mu _A}{𝖦_A^n(1)}(n=1,2,\mathrm{};\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}<\alpha <1)$$
(3.4)
In Eq. (3.3) the dots stand for a power series in $`(1z)^{\frac{1}{\alpha }1}`$ and the number of nonanalytic terms between the zeroth and the first power of $`1z`$ is equal to $`n_\alpha \frac{\alpha }{1\alpha }1`$. That is, $`n_\alpha `$ is zero for $`0<\alpha <\frac{1}{2}`$ and, as $`\alpha `$ goes up, jumps to $`1,2,3,\mathrm{}`$ at $`\alpha =\frac{1}{2},\frac{2}{3},\frac{3}{4},\mathrm{},`$ respectively. After Laplace inversion we get for the average $`\overline{M(t)}`$ in the limit $`t\mathrm{}`$ the explicit result
$$\overline{M(t)}=m_1t+\frac{B_\alpha }{\mathrm{\Gamma }(3\frac{1}{\alpha })}m_2t^{2\frac{1}{\alpha }}+\frac{B_\alpha ^2}{\mathrm{\Gamma }(4\frac{2}{\alpha })}m_3t^{3\frac{2}{\alpha }}+\mathrm{}+𝒪(1)$$
(3.5)
where the number of nonanalytic terms between the leading and the $`𝒪(1)`$ term is again equal to $`n_\alpha `$.
### 3.2 Correlations
In this subsection we consider two – possibly equal – observables $`M`$ and $`M^{}`$, represented by sets of coefficients $`\{\mu _A\}`$ and $`\{\mu _A^{}\}`$, respectively, and wish to study their correlation. The starting point is Eq. (2.6) for $`C_{MM^{}}(z)`$, in which we substitute Eq. (2.15). Whereas $`𝖦_A(z)`$ and $`𝖦_B(z)`$ tend to finite values in the limit $`z1`$, the Green function $`\widehat{G}(r,z)`$ vanishes in that limit when taken with $`\xi `$ fixed. This suggests that we expand in powers of $`\widehat{G}(r,z)`$,
$$\frac{1}{𝖦_{A(r+B)}(z)}\frac{1}{𝖦_A(z)}\frac{1}{𝖦_B(z)}=\underset{n=1}{\overset{\mathrm{}}{}}C_{AB}^{(n)}(z)\widehat{G}^n(r,z)+𝒪(\text{V}^2)$$
(3.6)
with coefficients
$$C_{AB}^{(n)}(z)=\{\begin{array}{cc}[𝖦_A(z)+𝖦_B(z)][𝖦_A(z)𝖦_B(z)]^{\frac{n}{2}1}\hfill & (n\text{ even})\hfill \\ 2[𝖦_A(z)𝖦_B(z)]^{\frac{n+1}{2}}\hfill & (n\text{ odd})\hfill \end{array}$$
(3.7)
The sum over space that occurs in Eq. (2.6) leads us to now consider the sums $`_r\widehat{G}^n(r,z)`$. From conservation of probability one finds that for $`n=1`$
$$\underset{r}{}\widehat{G}(r,z)=(1z)^1$$
(3.8)
For general $`n`$ the calculation of $`_r\widehat{G}^n(r,z)`$ is slightly more laborious; after substituting Eq. (2.16) for $`\widehat{G}`$ one finds by explicit expansion in powers of $`1z`$ that for $`z1`$ the sum on $`r`$ behaves as
$$\underset{r}{}\widehat{G}^n(r,z)F_{\alpha ,n}(1z)^{1+(n1)(\frac{1}{\alpha }1)}+𝒪(1)(0<\alpha <1;\alpha 1\frac{1}{n})$$
(3.9)
For $`n=1\frac{1}{1\alpha }`$ (with $`n=1,2,\mathrm{}`$) the sum on $`r`$ instead behaves as $`\mathrm{log}(1z)`$; this happens, in particular, for $`n=3`$ when $`\alpha =\frac{2}{3}`$, a case studied separately in Sec. 5.
We will now continue to consider the generic case. The nonanalytic term on the RHS of Eq. (3.9) dominates the $`𝒪(1)`$ term only for $`n<\frac{1}{1\alpha }`$. In that case (3.9) follows just from the scaling form (2.22) of $`\widehat{G}`$ and from the $`\xi 0`$ behavior (2.27) of $`F(\xi )`$. One then finds for the prefactor $`F_{\alpha ,n}`$ the expression
$$F_{\alpha ,n}=2_0^{\mathrm{}}\text{d}\xi F^n(\xi )(1\frac{1}{n}<\alpha <1)$$
(3.10)
For $`0<\alpha <1\frac{1}{n}`$ the expression for $`F_{\alpha ,n}`$ is different and will not be needed. Eq. (3.10) shows that the main contribution to the sum on $`r`$ comes from $`\xi 1`$, that is, from $`r(1z)^{\frac{1}{\alpha }}`$. For $`n=1`$ and $`n=2`$ the integral (3.10) yields the explicit results $`F_{\alpha ,1}=1`$, in agreement with Eq. (3.8), and $`F_{\alpha ,2}=(\frac{1}{\alpha }1)B_\alpha `$, respectively. For $`n>\frac{1}{1\alpha }`$ the sum on $`r`$ draws its main contribution from the short distance (nonscaling) regime $`r1`$, and is of $`𝒪(1)`$ for $`z1`$.
By successively substituting Eq. (3.7) in Eq. (3.6), neglecting the $`𝒪`$(V$`{}_{}{}^{2})`$ terms in that equation – which is justified in Appendix B – , then substituting Eq. (3.6) in Eq. (2.6), expanding $`𝖦_A(z)`$ and $`𝖦_B(z)`$ according to Eq. (3.2), and using the $`z1`$ behavior of $`_r\widehat{G}^n(r,z)`$ obtained in Eqs. (3.8) and (3.9) we find
$`C_{MM^{}}(z)=2(1z)^1m_1m_1^{}`$
$`(1z)^{\frac{1}{\alpha }2}B_\alpha \left(3\frac{1}{\alpha }\right)(m_1m_2^{}+m_2m_1^{})`$
$`(1z)^{\frac{2}{\alpha }3}\left[\left(4\frac{2}{\alpha }\right)B_\alpha ^2(m_1m_3^{}+m_3m_1^{}+m_2m_2^{})+2F_{\alpha ,3}m_2m_2^{}\right]`$
$`\mathrm{}+𝒪\left(1\right)`$ (3.11)
Here the dots stand for terms of order $`(1z)^{1+k(\frac{1}{\alpha }1)}`$, with $`k=3,4,\mathrm{};`$ and the $`m_n^{}`$ are related to $`M^{}`$ in the same way as the $`m_n`$ are to $`M`$. For $`t\mathrm{}`$ we therefore find by substituting (3.11) in Eq. (2.8) and Laplace inverting
$$\begin{array}{cc}& \overline{M(t)M^{}(t)}=m_1m_1^{}t^2\hfill \\ & +\frac{B_\alpha }{\mathrm{\Gamma }(3\frac{1}{\alpha })}(m_1m_2^{}+m_2m_1^{})t^{3\frac{1}{\alpha }}\hfill \\ & +\left[\frac{B_\alpha ^2}{\mathrm{\Gamma }(4\frac{2}{\alpha })}(m_1m_3^{}+m_3m_1^{}+m_2m_2^{})+\frac{2F_{\alpha ,3}}{\mathrm{\Gamma }(5\frac{2}{\alpha })}m_2m_2^{}\right]t^{4\frac{2}{\alpha }}\hfill \\ & +\mathrm{}+𝒪(t)\hfill \end{array}$$
(3.12)
The successive terms in the above series all have one power of $`t`$ more than the corresponding terms in the series (3.5) for $`\overline{M(t)}`$, and the number of nonanalytic terms between the leading and the $`𝒪(t)`$ term is once more equal to $`n_\alpha 1`$. The product $`\overline{M(t)}`$ $`\overline{M^{}(t)}`$, which follows from Eq. (3.5), now has to be subtracted from the series (3.2). This exactly cancels the terms in (3.2) proportional to $`t^2`$ and to $`t^{3\frac{1}{\alpha }}`$ but leaves those proportional to $`t^{4\frac{2}{\alpha }}`$ and of $`𝒪(t)`$. The $`t^{4\frac{2}{\alpha }}`$ terms are leading only if $`\frac{2}{3}<\alpha <1`$. Hence we find for the correlation between observables $`M`$ and $`M^{}`$ in the limit $`t\mathrm{}`$
$$\overline{\mathrm{\Delta }M(t)\mathrm{\Delta }M^{}(t)}_\alpha ^2m_2m_2^{}t^{4\frac{2}{\alpha }}$$
(3.13)
valid for $`\frac{2}{3}<\alpha <1`$, and in which
$$_\alpha ^2=\frac{2F_{\alpha ,3}}{\mathrm{\Gamma }(5\frac{2}{\alpha })}+\frac{B_\alpha ^2}{\mathrm{\Gamma }(4\frac{2}{\alpha })}\frac{B_\alpha ^2}{\mathrm{\Gamma }^2(3\frac{1}{\alpha })}$$
(3.14)
We have supposed here that $`m_2,m_2^{}0`$. The preceding analysis changes when either of these two coefficients vanishes. We do not know of any physically interesting examples where this happens, and do not pursue our analysis in this direction.
The borderline case $`\alpha =\frac{2}{3}`$ is considered in Sec. 5. In the interval $`0<\alpha <\frac{2}{3}`$ the calculation of the present section applies, but with the result that
$$\overline{\mathrm{\Delta }M(t)\mathrm{\Delta }M^{}(t)}\kappa _{MM^{}}t(0<\alpha <\frac{2}{3})$$
(3.15)
in which the coefficient $`\kappa _{MM^{}}`$ has contributions from the $`𝒪(1)`$ terms in Eq. (3.5) and the $`𝒪(t)`$ terms in Eq. (3.2), and does not factor into an $`M`$ and an $`M^{}`$ dependent constant. This difference between Eqs. (3.13) and (3.15) is crucial for the phenomenon of universality discussed in Sec. 6.
## 4 Riemann walk of exponent $`\mathrm{\hspace{0.17em}1}<\alpha <2`$
### 4.1 Averages
The calculation of $`\overline{M(t)}`$ starts again from the series (2.5) for $`C_M(z)`$. The calculation in the exponent regime $`1<\alpha <2`$ is different from that of the preceding section because now $`G_0(z)`$ diverges as $`z1`$. Since $`g_A(z)`$ remains finite for $`z1`$, this suggests that we use Eq. (2.12) and expand $`𝖦_A(z)`$ in powers of $`g_A(z)/G_0(z)`$. This yields
$$C_M(z)=\underset{A\mathrm{}}{}\mu _A\frac{1}{G_0(z)}\left[\mathrm{\hspace{0.17em}\hspace{0.17em}1}+\frac{g_A(z)}{G_0(z)}+\frac{g_A^2(z)}{G_0^2(z)}+𝒪(\frac{g_A^3}{G_0^3})\right]$$
(4.1)
We now substitute in Eq. (4.1) the expansion (2.22) for $`G_0(z)`$ and use that $`g_A(1)`$ is finite. The result is a power series in $`1z`$ in which there appear coefficients that we denote again by $`m_n`$ but that are defined for $`1\alpha <2`$ as
$$m_n[M]=\underset{A\mathrm{}}{}\mu _Ag_A^n(1)(n=0,1,2,\mathrm{};\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}1}\alpha <2)$$
(4.2)
It will turn out that we need only the leading term, which is
$$C_M(z)\{\begin{array}{c}A_\alpha ^1m_0(1z)^{1\frac{1}{\alpha }}(m_00)\hfill \\ A_\alpha ^2m_1(1z)^{2\frac{2}{\alpha }}(m_0=0,m_10)\hfill \end{array}$$
(4.3)
We pause to note that in the terminology of Sec. 2.1 the condition $`m_00`$ characterizes the bulk or ”S-like” observables, and the condition $`m_0=0`$ the surface or ”I-like” observables. This is the first equation where a difference appears between these two subclasses; in its analog, Eq. (3.5) of the preceding section, no such distinction appears.
Upon using Eq. (4.3) in Eq. (2.7) we obtain after an inverse Laplace transformation the asymptotic expansion of $`\overline{M(t)}`$ as $`t\mathrm{}`$,
$$\overline{M(t)}\{\begin{array}{c}[A_\alpha \mathrm{\Gamma }(1+\frac{1}{\alpha })]^1m_0t^{\frac{1}{\alpha }}(m_00)\hfill \\ [A_\alpha \mathrm{\Gamma }(\frac{2}{\alpha })]^1m_1t^{\frac{2}{\alpha }1}(m_0=0,m_10)\hfill \end{array}$$
(4.4)
where the dots indicate terms of lower order in $`t`$.
### 4.2 Correlations
For the calculation of the correlation $`\overline{M(t)M^{}(t)}`$ via Eqs. (2.8) and (2.6) we have to return again to expression (2.15) for $`1/𝖦_{A(r+B)}(z)`$, which is needed in Eq. (2.6). We use Eqs. (2.15), (2.12), and (2.30) to rewrite this quantity as
$$\frac{1}{𝖦_{A(r+B)}(z)}=\frac{1}{G_0(z)}\frac{2(1f(r,z))\frac{g_A(z)}{G_0(z)}\frac{g_B(z)}{G_0(z)}}{1f^2(r,z)\frac{g_A(z)}{G_0(z)}\frac{g_B(z)}{G_0(z)}+\frac{g_A(z)g_B(z)}{G_0^2(z)}}+𝒪(\text{V}^2)$$
(4.5)
The function $`f(r,z)`$ was defined in Eq. (2.30). In the scaling limit $`f(r,z),`$ $`g_A(z)`$, and $`g_B(z)`$ have finite limits, whereas $`G_0(z)`$ diverges. An expansion in inverse powers of $`G_0(z)`$ corresponds therefore to an expansion in ascending powers of $`1z`$. Writing for short $`f,g_A,g_B`$, and $`G_0`$ when $`f(r,z),`$$`g_A(z),`$$`g_B(z)`$, and $`G_0(z)`$ are meant, we find after a straightforward calculation
$`C_{MM^{}}(z)`$ $`=`$ $`{\displaystyle \underset{r=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{G_0}}{\displaystyle \frac{f}{1+f}}{\displaystyle \underset{A,B\mathrm{}}{}}\mu _A\mu _B^{}[\mathrm{\hspace{0.17em}\hspace{0.17em}2}+(2+f){\displaystyle \frac{g_A+g_B}{G_0}}`$ (4.6)
$`+{\displaystyle \frac{22ff^3}{(1f)(1+f)^2}}{\displaystyle \frac{g_A^2+g_B^2}{G_0^2}}`$
$`+{\displaystyle \frac{2}{(1f)(1+f)^2}}{\displaystyle \frac{g_Ag_B}{G_0^2}}+𝒪({\displaystyle \frac{g_A^3}{G_0^3}},{\displaystyle \frac{g_B^3}{G_0^3}})]`$
Let us write $`m_n^{}=m_n[M^{}]`$ for the coefficients that characterize the observable $`M^{}`$. The two distinct cases described by Eq. (4.4) now lead to the following possibilities.
Case (i): $`m_00`$ and $`m_0^{}0`$. In this case the leading term in the expression in brackets in Eq. (4.6) survives under the sum on $`A`$ and $`B`$.
Case (ii): $`m_0=0,m_10`$, and $`m_0^{}0`$. In this case in order to survive a term in the bracketed expression should contain at least one factor $`g_A(z)`$.
Case (iii): $`m_0=m_0^{}=0`$ but $`m_10`$ and $`m_1^{}0`$. In this case a term in order to survive must contain at least one factor $`g_A(z)`$ and one factor $`g_B(z)`$.
Upon using in each of these cases for $`G_0(z)`$the expansion (2.22), passing to the scaling limit, and writing $`f`$ for $`f(\xi )`$, we find that the result is
$$C_{MM^{}}(z)\{\begin{array}{c}A_\alpha ^1(1z)^{1\frac{2}{\alpha }}f_{00}m_0m_0^{}\hfill \\ A_\alpha ^1(1z)^{2\frac{3}{\alpha }}f_{10}m_1m_0^{}\hfill \\ A_\alpha ^1(1z)^{3\frac{4}{\alpha }}f_{11}m_1m_1^{}\hfill \end{array}$$
(4.7)
in the three cases (i), (ii), and (iii), respectively; here the coefficients $`f_k\mathrm{}`$ represent the integrals
$`f_{00}`$ $`=`$ $`4{\displaystyle _0^{\mathrm{}}}\text{d}\xi f(1+f)^1`$
$`f_{10}`$ $`=`$ $`2{\displaystyle _0^{\mathrm{}}}\text{d}\xi f(2+f)(1+f)^1`$ (4.8)
$`f_{11}`$ $`=`$ $`4{\displaystyle _0^{\mathrm{}}}\text{d}\xi f(1f)^1(1+f)^3`$
After substituting Eqs. (4.7) in Eq. (2.8) and carrying out the inverse Laplace transformation we find, in the limit $`t\mathrm{}`$,
$$\overline{M(t)M^{}(t)}\{\begin{array}{cc}\hfill \mathrm{\Gamma }^1(1+\frac{2}{\alpha })& A_\alpha ^1f_{00}m_0m_0^{}t^{\frac{2}{\alpha }}\hfill \\ \hfill \mathrm{\Gamma }^1(\frac{3}{\alpha })& A_\alpha ^2f_{10}m_1m_0^{}t^{\frac{3}{\alpha }1}\hfill \\ \hfill \mathrm{\Gamma }^1(1+\frac{4}{\alpha })& A_\alpha ^3f_{11}m_1m_1^{}t^{\frac{4}{\alpha }2}\hfill \end{array}$$
(4.9)
respectively, for the three cases distinguished above. Upon combining these results with those of Section 4.1 one obtains, for $`t\mathrm{}`$,
$$\overline{\mathrm{\Delta }M(t)\mathrm{\Delta }M^{}(t)}\{\begin{array}{c}_{\alpha 00}^2m_0m_0^{}t^{\frac{2}{\alpha }}\hfill \\ _{\alpha 10}^2m_1m_0^{}t^{\frac{3}{\alpha }1}\hfill \\ _{\alpha 11}^2m_1m_1^{}t^{\frac{4}{\alpha }2}\hfill \end{array}$$
(4.10)
in which the coefficients $`_{\alpha k\mathrm{}}`$ are given by
$$\begin{array}{c}_{\alpha 00}^2=A_\alpha ^1[f_{00}\mathrm{\Gamma }^1(1+\frac{2}{\alpha })A_\alpha ^1\mathrm{\Gamma }^2(1+\frac{1}{\alpha })]\hfill \\ _{\alpha 10}^2=A_\alpha ^2[f_{10}\mathrm{\Gamma }^1(\frac{3}{\alpha })A_\alpha ^1\mathrm{\Gamma }^1(1+\frac{1}{\alpha })\mathrm{\Gamma }^1(\frac{2}{\alpha })]\hfill \\ _{\alpha 11}^2=A_\alpha ^3[f_{11}\mathrm{\Gamma }^1(1+\frac{4}{\alpha })A_\alpha ^1\mathrm{\Gamma }^2(\frac{2}{\alpha })]\hfill \end{array}$$
(4.11)
in the three cases (i), (ii), and (iii) defined above, respectively. We recall that the $`f_k\mathrm{}`$ on the RHS of Eq. (4.11) are given by Eq. (4.8) as integrals on $`f(\xi )`$, with $`f(\xi )`$ in turn given by Eq. (2.26).
## 5 Riemann walk of exponents $`\alpha =\frac{2}{3}`$ and $`\alpha =1`$
### 5.1 Exponent $`\alpha =\frac{2}{3}`$
In this special case $`\overline{M(t)}`$ is still given by Eq. (3.5). However, the calculation of $`\overline{M(t)M^{}(t)}`$ has to be reconsidered, as signalled by the fact that $`F_{\alpha ,3}`$ in Eq. (3.11) diverges for $`\alpha \frac{2}{3}^+`$. In order to calculate $`_r\widehat{G}(r,z)`$ we cannot now use Eq. (3.9). Instead we replace $`\widehat{G}(r,z)`$ by its scaling form (2.22) but take into account that $`|\xi |`$ has a lower cutoff $`|\xi |\text{cst}\times (1z)^{\frac{2}{3}}`$. This gives
$`{\displaystyle \underset{r}{}}\widehat{G}^3(r,z)`$ $``$ $`2{\displaystyle _{\text{cst}\times (1z)^{2/3}}^{\mathrm{}}}\text{d}\xi F^3(\xi )`$ (5.1)
$``$ $`\frac{1}{2}𝒞^2\mathrm{log}(1z)^1+𝒪(1)(z1)`$
where in the second step we used Eq. (2.24) and found for the coefficient the value $`𝒞^2=2^{12}3^{11/2}\pi ^3\zeta ^3(\frac{3}{2})`$. In this case, due to the $`\mathrm{log}(1z)`$ in the equation above, $`\overline{M(t)M^{}(t)}`$ is larger than the product $`\overline{M(t)}`$ $`\overline{M^{}(t)}`$ by a factor $`\mathrm{log}t`$, and determines by itself alone the final result, which reads
$$\overline{\mathrm{\Delta }M(t)\mathrm{\Delta }M^{}(t)}𝒞^2m_2m_2^{}t\mathrm{log}t(t\mathrm{})$$
(5.2)
This $`t\mathrm{log}t`$ behavior is the same as in the well-known case of the simple random walk in spatial dimension $`d=3`$ .
### 5.2 Exponent $`\alpha =1`$
The case of Lévy exponent $`\alpha =1`$ is subtler than the others. Since it is closely analogous to the simple random walk in dimension $`d=2`$ , we will not present all steps in detail. Eq. (2.21) shows that for $`\alpha =1`$ the Green function in the origin, $`G_0(z)`$, diverges as $`z1`$. We can therefore expand $`C_M(z)`$ as a series in the same way as in Eq. (4.1). Since here again $`g_A(z)=g_A(1)+𝒪(1z)`$, and in view of the logarithmic behavior (2.21), this series now leads to an expansion of $`C_M(z)`$ in inverse powers of $`G_0(z)`$. If the first nonzero term is of order $`k+1`$, then we have explicitly
$$C_M(z)=m_kG_0^{k1}(z)m_{k+1}G_0^{k2}(z)m_{k+2}G_0^{k3}(z)+\mathrm{}$$
(5.3)
with the $`m_n`$ defined by Eq. (4.2). The cases of physical interest have $`k=0`$ (bulk observables) or $`k=1`$ (surface observables), but it will be notationally convenient to keep $`k`$ as a parameter. We will also refer to it as the order of $`M`$.
To find $`C_{MM^{}}(z)`$ we may still start from Eq. (4.5), but now the expansion of this equation runs differently. The reason is that for $`z1`$ at fixed $`\xi `$ the function $`f(r,z)`$ (defined by (2.30)) behaves as $`F(\xi )/G_0(z)`$ and so is of the same order as $`g_A(z)/G_0(z)`$ and $`g_B(z)/G_0(z)`$. We therefore have to perform a double expansion of the RHS of Eq. (4.5) in terms of on the one hand $`F/G_0`$ and on the other hand $`g_A/G_0`$ and $`g_B/G_0`$. The sum on $`r`$, which in the scaling limit becomes an integral on $`\xi `$, then leads to the appearance of coefficients $`F_{1,n}`$ defined as in Eq. (3.10) but with $`\alpha =1`$. Special cases are $`F_{1,1}=1`$ and $`F_{1,2}=\frac{1}{3}`$. Let $`m_k`$ and $`m_k^{}^{}`$ be the first nonzero coefficients in the expansions of $`C_M(z)`$ and $`C_M^{}(z)`$, respectively. Then we find for $`C_{MM^{}}(z)`$, retaining only the three leading order terms in the limit $`z1`$,
$`C_{MM^{}}(z)`$ $`{\displaystyle \frac{1}{(1z)G_0^{k+k^{}+2}(z)}}[2m_km_k^{}^{}`$
$``$ $`G_0^1(z)\left(\frac{1}{3}(k+k^{}+2)a_2m_km_k^{}^{}2(m_km_{k^{}+1}^{}+m_{k+1}m_k^{}^{})\right)`$
$`+`$ $`G_0^2(z)(2F_{1,3}(k+1)(k^{}+1)a_3m_km_k^{}^{}`$
$`\frac{1}{3}(k+k^{}+3)a_2(m_km_{k^{}+1}^{}+m_{k+1}m_k^{}^{})`$
$`+2(m_km_{k^{}+2}^{}+m_{k+1}m_{k^{}+1}^{}+m_{k+2}m_k^{}^{}))]`$ (5.4)
The inverse Laplace transforms of $`C_M(z)`$ and $`C_{MM^{}}(z)`$ may be found with the help of the explicit expression (2.21) for $`G_0(z)`$ and the integrals of Ref. . We state only the explicit result for $`\overline{M(t)}`$, which is, for $`t\mathrm{}`$,
$`\overline{M(t)}{\displaystyle \frac{3^{k+1}t}{\mathrm{log}^{k+1}ct}}[`$ $`m_k+{\displaystyle \frac{1}{\mathrm{log}ct}}\left((1\gamma )(k+1)m_k+3m_{k+1}\right)`$
$`+{\displaystyle \frac{1}{\mathrm{log}^2ct}}((1\frac{1}{12}\pi ^2\gamma +\frac{1}{2}\gamma ^2)(k+1)(k+2)m_k`$
$`3(1\gamma )(k+2)m_{k+1}+9m_{k+2})]`$ (5.5)
in which $`\gamma =0.577215\mathrm{}`$ denotes Euler’s constant. Both $`\overline{M(t)M^{}(t)}`$ and the product $`\overline{M(t)}`$ $`\overline{M^{}(t)}`$ the appear as $`t^2`$ times a power series in $`1/\mathrm{log}ct`$ of which the leading term is of order $`k+k^{}+2`$, and in which the three leading orders have to be retained. Upon carrying out the subtraction one finds that the two leading orders cancel and the correlation $`\overline{\mathrm{\Delta }M(t)\mathrm{\Delta }M^{}(t)}`$ appears to be proportional to $`t^2/\mathrm{log}^{k+k^{}+4}ct`$. Explicitly, as $`t\mathrm{}`$,
$$\overline{\mathrm{\Delta }M(t)\mathrm{\Delta }M^{}(t)}𝒜^2(k+1)(k^{}+1)m_km_k^{}^{}\frac{3^{k+k^{}+2}t^2}{\mathrm{log}^{k+k^{}+4}ct}(t\mathrm{})$$
(5.6)
in which the coefficient $`𝒜`$ is given by
$$𝒜^2=1+(F_{1,3}\frac{1}{6})\pi ^2$$
(5.7)
Numerical evaluation gives $`F_{1,3}=0.27415\mathrm{}`$, whence $`𝒜^2=2.0608\mathrm{}`$. Eq. (5.6) is the same as for the two-dimensional simple random walk , but with a different constant $`𝒜`$.
## 6 Universality of fluctuations
We consider in this section the normalized deviations from average
$$\theta _M(t)=\frac{\mathrm{\Delta }M(t)}{\overline{\mathrm{\Delta }M^2(t)}^{1/2}}$$
(6.1)
These random functions of time satisfy by construction
$$\overline{\theta _M(t)}=0,\overline{\theta _M^2(t)}=1$$
(6.2)
We consider now two arbitrary observables $`M`$ and $`M^{}`$. When $`\frac{2}{3}\alpha <1`$ we have from Eq. (6.1) together with either Eq. (5.2) or Eq. (3.13) that
$$\overline{\theta _M(t)\theta _M^{}(t)}=1$$
(6.3)
It then follows from Eqs. (6.2) and (6.3) that the difference $`\theta _M\theta _M^{}`$ is a random variable of zero average and zero variance. Such a random variable can only be itself equal to zero. We therefore deduce that, when $`\frac{2}{3}\alpha <1`$, in the limit $`t\mathrm{}`$ all $`\theta _M(t)`$ are equal to a single random variable, which we will call $`\mathrm{\Theta }_\alpha (t)`$, thus indicating explicitly its $`\alpha `$ dependence.
When $`1\alpha <2`$ we have for two observables $`M`$ and $`M^{}`$ whose orders, $`k`$ and $`k^{}`$, are equal from Eq. (6.1) and either Eq. (5.6) or Eq. (4.10) again the result (6.3). Hence, when $`1\alpha <2`$, in the limit $`t\mathrm{}`$ all $`\theta _M(t)`$ with $`k=0`$ are equal to a single random variable – that we will call $`\mathrm{\Theta }_{\alpha 0}(t)`$ –, and all $`\theta _M(t)`$ with $`k=1`$ are similarly equal to a single random variable – that we will call $`\mathrm{\Theta }_{\alpha 1}(t)`$. The variables $`\mathrm{\Theta }_\alpha (t)`$, $`\mathrm{\Theta }_{\alpha 0}(t)`$, and $`\mathrm{\Theta }_{\alpha 1}(t)`$ are universal in the sense that they are independent of the observables $`M`$ (but depend at most on their order).
In each of these cases the key ingredient necessary for arriving at Eq. (6.3) is the factorization of $`\overline{M(t)M^{}(t)}`$ into an $`M`$ and an $`M^{}`$ dependent part. This also explains why for $`\alpha <\frac{2}{3}`$ the same reasoning fails.
The cross correlation between $`\mathrm{\Theta }_{\alpha 0}(t)`$ and $`\mathrm{\Theta }_{\alpha 1}(t)`$ is easily found from the correlation between a $`\theta _M(t)`$ and a $`\theta _M^{}(t)`$ with $`k=0`$ and $`k^{}=1`$, and use of the second one of Eqs. (4.10). The answer is independent of the choice of $`M`$ and $`M^{}`$, as it had to be, and reads
$$\overline{\mathrm{\Theta }_{\alpha 0}(t)\mathrm{\Theta }_{\alpha 1}(t)}=_{\alpha 10}^2(_{\alpha 00}_{\alpha 11})^1$$
(6.4)
The coefficient ratio on the RHS of this equation depends only on the exponent $`\alpha `$ and must necessarily be less than unity. In the limit $`\alpha 1^+`$ it approaches unity and for $`\alpha <1`$ the distinction between surface and bulk observables is no longer reflected in the fluctuations.
Upon combining all these conclusions we get explicitly
$$\mathrm{\Delta }M(t)=\{\begin{array}{cc}m_2𝒞(t\mathrm{log}t)^{\frac{1}{2}}\mathrm{\Theta }_{\frac{2}{3}}(t)\hfill & (\alpha =\frac{2}{3})\hfill \\ m_2_\alpha t^{2\frac{1}{\alpha }}\mathrm{\Theta }_\alpha (t)\hfill & (\frac{2}{3}<\alpha <1)\hfill \\ m_k𝒜(k+1)t(\frac{1}{3}\mathrm{log}ct)^{k2}\mathrm{\Theta }_1(t)\hfill & (\alpha =1;k=0,1)\hfill \\ m_0_{\alpha 00}t^{\frac{1}{\alpha }}\mathrm{\Theta }_{\alpha 0}(t)\hfill & (1<\alpha <2;k=0)\hfill \\ m_1_{\alpha 11}t^{\frac{2}{\alpha }1}\mathrm{\Theta }_{\alpha 1}(t)\hfill & (1<\alpha <2;k=1)\hfill \end{array}$$
(6.5)
Here all $`M`$ dependence is contained in the coefficients $`m_n`$.
## 7 Conclusions
We have studied a large class of properties $`M(t)`$ of the support of the one-dimensional $`t`$ step Riemann walk. These include the number $`S(t)`$ of distinct sites visited, and the number $`I(t)`$ of sequences of visited sites. The $`M(t)`$ fall into two classes, the bulk or S-like properties, and the surface or I-like properties. The asymptotic laws found in the preceding sections for the averages, variances, and correlation of $`S(t)`$ and $`I(t)`$ have been summarized in Table I in the Introduction.
It appears from that table that in the exponent regime $`0<\alpha 1`$ the ratios $`\overline{\mathrm{\Delta }S^2(t)}^{1/2}/\overline{S(t)}`$ and $`\overline{\mathrm{\Delta }I^2(t)}^{1/2}/\overline{I(t)}`$ tend to zero when $`t\mathrm{}`$, which indicates that the distributions of $`S(t)`$ and of $`I(t)`$ become infinitely narrowly peaked around their average. Hence in this exponent regime the ratio $`s(t)\overline{S(t)}/\overline{I(t)}`$ represents the average number of sites per visited sequence. When $`\alpha `$ is strictly less than unity we have explicitly
$$\underset{t\mathrm{}}{lim}s(t)=\frac{m_1[S]}{m_1[I]}=\frac{\widehat{G}(0,1)\widehat{G}(1,1)}{\widehat{G}(0,1)+\widehat{G}(1,1)}(0<\alpha <1)$$
(7.1)
The first equality is based on Eq. (3.5) and in the second one we used the definition (3.4) of the $`m_n[M]`$ and the remarks at the end of Sec. 2.2. The finiteness of the result (7.1) means that in the large $`t`$ limit every new step of the walk creates a new visited sequence with a finite nonzero probability. This explains that in this exponent regime the asymptotic power laws do not distinguish between bulk and surface properties. For $`\alpha 1^{}`$ expression (7.1) diverges, and when $`\alpha =1`$ the ratio $`s(t)`$ increases logarithmically with $`t`$.
In the exponent regime $`1<\alpha <2`$ the appropriately scaled distributions of $`S(t)`$ and $`I(t)`$ are of finite width even in the limit $`t\mathrm{}`$. The support has an ”interior”, bulk and surface properties have different asymptotic power laws, and $`s(t)\mathrm{}`$ as $`t\mathrm{}`$. In the terminology of critical phenomena, this regime is fluctuation dominated. In this regime the universality of fluctuations holds in a slightly weaker but at least as interesting a sense as for $`0<\alpha 1`$. To describe the fluctuations, not a single but two universal stochastic variables are needed, one applying to the bulk and the other to the surface properties. These two variables become fully correlated in the limit $`\alpha 1^+`$.
Finally we remark that when $`2/\alpha `$ is equal to one of the integers $`2,3,4,\mathrm{}`$, the asymptotic laws of Table I coincide with the ones known to hold for the simple random walk on a lattice of dimension $`d=2/\alpha `$. Similarly, the universality properties for those $`\alpha `$ values have their analogs in the $`d`$-dimensional simple random walk. Hence the ”rule of the effective dimensionality”, which states the correspondence $`\alpha 2/d`$, applies to all properties that we have studied. Of course it must break down when the comparison between the Riemann walk and the simple random walk is refined sufficiently. Also, we have not considered the borderline value $`\alpha =2`$, which is special , and for which this rule fails.
## Acknowledgments
The authors acknowledge support from the French-Brazilian scientific cooperation project CAPES/COFECUB 229/97.
## Appendix A Relations for the inverse sums $`𝖦_A`$ and $`g_A`$
We collect here some elementary matrix algebra relations useful for dealing with the inverse sums $`𝖦_A`$ and $`g_A`$ occurring in the main text. The $`z`$ dependence of these quantities plays no role. The presentation and notation are independent of the body of the paper.
Let $`L`$ be an invertible $`\mathrm{}\times \mathrm{}`$ matrix. We define the ”inverse sum” $`(L)`$ by
$$^1(L)=\underset{i,j}{}L_{ij}^1$$
(A.1)
In the remainder $`\alpha ,\beta ,`$ and $`\gamma `$ will denote constants.
Property 1. Let $`J`$ be the $`\mathrm{}\times \mathrm{}`$ matrix with all $`J_{ij}=1`$, and let $`M`$ be an invertible $`\mathrm{}\times \mathrm{}`$ matrix. Let $`L=\alpha J+\gamma M`$. Then
$$(L)=\alpha +\gamma (M)$$
(A.2)
The proof of this relation is given in Ref. .
Property 2. Let $`M`$ and $`N`$ be invertible matrices of dimensions $`m\times m`$ and $`n\times n`$, respectively, and let $`L`$ be the block diagonal $`\mathrm{}\times \mathrm{}`$ matrix with blocks $`M`$ and $`N`$. Then
$$\frac{1}{(L)}=\frac{1}{(M)}+\frac{1}{(N)}$$
(A.3)
This follows directly from the definition (A.1). The calculation of $`(L)`$ for an $`\mathrm{}\times \mathrm{}`$ matrix may be reduced to an inversion problem of dimension less than $`\mathrm{}`$ also in certain cases where $`L`$ is not block diagonal, as shown below.
Property 3. Let $`J^{mn}`$ be the $`m\times n`$ matrix with all elements equal to 1. Let $`L`$ be $`\mathrm{}\times \mathrm{}`$ and of the form
$$L=\left(\begin{array}{cc}\gamma M& \beta J^{mn}\\ \beta J^{nm}& \gamma N\end{array}\right)$$
(A.4)
Then
$$(L)=\frac{\gamma ^2(M)(N)\beta ^2}{\gamma (M)+\gamma (N)2\beta }$$
(A.5)
To prove this we rewrite $`L`$ as $`L=\beta J+\stackrel{~}{L}`$, where $`J`$ is as before and where
$$\stackrel{~}{L}=\left(\begin{array}{cc}\gamma M\beta J^{mm}& 0\\ 0& \gamma N\beta J^{nn}\end{array}\right)$$
(A.6)
From Property 1 we have that $`(L)=\beta +(\stackrel{~}{L})`$, after which by applying Property 2 and once more Property 1, we obtain after some rearrangement Eq. (A.5). For $`\beta =0`$ Eq. (A.5) reduces to Property 2.
In this work the need for Properties 1 and 3 arises when the limit $`\gamma 0`$ has to be taken. For $`\gamma =0`$ the matrices $`L`$ that occur on the LHS of Eqs. (A.2) and (A.5) are no longer invertible, but these properties allow nevertheless $`(L)`$ to be calculated in that limit.
## Appendix B Corrections to scaling
In Eq. (3.6) we have neglected the $`𝒪`$(V<sup>2</sup>) terms that appear in Eq. (2.15). Since in the last step that led to Eq. (3.13) the leading order in $`1z`$ went down due to cancellations, we must now check that the $`𝒪`$(V<sup>2</sup>) terms remain subdominant. In this Appendix we will write Eq. (2.13) in the simplified notation $`\text{G}_{A(r+B)}=\text{G}+\text{W}`$ where G and W are the first and second matrix, respectively, on the RHS of Eq. (2.13). Upon writing the inverse G$`{}_{}{}^{1}{}_{A(r+B)}{}^{}`$ as a perturbation series in W and applying Eq. (2.4) one finds
$$𝖦_{A(r+B)}^1=\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}(1)^{\mathrm{}}\underset{c,c^{}}{}[\text{G}^1(\text{WG}^1)^{\mathrm{}}]_{cc^{}}$$
(B.1)
The $`\mathrm{}=0`$ term of this series is the term shown explicitly on the RHS of Eq. (2.15), and has been the object of study in Sec. 3.2. We will show here that the terms with $`\mathrm{}1`$ produce, in the scaling limit, only higher order corrections to the final result. To this end we first consider $`C_{MM^{}}(z)`$ defined by Eq. (2.6). Let $`R_{\mathrm{}}(z)`$ denote the contribution to $`C_{MM^{}}(z)`$ from the $`\mathrm{}`$th term in Eq. (B.1). In order to estimate the order in $`\mathrm{\hspace{0.17em}1}z`$ of $`R_{\mathrm{}}(z)`$ as $`z1`$ we first deduce from Eqs. (2.14) and (2.25) that in the scaling limit the matrix elements of W behave as $`V_{a,r+b}(1z)^{\frac{2}{\alpha }1}(ba)F^{}(\xi )`$, and that summing on $`r`$ amounts to applying $`(1z)^{\frac{1}{\alpha }}\text{d}\xi `$. This yields the asymptotic proportionality
$$R_{\mathrm{}}(z)(1z)^{\frac{1}{\alpha }}\text{d}\xi \frac{1}{G_0(z)}\left[\frac{(1z)^{\frac{2}{\alpha }1}F^{}(\xi )}{G_0(z)}\right]^{\mathrm{}}$$
(B.2)
where $`G_0(z)`$ represents the order in $`\mathrm{\hspace{0.17em}1}z`$ of the matrix G. When $`\mathrm{}`$ is odd, this integral vanishes by symmetry, which shows that the leading correction is of order V<sup>2</sup>, as anticipated. For $`\xi 0`$ we have, in virtue of Eq. (2.27), that $`F^{}(\xi )\xi ^{2+\alpha }`$. Hence the $`\xi `$ integral in Eq. (B.2) diverges in the origin for all $`\mathrm{}1`$ when $`0<\alpha <1`$. This signals that the main contribution comes from $`r`$ values near the origin. The order in $`\mathrm{\hspace{0.17em}1}z`$ of $`R_{\mathrm{}}(z)`$ may then be estimated by introducing in the integral the cutoff $`|\xi |\text{cst}\times (1z)^{\frac{1}{\alpha }}`$, which leads to $`R_{\mathrm{}}(z)(1z)^0`$. When $`\mathrm{}1`$ these additive corrections to $`C_{MM^{}}(z)`$ in Eq. (3.11) are negligible, therefore, with respect to the $`(1z)^{\frac{2}{\alpha }3}`$ term which, in the relevant exponent regime $`\frac{2}{3}<\alpha <1`$, determines the final result. |
warning/0003/hep-lat0003002.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Lattice QCD simulations are usually performed with periodic boundary conditions (BC). However, other type of BC may be important in certain cases. For example, a comparison between systems with different BC can be used to understand finite size-effects in lattice QCD. Some years ago<sup>?</sup> $`C`$-periodic BC were studied as an alternative to periodic conditions. Then they were considered with the general idea of studying the spontaneous symmetry breaking aspects of the QCD dynamics in a simple way.<sup>?</sup> In that work an analysis was done in the continuum, and it was shown that in pure gauge theory these BC break the $`Z(3)`$ symmetry explicitely, which has important consequences for the high-temperature deconfinement phase transition. These conditions are also useful in numerical lattice simulations of this transition. When quarks are present, $`C`$-periodic BC break both chiral and flavour symmetries.
These boundary conditions are also especially important when topological properties are relevant in the system under consideration. This is the case of the lattice studies of confinement through monopole condensation. Recently the role of monopoles in connection with colour confinement has been evidentiated in SU(2) and SU(3) gluodynamics,<sup>?</sup> for which a disorder parameter based on the magnetic U(1) symmetry has been constructed and studied by Monte Carlo techniques. The disorder parameter is the vacuum expectation value (vev) of a disorder operator, which is an operator that creates a magnetic monopole in the gauge configuration. The definition of this disorder operator requires $`C`$-periodic BC in the time direction. For the pure gauge case, this means that the links at time $`t+N_t`$, where $`N_t`$ is the temporal extension of the lattice, are the complex conjugate of the links at time $`t`$. The effect on the simulation algorithm is a simple redefinition of the staples containing links that pierce the temporal boundary. The natural extension of the procedure used for the pure gauge case to full QCD requires the implementation of $`C`$-periodic BC in the presence of fermions. In particular, we will be concerned with the case of staggered fermions. $`C`$-periodic BC modify the fermionic matrix, and many proofs of properties used for the setup of standard simulation algorithms no longer hold.
$`C`$-periodic BC in the continuum are defined by the action of the charge conjugation operator $`C`$ on the fields.<sup>?</sup> However, lattice fermions are different from fermions in the continuum. In particular, $`C`$ is not a symmetry of the lattice action with staggered fermions. It also breaks translation invariance for finite lattice spacing. However, there is a discrete symmetry of the staggered fermion action,
$$S_f=\underset{i,\mu }{}\left[\frac{1}{2}\eta _{i,\mu }\left(\overline{\psi }_iU_{i,\mu }\psi _{i+\mu }\overline{\psi }_{i+\mu }U_{i,\mu }^{}\psi _i\right)+m\overline{\psi }_i\psi _i\right]$$
(1.1)
(here $`i`$ indicates the lattice point, $`U_{i,\mu }`$ is the SU(3) matrix associated with the link leaving the $`i`$-th lattice point in the $`\mu `$ direction, $`\psi _i`$ is the staggered fermion field at the point $`i`$, and $`\eta _{i,\mu }`$ the usual staggered fermion phase), which corresponds to charge conjugation in the gluon sector but which in the continuum limit also contains a flavour transformation.<sup>?,?</sup> We will call $`C^{}`$ this symmetry of the lattice action:
$${}_{}{}^{C^{}}U_{i,\mu }^{}=U_{i,\mu }^{},^C^{}\psi _i=ϵ_i\overline{\psi }_i^\mathrm{T},^C^{}\overline{\psi }_i=\psi _i^\mathrm{T}ϵ_i,$$
(1.2)
where $`ϵ_i=(1)^{x_i+y_i+z_i+t_i}`$, $`(x_i,y_i,z_i,t_i)`$ being the lattice coordinates of point $`i`$, the “T” represents the traspose operation, and we will use $`\psi ^{}\overline{\psi }^\mathrm{T}`$. Translation invariance implies that a BC must correspond to a symmetry of the action. $`C^{}`$-BC are defined as the boundary conditions corresponding to the symmetry (1.2):
$$\mathrm{\Phi }_{i+N}=^C^{}\mathrm{\Phi }_i,$$
(1.3)
where $`\mathrm{\Phi }`$ is a field, $`U_\mu `$ or $`\psi `$, and $`N`$ is the number of lattice points in the direction in which we use this boundary condition.
In the chiral limit the lattice action has a $`U(1)_\mathrm{E}U(1)_\mathrm{O}`$ chiral symmetry of independent rotations on $`(x+y+z+t)`$-even and -odd lattice points. $`C^{}`$-BC break explicitely this symmetry to $`U(1)_{\mathrm{E}=\mathrm{O}^{}}`$. Baryon number $`U(1)_{\mathrm{E}=\mathrm{O}}`$ is also broken explicitely down to $`Z(2)_{\mathrm{E}=\mathrm{O}}`$.<sup>?</sup>
In this paper we will have in mind the physical problem mentioned above: computation of the vev of a monopole creation operator in lattice QCD. This means that we will consider imposing $`C^{}`$-BC in the time direction, and periodic BC in the spatial directions. However this is just to fix the notation in what follows; the $`C^{}`$ conditions could in fact be assumed in any direction. The purpose of the paper is to show the theoretical framework to be used in a lattice simulation with these BC (section 2) and how usual algorithms<sup>?</sup> need to be modified (section 3). Our conclusions are summarized in section 4.
## 3 Mathematical description
Let us consider the partition function of lattice QCD with staggered fermions
$$Z=(𝒟U)(𝒟\psi 𝒟\overline{\psi })e^{S_gS_f},$$
(3.1)
where $`S_g(U)`$ is the Wilson action for the pure gauge sector, and $`S_f`$ is given by Eq. (1.1). The fermionic variables can be integrated out to give
$$Z=(𝒟U)e^{S_g(U)}detM(U),$$
(3.2)
where<sup>*</sup><sup>*</sup>*We have absorbed the staggered phases by a redefinition of the link matrices $`U`$.
$$M(U)_{i,j}=m\delta _{i,j}+\underset{\mu }{}\frac{1}{2}(U_{i,\mu }\delta _{i,j\mu }U_{i\mu ,\mu }^{}\delta _{i,j+\mu })$$
(3.3)
is the fermionic matrix, and periodic BC are assumed. Using this matrix notation, we can write
$$S_f=\overline{\psi }M\psi =\frac{1}{2}\left[\left(\overline{\psi }M\psi \right)+\left(\overline{\psi }M\psi \right)^\mathrm{T}\right],$$
(3.4)
since $`S_f`$ is a number. We can use the new variable $`\mathrm{\Psi }`$ defined as the column vector formed by $`\psi `$ and $`\psi ^{}`$, so that the fermionic integral is
$$(𝒟\psi )(𝒟\overline{\psi })e^{S_f}=(𝒟\mathrm{\Psi })e^{\frac{1}{2}\mathrm{\Psi }^\mathrm{T}A\mathrm{\Psi }}=\mathrm{Pf}(A),$$
(3.5)
$$\mathrm{\Psi }^\mathrm{T}A\mathrm{\Psi }\left(\begin{array}{cc}\psi ^\mathrm{T}& \overline{\psi }\end{array}\right)A\left(\begin{array}{c}\psi \hfill \\ \psi ^{}\hfill \end{array}\right),$$
(3.6)
with
$$A=\left(\begin{array}{cc}0& M^\mathrm{T}\\ M& 0\end{array}\right),$$
(3.7)
and we have introduced the Pfaffian of the matrix $`A`$, $`\mathrm{Pf}(A)`$. It is well known that<sup>?</sup>
$$\mathrm{Pf}^2(A)=detA.$$
(3.8)
Now let us consider $`C^{}`$-BC in the time direction. Then, following Eqs. (1.2) and (1.3), $`\psi _{N_t}=ϵ_0\psi _0^{}`$, and $`\overline{\psi }_{N_t}=ϵ_0\psi _0^\mathrm{T}`$, where we have written in the subscript the temporal coordinate and omitted the spatial coordinates. Now Eq. (3.4) gives, for the terms connecting the slices $`N_t1`$ and $`N_t`$,
$`{\displaystyle \frac{1}{2}}\left(\overline{\psi }_{N_t1}U_{N_t1,t}\psi _0^{}+\psi _0^\mathrm{T}U_{N_t1,t}^{}\psi _{N_t1}\right)ϵ_0`$ (3.9)
$``$ $`{\displaystyle \frac{1}{2}}\left(\overline{\psi }_0U_{N_t1,t}^\mathrm{T}\psi _{N_t1}^{}+\psi _{N_t1}^\mathrm{T}U_{N_t1,t}^{}\psi _0\right)ϵ_0.`$
In this way the matrix $`A`$ of Eq. (3.7) is substituted by
$$A=\left(\begin{array}{cc}B& \stackrel{~}{M}^\mathrm{T}\\ \stackrel{~}{M}& B^{}\end{array}\right),$$
(3.10)
where $`B`$ satisfies the properties:
$$B^{}=B^{}B=B^\mathrm{T}$$
(3.11)
and $`\stackrel{~}{M}`$ is the fermionic matrix Eq. (3.3) apart from the terms connecting the slices $`N_t1`$ and $`N_t`$, which have gone to the matrices $`B`$ and $`B^{}`$.
Eq. (3.5) is still valid and we are interested in calculating $`\mathrm{Pf}(A)`$.
### 3.1 Pseudofermionic variables
The usual approach<sup>?</sup> to the simulation of theories with dynamical fermions is to rewrite the determinant of the fermionic matrix in Eq. (3.2) using that
$$det(M^{}M)𝒟\varphi ^{}𝒟\varphi \mathrm{exp}[\varphi ^{}(M^{}M)^1\varphi ],$$
(3.12)
where $`\varphi `$ is a complex bosonic field with the same quantum numbers as the Grassmann field. One introduces the matrix $`(M^{}M)^1`$, instead of $`M^1`$, so that the pseudofermionic fields can be generated using a simple heatbath method. Since $`detM`$ is a real number, $`det(M^{}M)=(detM)^2`$. So, actually, this corresponds to a double number of flavours with respect to the original theory. However, the matrix $`M^{}M`$ has two important properties: it has no matrix elements connecting even and odd lattice sites, and the determinants of its submatrices on the even and odd sites are equal.<sup>?</sup> Therefore one can avoid the redoubling of flavours by defining the pseudofermionic field only on even lattice sites.
A remarkable difference between the partition function with periodic BC and the partition function with $`C^{}`$-BC is that in the latter case the determinant of $`M`$, Eq. (3.2), has to be replaced by the Pfaffian of $`A`$, i.e. $`\pm \sqrt{detA}`$. Because of the square root, the usual trick of introducing pseudofermionic fields and rewriting this factor as the integral over these fields can only be applied if the number of continuum fermion flavours is such that the square root cancels. Moreover, in order to have a positive-definite integration measure, we need that the sign in front of this factor be $`+`$. Both conditions are satisfied if the number of continuum flavours is a multiple of eight. This generates a further unavoidable redoubling. Until Sect. 3.2, we will be concerned with the numerical simulation of $`detA`$, which is equivalent to simulating a double number of fermion flavours in the continuum limit (eight instead of four). In Sect. 3.2 we will discuss how to deal with the usual case of four staggered fermion flavours.
Since the case of a system with eight fermion flavours in the continuum limit and $`C^{}`$-BC and the case of a system with four fermion flavours in the continuum limit and periodic BC are similar, we would like to follow in the former case the standard procedure to obtain the determinant of the matrix (3.10). First, in the Appendix A it is shown that $`detA`$ is a real number. So we can also in this case use the matrix $`A^{}A`$ to introduce the pseudofermionic field, which will have now twice the number of components as in the usual case. Second, we will now see that the matrix $`A^{}A`$ does not connect even and odd lattice sites.
Using the form of the fermionic matrix (3.3) and the definitions of $`B`$ and $`\stackrel{~}{M}`$ given by Eqs. (3.9) and (3.10), we can split the blocks of the matrix $`A`$ in even and odd lattice sites:
$$A=\left(\begin{array}{cccc}0& \frac{1}{2}B_{\mathrm{eo}}& m& \frac{1}{2}D_{\mathrm{oe}}^\mathrm{T}\\ \frac{1}{2}B_{\mathrm{oe}}& 0& \frac{1}{2}D_{\mathrm{eo}}^\mathrm{T}& m\\ m& \frac{1}{2}D_{\mathrm{eo}}& 0& \frac{1}{2}B_{\mathrm{eo}}^{}\\ \frac{1}{2}D_{\mathrm{oe}}& m& \frac{1}{2}B_{\mathrm{oe}}^{}& 0\end{array}\right).$$
(3.13)
From Eqs. (3.3) and (3.11) it is easy to see that
$`D_{\mathrm{oe}}^{}=D_{\mathrm{eo}}`$ $`D_{\mathrm{oe}}^\mathrm{T}=D_{\mathrm{eo}}^{}`$ (3.14)
$`B_{\mathrm{oe}}^{}=B_{\mathrm{eo}}^{}`$ $`B_{\mathrm{oe}}^\mathrm{T}=B_{\mathrm{eo}}.`$ (3.15)
In the expression (3.13), the new blocks divide the matrix in rows and columns identified by the pairs $`[0,\mathrm{e}]`$, $`[0,\mathrm{o}]`$, $`[1,\mathrm{e}]`$ and $`[1,\mathrm{o}]`$, where the number indicates the blocks defined in Eq. (3.10) and the letter the even/odd subblock. Reorganizing the rows and columns to $`[0,\mathrm{e}]`$, $`[1,\mathrm{e}]`$, $`[0,\mathrm{o}]`$ and $`[1,\mathrm{o}]`$ (the determinant does not change), we rewrite the matrix A in the following form:
$$A=\left(\begin{array}{cc}\mu & \frac{1}{2}A_{\mathrm{eo}}\\ \frac{1}{2}A_{\mathrm{oe}}& \mu \end{array}\right),$$
(3.16)
where
$$\mu =\left(\begin{array}{cc}0& m\\ m& 0\end{array}\right),$$
(3.17)
$$A_{\mathrm{eo}}=\left(\begin{array}{cc}B_{\mathrm{eo}}& D_{\mathrm{eo}}^{}\\ D_{\mathrm{eo}}& B_{\mathrm{eo}}^{}\end{array}\right),$$
(3.18)
and $`A_{\mathrm{oe}}`$ has the same form as $`A_{\mathrm{eo}}`$ with the exchange $`\mathrm{e}\mathrm{o}`$. Eqs. (3.14) and (3.15) give
$$A_{\mathrm{oe}}^{}=A_{\mathrm{eo}}^{}.$$
(3.19)
From Eqs. (3.3) and (3.9) we obtain the explicit form of the matrices $`D_{\mathrm{eo}}`$ and $`B_{\mathrm{eo}}`$:
$$(D_{\mathrm{eo}})_{j,k}=\underset{\mu }{}(U_{j,\mu }\delta _{k,j+\mu }U_{j\mu ,\mu }^{}\delta _{k,j\mu }),$$
(3.20)
where $`j`$ is an even site and the link between $`j`$ and $`k`$ does not connect the $`t=N_t1`$ and $`t=0`$ time slices (otherwise $`(D_{\mathrm{eo}})_{j,k}=0`$), and
$$\begin{array}{cc}(B_{\mathrm{eo}})_{j,k}=U_{j,t}^{}\hfill & \hfill j=N_t1,k=0\\ (B_{\mathrm{eo}})_{j,k}=U_{k,t}^{}\hfill & \hfill j=0,k=N_t1\\ (B_{\mathrm{eo}})_{j,k}=0\hfill & \hfill \text{in every other case}\end{array}\}$$
(3.21)
where $`j=N_t1`$ or $`j=0`$ means that the temporal coordinate of site $`j`$ is $`N_t1`$ or zero, respectively. When $`j`$ is an odd site, $`(D_{\mathrm{oe}})_{j,k}`$ has the same expression as Eq. (3.20), and
$$\begin{array}{cc}(B_{\mathrm{oe}})_{j,k}=U_{j,t}^{}\hfill & \hfill j=N_t1,k=0\\ (B_{\mathrm{oe}})_{j,k}=U_{k,t}^{}\hfill & \hfill j=0,k=N_t1\\ (B_{\mathrm{oe}})_{j,k}=0\hfill & \hfill \text{in every other case}\end{array}\}$$
(3.22)
These expressions completely determine every element of the matrices $`A_{\mathrm{eo}}`$ and $`A_{\mathrm{oe}}`$.
We now compute
$$A^{}A=\left(\begin{array}{cc}\mu ^2+\frac{1}{4}A_{\mathrm{oe}}^{}A_{\mathrm{oe}}& \frac{1}{2}(\mu A_{\mathrm{eo}}A_{\mathrm{oe}}^{}\mu )\\ \frac{1}{2}(A_{\mathrm{eo}}^{}\mu \mu A_{\mathrm{oe}})& \mu ^2+\frac{1}{4}A_{\mathrm{eo}}^{}A_{\mathrm{eo}}\end{array}\right).$$
(3.23)
But using the properties (3.14) and (3.15) it is direct to see that the blocks out of the diagonal in Eq. (3.23) are zero. Then, using Eq. (3.19),
$$A^{}A=\left(\begin{array}{cc}\mu ^2\frac{1}{4}A_{\mathrm{eo}}^{}A_{\mathrm{oe}}& 0\\ 0& \mu ^2\frac{1}{4}A_{\mathrm{oe}}^{}A_{\mathrm{eo}}\end{array}\right)$$
(3.24)
and, as a result, we see that the matrix $`A^{}A`$ connects only lattice points of the same parity.
In order to use the same trick to avoid the flavour doubling produced by the introduction of $`A^{}A`$, that is, to define the pseudofermionic field on even sites only, we need to show that, also in this case, the determinant of the even and odd parts in the matrix (3.24) are equal. This is done in the following subsection.
### 3.2 Reducing the number of flavours: even-odd partitioning
Let us write $`KA^{}A`$. Then we have
$$K=\left(\begin{array}{cc}m^2I_2\frac{1}{4}A_{\mathrm{eo}}^{}A_{\mathrm{oe}}& 0\\ 0& m^2I_2\frac{1}{4}A_{\mathrm{oe}}^{}A_{\mathrm{eo}}\end{array}\right)\left(\begin{array}{cc}K_\mathrm{e}& 0\\ 0& K_\mathrm{o}\end{array}\right),$$
(3.25)
where $`I_2`$ is the $`2\times 2`$ identity matrix. We need to show that $`detK_\mathrm{e}=detK_\mathrm{o}`$. The strategy will be to show that both $`detK_\mathrm{e}`$ and $`detK_\mathrm{o}`$ are equal to $`detA`$. To this aim, let us consider the matrix $`A`$ written in the e–o partitioned form, Eq. (3.16),
$$A=\left(\begin{array}{cc}\mu & \frac{1}{2}A_{\mathrm{eo}}\\ \frac{1}{2}A_{\mathrm{oe}}& \mu \end{array}\right),$$
(3.26)
where
$$\mu =\left(\begin{array}{cc}0& m\\ m& 0\end{array}\right),$$
(3.27)
$$A_{\mathrm{eo}}=\left(\begin{array}{cc}B_{\mathrm{eo}}& D_{\mathrm{eo}}^{}\\ D_{\mathrm{eo}}& B_{\mathrm{eo}}^{}\end{array}\right),$$
(3.28)
and
$$A_{\mathrm{oe}}=\left(\begin{array}{cc}B_{\mathrm{oe}}& D_{\mathrm{oe}}^{}\\ D_{\mathrm{oe}}& B_{\mathrm{oe}}^{}\end{array}\right).$$
(3.29)
We will make use of a general property of the determinant of a square $`2N\times 2N`$ block matrix,
$$det\left(\begin{array}{cc}X& Y\\ W& Z\end{array}\right)=det(X)det(ZWX^1Y),$$
(3.30)
where $`X,Y,W,Z`$ are $`N\times N`$ square matrices and $`X`$ is invertible. Applying this equality to the computation of $`detA`$, we obtain
$$detA=det(\mu )det\left(\mu \frac{1}{4}A_{\mathrm{oe}}\mu ^1A_{\mathrm{eo}}\right)=det\left(\mu ^2\frac{1}{4}\mu A_{\mathrm{oe}}\mu ^1A_{\mathrm{eo}}\right).$$
(3.31)
Using the explicit form of $`\mu `$ and $`A_{\mathrm{oe}}`$ it is easily shown that
$$\mu A_{\mathrm{oe}}=A_{\mathrm{oe}}^{}\mu .$$
(3.32)
Inserting this equality in Eq. (3.31), it follows that
$$detA=det\left(\mu ^2+\frac{1}{4}A_{\mathrm{oe}}^{}A_{\mathrm{eo}}\right)=detK_\mathrm{o},$$
(3.33)
where the last equality is implied by the fact that $`N`$ is an even number.
We can compute again $`detA`$ using the same procedure after exchanging even and odd variables in $`A`$. Using a property analogous to that in Eq. (3.32), namely $`\mu A_{\mathrm{eo}}=A_{\mathrm{eo}}^{}\mu `$, it is then easily shown that
$$detA=det\left(\mu ^2+\frac{1}{4}A_{\mathrm{eo}}^{}A_{\mathrm{oe}}\right)=detK_\mathrm{e}.$$
(3.34)
This is a proof that $`detK_\mathrm{e}=detK_\mathrm{o}`$.
We have shown that $`detK_\mathrm{e}detK_\mathrm{o}=(detA)^2`$. On the other hand it is also true that $`detK_\mathrm{e}detK_\mathrm{o}=det(A^{}A)=detA^{}detA`$, therefore we have also obtained a proof that $`detA=detA^{}`$ alternative to that given in Appendix A. Moreover, we have that
$$detA=detK_\mathrm{e}=det\left(m^2I_2+\frac{1}{4}A_{\mathrm{oe}}^{}A_{\mathrm{oe}}\right)>0,$$
(3.35)
since both $`m^2I_2`$ and $`A_{\mathrm{oe}}^{}A_{\mathrm{oe}}`$ are positive definite matrices. Therefore $`detA`$ is a positive number.
We notice that this method can be easily applied also to the standard case, thus providing a proof alternative to those presented in Ref. 8. This is shown in detail in Appendix B.
## 4 Hybrid Monte Carlo Implementation
We will show now how the standard Hybrid Monte Carlo (HMC) algorithm<sup>?</sup> needs to be modified to incorporate $`C^{}`$-BC. In this algorithm one introduces fictitious momenta, conjugate variables of the links, as dynamical variables, and makes fields evolve with a mixed dynamics, in which deterministic and stochastic steps are alternated in a prescribed way. In the deterministic part of the algorithm, the system follows the equations of motion derived from the Hamiltonian of the (4+1)-dimensional system. These equations give the evolution of the fields in the fictitious time $`\tau `$. The equation of motion for the matrix $`U_{j,\mu }`$ is
$$\dot{U}_{j,\mu }=iH_{j,\mu }U_{j,\mu },$$
(4.1)
where the conjugate momentum $`H_{j,\mu }`$ is a traceless Hermitian matrix, and $`\dot{U}`$ is the derivative of $`U`$ with respect to $`\tau `$. The equations of motion for the momenta $`H`$ are obtained by imposing that the Hamiltonian be constant. The integration of these equations of motion is carried out numerically after discretization of $`\tau `$. Usually the temporal step is of order $`10^2`$$`10^3`$ and the configuration space is sampled with a total length of the trajectory of order 1. The stochastic part of the algorithm consists in the generation of new momenta and new pseudofermionic variables according to their probability distributions at the beginning of each trajectory. Moreover, at the end of each trajectory a Metropolis accept-reject step is performed, which makes the algorithm exact.
### 4.1 HMC algorithm with $`C^{}`$ boundary conditions
Once introduced the pseudofermionic fields, defined only on even sites, and the auxiliary momenta fields, the partition function of the system is
$$Z=(𝒟U𝒟\varphi ^{}𝒟\varphi 𝒟H)e^{},$$
(4.2)
with
$$=\frac{1}{2}\underset{j,\mu }{}\mathrm{tr}H_{j,\mu }^2+S_g+\varphi ^{}(A^{}A)^1\varphi .$$
(4.3)
To obtain an equation of motion for $`H`$ we require that $``$ be a constant of motion, that is, $`\dot{}=0`$. The tricky part in this differentiation is in the fermionic term. The derivative of $`_f=\varphi ^{}(A^{}A)^1\varphi `$ is
$`\dot{}_f`$ $`=`$ $`{\displaystyle \underset{j,\mu }{}}\varphi ^{}(A^{}A)^1i[{\displaystyle \frac{A^{}}{U_{j,\mu }}}H_{j,\mu }U_{j,\mu }A+{\displaystyle \frac{A^{}}{U_{j,\mu }^\mathrm{T}}}U_{j,\mu }^\mathrm{T}H_{j,\mu }^\mathrm{T}A+`$ (4.4)
$`A^{}{\displaystyle \frac{A}{U_{j,\mu }}}H_{j,\mu }U_{j,\mu }+A^{}{\displaystyle \frac{A}{U_{j,\mu }^\mathrm{T}}}U_{j,\mu }^\mathrm{T}H_{j,\mu }^\mathrm{T}H_{j,\mu }^\mathrm{T}U_{j,\mu }^{}{\displaystyle \frac{A^{}}{U_{j,\mu }^{}}}A`$
$`U_{j,\mu }^{}H_{j,\mu }{\displaystyle \frac{A^{}}{U_{j,\mu }^{}}}AA^{}H_{j,\mu }^\mathrm{T}U_{j,\mu }^{}{\displaystyle \frac{A}{U_{j,\mu }^{}}}A^{}U_{j,\mu }^{}H_{j,\mu }{\displaystyle \frac{A}{U_{j,\mu }^{}}}]`$
$`\times (A^{}A)^1\varphi ,`$
where we have made use of Eq. (4.1) and the fact that $`H`$ is Hermitian. We have also used that $`A`$ is linear in $`U`$, $`U^\mathrm{T}`$, $`U^{}`$ and $`U^{}`$, and then the partial derivatives in the previous expression commute with $`H`$ and $`U`$.
It is convenient to introduce the operator
$$P_{ij}=X_iX_j^{},$$
(4.5)
where
$$X=(A^{}A)^1\varphi .$$
(4.6)
It is easily seen that $`P`$ is Hermitian:
$$(P^{})_{ij}=(P_{ji})^{}=X_iX_j^{}=P_{ij},$$
(4.7)
and, being $`\varphi `$ defined only on even sites, $`P_{ij}`$ is taken to be zero unless $`i`$ and $`j`$ are both even sites. On the other hand, we have that
$`\left({\displaystyle \frac{A^{}}{U}}\right)^{}={\displaystyle \frac{A}{U^{}}}`$ $`\left({\displaystyle \frac{A}{U}}\right)^{}={\displaystyle \frac{A^{}}{U^{}}}`$
$`\left({\displaystyle \frac{A^{}}{U^\mathrm{T}}}\right)^{}={\displaystyle \frac{A}{U^{}}}`$ $`\left({\displaystyle \frac{A}{U^\mathrm{T}}}\right)^{}={\displaystyle \frac{A^{}}{U^{}}}`$ (4.8)
and then we can write
$`\dot{}_f`$ $`=`$ $`{\displaystyle \underset{j,\mu }{}}\mathrm{tr}[iH_{j,\mu }(U_{j,\mu }AP{\displaystyle \frac{A^{}}{U_{j,\mu }}}+U_{j,\mu }PA^{}{\displaystyle \frac{A}{U_{j,\mu }}})+`$ (4.9)
$`iH_{j,\mu }^\mathrm{T}(AP{\displaystyle \frac{A^{}}{U_{j,\mu }^\mathrm{T}}}U_{j,\mu }^\mathrm{T}+PA^{}{\displaystyle \frac{A}{U_{j,\mu }^\mathrm{T}}}U_{j,\mu }^\mathrm{T})+\mathrm{H}.\mathrm{c}.],`$
where H.c. means the Hermitian conjugate and we have used the cyclic property of the trace operation.
The trace in Eq. (4.9) is taken over both color and site indices. Taking this into account and using the cyclic property again, we have
$`\dot{}_f`$ $`=`$ $`{\displaystyle \underset{j,\mu }{}}\mathrm{tr}(H_{j,\mu }\{i[U_{j,\mu }\left(AP{\displaystyle \frac{A^{}}{U_{j,\mu }}}\right)+U_{j,\mu }\left(PA^{}{\displaystyle \frac{A}{U_{j,\mu }}}\right)+`$ (4.10)
$`U_{j,\mu }\left(AP{\displaystyle \frac{A^{}}{U_{j,\mu }^\mathrm{T}}}\right)^\mathrm{T}+U_{j,\mu }\left(PA^{}{\displaystyle \frac{A}{U_{j,\mu }^\mathrm{T}}}\right)^\mathrm{T}]+\mathrm{H}.\mathrm{c}.\}).`$
The following step is to calculate the derivatives appearing in Eq. (4.10). From Eqs. (3.16) and (3.19), we have that
$$A^{}=A^{}.$$
(4.11)
We write again $`A`$ and $`A^{}`$ in terms of subblocks, identifying the rows and columns with the pairs $`[0,\mathrm{e}]`$, $`[1,\mathrm{e}]`$, $`[0,\mathrm{o}]`$ and $`[1,\mathrm{o}]`$:
$$A=\left(\begin{array}{cccc}0& m& \frac{1}{2}B_{\mathrm{eo}}& \frac{1}{2}D_{\mathrm{eo}}^{}\\ m& 0& \frac{1}{2}D_{\mathrm{eo}}& \frac{1}{2}B_{\mathrm{eo}}^{}\\ \frac{1}{2}B_{\mathrm{oe}}& \frac{1}{2}D_{\mathrm{oe}}^{}& 0& m\\ \frac{1}{2}D_{\mathrm{oe}}& \frac{1}{2}B_{\mathrm{oe}}^{}& m& 0\end{array}\right),$$
(4.12)
$$A^{}=\left(\begin{array}{cccc}0& m& \frac{1}{2}B_{\mathrm{eo}}^{}& \frac{1}{2}D_{\mathrm{eo}}\\ m& 0& \frac{1}{2}D_{\mathrm{eo}}^{}& \frac{1}{2}B_{\mathrm{eo}}\\ \frac{1}{2}B_{\mathrm{oe}}^{}& \frac{1}{2}D_{\mathrm{oe}}& 0& m\\ \frac{1}{2}D_{\mathrm{oe}}^{}& \frac{1}{2}B_{\mathrm{oe}}& m& 0\end{array}\right).$$
(4.13)
Once written everything in terms of $`B`$ and $`D`$ we can easily see where $`U`$ and $`U^\mathrm{T}`$ appear, from Eqs. (3.20)–(3.22), and compute the derivatives. The result is given in Table 4.1.
Now it is easy to calculate the different terms in Eq. (4.10). Let us do, as an example, the first one. The trace over lattice site indices affects only the expressions between parentheses. From Table 4.1, we see that, for $`j`$ even, and $`jN_t1`$ or $`\mu t`$, the first one gives
$`\mathrm{tr}\left(AP{\displaystyle \frac{A^{}}{U_{j,\mu }}}\right)`$ $`=`$ $`A_{k_2,k}X_kX_{k_1}^{}\left({\displaystyle \frac{A^{}}{U_{j,\mu }}}\right)_{k_1,k_2}={\displaystyle \frac{1}{2}}A_{[1,\mathrm{o}]j+\mu ,k}X_kX_{[0,\mathrm{e}]j}^{}=`$ (4.14)
$`\left({\displaystyle \frac{1}{4}}(D_{\mathrm{oe}})_{j+\mu ,k}X_k^0+{\displaystyle \frac{1}{4}}(B_{\mathrm{oe}}^{})_{j+\mu ,k}X_k^1\right)(X^{})_j^0=`$
$`{\displaystyle \frac{1}{4}}((A_{\mathrm{oe}})_{j+\mu ,k}X_k)^1(X^{})_j^0={\displaystyle \frac{1}{4}}(A_{\mathrm{oe}}X)_{j+\mu }^1(X^{})_j^0,`$
where we have used Eqs. (4.12) and (3.18), and called $`X_k^0`$ the first three components of the vector $`X_k`$, and $`X_k^1`$ the second three components. Notice that there is not a mass term because it is diagonal, which would imply that $`k=j+\mu `$, that is, $`k`$ would be odd, and then $`X_k=0`$. For $`j`$ even and $`j=N_t1`$, $`\mu =t`$, an analogous calculation gives
$$\mathrm{tr}\left(AP\frac{A^{}}{U_{j,\mu }}\right)=\frac{1}{4}(A_{\mathrm{oe}}X)_{j+\mu }^0(X^{})_j^0.$$
(4.15)
This term does not contribute when $`j`$ is odd, because in this case $`(X^{})_j^0`$ is zero ($`X`$ is defined on even sites only).
Computing every other contribution in Eq. (4.10), the final result can be written
$$\dot{}_f=\underset{j,\mu }{}\mathrm{tr}\left\{H_{j,\mu }\left(iF_{j,\mu }iF_{j,\mu }^{}\right)\right\},$$
(4.16)
where $`F_{j,\mu }`$ takes a different form depending on $`j`$ and $`\mu `$.
For $`j`$ even and $`(j,\mu )(N_t1,t)`$,
$$F_{j,\mu }=U_{j,\mu }\left(\frac{1}{4}(A_{\mathrm{oe}}X)_{j+\mu }^1(X^{})_j^0\left(\frac{1}{4}X_j^1((A_{\mathrm{oe}}X)^{})_{j+\mu }^0\right)^\mathrm{T}\right);$$
(4.17)
for $`j`$ even and $`(j,\mu )=(N_t1,t)`$,
$$F_{j,\mu }=U_{j,\mu }\left(\frac{1}{4}(A_{\mathrm{oe}}X)_{j+\mu }^0(X^{})_j^0+\left(\frac{1}{4}X_j^1((A_{\mathrm{oe}}X)^{})_{j+\mu }^1\right)^\mathrm{T}\right);$$
(4.18)
for $`j`$ odd and $`(j,\mu )(N_t1,t)`$,
$$F_{j,\mu }=U_{j,\mu }\left(\frac{1}{4}X_{j+\mu }^0((A_{\mathrm{oe}}X)^{})_j^1+\left(\frac{1}{4}(A_{\mathrm{oe}}X)_j^0(X^{})_{j+\mu }^1\right)^\mathrm{T}\right);$$
(4.19)
for $`j`$ odd and $`(j,\mu )=(N_t1,t)`$,
$$F_{j,\mu }=U_{j,\mu }\left(\frac{1}{4}X_{j+\mu }^1((A_{\mathrm{oe}}X)^{})_j^1\left(\frac{1}{4}(A_{\mathrm{oe}}X)_j^0(X^{})_{j+\mu }^0\right)^\mathrm{T}\right).$$
(4.20)
Since we determined completely the elements of the matrices $`A_{\mathrm{eo}}`$ and $`A_{\mathrm{oe}}`$ in section 2.1, these expressions allow us to obtain $`F_{j,\mu }`$ as a function of the links $`U`$ for every $`j`$ and $`\mu `$.
Coming back to Eq. (4.3), the condition $`\dot{}=0`$ gives
$$\dot{}=0=\underset{j,\mu }{}\mathrm{tr}(\dot{H}_{j,\mu }H_{j,\mu }+\frac{\beta }{6}(iH_{j,\mu }U_{j,\mu }V_{j,\mu }+\mathrm{H}.\mathrm{c}.)(iH_{j,\mu }F_{j,\mu }+\mathrm{H}.\mathrm{c}.)),$$
(4.21)
where $`V_{j,\mu }`$ is the sum of staples, or products of the other three matrices in the plaquettes containing $`U_{j,\mu }`$, and arises from the differentiation of the Wilson action $`S_g`$. Of course, the staples at the border of the lattice contain links defined by the boundary conditions ($`C^{}`$ or periodic, depending on the direction). The final solution for $`\dot{H}_{j,\mu }`$ is
$$i\dot{H}_{j,\mu }=\left[\frac{\beta }{3}U_{j,\mu }V_{j,\mu }2F_{j,\mu }\right]_{\mathrm{T}A},$$
(4.22)
where the subscript TA indicates the traceless anti-Hermitian part of the matrix:
$$Q_{\mathrm{T}A}=\frac{1}{2}(QQ^{})\frac{1}{6}\mathrm{tr}(QQ^{}).$$
(4.23)
### 4.2 Reducing the number of flavours: the Hybrid algorithm
Because of Eq. (3.8), the HMC algorithm that we have just described simulates eight fermion flavours in the continuum. In order to come back to four flavours, we note that
$$\mathrm{Pf}(A)=\pm (detA)^{1/2}$$
(4.24)
(we saw in Eq. (3.35) that $`detA`$ is a positive number).
One can simulate $`(detA)^{1/2}`$ by reverting to an approximate algorithm. A suitable choice is the $`R`$ algorithm,<sup>?</sup> in which discretization errors in the molecular dynamics part are of $`𝒪(\mathrm{\Delta }\tau ^2)`$.
On a lattice closed with $`C^{}`$-BC, the full QCD action in the presence of $`N_f`$ families of degenerate continuum fermions can be written as follows:
$$Z=(𝒟U)[det(A^{}A)_\mathrm{e}]^{N_f/8}e^{S_g(U)}=(𝒟U)e^{S_g(U)+\frac{N_f}{8}\mathrm{tr}\mathrm{log}(A^{}A)_\mathrm{e}},$$
(4.25)
where $`(A^{}A)_\mathrm{e}`$ is the restriction of $`A^{}A`$ to the even lattice sites.
Since the implementation of the $`R`$ algorithm in the present case reduces to simple adaptation of a standard technique to the system described by the equation of motion obtained in the previous subsection, we do not elaborate further on this point. However, we still have the “sign problem” of Eq. (4.24). The sign of the Pfaffian can be included by reweighting the expectation values according to
$$𝒪=\frac{𝒪\mathrm{signPf}(A)_+}{\mathrm{signPf}(A)_+},$$
(4.26)
where $`\mathrm{}_+`$ means that the expectation values have been obtained simulating $`+(detA)^{1/2}`$. One then needs to monitorize the sign of the Pfaffian depending on the gauge configurations. Some techniques to do that are explained in Ref. 9.
## 5 Conclusions
$`C^{}`$ boundary conditions are interesting to study some spontaneous symmetry breaking aspects of QCD. They are relevant when one analyses confinement through monopole condensation. We have shown in this work how these boundary conditions can be implemented to carry out a lattice simulation of full QCD with staggered fermions. We have proved that the common even-odd trick used to avoid the fermion redoubling produced by the introduction of the pseudofermionic field can be applied also to this case. However, there is an additional redoubling which forces to work with a minimum number of eight flavours with the usual Hybrid Monte Carlo algorithm, which we have adapted to this case. An alternative to avoid that is to consider a non-exact algorithm, which can be applied to any number of flavours.
These algorithms have been implemented and are presently running on an APE Quadrics machine to explore monopole condensation in full QCD.
Acknowledgments
This work has been partially supported by EU TMR program ERBFMRX-CT97-0122, Italian MURST and PPARC Grant PPA/G/0/1998/00567. We thank valuable comments from Isabel Campos, Simon Hands, Pilar Hernández and Giampiero Paffuti.
References
Appendix A
We will show that $`detA`$ is a real number, or, equivalently, that $`detA^{}=detA`$. This property is required in order that one can use the matrix $`A^{}A`$ to introduce the pseudofermionic field.
Let us consider $`A`$ in the form
$$A=\left(\begin{array}{cc}B& mD^\mathrm{T}\\ m+D& B^{}\end{array}\right),$$
(A.1)
where $`D_{i,j}=\stackrel{~}{M}_{i,j}m\delta _{i,j}`$.
Using the general property of the determinant of a square block matrix reported in Eq. (3.30) and exchangingThis corresponds to an even number of row exchanges, since $`A`$ has dimension $`2N`$ and $`N`$, being the dimension of the fermion matrix with periodic boundary conditions, is always an even number. Therefore $`detA`$ does not change under this operation. the two rows of blocks of $`A`$ we can rewrite
$`detA`$ $`=`$ $`det\left(\begin{array}{cc}m+D& B^{}\\ B& mD^\mathrm{T}\end{array}\right)`$ (A.4)
$`=`$ $`det(m+D)det(mD^\mathrm{T}B(D+m)^1B^{}).`$ (A.5)
Remembering that $`D^{}=D`$, we can write $`A^{}`$ in the form
$$A^{}=\left(\begin{array}{cc}B^{}& mD\\ m+D^\mathrm{T}& B\end{array}\right),$$
(A.6)
so that, exchanging the two columns of blocks in $`A^{}`$ and using again Eq. (3.30), we can write
$`detA^{}`$ $`=`$ $`det\left(\begin{array}{cc}mD& B^{}\\ B& m+D^\mathrm{T}\end{array}\right)`$ (A.9)
$`=`$ $`det(mD)det(m+D^\mathrm{T}B(D+m)^1B^{}).`$ (A.10)
After extracting a factor $`(1)^Nm^{2N}=m^{2N}`$ from both $`detA`$ and $`detA^{}`$ and defining $`\alpha =1/m`$, we can write
$`detA`$ $`=`$ $`m^{2N}det(1+\alpha D)det(1+\alpha D^\mathrm{T}+\alpha ^2B(1+\alpha D)^1B^{})`$
$`detA^{}`$ $`=`$ $`m^{2N}det(1\alpha D)det(1\alpha D^\mathrm{T}+\alpha ^2B(1\alpha D)^1B^{})`$ (A.11)
It clearly appears from Eq. (A.11) that $`detA^{}`$ is obtained from $`detA`$ by changing the sign of $`\alpha `$. Therefore, in order to show that $`detA=detA^{}`$, it is sufficient to show that $`detA`$ is an even function of $`\alpha `$.
Let us consider $`det(1+\alpha D)`$ at first. We can expand the determinant as follows:
$`det(1+\alpha D)`$ $`=`$ $`\mathrm{exp}\left(\mathrm{tr}\mathrm{ln}(1+\alpha D)\right)`$ (A.12)
$`=`$ $`\mathrm{exp}\left(\mathrm{tr}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^k}{k}}(\alpha D)^k\right).`$
It is easy to see that the trace of the product of an odd number of $`D`$ matrices is zero. Indeed $`D`$ only connects nearest neighbour lattice sites, so it is not possible to connect a site to itself using the product of an odd number of $`D`$ matrices. Therefore only even powers of $`\alpha `$ appear in the expansion in Eq. (A.12) and $`det(1+\alpha D)`$ is an even function of $`\alpha `$.
Let us consider now $`det(1+\alpha D^\mathrm{T}+\alpha ^2B(1+\alpha D)^1B^{})`$, which we rewrite as $`det(1+P(\alpha ))`$, where
$$P(\alpha )=\alpha D^\mathrm{T}+\alpha ^2B\left(\underset{k=0}{\overset{\mathrm{}}{}}(1)^k\alpha ^kD^k\right)B^{}.$$
(A.13)
We notice that the matrix $`P(\alpha )`$ is expressed as series expansion in $`\alpha `$, where the coefficient of the $`k`$-th term is a homogeneous polynomial of degree $`k`$ in the matrices $`B,B^{},D`$ and $`D^\mathrm{T}`$. Therefore, expanding again the determinant as
$`det(1+P(\alpha ))`$ $`=`$ $`\mathrm{exp}\left(\mathrm{tr}\mathrm{ln}(1+P(\alpha ))\right)`$ (A.14)
$`=`$ $`\mathrm{exp}\left(\mathrm{tr}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^k}{k}}P(\alpha )^k\right),`$
we see that $`det(1+P(\alpha ))`$ can be expanded as a power series in $`\alpha `$ and that the coefficient of the $`k`$-th term is the trace of a homogeneous polynomial of degree $`k`$ in the matrices $`B,B^{},D`$ and $`D^\mathrm{T}`$: since these matrices only connect nearest neighbour sites, the trace is zero for $`k`$ odd. Therefore also in this case the determinant is an even function of $`\alpha `$.
We conclude that $`detA`$, being the product of even functions of $`\alpha `$, is also an even function of $`\alpha `$, and therefore, from Eq. (A.11), $`detA=detA^{}`$.
Appendix B
We will give here a proof, alternative to that presented in Ref. 8, of the equality of the determinants of the submatrices on even and odd sites of $`M^{}M`$ in the case with standard boundary conditions.
In the standard case the fermion matrix $`M`$, defined in Eq. (3.3), has the following form
$$M=\left(\begin{array}{cc}m& \frac{1}{2}D_{\mathrm{eo}}\\ \frac{1}{2}D_{\mathrm{oe}}& m\end{array}\right),$$
(B.1)
and it is easily shown that $`D_{\mathrm{eo}}^{}=D_{\mathrm{oe}}`$. Using this last property it follows that
$$M^{}M=\left(\begin{array}{cc}m^2\frac{1}{4}D_{\mathrm{eo}}D_{\mathrm{oe}}& 0\\ 0& m^2\frac{1}{4}D_{\mathrm{oe}}D_{\mathrm{eo}}\end{array}\right).$$
(B.2)
Using the decomposition of the determinant given in Eq. (3.30), it is easy to show that
$$detM=det\left(m^2\frac{1}{4}D_{\mathrm{eo}}D_{\mathrm{oe}}\right).$$
(B.3)
On the other hand, if we exchange even and odd variables in $`M`$ before applying Eq. (3.30), we obtain
$$detM=det\left(m^2\frac{1}{4}D_{\mathrm{oe}}D_{\mathrm{eo}}\right),$$
(B.4)
and therefore
$$det\left(m^2\frac{1}{4}D_{\mathrm{oe}}D_{\mathrm{eo}}\right)=det\left(m^2\frac{1}{4}D_{\mathrm{eo}}D_{\mathrm{oe}}\right).$$
(B.5) |
warning/0003/cond-mat0003326.html | ar5iv | text | # Swelling-Collapse Transition of Self-Attracting Walks
## Abstract
We study the structural properties of self-attracting walks in $`d`$ dimensions using scaling arguments and Monte Carlo simulations. We find evidence for a transition analogous to the $`\mathrm{\Theta }`$ transition of polymers. Above a critical attractive interaction $`u_\mathrm{c}`$, the walk collapses and the exponents $`\nu `$ and $`k`$, characterising the scaling with time $`t`$ of the mean square end-to-end distance $`R^2t^{2\nu }`$ and the average number of visited sites $`St^k`$, are universal and given by $`\nu =1/(d+1)`$ and $`k=d/(d+1)`$. Below $`u_\mathrm{c}`$, the walk swells and the exponents are as with no interaction, i.e. $`\nu =1/2`$ for all $`d`$, $`k=1/2`$ for $`d=1`$ and $`k=1`$ for $`d2`$. At $`u_\mathrm{c}`$, the exponents are found to be in a different universality class.
In recent years different models of random walks with memory or interaction have been studied. They can be divided in static and dynamic models, for an overview we refer to the papers of Duxbury and Queiroz and Oettinger . Most efforts concentrated on models with repulsive interactions, in particular self-avoiding walks (SAW), which have been found useful for investigating polymers in dilute solution. When an attraction term $`\mathrm{exp}(A/T)`$, $`A<0`$, is included, the SAW model reveals a swelling-collapse transition at the ‘$`\mathrm{\Theta }`$ point’ $`T=\mathrm{\Theta }`$ . In contrast, the likewise challenging case of random walks with a similar attractive interaction, but without repulsion, has been less understood. This problem was solved only for one dimension, while in higher dimensions the results are highly controversial. Our numerical and analytical study of attractive random walks suggests that there exists a swelling-collapse transition, too, that is analogous to the $`\mathrm{\Theta }`$ transition in polymers.
We focus on the dynamic model of self-attracting walks (SATW, where a random walker jumps with probability $`p\mathrm{exp}(nu)`$ to a nearest neighbor site, with $`n=1`$ for already visited sites and $`n=0`$ for not visited sites. The interaction parameter $`u`$ is equivalent to $`A/T`$ for linear polymers. For $`u>0`$, the walk is attracted to its own trajectory . The structural behavior of the walk can be characterized by the mean square end-to-end distance $`R^2(t)`$ and the average number of visited sites $`S(t)`$. It is expected that these quantities scale with time $`t`$ as
$`R^2(t)t^{2\nu }(1\text{a})`$ (1)
and (2)
$`S(t)t^k.(3\text{b})`$ (3)
Earlier analyses for the SATW in two and three dimensions were not conclusive and the numerical data have been controversially interpreted . While Sapozhnikov considered the possibility of the existence of a critical attraction $`u_\mathrm{c}`$ (but his numerical results were not conclusive), Lee and Reis argue strongly against the existence of $`u_\mathrm{c}`$, since they find $`\nu `$ and $`k`$ continuously decreasing with $`u`$.
In this Letter we present scaling arguments and extensive numerical simulations for $`R^2`$ and $`S`$ that strongly suggest the existence of a critical attraction $`u_\mathrm{c}`$ in $`d2`$, with three different universality classes for $`u>u_\mathrm{c}`$, $`u<u_\mathrm{c}`$ and $`u=u_\mathrm{c}`$. Below $`u_\mathrm{c}`$, the SATW is in the universality class of random walks, with $`\nu =1/2`$ and $`k=1`$. Above $`u_\mathrm{c}`$, the SATW collapses and the exponents change to $`\nu =1/(d+1)`$ and $`k=d/(d+1)`$. At the critical point, the exponents are $`\nu _\mathrm{c}=0.40\pm 0.01`$ and $`k_\mathrm{c}=0.80\pm 0.01`$ in $`d=2`$ and $`\nu _\mathrm{c}=0.32\pm 0.01`$ and $`k_\mathrm{c}=0.91\pm 0.03`$ in $`d=3`$ . The existence of $`u_\mathrm{c}`$ is in striking similarity to the ‘$`\mathrm{\Theta }`$ point’ phenomenon of linear polymers where three different universality classes for $`T>\mathrm{\Theta }`$, $`T=\mathrm{\Theta }`$ and $`T<\mathrm{\Theta }`$ exist.
We used Monte-Carlo simulations to study $`R^2(t)`$ and $`S(t)`$. Figure 1 shows representative results of $`R^2(t)`$, for several values of $`u`$ in $`d=3`$. For large values of $`u`$, the curves bend down towards the slope of $`2\nu 0.5`$, while for small values of $`u`$, the curves bend up towards the slope of $`2\nu 1`$. At some intermediate critical value $`u_\mathrm{c}1.9`$, the slope is approximately $`2\nu _\mathrm{c}0.64`$. The mean number of visited sites $`S(t)`$ shows similar behavior, with $`k1`$ below $`u_\mathrm{c}`$, $`k_\mathrm{c}0.91`$ at $`u_\mathrm{c}`$, and $`k0.75`$ above $`u_\mathrm{c}`$. Figure 2a summarizes the asymptotic exponents $`\nu `$ and $`k`$ as a function of $`u`$ in $`d=3`$. We obtained similar results in $`d=2`$, the asymptotic values of $`\nu `$ and $`k`$ are presented in Fig. 2b. In Table I the values of the exponents are summarized and compared with the analogous known exponents for the $`\mathrm{\Theta }`$ transition in linear polymers.
In the following we present analytical arguments for the exponents above criticaltity, which can explain our numerical findings. We assume that for sufficiently strong attraction $`u>u_\mathrm{c}`$ the grown clusters are compact, so that the average number of visited sites scales with the rms displacement $`R(t)R^2(t)^{1/2}`$ as
$$S(t)R(t)^d,uu_\mathrm{c}.$$
(5)
Comparing Eq. (1) and Eq. (5) yields
$$k=\nu d,uu_\mathrm{c}.$$
(6)
For sufficiently strong attraction it takes a very long time for the walker to jump to an unvisited site. Before doing this, the walker diffuses around on the visited sites, being located with equal probability on any of the cluster sites. Hence the mean cluster growth rate is proportional to the ratio between the number of boundary sites and the total number of the cluster sites :
$$\frac{\mathrm{d}S}{\mathrm{d}t}\frac{R^{d1}}{R^d}t^\nu .$$
(7)
Thus $`St^{\nu +1}`$. Combining this result with Eq. (1b) and (6), we obtain
$`\nu ={\displaystyle \frac{1}{d+1}}(8\text{a})`$ (8)
and (9)
$`k={\displaystyle \frac{d}{d+1}}(10\text{b})`$ (10)
for $`uu_\mathrm{c}`$.
Because of universality we assume that these results, which are in agreement with the exact values $`\nu =1/2`$ and $`k=1/2`$ in $`d=1`$ and are supported by our extensive Monte Carlo simulations in $`d=2`$ and $`d=3`$, are valid for all $`u>u_\mathrm{c}`$. Indeed, Fig. 2 suggests that the predictions for $`u>u_\mathrm{c}`$ (Eq. (8)) are approached asymptotically. We like to note that in $`d=2`$ the relation $`k=\nu d`$ also holds for $`uu_\mathrm{c}`$, while in $`d=3`$ the numerical results yield $`k<\nu d`$ for $`u<u_\mathrm{c}`$. Since the mass of the generated clusters scales like $`MSR^{k/\nu }`$, $`k/\nu `$ corresponds to the fractal dimension $`d_\mathrm{f}`$ of the cluster. In $`d=2`$ the clusters are compact for all $`u`$ as $`k/\nu =d_\mathrm{f}=d`$. In $`d=3`$ they are compact for $`u>u_\mathrm{c}`$, while for $`u<u_\mathrm{c}`$, the fractal dimension of clusters generated by simple random walks $`d_\mathrm{f}=2<d`$ is obtained. At the criticality, we find $`d_\mathrm{f}=2.84\pm 0.25`$, but we cannot rule out the possibility that $`d_\mathrm{f}=d`$.
To understand the behavior in the critical regime we suggest the following scaling approach. Guided by Fig. 1, we assume that there exists a crossover time $`t_\xi `$ below which the exponent $`\nu `$ is close to $`\nu _\mathrm{c}`$ and above which $`\nu `$ approaches $`1/2`$ for $`u<u_\mathrm{c}`$ and $`1/(d+1)`$ for $`u>u_\mathrm{c}`$. This suggests the following scaling relations:
$`R(t)t^{\nu _\mathrm{c}}f_\pm (t/t_\xi )(12\text{a})`$ (12)
and (13)
$`S(t)t^{k_\mathrm{c}}g_\pm (t/t_\xi ),(14\text{b})`$ (14)
where (15)
$`t_\xi =|uu_\mathrm{c}|^\alpha .(16\text{c})`$ (16)
The plus sign refers to $`u>u_\mathrm{c}`$, the minus sign to $`u<u_\mathrm{c}`$, and the exponent $`\alpha `$ has to be determined numerically. As $`t_\xi `$ is assumed to be the only relevant time scale, the scaling functions bridge the short time and the long time regime. To match both regimes, we require that $`f_\pm (x)=\mathrm{const}`$ for $`x1`$ ($`tt_\xi `$), and $`f_+(x)x^{1/(d+1)\nu _\mathrm{c}}`$, $`f_{}(x)x^{1/2\nu _\mathrm{c}}`$ for $`x1`$. Analogous results are expected for $`g_\pm (x)`$, with $`g_\pm (x)=\mathrm{const}`$ for $`x1`$, and $`g_+(x)x^{d/(d+1)k_\mathrm{c}}`$, $`g_{}(x)x^{1k_\mathrm{c}}`$ for $`x1`$.
To test the scaling theory and to determine the exponent $`\alpha `$ we plotted $`R^2(t)/t_\xi ^{2\nu _\mathrm{c}}`$ and $`S(t)/t_\xi ^{k_\mathrm{c}}`$ as functions of $`t/t_\xi `$ for several values of $`\alpha `$ in $`d=2`$ and $`d=3`$. We obtained the best data collapse for $`\alpha =5.0\pm 0.5`$ in $`d=3`$ and $`\alpha =7\pm 1`$ in $`d=2`$, which are shown in Fig. 3a and 3b, respectively. The excellent data collapse strongly supports the above scaling assumptions.
We would like to thank Dmitry Malykhanov for the assistance with the simulations. Financial support from the German-Israeli Foundation (GIF), the Minerva Center for Mesoscopics, Fractals, and Neural Networks and the Deutsche Forschungsgemeinschaft is gratefully acknowledged. |
warning/0003/hep-th0003052.html | ar5iv | text | # Brane New World
## 1 Introduction
Randall and Sundrum (RS) have suggested that four dimensional gravity may be recovered in the presence of an infinite fifth dimension provided that we live on a domain wall embedded in anti-de Sitter space (AdS). Their linearized analysis showed that there is a massless bound state of the graviton associated with such a wall as well as a continuum of massive Kaluza-Klein modes. More recently, linearized analyses have examined the spacetime produced by matter on the domain wall and concluded that it is in close agreement with four dimensional Einstein gravity .
RS used horospherical coordinates based on slicing AdS into flat hypersurfaces. These horospherical coordinates break down at the horizons shown in figure 1.
An issue that has not received much attention so far is the role of boundary conditions at these Cauchy horizons in AdS. With stationary perturbations, one can impose the boundary conditions that the horizons remain regular. Indeed, without this boundary condition the solution for stationary perturbations is not well defined. Even for non-perturbative departures from the RS solution, like black holes, one can impose the boundary condition that the AdS horizons remain regular . Non-stationary perturbations on the domain wall, however, will give rise to gravitational waves that cross the horizons. This will tend to focus the null geodesic generators of the horizon, which will mean that they will intersect each other on some caustic. Beyond the caustic, the null geodesics will not lie in the horizon. However, null geodesic generators of the future event horizon cannot have a future endpoint and so the endpoint must lie to the past. We conclude that if the past and future horizons remain non-singular when perturbed<sup>1</sup><sup>1</sup>1 It has been shown that the KK modes of RS give rise to singular horizons . (as required for a well-defined boundary condition) then they must intersect at a finite distance from the wall. By contrast, the past and future horizons don’t intersect in the RS ground state but go off to infinity in AdS.
The RS horizons are like the horizons of extreme black holes. When considering perturbations of black holes, one normally assumes that radiation can flow across the future horizon but that nothing comes out of the past horizon. This is because the past horizon isn’t really there, and should be replaced by the collapse that formed the black hole. To justify a similar boundary condition on the Randall-Sundrum past horizon, one needs to consider the initial conditions of the universe.
The main contender for a theory of initial conditions is the “no boundary” proposal<sup>2</sup><sup>2</sup>2 Other approaches to quantum cosmology in the RS model have been discussed in . Boundary conditions motivated by a Euclidean approach were also used in for a flat domain wall. that the quantum state of the universe is given by a Euclidean path integral over compact metrics. The simplest way to implement this proposal for the Randall Sundrum idea is to take the Euclidean version of the wall to be a four sphere at which two balls of $`AdS_5`$ are joined together. In other words, take two balls in $`AdS_5`$, and glue them together along their four sphere boundaries. The result is topologically a five sphere, with a delta function of curvature on a four dimensional domain wall separating the two hemispheres. If one analytically continues to Lorentzian signature, one obtains a four dimensional de Sitter hyperboloid, embedded in Lorentzian anti de Sitter space, as shown in figure 2.
The past and future RS horizons, are replaced by the past and future light cones of the points at the centres of the two balls. Note that the past and future horizons now intersect each other and are non extreme, which means they are stable to small perturbations. A perfectly spherical Euclidean domain wall will give rise to a four dimensional Lorentzian universe that expands forever in an inflationary manner<sup>3</sup><sup>3</sup>3 Such inflationary brane-world solutions have been studied in . For a discussion of other cosmological aspects of the RS model, see and references therein..
In order for a spherical domain wall solution to exist, the tension of the wall must be larger than the value assumed by RS, who had a flat domain wall. We shall assume that matter on the wall increases its effective tension, permitting a spherical solution. In section 3, we consider a strongly coupled large $`N`$ CFT on the domain wall. On a spherical domain wall, the conformal anomaly of the CFT increases the effective tension of the domain wall, making the spherical solution possible. The Lorentzian geometry is a de Sitter universe with the conformal anomaly driving inflation<sup>4</sup><sup>4</sup>4 A similar idea was recently discussed within the context of renormalization group flow in the AdS/CFT correspondence . However, in the case the CFT was the CFT dual to the bulk AdS geometry, not a new CFT living on the domain wall. , an idea introduced long ago by Starobinsky .
The no boundary proposal allows one to calculate unambiguously the graviton correlator on the domain wall. In particular, the Euclidean path integral itself uniquely specifies the allowed fluctuation modes, because perturbations that have infinite Euclidean action are suppressed in the path integral. Therefore, in this framework, there is no need to impose by hand an additional, external prescription for the vacuum state for each perturbation mode. In addition, the AdS/CFT correspondence allows a fully quantum mechanical treatment of the CFT, in contrast with the usual classical treatment of matter fields in inflationary cosmology.
Finally, we analytically continue the Euclidean correlator into the Lorentzian region, where it describes the spectrum of quantum mechanical vacuum fluctuations of the graviton field on an inflating domain wall with conformally invariant matter living on it. We find that the quantum loops of the large $`N`$ CFT give spacetime a rigidity that strongly suppresses metric fluctuations on small scales. Since any matter would be expected to behave like a CFT at small scales, this result probably extends to any inflationary model with sufficiently many matter fields. It has long been known that matter loops lead to short distance modifications of gravity. Our work shows that these modifications can lead to observable consequences in an inflationary scenario.
Although we have carried out our calculations for the RS model, we shall show that results for four dimensional Einstein gravity coupled to the CFT can be recovered by taking the domain wall to be large compared with the AdS scale. Thus our conclusion that metric fluctuations are suppressed holds independently of the RS scenario.
The spherical domain wall considered in this paper analytically continues to a Lorentzian de Sitter universe that inflates forever. However, Starobinsky showed that the conformal anomaly driven de Sitter phase is unstable to evolution into a matter dominated universe. If such a solution could be obtained from a Euclidean instanton then it would have an $`O(4)`$ symmetry group, rather than the $`O(5)`$ symmetry of a spherical instanton. We shall study such models for both the RS model and four dimensional Einstein gravity in a separate paper.
The AdS/CFT correspondence provides an explanation of the RS behaviour<sup>5</sup><sup>5</sup>5 This was first pointed out in unpublished remarks of Maldacena and Witten. . It relates the RS model to an equivalent four dimensional theory consisting of general relativity coupled to a strongly interacting conformal field theory and a logarithmic correction. Under certain circumstances, the effects of the CFT and logarithmic term are negligible and pure gravity is recovered. We review this correspondence in section 2.
In section 3 we present our calculation of the graviton correlator on the instanton and demonstrate how the result is continued to Lorentzian signature. Section 4 contains our conclusions and some speculations. This paper also includes two appendices which contain technical details that we have omitted from the text.
## 2 Randall-Sundrum from AdS/CFT
The AdS/CFT correspondence relates IIB supergravity theory in $`AdS_5\times S^5`$ to a $`𝒩=4`$ $`U(N)`$ superconformal field theory. If $`g_{YM}`$ is the coupling constant of this theory then the ’t Hooft parameter is defined to be $`\lambda =g_{YM}^2N`$. The CFT parameters are related to the supergravity parameters by
$$l=\lambda ^{1/4}l_s,$$
(2.1)
$$\frac{l^3}{G}=\frac{2N^2}{\pi },$$
(2.2)
where $`l_s`$ is the string length, $`l`$ the AdS radius and $`G`$ the five dimensional Newton constant. Note that $`\lambda `$ and $`N`$ must be large in order for stringy effects to be small. The CFT lives on the conformal boundary of $`AdS_5`$. The correspondence takes the following form:
$$Z[𝐡]d[𝐠]\mathrm{exp}(S_{grav}[𝐠])=d[\varphi ]\mathrm{exp}(S_{CFT}[\varphi ;𝐡])\mathrm{exp}(W_{CFT}[𝐡]),$$
(2.3)
here $`Z[𝐡]`$ denotes the supergravity partition function in $`AdS_5`$. This is given by a path integral over all metrics in $`AdS_5`$ which induce a given conformal equivalence class of metrics $`𝐡`$ on the conformal boundary of $`AdS_5`$. The correspondence relates this to the generating functional $`W_{CFT}`$ of connected Green’s functions for the CFT on this boundary. This functional is given by a path integral over the fields of the CFT, denoted schematically by $`\varphi `$. Other fields of the supergravity theory can be included on the left hand side; these act as sources for operators of the CFT on the right hand side.
A problem with equation 2.3 as it stands is that the usual gravitational action in AdS is divergent, rendering the path integral ill-defined. A procedure for solving this problem was developed in . First one brings the boundary into a finite radius. Next one adds a finite number of counterterms to the action in order to render it finite as the boundary is moved back off to infinity. These counterterms can be expressed solely in terms of the geometry of the boundary. The total gravitational action for $`AdS_{d+1}`$ becomes
$$S_{grav}=S_{EH}+S_{GH}+S_1+S_2+\mathrm{}.$$
(2.4)
The first term is the usual Einstein-Hilbert action<sup>6</sup><sup>6</sup>6We use a positive signature metric and a curvature convention for which a sphere has positive Ricci scalar. with a negative cosmological constant:
$$S_{EH}=\frac{1}{16\pi G}d^{d+1}x\sqrt{g}\left(R+\frac{d(d1)}{l^2}\right)$$
(2.5)
the overall minus sign arises because we are considering a Euclidean theory. The second term in the action is the Gibbons-Hawking boundary term, which is necessary for a well-defined variational problem :
$$S_{GH}=\frac{1}{8\pi G}d^dx\sqrt{h}K,$$
(2.6)
where $`K`$ is the trace of the extrinsic curvature of the boundary<sup>7</sup><sup>7</sup>7Our convention is the following. Let $`n`$ denotes the outward unit normal to the boundary. The extrinsic curvature is defined as $`K_{\mu \nu }=h_\mu ^\rho h_\nu ^\sigma _\rho n_\sigma `$, where $`h_\mu ^\nu =\delta _\mu ^\nu n_\mu n^\nu `$ projects quantities onto the boundary. and $`h`$ the determinant of the induced metric. The first two counterterms are given by the following (we use the results of rotated to Euclidean signature)
$$S_1=\frac{d1}{8\pi Gl}d^dx\sqrt{h},$$
(2.7)
$$S_2=\frac{l}{16\pi G(d2)}d^dx\sqrt{h}R,$$
(2.8)
where $`R`$ now refers to the Ricci scalar of the boundary metric. The third counterterm is
$$S_3=\frac{l^3}{16\pi G(d2)^2(d4)}d^dx\sqrt{h}\left(R_{ij}R^{ij}\frac{d}{4(d1)}R^2\right),$$
(2.9)
where $`R_{ij}`$ is the Ricci tensor of the boundary metric and boundary indices $`i,j`$ are raised and lowered with the boundary metric $`h_{ij}`$. This expression is ill-defined for $`d=4`$, which is the case of most interest to us. With just the first two counterterms, the gravitational action exhibits logarithmic divergences so a third term is needed. This term cannot be written solely in terms of a polynomial in scalar invariants of the induced metric and curvature tensors; it makes explicit reference to the cut-off (i.e. the finite radius to which the boundary is brought before taking the limit in which it tends to infinity). The form of this term is the same as 2.9 with the divergent factor of $`1/(d4)`$ replaced by $`\mathrm{log}(R/\rho )`$, where $`R`$ measure the boundary radius and $`\rho `$ is some finite renormalization length scale.
Following , we can now use the AdS/CFT correspondence to explain the behaviour discovered by Randall and Sundrum. The (Euclidean) RS model has the following action:
$$S_{RS}=S_{EH}+S_{GH}+2S_1+S_m.$$
(2.10)
Here $`2S_1`$ is the action of a domain wall with tension $`(d1)/(4\pi Gl)`$. The final term is the action for any matter present on the domain wall. The domain wall tension can cancel the effect of the bulk cosmological constant to produce a flat domain wall. However, we are interested in a spherical domain wall so we assume that the matter on the wall gives an extra contribution to the effective tension. We shall discuss a specific candidate for the matter on the wall later on. The wall separates two balls $`B_1`$ and $`B_2`$ of $`AdS`$.
We want to study quantum fluctuations of the metric on the domain wall. Let $`𝐠_0`$ denote the five dimensional background metric we have just described and $`𝐡_0`$ the metric it induces on the wall. Let $`𝐡`$ denote a metric perturbation on the wall. If we wish to calculate correlators of $`𝐡`$ on the domain wall then we are interested in a path integral of the form<sup>8</sup><sup>8</sup>8 In principle, we should worry about gauge fixing and ghost contributions to the gravitational action. A convenient gauge to use in the bulk is transverse traceless gauge. We shall only deal with metric perturbations that also appear transverse and traceless on the domain wall. The gauge fixing terms vanish for such perturbations and the ghosts only couple to these perturbations at higher orders.
$$h_{ij}(x)h_{i^{}j^{}}(x^{})=d[𝐡]Z[𝐡]h_{ij}(x)h_{i^{}j^{}}(x^{}),$$
(2.11)
where
$`Z[𝐡]`$ $`=`$ $`{\displaystyle _{B_1B_2}}d[\delta 𝐠]d[\varphi ]\mathrm{exp}(S_{RS}[𝐠_0+\delta 𝐠])`$
$`=`$ $`\mathrm{exp}(2S_1[𝐡_0+𝐡])`$
$`\times {\displaystyle _{B_1B_2}}d[\delta 𝐠]d[\varphi ]\mathrm{exp}(S_{EH}[𝐠_0+\delta 𝐠]S_{GH}[𝐠_0+\delta 𝐠]S_m[\varphi ;𝐡_0+𝐡]),`$
$`\delta 𝐠`$ denotes a metric perturbation in the bulk that approaches $`𝐡`$ on the boundary and $`\varphi `$ denotes the matter fields on the domain wall. The integrals in the two balls are independent so we can replace the path integral by
$`Z[𝐡]`$ $`=`$ $`\mathrm{exp}(2S_1[𝐡_0+𝐡])\left({\displaystyle _B}d[\delta 𝐠]\mathrm{exp}(S_{EH}[𝐠_0+\delta 𝐠]S_{GH}[𝐠_0+\delta 𝐠])\right)^2`$ (2.13)
$`\times {\displaystyle }d[\varphi ]\mathrm{exp}(S_m[\varphi ;𝐡_0+𝐡]),`$
where $`B`$ denotes either ball. We now take $`d=4`$ and use the AdS/CFT correspondence 2.3 to replace the path integral over $`\delta 𝐠`$ by the generating functional for a conformal field theory:
$`{\displaystyle _B}`$ $`d[\delta 𝐠]\mathrm{exp}(S_{EH}[𝐠_0+\delta 𝐠]S_{GH}[𝐠_0+\delta 𝐠])=`$ (2.14)
$`\mathrm{exp}(W_{RS}[𝐡_0+𝐡]+S_1[𝐡_0+𝐡]+S_2[𝐡_0+𝐡]+S_3[𝐡_0+𝐡]),`$
we shall refer to this CFT as the RS CFT since it arises as the dual of the RS geometry. It has gauge group $`U(N_{RS})`$, where $`N_{RS}`$ is given by equation 2.2. Strictly speaking, we are using an extended form of the AdS/CFT conjecture, which asserts that supergravity theory in a finite region of AdS is dual to a CFT on the boundary of that region with an ultraviolet cut-off related to the radius of the boundary<sup>9</sup><sup>9</sup>9Evidence in support of this extended version of the duality was given in .. The path integral for the metric perturbation becomes
$$Z[𝐡]=\mathrm{exp}(2W_{RS}[𝐡_0+𝐡]+2S_2[𝐡_0+𝐡]+2S_3[𝐡_0+𝐡])d[\varphi ]\mathrm{exp}(S_m[\varphi ;𝐡_0+𝐡]).$$
(2.15)
The RS model has been replaced by a CFT and a coupling to matter fields and the domain wall metric given by the action
$$2S_2[𝐡_0+𝐡]2S_3[𝐡_0+𝐡]+S_m[\varphi ;𝐡_0+𝐡].$$
(2.16)
The remarkable feature of this expression is that the term $`2S_2`$ is precisely the (Euclidean) Einstein-Hilbert action for four dimensional gravity with a Newton constant given by the RS value
$$G_4=G/l.$$
(2.17)
Therefore the RS model is equivalent to four dimensional gravity coupled to a CFT with corrections to gravity coming from the third counter term. This explains why gravity is trapped to the domain wall.
At first sight this appears rather amazing. We started off with a quite complicated five dimensional system and have argued that it is dual to four dimensional Einstein gravity with some corrections and matter fields. However in order to use this description, we have to know how to calculate with the RS CFT. At present, the only way we know of doing this is via AdS/CFT, i.e., going back to the five dimensional description. The point of the AdS/CFT argument is to explain why the RS “alternative to compactification” works and also to explain the origin of the corrections to Einstein gravity in the RS model. Note that if the matter on the domain wall dominates the RS CFT and the third counterterm then these can be neglected and a purely four dimensional description is adequate.
## 3 CFT on the Domain Wall
### 3.1 Introduction
Long ago, Starobinsky studied the cosmology of a universe containing conformally coupled matter . CFTs generally exhibit a conformal anomaly when coupled to gravity (for a review, see ). Starobinsky gave a de Sitter solution in which the anomaly provides the cosmological constant. By analyzing homogeneous perturbations of this model, he showed that the de Sitter phase is unstable but could be long lived, eventually decaying to a FRW cosmology.
In this section we will consider the RS analogue of Starobinsky’s model by putting a CFT on the domain wall. On a spherical domain wall, the conformal anomaly provides the extra tension required to satisfy the Israel equations. It is appealing to choose the new CFT to be a $`𝒩=4`$ superconformal field theory because then the AdS/CFT correspondence makes calculations relatively easy<sup>10</sup><sup>10</sup>10 We emphasize that this use of the AdS/CFT correspondence is independent of the use described above because this new CFT is unrelated to the RS CFT.. This requires that the CFT is strongly coupled, in contrast with Starobinsky’s analysis<sup>11</sup><sup>11</sup>11Note that the conformal anomaly is the same at strong and weak coupling so any differences arising from strong coupling can only show up when we perturb the system..
Our five dimensional (Euclidean) action is the following:
$$S=S_{EH}+S_{GH}+2S_1+W_{CFT}.$$
(3.1)
We seek a solution in which two balls of $`AdS_5`$ are separted by a spherical domain wall. Inside each ball, the metric can be written
$$ds^2=l^2(dy^2+\mathrm{sinh}^2yd\mathrm{\Omega }_d^2),$$
(3.2)
with $`0yy_0`$. The domain wall is at $`y=y_0`$ and has radius
$$R=l\mathrm{sinh}y_0.$$
(3.3)
The effective tension of the domain wall is given by the Israel equations as
$$\sigma _{eff}=\frac{3}{4\pi Gl}\mathrm{coth}y_0.$$
(3.4)
The actual tension of the domain wall is
$$\sigma =\frac{3}{4\pi Gl}.$$
(3.5)
We therefore need a contribution to the effective tension from the CFT. This is provided by the conformal anomaly, which takes the value
$$T=\frac{3N^2}{8\pi ^2R^4},$$
(3.6)
This contributes an effective tension $`T/4`$. We can now obtain an equation for the radius of the domain wall:
$$\frac{R^3}{l^3}\sqrt{\frac{R^2}{l^2}+1}=\frac{N^2G}{8\pi l^3}+\frac{R^4}{l^4}.$$
(3.7)
It is easy to see that this has a unique positive solution for $`R`$. We shall derive this equation directly from the action in subsection 3.3.
We are particularly interested in how perturbations of this model would appear to inhabitants of the domain wall. Thus we are interested in metric perturbations on the sphere
$$ds^2=(R^2\widehat{\gamma }_{ij}+h_{ij})dx^idx^j.$$
(3.8)
Here $`\widehat{\gamma }_{ij}`$ is the metric on a unit $`d`$-sphere. We shall only consider tensor perturbations, for which $`h_{ij}`$ is transverse and traceless with respect to $`\widehat{\gamma }_{ij}`$. In order to calculate correlators of the metric perturbation, we need to know the action to second order in the perturbation. The most difficult part here is obtaining $`W_{CFT}`$ to second order. This is the subject of the next subsection.
### 3.2 CFT Generating Function
We want to work out the effect of the perturbation on the CFT on the sphere. To do this we use AdS/CFT. Introduce a fictional AdS region that fills in the sphere. Let $`\overline{l},\overline{G}`$ be the AdS radius and Newton constant of this region. We emphasize that this region has nothing to do with the regions of AdS that “really” lie inside the sphere in the RS scenario. This new AdS region is bounded by the sphere. If we take $`\overline{l}`$ to zero then the sphere is effectively at infinity in AdS so we can use AdS/CFT to calculate the generating functional of the CFT on the sphere. In other words, $`\overline{l}`$ is acting like a cut-off in the CFT and taking it to zero corresponds to removing the cut-off. However the relation
$$\frac{\overline{l}^3}{\overline{G}}=\frac{2N^2}{\pi },$$
(3.9)
implies that if $`\overline{l}`$ is taken to zero then we must also take $`\overline{G}`$ to zero since $`N`$ is fixed (and large).
For the unperturbed sphere, the metric in the new AdS region is
$$ds^2=\overline{l}^2(dy^2+\mathrm{sinh}^2y\widehat{\gamma }_{ij}dx^idx^j),$$
(3.10)
and the sphere is at $`y=y_0`$ given by $`R=\overline{l}\mathrm{sinh}y_0`$. Note that $`y_0\mathrm{}`$ as $`\overline{l}0`$ since $`R`$ is fixed. In order to use AdS/CFT for the perturbed sphere, we need to know how the perturbation extends into the bulk. This is done by solving the linearized Einstein equations. It is always possible to choose a gauge in which the bulk metric perturbation takes the form
$$h_{ij}(y,x)dx^idx^j,$$
(3.11)
where $`h_{ij}`$ is transverse and traceless with respect to the metric on the spherical spatial sections:
$$\widehat{\gamma }^{ij}(x)h_{ij}(y,x)=\widehat{}^ih_{ij}(y,x)=0,$$
(3.12)
with $`\widehat{}`$ denoting the covariant derivative defined by the metric $`\widehat{\gamma }_{ij}`$. Since we are only dealing with tensor perturbations, this choice of gauge is consistent with the boundary sitting at constant $`y`$. If scalar metric perturbations were included then we would have to take account of a perturbation in the position of the boundary. These issues are discussed in detail in Appendix A.
The linearized Einstein equations in the bulk are (for any dimension)
$$^2h_{\mu \nu }=\frac{2}{\overline{l}^2}h_{\mu \nu },$$
(3.13)
where $`\mu ,\nu `$ are $`d+1`$ dimensional indices. It is convenient to expand the metric perturbation in terms of tensor spherical harmonics $`H_{ij}^{(p)}(x)`$. These obey
$$\widehat{\gamma }^{ij}H_{ij}^{(p)}(x)=\widehat{}^iH_{ij}^{(p)}(x)=0,$$
(3.14)
and they are tensor eigenfunctions of the Laplacian:
$$\widehat{}^2H_{ij}^{(p)}=\left(2p(p+d1)\right)H_{ij}^{(p)},$$
(3.15)
where $`p=2,3,\mathrm{}`$. We have suppressed extra labels $`k,l,m,\mathrm{}`$ on these harmonics. The harmonics are orthonormal with respect to the obvious inner product. See Appendix B and for more details of their properties. The metric perturbation can be written as a sum of separable perturbations of the form
$$h_{ij}(y,x)=f_p(y)H_{ij}^{(p)}(x).$$
(3.16)
Substituting this into equation 3.13 gives
$$f_p^{\prime \prime }(y)+(d4)\mathrm{coth}yf_p^{}(y)(2(d2)+(p(p+d1)+2(d3))\mathrm{cosech}^2y)f_p(y)=0.$$
(3.17)
The roots of the indicial equation are $`p+2`$ and $`pd+3`$, yielding two linearly independent solutions for each $`p`$. In order to compute the generating functional $`W_{CFT}`$ we have to calculate the Euclidean action of these solutions. However, because the latter solution goes as $`y^{(p+d3)}`$ at the origin $`y=0`$ of the instanton, the corresponding fluctuation modes have infinite Euclidean action<sup>12</sup><sup>12</sup>12This can be seen by surrounding the origin by a small sphere $`y=ϵ`$ and calculating the surface terms in the action that arise on this sphere. They are the same as the surface terms in equations 3.25 and 3.26 below, which are obviously divergent for the modes in question.. Hence they are suppressed in the path integral. Therefore, in contrast to other methods where one requires a (rather ad hoc) prescription for the vacuum state of each perturbation mode, there is no need to impose boundary conditions by hand in our approach: the Euclidean path integral defines its own boundary conditions, which automatically gives a unique Green function. The path integral unambiguously specifies the allowed fluctuation modes as those which vanish at $`y=0`$. Note that boundary conditions at the origin in Euclidean space replace the need for boundary conditions at the horizon in Lorentzian space. The solution regular at $`y=0`$ is given by
$$f_p(y)=\frac{\mathrm{sinh}^{p+2}y}{\mathrm{cosh}^py}F(p/2,(p+1)/2,p+(d+1)/2,\mathrm{tanh}^2y).$$
(3.18)
This solution can also be written in terms of associated Legendre functions:
$$f_p(y)(\mathrm{sinh}y)^{(5d)/2}P_{(d+1)/2}^{(p+(d1)/2)}(\mathrm{cosh}y)(\mathrm{sinh}y)^{(4d)/2}Q_{p+(d2)/2}^{d/2}(\mathrm{coth}y),$$
(3.19)
and the latter can be related to Legendre functions if $`d/2`$ is an integer, using
$$Q_\nu ^m(z)=(z^21)^{m/2}\frac{d^mQ_\nu }{dz^m}.$$
(3.20)
The full solution for the metric perturbation is
$$h_{ij}(y,x)=\underset{p}{}\frac{f_p(y)}{f_p(y_0)}H_{ij}^{(p)}(x)d^dx^{}\sqrt{\widehat{\gamma }}h^{kl}(x^{})H_{kl}^{(p)}(x^{}).$$
(3.21)
We have a solution for the metric perturbation throughout the bulk region. The AdS/CFT correspondence can now be used to give the generating functional of the CFT on the perturbed sphere:
$$W_{CFT}=S_{EH}+S_{GH}+S_1+S_2+\mathrm{}.$$
(3.22)
We shall give the terms on the right hand side for $`d=4`$.
The Einstein-Hilbert action with cosmological constant is
$$S_{EH}=\frac{1}{16\pi \overline{G}}d^5x\sqrt{g}\left(R+\frac{12}{\overline{l}^2}\right),$$
(3.23)
and perturbing this gives
$`S_{bulk}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi \overline{G}}}{\displaystyle d^5x\sqrt{g}\left(\frac{8}{\overline{l}^2}+\frac{1}{4}h^{\mu \nu }^2h_{\mu \nu }+\frac{1}{2\overline{l}^2}h^{\mu \nu }h_{\mu \nu }\right)}`$ (3.24)
$``$ $`{\displaystyle \frac{1}{16\pi \overline{G}}}{\displaystyle d^4x\sqrt{\gamma }\left(\frac{1}{2}n^\mu h^{\nu \rho }_\nu h_{\mu \rho }+\frac{3}{4}h_{\nu \rho }n^\mu _\mu h^{\nu \rho }\right)},`$
where Greek indices are five dimensional and we are raising and lowering with the unperturbed five dimensional metric. $`n=ldy`$ is the unit normal to the boundary and $``$ is the covariant derivative defined with the unperturbed bulk metric. $`\gamma _{ij}=R^2\widehat{\gamma }_{ij}`$ is the unperturbed boundary metric. It is important to keep track of all the boundary terms arising from integration by parts. Evaluating on shell gives
$$S_{EH}=\frac{\overline{l}^3}{2\pi \overline{G}}d^4x\sqrt{\widehat{\gamma }}_0^{y_0}𝑑y\mathrm{sinh}^4y\frac{\overline{l}^3}{16\pi \overline{G}}d^4x\sqrt{\widehat{\gamma }}\left(\frac{3}{4\overline{l}^4}h^{ij}_yh_{ij}\frac{\mathrm{coth}y_0}{\overline{l}^4}h^{ij}h_{ij}\right).$$
(3.25)
where we are now raising and lowering with $`\widehat{\gamma }_{ij}`$. The Gibbons-Hawking term is
$$S_{GH}=\frac{\overline{l}^3}{2\pi \overline{G}}d^4x\sqrt{\widehat{\gamma }}\left(\mathrm{sinh}^3y_0\mathrm{cosh}y_0\frac{1}{8\overline{l}^4}h^{ij}_yh_{ij}\right).$$
(3.26)
The first counter term is
$`S_1`$ $`=`$ $`{\displaystyle \frac{3}{8\pi \widehat{G}\overline{l}}}{\displaystyle d^4x\sqrt{\gamma }}`$ (3.27)
$`=`$ $`{\displaystyle \frac{3\overline{l}^3}{8\pi \overline{G}}}{\displaystyle d^4x\sqrt{\widehat{\gamma }}\left(\mathrm{sinh}^4y_0\frac{1}{4\overline{l}^4}h^{ij}h_{ij}\right)}.`$
The second counter term is
$`S_2`$ $`=`$ $`{\displaystyle \frac{\overline{l}}{32\pi \overline{G}}}{\displaystyle d^4x\sqrt{\gamma }R}`$ (3.28)
$`=`$ $`{\displaystyle \frac{\overline{l}^3}{32\pi \overline{G}}}{\displaystyle d^4x\sqrt{\widehat{\gamma }}\left(12\mathrm{sinh}^2y_0\frac{2}{\overline{l}^4\mathrm{sinh}^2y_0}h^{ij}h_{ij}+\frac{1}{4\overline{l}^4\mathrm{sinh}^2y_0}h^{ij}\widehat{}^2h_{ij}\right)}.`$
Thus with only two counter terms we would have
$`W_{CFT}={\displaystyle \frac{3N^2\mathrm{\Omega }_4}{8\pi ^2}}\mathrm{log}{\displaystyle \frac{R}{\overline{l}}}`$ $``$ $`{\displaystyle \frac{\overline{l}^3}{16\pi \overline{G}}}{\displaystyle }d^4x\sqrt{\widehat{\gamma }}({\displaystyle \frac{1}{4\overline{l}^4}}h^{ij}_yh_{ij}+{\displaystyle \frac{1}{\overline{l}^4}}h^{ij}h_{ij}({\displaystyle \frac{3}{2}}\sqrt{1+{\displaystyle \frac{\overline{l}^2}{R^2}}})`$ (3.29)
$`+`$ $`{\displaystyle \frac{1}{\overline{l}^2R^2}}h^{ij}h_{ij}{\displaystyle \frac{1}{8\overline{l}^2R^2}}h^{ij}\widehat{}^2h_{ij}).`$
$`\mathrm{\Omega }_4`$ is the area of a unit four-sphere and we have used equation 3.9. The expansion of $`_yh_{ij}`$ at $`y=y_0`$ is obtained from
$$_yh_{ij}=\underset{p}{}\frac{f_p^{}(y_0)}{f_p(y_0)}H_{ij}^{(p)}(x)d^4x^{}\sqrt{\widehat{\gamma }}h^{kl}(x^{})H_{kl}^{(p)}(x^{})$$
(3.30)
and
$`{\displaystyle \frac{f_p^{}(y_0)}{f_p(y_0)}}`$ $`=`$ $`2+{\displaystyle \frac{\overline{l}^2}{2R^2}}(p+1)(p+2)+p(p+1)(p+2)(p+3){\displaystyle \frac{\overline{l}^4}{4R^4}}\mathrm{log}(\overline{l}/R)+{\displaystyle \frac{\overline{l}^4}{8R^4}}[p^4+2p^3`$ (3.31)
$``$ $`5p^210p2p(p+1)(p+2)(p+3)(\psi (1)+\psi (2)\psi (p/2+2)\psi (p/2+5/2))]`$
$`+`$ $`𝒪\left({\displaystyle \frac{\overline{l}^6}{R^6}}\mathrm{log}(\overline{l}/R)\right).`$
The psi function is defined by $`\psi (z)=\mathrm{\Gamma }^{}(z)/\mathrm{\Gamma }(z)`$. Substituting into the action we find that the divergences as $`\overline{l}0`$ cancel at order $`R^4/\overline{l}^4`$ and $`R^2/\overline{l}^2`$. The term of order $`\overline{l}^4/R^4`$ in the above expansion makes a contribution to the finite part of the action (along with a term from the square root in equation 3.29):
$`W_{CFT}`$ $`=`$ $`{\displaystyle \frac{3N^2\mathrm{\Omega }_4}{8\pi ^2}}\mathrm{log}{\displaystyle \frac{R}{\overline{l}}}`$
$`+`$ $`{\displaystyle \frac{N^2}{256\pi ^2R^4}}{\displaystyle \underset{p}{}}\left({\displaystyle d^4x^{}\sqrt{\widehat{\gamma }}h^{kl}(x^{})H_{kl}^{(p)}(x^{})}\right)^2\left(2p(p+1)(p+2)(p+3)\mathrm{log}(\overline{l}/R)+\mathrm{\Psi }(p)\right),`$
where
$`\mathrm{\Psi }(p)`$ $`=`$ $`p(p+1)(p+2)(p+3)\left[\psi (p/2+5/2)+\psi (p/2+2)\psi (2)\psi (1)\right]`$ (3.33)
$`+p^4+2p^35p^210p6.`$
To cancel the logarithmic divergences as $`\overline{l}0`$, we have to introduce a length scale $`\rho `$ defined by $`\overline{l}=ϵ\rho `$ and add a counter term proportional to $`\mathrm{log}ϵ`$ to cancel the divergence as $`ϵ`$ tends to zero. The counter term is
$`S_3`$ $`=`$ $`{\displaystyle \frac{\overline{l}^3}{64\pi \overline{G}}}\mathrm{log}ϵ{\displaystyle d^4x\sqrt{\gamma }\left(\gamma ^{ik}\gamma ^{jl}R_{ij}R_{kl}\frac{1}{3}R^2\right)}`$ (3.34)
$`=`$ $`{\displaystyle \frac{\overline{l}^3}{64\pi \overline{G}}}\mathrm{log}ϵ{\displaystyle d^4x\sqrt{\widehat{\gamma }}\left(12+\frac{1}{R^4}\left[2h^{ij}h_{ij}\frac{3}{2}h^{ij}\widehat{}^2h_{ij}+\frac{1}{4}h^{ij}\widehat{}^4h_{ij}\right]\right)}.`$
This term does indeed cancel the logarithmic divergence, leaving us with
$`W_{CFT}`$ $`=`$ $`{\displaystyle \frac{3N^2\mathrm{\Omega }_4}{8\pi ^2}}\mathrm{log}{\displaystyle \frac{R}{\rho }}`$
$`+`$ $`{\displaystyle \frac{N^2}{256\pi ^2R^4}}{\displaystyle \underset{p}{}}\left({\displaystyle d^4x^{}\sqrt{\widehat{\gamma }}h^{kl}(x^{})H_{kl}^{(p)}(x^{})}\right)^2\left(2p(p+1)(p+2)(p+3)\mathrm{log}(\rho /R)+\mathrm{\Psi }(p)\right)`$
Note that varying $`W_{CFT}`$ twice with respect to $`h_{ij}`$ yields the expression for the transverse traceless part of the correlator $`T_{ij}(x)T_{i^{}j^{}}(x^{})`$ on a round four sphere. At large $`p`$, this behaves like $`p^4\mathrm{log}p`$, as expected from the flat space result . In fact this correlator can be determined in closed form solely from the trace anomaly and symmetry considerations<sup>13</sup><sup>13</sup>13 See for a general discussion of such correlators on maximally symmetric spaces.. However, we shall be be interested in calculating cosmologically observable effects, for which our mode expansion is more useful.
### 3.3 The Total Action.
Recall that our five dimensional action is
$$S=S_{EH}+S_{GH}+2S_1+W_{CFT}.$$
(3.36)
In order to calculated correlators of the metric, we need to evaluate the path integral
$`Z[𝐡]`$ $`=`$ $`{\displaystyle _{B_1B_2}}d[\delta 𝐠]\mathrm{exp}(S)`$
$`=`$ $`\mathrm{exp}(2S_1[𝐡_0+𝐡]W_{CFT}[𝐡_0+𝐡])\left({\displaystyle _B}d[\delta 𝐠]\mathrm{exp}(S_{EH}[𝐠_0+\delta 𝐠]S_{GH}[𝐠_0+\delta 𝐠])\right)^2.`$
Here $`𝐠_0`$ and $`𝐡_0`$ refer to the unperturbed background metrics in the bulk and on the wall respectively and $`𝐡`$ denotes the metric perturbation on the wall. Many of the terms required here can be obtained from results in the previous section by simply replacing $`\overline{l}`$ and $`\overline{G}`$ with $`l`$ and $`G`$. For example, from equation 3.27 we obtain
$$S_1[𝐡_0+𝐡]=\frac{3l^3}{8\pi G}d^4x\sqrt{\widehat{g}}\left(\mathrm{sinh}^4y_0\frac{1}{4l^4}\right),$$
(3.38)
where $`y_0`$ is defined by $`R=l\mathrm{sinh}y_0`$. The path integral over $`\delta 𝐠`$ is performed by splitting it into a classical and quantum part:
$$\delta 𝐠=𝐡+𝐡^{},$$
(3.39)
where the boundary perturbation $`𝐡`$ is extended into the bulk using the linearized Einstein equations and the requirement of finite Euclidean action, i.e., $`𝐡`$ is given in the bulk by equation 3.21. $`𝐡^{}`$ denotes a quantum fluctuation that vanishes at the domain wall. The gravitational action splits into separate contributions from the classical and quantum parts:
$$S_{EH}+S_{GH}=S_0[𝐡]+S^{}[𝐡^{}],$$
(3.40)
where $`S_0`$ can be read off from equations 3.25 and 3.26 as
$$S_0=\frac{3l^3\mathrm{\Omega }_4}{2\pi G}_0^{y_0}𝑑y\mathrm{sinh}^2y_0\mathrm{cosh}^2y_0+\frac{l^3}{16\pi G}d^4x\sqrt{\widehat{\gamma }}\left(\frac{1}{4l^4}h^{ij}_yh_{ij}+\frac{\mathrm{coth}y_0}{l^4}h^{ij}h_{ij}\right),$$
(3.41)
Note that $`S^{}`$ cannot be converted to a surface term since $`𝐡^{}`$ does not satisfy the Einstein equations. We shall not need the explicit form for $`S^{}`$ since the path integral over $`𝐡^{}`$ just contributes a factor of some determinant $`Z_0`$ to $`Z[𝐡]`$. We obtain
$$Z[𝐡]=Z_0\mathrm{exp}(2S_0[𝐡_0+𝐡]2S_1[𝐡_0+𝐡]W_{CFT}[𝐡_0+𝐡]).$$
(3.42)
The exponent is given by
$`2S_0+2S_1`$ $`+`$ $`W_{CFT}={\displaystyle \frac{3l^3\mathrm{\Omega }_4}{\pi G}}{\displaystyle _0^{y_0}}𝑑y\mathrm{sinh}^2y\mathrm{cosh}^2y+{\displaystyle \frac{3\mathrm{\Omega }_4R^4}{4\pi Gl}}+{\displaystyle \frac{3N^2\mathrm{\Omega }_4}{8\pi ^2}}\mathrm{log}{\displaystyle \frac{R}{\rho }}`$ (3.43)
$`+`$ $`{\displaystyle \frac{1}{l^4}}{\displaystyle \underset{p}{}}\left({\displaystyle }d^4x^{}\sqrt{\widehat{\gamma }}h^{kl}(x^{})H_{kl}^{(p)}(x^{})\right)^2[{\displaystyle \frac{l^3}{32\pi G}}({\displaystyle \frac{f_p^{}(y_0)}{f_p(y_0)}}+4\mathrm{coth}y_06)`$
$`+`$ $`{\displaystyle \frac{N^2}{256\pi ^2\mathrm{sinh}^4y_0}}(2p(p+1)(p+2)(p+3)\mathrm{log}(\rho /R)+\mathrm{\Psi }(p))].`$
We have kept the unperturbed action in order to demonstrate how the conformal anomaly arises: it is simply the coefficient of the $`\mathrm{log}(R/\rho )`$ term divided by the area $`\mathrm{\Omega }_4R^4`$ of the sphere. If we set the metric perturbation to zero and vary $`R`$ in equation 3.43 (using $`R=l\mathrm{sinh}y_0`$) then we reproduce equation 3.7.
Having calculated $`R`$, we can now choose a convenient value for the renormalization scale $`\rho `$. If we were dealing purely with the CFT then we could keep $`\rho `$ arbitrary. However, since the third counter term (equation 3.34) involves the square of the Weyl tensor (the integrand is proportional to the difference of the Euler density and the square of the Weyl tensor), we can expect pathologies to arise if this term is present when we couple the CFT to gravity. In other words, when coupled to gravity, different choices of $`\rho `$ lead to different theories. We shall choose the value $`\rho =R`$ so that the third counter term exactly cancels the divergence in the CFT, with no finite remainder and hence no residual curvature squared terms in the action.
The (Euclidean) graviton correlator can be read off from the action as
$$h_{ij}(x)h_{i^{}j^{}}(x^{})=\frac{128\pi ^2R^4}{N^2}\underset{p=2}{\overset{\mathrm{}}{}}W_{iji^{}j^{}}^{(p)}(x,x^{})F(p,y_0)^1$$
(3.44)
where we have eliminated $`l^3/G`$ using equation 3.7. The function $`F(p,y_0)`$ is given by
$$F(p,y_0)=e^{y_0}\mathrm{sinh}y_0\left(\frac{f_p^{}(y_0)}{f_p(y_0)}+4\mathrm{coth}y_06\right)+\mathrm{\Psi }(p),$$
(3.45)
and the bitensor $`W_{iji^{}j^{}}^{(p)}(x,x^{})`$ is defined as
$$W_{iji^{}j^{}}^{(p)}(x,x^{})=\underset{k,l,m,\mathrm{}}{}H_{ij}^{(p)}(x)H_{i^{}j^{}}^{(p)}(x^{}),$$
(3.46)
with the sum running over all the suppressed labels $`k,l,m,\mathrm{}`$ of the tensor harmonics.
The appearance of $`N^2`$ in the denominator in equation 3.44 suggests that the CFT suppresses metric perturbations on all scales. This is misleading because $`R`$ also depends on $`N`$. The function $`F(p,y_0)`$ has the following limiting forms for large and small radius:
$$\underset{y_0\mathrm{}}{lim}F(p,y_0)=\mathrm{\Psi }(p)+p^2+3p+6,$$
(3.47)
$$\underset{y_00}{lim}F(p,y_0)=\mathrm{\Psi }(p)+p+6.$$
(3.48)
$`F(p,y_0)`$ has poles at $`p=4,5,6,\mathrm{}`$ with zeros between each pair of negative integers starting at $`3,4`$. When we analytically continue to Lorentzian signature, we shall be particularly interested in zeros lying in the range $`p3/2`$. There is one such zero exactly at $`p=0`$, another near $`p=0`$ and a third near $`p=3/2`$. For large radius, these extra zeros are at $`p0.054`$ and $`p1.48`$ while for small radius they are at $`p0.094`$ and $`p1.60`$. For intermediate radius they lie between these values, with the zeros crossing through $`3/2`$ and $`0`$ at $`y_00.632`$ and $`y_01.32`$ respectively.
### 3.4 Comparison With Four Dimensional Gravity.
We discussed in section 2 how the RS scenario reprodues the predictions of four dimensional gravity when the effects of matter on the domain wall dominates the effects of the RS CFT. In our case we have a CFT on the domain wall. This has action proportional to $`N^2`$. The RS CFT is a similar CFT (but with a cut-off) and therefore has action proportional to $`N_{RS}^2`$. Hence we can neglect it when $`NN_{RS}`$. The logarithmic counterterm is also proportional to $`N_{RS}^2`$ and therefore also negligible. We therefore expect the predictions of four dimensional gravity to be recovered when $`NN_{RS}`$. We shall now demonstrate this explicitly.
First consider the radius $`R`$ of the domain wall given by equation 3.7. It is convenient to write this in terms of the rank $`N_{RS}`$ of the RS CFT (given by $`l^3/G=2N_{RS}^2/\pi `$)
$$\frac{R^3}{l^3}\sqrt{\frac{R^2}{l^2}+1}=\frac{N^2}{16N_{RS}^2}+\frac{R^4}{l^4}.$$
(3.49)
If we assume $`NN_{RS}1`$ then the solution is
$$\frac{R}{l}=\frac{N}{2\sqrt{2}N_{RS}}\left[1+\frac{N_{RS}^2}{N^2}+𝒪(N_{RS}^4/N^4)\right].$$
(3.50)
Note that this implies $`Rl`$, i.e., the domain wall is large compared with the anti-de Sitter length scale.
Now let’s turn to a four dimensional description in which we are considering a four sphere with no interior. The only matter present is the CFT. The metric is simply
$$ds^2=R_4^2\widehat{\gamma }_{ij}dx^idx^j,$$
(3.51)
where $`R_4`$ remains to be determined. The action is the four dimensional Einstein-Hilbert action (without cosmological constant) together with $`W_{CFT}`$. There is no Gibbons-Hawking term because there is no boundary. Without a metric perturbation, the action is simply
$$S=\frac{1}{16\pi G_4}d^4x\sqrt{\gamma }R+W_{CFT}=\frac{3\mathrm{\Omega }_4R_4^2}{4\pi G_4}+\frac{3N^2\mathrm{\Omega }_4}{8\pi ^2}\mathrm{log}\frac{R_4}{\rho }.$$
(3.52)
where $`G_4`$ is the four dimensional Newton constant. We want to calculate the value of $`R_4`$ so we can’t choose $`\rho =R_4`$ yet. Varying $`R_4`$ gives
$$R_4^2=\frac{N^2G_4}{4\pi },$$
(3.53)
and $`N`$ is large hence $`R_4`$ is much greater than the four dimensional Planck length. Substituting $`G_4=G_5/l`$, this reproduces the leading order value for $`R`$ found above from the five dimensional calculation.
We can now go further and include the metric perturbation. The perturbed four dimensional Einstein-Hilbert action is
$$S_{EH}^{(4)}=\frac{1}{16\pi G_4}d^4x\sqrt{\widehat{\gamma }}\left(12R_4^2\frac{2}{R_4^2}h^{ij}h_{ij}+\frac{1}{4R_4^2}h^{ij}\widehat{}^2h_{ij}\right).$$
(3.54)
Adding the perturbed CFT gives
$`S`$ $`=`$ $`{\displaystyle \frac{3N^2\mathrm{\Omega }_4}{16\pi ^2}}+{\displaystyle \frac{3N^2\mathrm{\Omega }_4}{8\pi ^2}}\mathrm{log}{\displaystyle \frac{R_4}{\rho }}+{\displaystyle \underset{p}{}}\left({\displaystyle }d^4x^{}\sqrt{\widehat{\gamma }}h^{kl}(x^{})H_{kl}^{(p)}(x^{})\right)^2[{\displaystyle \frac{1}{64\pi G_4R_4^2}}(p^2+3p+6)`$ (3.55)
$`+`$ $`{\displaystyle \frac{N^2}{256\pi ^2R_4^4}}(2p(p+1)(p+2)(p+3)\mathrm{log}(\rho /R_4)+\mathrm{\Psi }(p))].`$
Setting $`\rho =R_4`$, we find that the graviton correlator for a four dimensional universe containing the CFT is
$$h_{ij}(x)h_{i^{}j^{}}(x^{})=8N^2G_4^2\underset{p=2}{\overset{\mathrm{}}{}}W_{iji^{}j^{}}^{(p)}(x,x^{})\left[p^2+3p+6+\mathrm{\Psi }(p)\right]^1.$$
(3.56)
This can be compared with the expression obtained from the five dimensional calculation, which can be written
$`h_{ij}(x)h_{i^{}j^{}}(x^{})`$ $`=`$ $`{\displaystyle \frac{8N^2G^2}{l^2}}[1+𝒪(N_{RS}^2/N^2)]{\displaystyle \underset{p=2}{\overset{\mathrm{}}{}}}W_{iji^{}j^{}}^{(p)}(x,x^{})[p^2+3p+6+\mathrm{\Psi }(p)`$
$`+`$ $`4p(p+1)(p+2)(p+3)(N_{RS}^2/N^2)\mathrm{log}(N_{RS}/N)+𝒪(N_{RS}^2/N^2)]^1.`$
We have expanded in terms of
$$\frac{N_{RS}^2}{N^2}=\frac{\pi l^3}{2N^2G}.$$
(3.58)
The four and five dimensional expressions clearly agree (for $`G_4=G/l`$) when $`NN_{RS}`$, i.e., $`Rl`$. There are corrections of order $`(N_{RS}^2/N^2)\mathrm{log}(N_{RS}/N)`$ coming from the RS CFT and the logarithmic counter term. In fact, these corrections can be absorbed into the renormalization of the CFT on the domain wall if, instead of choosing $`\rho =R`$, we choose
$$\rho =R\left(1\frac{2N_{RS}^2}{N^2}\mathrm{log}(N_{RS}/N)\right).$$
(3.59)
The corrections to the four dimensional expression are then of order $`N_{RS}^2/N^2`$. We shall not give these correction terms explicitly although they are easily obtained from the exact result 3.44.
### 3.5 Lorentzian Correlator.
In this subsection we shall show how the Euclidean correlator calculated above is analytically continued to give a correlator for Lorentzian signature. We have put many of the details in Appendix B but the analysis is still rather technical so the reader may wish to skip to the final result, which is given in equation 3.66. The techniques used here were developed in .
Let us first introduce a new label $`p^{}=i(p+3/2)`$, so that on the four sphere
$$\widehat{}^2H_{ij}^{(p^{})}=\lambda _p^{}H_{ij}^{(p^{})},$$
(3.60)
where $`p^{}=7i/2,9i/2,\mathrm{}`$ and
$$\lambda _p^{}=(p^2+17/4).$$
(3.61)
Recall that there are extra labels on the tensor harmonics that we have suppressed. The set of rank-two tensor eigenmodes on $`S^4`$ forms a representation of the symmetry group of the manifold. Hence the sum (equation B.2) of the degenerate eigenfunctions with eigenvalue $`\lambda _p^{}`$ defines a maximally symmetric bitensor $`W_{(p^{})i^{}j^{}}^{ij}(\mu (\mathrm{\Omega },\mathrm{\Omega }^{}))`$, where $`\mu (\mathrm{\Omega },\mathrm{\Omega }^{})`$ is the distance along the shortest geodesic between the points with polar angles $`\mathrm{\Omega }`$ and $`\mathrm{\Omega }^{}`$. The expression of the bitensor in terms of a set of fundamental bitensors with $`\mu `$-dependent coefficient functions together with the relation between the bitensors on $`S^4`$ and Lorentzian de Sitter space are obtained in Appendix B.
The motivation for the unusual labelling is that, as demonstrated in Appendix B, in terms of the label $`p^{}`$ the bitensor on $`S^4`$ has exactly the same formal expression as the corresponding bitensor on Lorentzian de Sitter space. This property will enable us to analytically continue the Euclidean correlator into the Lorentzian region without Fourier decomposing it. In other words, instead of imposing by hand a prescription for the vacuum state of the graviton on each mode separately and propagating the individual modes into the Lorentzian region, we compute the two-point tensor correlator in real space, directly from the no boundary path integral. Since the path integral unambiguously specifies the allowed fluctuation modes as those which vanish at the origin of the instanton (see discussion in subsection 3.2), this automatically gives a unique Euclidean correlator. The technical advantage of our method is that dealing directly with the real space correlator makes the derivation independent of the gauge ambiguities involved in the mode decomposition .
We begin by continuing the graviton correlator (equation 3.44) obtained via the five dimensional calculation. The analytic continuation of the correlator for four dimensional gravity (equation 3.56) is completely analogous. In terms of the new label $`p^{}`$, the Euclidean correlator 3.44 between two points on the wall is given by
$$h_{ij}(\mathrm{\Omega })h_{i^{}j^{}}(\mathrm{\Omega }^{})=\frac{128\pi ^2R^4}{N^2}\underset{p^{}=7i/2}{\overset{i\mathrm{}}{}}W_{iji^{}j^{}}^{(p^{})}(\mu )G(p^{},y_0)^1$$
(3.62)
where
$`G(p^{},y_0)`$ $`=`$ $`F(ip^{}3/2,y_0)`$
$`=`$ $`e^{y_0}\mathrm{sinh}y_0({\displaystyle \frac{g_p^{}^{}(y_0)}{g_p^{}(y_0)}}+4\mathrm{coth}y_06)+(p^44ip^3+p^2/25ip^{}63/16`$
$`+`$ $`(p^2+1/4)(p^2+9/4)[\psi (ip^{}/2+5/4)+\psi (ip^{}/2+7/4)\psi (1)\psi (2)]).`$
with $`g_p^{}(y)=Q_{ip^{}1/2}^2(\mathrm{coth}y)`$, which follows from eq. 3.19. The function $`G(p^{},y_0)`$ is real and positive for all values of $`p^{}`$ in the sum and for arbitrary $`y_00`$.
We have the Euclidean correlator defined as an infinite sum. However, the eigenspace of the Laplacian on de Sitter space suggests that the Lorentzian propagator is most naturally expressed as an integral over real $`p^{}`$. We must therefore first analytically continue our result from imaginary to real $`p^{}`$. The coefficient $`G(p^{},y_0)^1`$ of the bitensor is analytic in the upper half complex $`p^{}`$-plane, apart from three simple poles on the imaginary axis. One of them is always at $`p^{}=3i/2`$, regardless of the radius of the sphere. Let the position of the remaining two poles be written $`p_k^{}=i\mathrm{\Lambda }_k(y_0)`$. If we take the radius of the domain wall to be large compared with the AdS scale (which is necessary for corrections to four dimensional Einstein gravity to be small) then<sup>14</sup><sup>14</sup>14If we decrease the radius of the domain wall, then the poles move away from each other. Their behaviour follows from the discussion below equations 3.47 and 3.48. For $`y_00.632`$, $`\mathrm{\Lambda }_1`$ becomes slightly smaller than zero while for $`y_01.32`$, $`\mathrm{\Lambda }_2`$ becomes slightly greater than $`3/2`$. $`0<\mathrm{\Lambda }_k3/2`$, with $`\mathrm{\Lambda }_10`$ and $`\mathrm{\Lambda }_23/2`$. Since $`G(p^{},y_0)`$ is real on the imaginary $`p^{}`$-axis, the residues at these poles are purely imaginary. In order to extend the correlator into the complex $`p^{}`$-plane, we must also understand the continuation of the bitensor itself. As shown in Appendix B, the condition of regularity at opposite points on the four sphere imposed by the completeness relation (equation B.4) is sufficient to uniquely specify the analytic continuation of $`W_{iji^{}j^{}}^{(p^{})}(\mu )`$ into the complex $`p^{}`$-plane. The extended bitensor is defined by equations B.6, B.12 and B.17.
Now we are able to write the sum in equation 3.62 as an integral along a contour $`𝒞_1`$ encircling the points $`p^{}=7i/2,9i/2,..ni/2`$, where $`n`$ tends to infinity. This yields
$$h_{ij}(\mathrm{\Omega })h_{i^{}j^{}}(\mathrm{\Omega }^{})=\frac{i64\pi ^2R^4}{N^2}_{𝒞_1}𝑑p^{}\mathrm{tanh}p^{}\pi W_{iji^{}j^{}}^{(p^{})}(\mu )G(p^{},y_0)^1.$$
(3.64)
Since we know the analytic properties of the integrand in the upper half complex $`p^{}`$-plane, we can distort the contour for the $`p^{}`$ integral to run along the real axis. At large imaginary $`p^{}`$ the integrand decays and the contribution vanishes in the large $`n`$ limit. However as we deform the contour towards the real axis, we encounter three extra poles in the $`\mathrm{cosh}p^{}\pi `$ factor, the pole at $`p^{}=3i/2`$ becoming a double pole due to the simple zero of $`G(p^{},y_0)`$. In addition, we have to take in account the two poles of $`G(p^{},y_0)^1`$ at $`p^{}=i\mathrm{\Lambda }_k`$.
For the $`p^{}=5i/2`$ pole, it follows from the normalization of the tensor harmonics that $`W_{iji^{}j^{}}^{(5i/2)}=0`$. Indirectly, this is a consequence of the fact that spin-2 perturbations do not have a dipole or monopole component. The meaning of the remaining two poles of the $`\mathrm{tanh}p^{}\pi `$ factor has been extensively discussed in , where the continuation is described of the two-point tensor fluctuation correlator from a four dimensional $`O(5)`$ instanton into open de Sitter space. They represent non-physical contributions to the graviton propagator, arising from the different nature of tensor harmonics on $`S^4`$ and on Lorentzian de Sitter space. In fact, a degeneracy appears between $`p_t^{}=3i/2`$ and $`p_t^{}=i/2`$ tensor harmonics and respectively $`p_v^{}=5i/2`$ vector harmonics and $`p_s^{}=5i/2`$ scalar harmonics on $`S^4`$. More precisely, the tensor harmonics that constitute the bitensors $`W_{(3i/2)}^{iji^{}j^{}}`$ and $`W_{(i/2)}^{iji^{}j^{}}`$ can be constructed from a vector (scalar) quantity. Consequently, the contribution to the correlator from the former pole is pure gauge, while the latter eigenmode should really be treated as a scalar perturbation, using the perturbed scalar action. Henceforth we shall exclude them from the tensor spectrum. This leaves us with the poles of $`G(p^{},y_0)`$ at $`p^{}=i\mathrm{\Lambda }_k`$. If we deform the contour towards the real axis, we must compensate for them by subtracting their residues from the integral. We will see that these residues correspond to discrete “supercurvature” modes in the Lorentzian tensor correlator.
The contribution from the closing of the contour in the upper half $`p^{}`$-plane vanishes. Hence our final result for the Euclidean correlator reads
$`h_{ij}(\mathrm{\Omega })h_{i^{}j^{}}(\mathrm{\Omega }^{})`$ $`=`$ $`{\displaystyle \frac{i64\pi ^2R^4}{N^2}}[{\displaystyle _{\mathrm{}}^+\mathrm{}}dp^{}\mathrm{tanh}p^{}\pi W_{iji^{}j^{}}^{(p^{})}(\mu )G(p^{},y_0)^1`$ (3.65)
$`+2\pi {\displaystyle \underset{k=1}{\overset{2}{}}}\mathrm{tan}\mathrm{\Lambda }_k\pi W_{iji^{}j^{}}^{(i\mathrm{\Lambda }_k)}(\mu )\mathrm{𝐑𝐞𝐬}(G(p^{},y_0)^1;i\mathrm{\Lambda }_k)].`$
The analytic continuation from a four sphere into Lorentzian closed de Sitter space is given by setting the polar angle $`\mathrm{\Omega }=\pi /2it`$. Without loss of generality we may take $`\mu =\mathrm{\Omega }`$, and $`\mu `$ then continues to $`\pi /2it`$. We then obtain the correlator in de Sitter space where one point has been chosen as the origin of the time coordinate.
The continuation of the bitensor $`W_{iji^{}j^{}}^{(p^{})}(\mu )`$ is given in Appendix B. An extra subtlety arises if one wants to identify the continued bitensor with the usual sum of tensor harmonics on de Sitter space. It turns out that in order to do so, one must extract a factor $`ie^{p\pi }/\mathrm{sinh}p^{}\pi `$ from its coefficient functions<sup>15</sup><sup>15</sup>15The underlying reason is that there exist two independent bitensors of the form defined by equations B.6 and B.12. Under the integral in the Lorentzian correlator, they are related by the factor $`ie^{p\pi }/\mathrm{sinh}p^{}\pi `$. It follows from the continuation of the completeness relation (equation B.4) that the sum of degenerate tensor harmonics on de Sitter space equals the second independent bitensor, rather then the bitensor that we obtain by continuation from $`S^4`$. Therefore, in order to express the Lorentzian two-point tensor correlator in terms of tensor harmonics, we must extract this factor from the bitensor. We refer the interested reader to the Appendix for the details.. We denote the final form of the bitensor by $`W_{iji^{}j^{}}^{L(p^{})}(\mu (x,x^{}))`$, which is defined in the Appendix, equations B.6, B.12 and B.24.
The extra factor $`ie^{p\pi }/\mathrm{sinh}p^{}\pi `$ combines with the factor $`i\mathrm{tanh}p^{}\pi `$ in the integrand to $`e^{p^{}\pi }/\mathrm{cosh}p^{}\pi `$. Furthermore, since $`G(p^{},y_0)=\overline{G}(p^{},y_0)`$, we can rewrite the correlator as an integral from $`0`$ to $`\mathrm{}`$. We finally obtain the Lorentzian tensor Feynman (time-ordered) correlator,
$`h_{ij}(x)h_{i^{}j^{}}(x^{})`$ $`=`$ $`{\displaystyle \frac{128\pi ^2R^4}{N^2}}[{\displaystyle _0^+\mathrm{}}dp^{}\mathrm{tanh}p^{}\pi W_{iji^{}j^{}}^{L(p^{})}(\mu )\mathrm{}(G(p^{},y_0)^1)`$ (3.66)
$`+\pi {\displaystyle \underset{k=1}{\overset{2}{}}}\mathrm{tan}\mathrm{\Lambda }_k\pi W_{iji^{}j^{}}^{L(i\mathrm{\Lambda }_k)}(\mu )\mathrm{𝐑𝐞𝐬}(G(p^{},y_0)^1;i\mathrm{\Lambda }_k)]`$
$`+i{\displaystyle \frac{128\pi ^2R^4}{N^2}}[{\displaystyle _0^+\mathrm{}}dp^{}W_{iji^{}j^{}}^{L(p^{})}(\mu )\mathrm{}(G(p^{},y_0)^1)`$
$`\pi {\displaystyle \underset{k=1}{\overset{2}{}}}W_{iji^{}j^{}}^{L(i\mathrm{\Lambda }_k)}(\mu )\mathrm{𝐑𝐞𝐬}(G(p^{},y_0)^1;i\mathrm{\Lambda }_k)].`$
In this integral the bitensor $`W_{iji^{}j^{}}^{L(p^{})}(\mu (x,x^{}))`$ may be written as the sum of the degenerate rank-two tensor harmonics on closed de Sitter space with eigenvalue $`\lambda _p^{}=(p_{}^{}{}_{}{}^{2}+17/4)`$ of the Laplacian. Note that the normalization factor $`\stackrel{~}{Q}_p^{}=p^{}(4p^2+25)/48\pi ^2`$ of the bitensor is imaginary at $`p^{}=i\mathrm{\Lambda }_k`$ and the residues of $`G^1`$ are also imaginary, so the quantities in square brackets are all real. Both integrands in equation 3.66 vanish as $`p^{}0`$, so the correlator is well-behaved in the infrared.
For cosmological applications, one is usually interested in the expectation of some quantity squared, like the microwave background multipole moments. For this purpose, all that matters is the symmetrized correlator, which is just the real part of the Feynman correlator.
Gravitational waves provide an extra source of time-dependence in the background in which the cosmic microwave background photons propagate. In particular, the contribution of gravitational waves to the CMB anisotropy is given by the integral in the Sachs-Wolfe formula, which is basically the integral along the photon trajectory of the time derivative of the tensor perturbation. Hence the resulting microwave multipole moments $`𝒞_l`$ can be directly determined from the graviton correlator.
We can therefore understand the effect of the strongly coupled CFT on the microwave fluctuation spectrum by comparing our result 3.66 with the transverse traceless part of the graviton propagator in four-dimensional de Sitter spacetime . On the four-sphere, this is easily obtained by varying the Einstein-Hilbert action with a cosmological constant. In terms of the bitensor, this yields
$`h_{ij}(\mathrm{\Omega })h_{i^{}j^{}}(\mathrm{\Omega }^{})`$ $`=`$ $`32\pi G_4R^2{\displaystyle \underset{p^{}=7i/2}{\overset{i\mathrm{}}{}}}{\displaystyle \frac{W_{iji^{}j^{}}^{(p^{})}(\mu (\mathrm{\Omega },\mathrm{\Omega }^{}))}{\lambda _p^{}2}},`$ (3.67)
which continues to
$`h_{ij}(x)h_{i^{}j^{}}(x^{})`$ $`=`$ $`32\pi G_4R^2{\displaystyle _0^+\mathrm{}}{\displaystyle \frac{dp^{}}{\lambda _p^{}2}}W_{iji^{}j^{}}^{L(p^{})}(\mu (x,x^{})).`$ (3.68)
This can be compared with equation 3.66. Note that (apart from the pole at $`p^{}=3i/2`$ corresponding to the gauge mode mentioned before) there are no supercurvature modes. We defer a detailed discussion of the effect of the CFT on the tensor perturbation spectrum in de Sitter space to the next section.
## 4 Conclusion
We have studied a Randall-Sundrum cosmological scenario consisting of a domain wall in anti-de Sitter space with a large $`N`$ conformal field theory living on the wall. The confomal anomaly of the CFT provides an effective tension which leads to a de Sitter geometry for the domain wall. We have computed the spectrum of quantum mechanical vacuum fluctuations of the graviton field on the domain wall, according to Euclidean no boundary initial conditions. The Euclidean path integral unambiguously specifies the tensor correlator with no additional assumptions. This is the first calculation of quantum fluctuations for RS cosmology.
In the usual inflationary models, one considers the classical action for a single scalar field. In that context, it is consistent to neglect quantum matter loops, on the grounds that they are small. On the other hand, in this paper we have studied a strongly coupled large $`N`$ CFT living on the domain wall, for which quantum loops of matter are important. By using the AdS/CFT correspondence, we have performed a fully quantum mechanical treatment of this CFT. The most notable effect of the large $`N`$ CFT on the tensor spectrum is that it suppresses small scale fluctuations on the microwave sky. It can be seen from equation 3.66 that the CFT yields a $`(p^4\mathrm{ln}p^{})^1`$ behaviour for the graviton propagator at large $`p^{}`$ (in agreement with the flat space results of ), instead of the usual $`p^2`$ falloff (equation 3.68). In other words, quantum loops of the CFT give spacetime a rigidity that strongly suppresses metric fluctuations on small scales. Note that this is true independently of how the de Sitter geometry arises, i.e. it is also true for four dimensional Einstein gravity. In addition, the coupling of the CFT to tensor perturbatons gives rise to two additional discrete modes in the tensor spectrum. Although this is a novel feature in the context of inflationary tensor perturbations, it is not surprising. In conventional open inflationary scenarios for instance, the coupling of scalar field fluctuations with scalar metric perturbations introduces a supercurvature mode with an eigenvalue of the Laplacian close to the discrete de Sitter gauge mode . The former discrete mode at $`p^{}=i\mathrm{\Lambda }_13i/2`$ in equation 3.66 is nothing else than the analogue of this well known supercurvature mode in the scalar fluctuation spectrum. The second mode has an eigenvalue $`p^{}=i\mathrm{\Lambda }_20`$. Its interpretation is less clear, but it is clearly an effect of the matter on the domain wall. However it hardly contributes to the correlator because $`\mathrm{tan}\mathrm{\Lambda }_2\pi `$ is very small.
The effect of the CFT on large scales is more difficult to quantify because of the complicated $`p^{}`$-dependence of the tensor correlator (equation 3.66) in the low-$`p^{}`$ regime. Generally speaking, however, long-wavelength tensor correlations in closed (or open) models for inflation are very sensitive to the details of the underlying theory, as well as to the boundary conditions at the instanton. Since tensor fluctuations do give a substantial contribution to the large scale CMB anisotropies, this may provide an additional way to observationally distinguish different inflationary scenarios .
Most matter fields can be expected to behave like a CFT at small scales. Furthermore, fundamental theories such as string theory predict the existence of a large number of matter fields. Therefore, our results based on a quantum treatment of a large $`N`$ CFT may be accurate at small scales for any matter. If this is the case then our result shows that tensor perturbations at small angular scales are much smaller than predicted by calculations that neglect quantum effects of matter fields.
Acknowledgments
It is a pleasure to thank Steven Gratton, Hugh Osborn and Neil Turok for useful discussions.
## Appendix A Choice of Gauge
This appendix demonstrates how a metric perturbation on the boundary of a ball of AdS is decomposed into vector, scalar and tensor components.
Consider a ball of perturbed AdS with a spherical boundary. Let $`\overline{l}`$ be the AdS length scale. Gaussian normal coordinates are introduced by defining $`\overline{l}y`$ to be the geodesic distance of a point from the origin. The surfaces of constant $`y`$ are spheres on which we introduce coordinates $`x^i`$. In these coordinates the metric takes the form
$$ds^2=\overline{l}^2(dy^2+\mathrm{sinh}^2y\widehat{\gamma }_{ij}(x)dx^idx^j)+h_{ij}(y,x)dx^idx^j.$$
(A.1)
The ball of AdS has been perturbed, so the boundary will be at a position $`y=y_0+\xi (x)`$.
Let the induced metric perturbation on the boundary be $`\widehat{h}_{ij}(x)`$. This can be decomposed into scalar, vector and tensor perturbations with respect to the round metric on the sphere :
$$\widehat{h}_{ij}(x)=\widehat{\theta }_{ij}+2\widehat{}_{(i}\widehat{\chi }_{j)}+\widehat{}_i\widehat{}_j\widehat{\varphi }+\widehat{\gamma }_{ij}\widehat{\psi },$$
(A.2)
where we use hats to denote quantities defined on the sphere (i.e. quantities that depend only on $`x`$). $`\widehat{\theta }_{ij}`$ is a transverse traceless tensor on the sphere and $`\widehat{\chi }_i`$ is a transverse vector on the sphere. $`\widehat{\varphi }`$ and $`\widehat{\psi }`$ are scalars on the sphere. $`\widehat{\chi }_i`$ and $`\widehat{\varphi }`$ can be gauged away by infinitesimal coordinate transformations on the sphere of the form $`x^i=\stackrel{~}{x}^i\eta ^i(\stackrel{~}{x})^i\eta (\stackrel{~}{x})`$ where $`\eta ^i`$ is transverse. Therefore we shall assume that $`\widehat{\chi }`$ and $`\widehat{\varphi }`$ vanish. Note that it is not possible to gauge away $`\widehat{\psi }`$ or $`\xi `$. This paper only deals with tensor perturbations so we shall assume that the scalars $`\widehat{\psi }`$ and $`\xi `$ are vanishing. The induced metric perturbation is then transverse and traceless and can be extended into the bulk as described in section 3. The scalars will be discussed in our next paper.
## Appendix B Maximally Symmetric Bitensors.
A maximally symmetric bitensor $`T`$ is one for which $`\sigma ^{}T=0`$ for any isometry $`\sigma `$ of the maximally symmetric manifold. Any maximally symmetric bitensor may be expanded in terms of a complete set of fundamental maximally symmetric bitensors with the correct index symmetries. For instance
$`T_{iji^{}j^{}}`$ $`=`$ $`t_1(\mu )g_{ij}^{}g_{i^{}j^{}}^{}+t_2(\mu )n_{(i}^{}g_{j)(i^{}}^{}n_{j^{})}^{}+t_3(\mu )\left[g_{ii^{}}^{}g_{jj^{}}^{}+g_{ji^{}}^{}g_{ij^{}}^{}\right]`$ (B.1)
$`+t_4(\mu )n_i^{}n_j^{}n_i^{}^{}n_j^{}^{}+t_5(\mu )\left[g_{ij}^{}n_i^{}^{}n_j^{}^{}+n_i^{}n_j^{}g_{i^{}j^{}}^{}\right].`$
The coefficient functions $`t_j(\mu )`$ depend only on the distance $`\mu (\mathrm{\Omega },\mathrm{\Omega }^{})`$ along the shortest geodesic from the point $`\mathrm{\Omega }`$ to the point $`\mathrm{\Omega }^{}`$. $`n_i^{}^{}(\mathrm{\Omega },\mathrm{\Omega }^{})`$ and $`n_i^{}(\mathrm{\Omega },\mathrm{\Omega }^{})`$ are unit tangent vectors to the geodesics joining $`\mathrm{\Omega }`$ and $`\mathrm{\Omega }^{}`$ and $`g_{ij^{}}(\mathrm{\Omega },\mathrm{\Omega }^{})`$ is the parallel propagator along the geodesic, i.e., $`V^ig_i^j^{}`$ is the vector at $`\mathrm{\Omega }^{}`$ obtained by parallel transport of $`V^i`$ along the geodesic from $`\mathrm{\Omega }`$ to $`\mathrm{\Omega }^{}`$ .
The set of tensor eigenmodes on $`S^4`$ (or on de Sitter space) forms a representation of the symmetry group of the manifold. It follows in particular that their sum over the parity states $`𝒫=\{e,o\}`$ and the quantum numbers $`k,l`$ and $`m`$ on the three sphere defines a maximally symmetric bitensor on $`S^4`$ (or dS space) :
$$W_{(p^{})i^{}j^{}}^{ij}(\mu )=\underset{𝒫klm}{}q_{𝒫klm}^{(p^{})ij}(\mathrm{\Omega })q_{i^{}j^{}}^{(p^{})𝒫klm}(\mathrm{\Omega }^{})^{}.$$
(B.2)
On $`S^4`$ the label $`p^{}`$ takes the value $`7i/2,9i/2,..`$. It is related to a real label $`p`$ by $`p^{}=i(p+3/2)`$. The ranges of the other labels are then $`0kp`$, $`0lk`$ and $`lml`$. On de Sitter space there is a continuum of eigenvalues $`p^{}[0,\mathrm{})`$. We will assume from now on that the eigenmodes are normalized by the condition
$$\sqrt{\gamma }d^4\mathrm{\Omega }q_{𝒫klm}^{(p^{})ij}q_{𝒫^{}k^{}l^{}m^{}ij}^{(p^{\prime \prime })}=\delta ^{p^{}p^{\prime \prime }}\delta _{𝒫𝒫^{}}\delta _{ll^{}}\delta _{mm^{}}$$
(B.3)
The completeness relation on the four sphere may then be written as
$$\gamma ^{\frac{1}{2}}\delta _{i^{}j^{}}^{ij}(\mathrm{\Omega }\mathrm{\Omega }^{})=\underset{p^{}=7i/2}{\overset{+i\mathrm{}}{}}W_{(p^{})i^{}j^{}}^{ij}(\mu (\mathrm{\Omega },\mathrm{\Omega }^{})).$$
(B.4)
Explicit formulae for the components of these tensors may be found in . In this Appendix we will determine $`W_{iji^{}j^{}}^{(p^{})}(\mu )`$ simultaneously on the four sphere and de Sitter space. The construction of the analogous bitensor on $`S^3`$ and $`H^3`$ is given in and their relation is described in .
The bitensor $`W_{(p^{})i^{}j^{}}^{ij}(\mu )`$ has some additional properties arising from its construction in terms of the transverse and traceless tensor harmonics $`q_{ij}^{(p)𝒫klm}`$. The tracelessness of $`W_{iji^{}j^{}}^{(p^{})}`$ allows one to eliminate two of the coefficient functions in equation B.1. It may then be written as
$`W_{iji^{}j^{}}^{(p^{})}(\mu )`$ $`=`$ $`w_1^{(p^{})}\left[g_{ij}^{}4n_i^{}n_j^{}\right]\left[g_{i^{}j^{}}^{}4n_i^{}^{}n_j^{}^{}\right]+w_2^{(p^{})}\left[4n_{(i}^{}g_{j)(i^{}}^{}n_{j^{})}^{}+4n_i^{}n_j^{}n_i^{}^{}n_j^{}^{}\right]`$ (B.6)
$`+w_3^{(p^{})}\left[g_{ii^{}}^{}g_{jj^{}}^{}+g_{ji^{}}^{}g_{ij^{}}^{}2n_i^{}g_{i^{}j^{}}^{}n_j^{}2n_i^{}^{}g_{ij}^{}n_j^{}^{}+8n_i^{}n_j^{}n_i^{}^{}n_j^{}^{}\right]`$
This expression is traceless on either the index pair $`ij`$ or $`i^{}j^{}`$. The requirement that the bitensor be transverse $`^iW_{iji^{}j^{}}^{(p^{})}=0`$ and the eigenvalue condition $`(^2\lambda _p^{})W_{(p^{})}^{iji^{}j^{}}=0`$ impose additional constraints on the remaining coefficient functions $`w_j^{(p^{})}(\mu )`$. To solve these constraint equations it is convenient to introduce the new variables on $`S^4`$ (in de Sitter space, $`\mu `$ is replaced by $`\pi /2i\stackrel{~}{\mu }`$)
$$\{\begin{array}{ccc}\alpha (\mu )\hfill & =\hfill & w_1^{(p)}(\mu )+\frac{2}{3}w_3^{(p)}(\mu )\hfill \\ \beta (\mu )\hfill & =\hfill & \frac{8}{(\lambda _p+8)\mathrm{sin}\mu }\frac{d\alpha (\mu )}{d\mu }\hfill \end{array}$$
(B.7)
In terms of a new argument $`z=\mathrm{cos}^2(\mu /2)`$ (or its continuation on de Sitter space) the transversality and eigenvalue conditions imply for $`\alpha (z)`$
$$z(1z)\frac{d^2\alpha (z)}{d^2z}+\left[48z\right]\frac{d\alpha (z)}{dz}=(\lambda _p^{}+8)\alpha (z)$$
(B.8)
and then for the coefficient functions
$`\{\begin{array}{ccc}w_1\hfill & =\hfill & \frac{6}{5}\left[(\lambda _p^{}+28)z(1z)45/6\right]\alpha (z)+\frac{6}{20}\left[(\lambda _p^{}+8)z(1z)(12z)\right]\beta (z)\hfill \\ w_2\hfill & =\hfill & \frac{9}{5}\left[(\lambda _p^{}+28)z(1z)+\frac{20}{3}(1z)\frac{20}{6}\right]\alpha (z)\frac{6}{20}\left[(\lambda _p^{}+8)z(1z)(43z)\right]\beta (z)\hfill \\ w_3\hfill & =\hfill & \frac{9}{5}\left[(\lambda _p^{}+28)z(1z)40/6\right]\alpha (z)\frac{9}{20}\left[(\lambda _p^{}+8)z(1z)(12z)\right]\beta (z)\hfill \end{array}`$ (B.12)
with $`\lambda _p^{}=(p^2+17/4)`$.
Notice that equation B.8 is precisely the hypergeometric differential equation, which has a pair of independent solutions $`\alpha (z)`$ and $`\alpha (1z)`$ where
$$\alpha (z)=Q_p^{}{}_{2}{}^{}F_{1}^{}(7/2+ip^{},7/2ip^{},4,z)$$
(B.13)
$`Q_p^{}`$ is a constant. The solution for $`\beta (z)`$ follows from equation B.7 and is given by
$$\beta (z)=Q_p^{}{}_{2}{}^{}F_{1}^{}(9/2ip^{},9/2+ip^{},5,z).$$
(B.14)
We will determine below which solution corresponds to the bitensor defined by B.2.
Our discussion so far applies to either $`S^4`$ or de Sitter space. We now specialize to the case of $`S^4`$ and will later obtain results for de Sitter space by analytic continuation. The hypergeometric functions on $`S^4`$ may be expressed in terms of Legendre polynomials in $`\mathrm{cos}\mu `$ (eq. \[15.4.19\] in ),
$`\{\begin{array}{cc}\alpha (\mu )\hfill & =Q_p^{}\mathrm{\Gamma }(4)2^3(\mathrm{sin}\mu )^3P_{1/2+ip^{}}^3(\mathrm{cos}\mu ),\hfill \\ \beta (\mu )\hfill & =Q_p^{}\mathrm{\Gamma }(5)2^4(\mathrm{sin}\mu )^4P_{1/2+ip^{}}^4(\mathrm{cos}\mu ).\hfill \end{array}`$ (B.17)
The solutions for $`\alpha (z)`$ and $`\beta (z)`$ are singular at $`z=1`$ (i.e. for coincident points on $`S^4`$) for generic values of $`p^{}`$. However, for the values of $`p^{}`$ corresponding to the eigenvalues of the Laplacian on $`S^4`$, they are regular everywhere on $`S^4`$. Similarly, $`\alpha (1z)`$ and $`\beta (1z)`$ are generically singular for antipodal points on $`S^4`$ and regular for these special values of $`p^{}`$. For these special values, $`\alpha (z)`$ and $`\alpha (1z)`$ are no longer linearly independent but related by a factor of $`(1)^{(n+1)/2}`$ where $`n=2ip^{}=7,9,11,\mathrm{}`$. This follows from the relation (eq.\[8.2.3\] in )
$$P_\nu ^\mu (z)=e^{i\nu \pi }P_\nu ^\mu (z)\frac{2}{\pi }e^{i\mu \pi }\mathrm{sin}[\pi (\nu +\mu )]Q_\nu ^\mu (z),$$
(B.18)
where the second term vanishes for $`p^{}=7i/2,9i/2,\mathrm{}`$. In fact, the hypergeometric series terminates for these values of $`p^{}`$ and the hypergeometric functions reduce to Gegenbauer polynomials $`C_{n7/2}^{(7/2)}(12z)`$. We have a choice between using $`\alpha (z)`$ and $`\alpha (1z)`$ in the bitensor for these values of $`p^{}`$. However, to obtain the Lorentzian correlator, we had to express the discrete sum 3.62 as a contour integral. Since the Euclidean correlator obeys a differential equation with a delta function source at $`\mu =0`$, we must maintain regularity of the integrand at $`\mu =\pi `$ when extending the bitensor in the complex $`p^{}`$-plane. In other words, for generic $`p^{}`$, we need to work with the solution $`\alpha (z)`$, rather then $`\alpha (1z)`$. We shall therefore choose $`\alpha (z)`$, since this is the solution that we will analytically continue.
The above conditions leave the overall normalisation of the bitensor undetermined. To fix the normalisation constant $`Q_p^{}`$, consider the biscalar quantity
$`g^{ii^{}}g^{jj^{}}W_{iji^{}j^{}}^{(p^{})}(\mu )=12w_1^{(p^{})}6w_2^{(p^{})}+24w_3^{(p^{})}`$ (B.19)
In the coincident limit $`\mathrm{\Omega }\mathrm{\Omega }^{}`$ and $`z1`$ this yields
$$W_{ij}^{(p^{})ij}(\mathrm{\Omega },\mathrm{\Omega })=\underset{𝒫klm}{}q_{ij}^{(p^{})𝒫klm}(\mathrm{\Omega })q^{(p^{})𝒫lmij}(\mathrm{\Omega })^{}=72\alpha (1).$$
(B.20)
Since $`F(0)=1`$ we have $`\alpha (1)=Q_p^{}(1)^{(1+n)/2}`$. By integrating over the four-sphere and using the normalization condition B.3 on the tensor harmonics one obtains, for $`n=2ip^{}=7,9,11,\mathrm{}`$
$$Q_p^{}=\frac{ip^{}(4p^2+25)}{48\pi ^2(1)^{(1+n)/2}}=\frac{p^{}(4p^2+25)}{48\pi ^2\mathrm{sinh}p^{}\pi }.$$
(B.21)
We conclude that the properties of the bitensor appearing in the tensor correlator completely determine its form. Notice that in terms of the label $`p^{}`$ we have obtained a unified functional description of the bitensor on $`S^4`$ and de Sitter space. However, its explicit form is very different in the two cases because the label $`p^{}`$ takes on different values. It is precisely this description that has enabled us in section 3 to analytically continue the correlator from the Euclidean instanton into de Sitter space without Fourier decomposing it. We shall conclude this Appendix by describing in detail the subtleties of this analytic continuation at the level of the bitensor.
To perform the continuation to de Sitter space we note that the geodesic separation $`\mu `$ on $`S^4`$ continues to $`\pi /2it`$, so $`z=\frac{1}{2}(1+i\mathrm{sinh}t)`$ on de Sitter space. The continuation of the hypergeometric functions (B.17) yields
$`\{\begin{array}{cc}\alpha (z)\hfill & =\mathrm{\Gamma }(4)2^3(\mathrm{cosh}t)^3P_{1/2+ip^{}}^3(i\mathrm{sinh}t),\hfill \\ \beta (z)\hfill & =\mathrm{\Gamma }(5)2^4(\mathrm{cosh}t)^4P_{1/2+ip^{}}^4(i\mathrm{sinh}t).\hfill \end{array}`$ (B.24)
However, an extra subtlety arises if one wants to identify the continued bitensor with the usual sum of tensor harmonics on de Sitter space. In particular, in order for the bitensor to correspond to the usual sum of rank-two tensor harmonics on the real $`p^{}`$-axis, one must choose the second solution $`\alpha (1z)`$ to the hypergeometric equation, rather then $`\alpha (z)`$ that enters in the continued bitensor. This is easily seen by performing the continuation on the completeness relation (equation B.4), which should continue to an integral over $`p^{}`$ from $`0`$ to $`\mathrm{}`$ of the Lorentzian bitensor, defined as the sum (B.2) over the degenerate tensor harmonics on de Sitter space. Writing (B.4) as a contour integral and continuing to Lorentzian de Sitter space yields
$$g^{\frac{1}{2}}\delta _{i^{}j^{}}^{ij}(xx^{})=_{\mathrm{}}^+\mathrm{}𝑑p^{}\mathrm{tanh}p^{}\pi W_{(p^{})i^{}j^{}}^{ij}(\mu (x,x^{})).$$
(B.25)
Clearly this is not the correct completeness relation according to the equivalent definition (B.2) of the bitensor on de Sitter space. But let us relate the continued bitensor in (B.25) to the independent bitensor in which the solutions $`\alpha (1z)`$ enter. This can be done by applying (B.18) to the Legendre polynomials in (B.24). By closing the contour in the upper half $`p^{}`$-plane, one sees there is no contribution to the integral (and indeed to the tensor correlator!) from the second term in equation B.18, because its prefactor cancels the $`\mathrm{cosh}^1(p^{}\pi )`$-factor in (B.25), making the integrand analytic in the upper half $`p^{}`$-plane (up to gauge modes). Hence, under the integral both solutions are simply related by the factor $`ie^{p\pi }`$. In addition one needs to extract the $`\mathrm{sinh}^1p^{}\pi `$-factor<sup>16</sup><sup>16</sup>16Remember that $`Q_p^{}`$ gained the factor $`\mathrm{sinh}^1p^{}\pi `$ because we have chosen the solution $`\alpha (z)`$ on the four sphere. The correct normalisation constant for the independent bitensor, obtained from the normalisation condition on the tensor harmonics, is then $`\stackrel{~}{Q}_p^{}=\mathrm{sinh}p^{}\pi Q_p^{}`$. from $`Q_p^{}`$. The completeness relation then becomes,
$$g^{\frac{1}{2}}\delta _{i^{}j^{}}^{ij}(xx^{})=_0^+\mathrm{}𝑑p^{}W_{L(p^{})i^{}j^{}}^{ij}(\mu (x,x^{})),$$
(B.26)
and the hypergeometric functions $`\alpha (1z)`$ and $`\beta (1z)`$ that constitute the final bitensor $`W_{L(p^{})i^{}j^{}}^{ij}(\mu (x,x^{}))`$ are given by
$`\{\begin{array}{cc}\alpha (1z)\hfill & =\stackrel{~}{Q}_p^{}\mathrm{\Gamma }(4)2^3(\mathrm{cosh}t)^3P_{1/2+ip^{}}^3(i\mathrm{sinh}t),\hfill \\ \beta (1z)\hfill & =\stackrel{~}{Q}_p^{}\mathrm{\Gamma }(5)2^4(\mathrm{cosh}t)^4P_{1/2+ip^{}}^4(i\mathrm{sinh}t),\hfill \end{array}`$ (B.29)
with $`\stackrel{~}{Q}_p^{}=p^{}(4p^2+25)/48\pi ^2`$.
On the real $`p^{}`$-axis, $`W_{iji^{}j^{}}^{L(p^{})}(\mu )`$ equals the sum (B.2) of the degenerate rank-two tensor harmonics on closed de Sitter space with eigenvalue $`\lambda _p^{}=(p^2+17/4)`$ of the Laplacian. |
warning/0003/hep-ph0003041.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The role of soft color partons in the high energy hadron interactions is the mostly intriguing modern problem of particles physics. So, the collective phenomena and symmetry breaking in the non-Abelian gauge theories, confinement of colored charges and the infrared divergences of the pQCD are the phenomena just of the soft color particles domain.
It seems natural that the very high multiplicity (VHM) hadron interaction, where the energy of created particles is small, should be sensitive to the soft color particles densities. Indeed, the aim of this paper is to show that even in the hard by definition deep inelastic scattering (DIS) the soft color particles role becomes important in the VHM region.
The standard (mostly popular) hadron theory considers perturbative QCD (pQCD) at small distances (in the scale of $`\mathrm{\Lambda }0.2Gev`$) as the exact theory. This statement is confirmed by a number of experiments, e.g. the DIS data, the QCD jets observation. But one should have in mind that the pQCD has finite range of validity since the non-perturbative effects should be taken into account at distances larger then $`1/\mathrm{\Lambda }`$.
It is natural to assume building the complete theory that at large distances the non-perturbative effects exceed the perturbative ones<sup>1</sup><sup>1</sup>1The corresponding formalism was described e.g. in .. In result the pQCD contributions becomes $`screened`$ by the non-perturbative interactions.
But exist another, more natural, possibility. It should be noted here that the pQCD running coupling constant $`\alpha _s(q^2)=1/b\mathrm{ln}(q^2/\mathrm{\Lambda }^2)`$ becomes infinite at $`q^2=\mathrm{\Lambda }^2`$ and we do not know what happens with pQCD if $`q^2<\mathrm{\Lambda }^2`$. There is a suspicion that at $`q^2\mathrm{\Lambda }^2`$ the properties of theory changed so drastically (being defined on new vacuum) that even the notions of pQCD is $`disappeared`$. This means that pQCD should be truncated from below on the ‘fundamental’ scale $`\mathrm{\Lambda }`$. It seems evident that this infrared cut-off should influence on the soft hadrons emission in the VHM domain.
So, it is important try to rise predictability of pQCD in the ‘forbidden’ area of large distances. For this purpose one should split experimentally the perturbative and non-perturbative effects at the large distances to check out quantitatively the role of soft color partons. One must realize for this purpose highly unusual condition that the non-perturbative effects must be negligible even if the distance among color charges is high.
The non-perturbative effects lead to the strong polarization of QCD vacuum and, in result, to the color charge confinement. This vacuum is unstable against creation of real quark-antiquark $`(q\overline{q})`$ pairs if the distance among charges became large. This pairs are captured into the colorless hadrons and just emission of this ‘vacuum’ hadrons is the mostly important non-perturbative effect. Therefore, generally, the number of hadrons $`n`$, produced even in the hard DIS process, $`nn_p`$, where $`n_p`$ is the number of $`\overline{q}q`$ pairs created ‘perturbatively’.
Notice, if the kinetic energy of colored partons is small, i.e. is comparable with hadron masses, creation of ‘vacuum’ hadrons should be negligible. Just this is the VHM process kinematics: because of the energy-momentum conversation law, produced (final-state) partons can not have high relative momentum and, if they was created at small distances, the production of ‘vacuum’ hadrons will be negligible (or did not play important role). Therefore, if the ‘vacuum’ channel is negligible, the pQCD contributions only should be considered .
The aim of this article is to show that $`n\mathrm{}`$ unavoidably leads to ‘low-$`x`$’ domain.
## 2 DIS kinematics in the VHM domain
To describe the hadron production in pQCD terms the parton-hadron duality is assumed. This means that the multiplicity, momentum etc. distributions of hadron and colored partons are the same. This reduce the problem practically on the level of QED.
Let us consider now $`n`$ particles (gluons) creation the DIS. We would like to calculate $`D_{ab}(x,q^2;n)`$, where
$$\underset{n}{}D_{ab}(x,q^2;n)=D_{ab}(x,q^2).$$
(2.1)
As usual, let $`D_{ab}(x,q^2)`$ be the probability to find parton $`b`$ with virtuality $`q^2<0`$ in the parton $`a`$ of $`\lambda `$ virtuality, $`\lambda >>\mathrm{\Lambda }`$ and $`\alpha _s(\lambda )<<1`$. We always may chose $`q^2`$ and $`x`$ so that the leading logarithm approximation (LLA) will be acceptable. One should assume also that $`(1/x)>>1`$ to have the phase space, into which the particles are produced, sufficiently large.
Then $`D_{ab}(x,q^2)`$ is described by ladder diagrams. From qualitative point of view this means approximation of Markovian process of random walk over coordinate $`\mathrm{ln}(1/x)`$ and time is $`\mathrm{ln}\mathrm{ln}|q^2|`$. LLA means that the ‘mobility’ $`\mathrm{ln}(1/x)/\mathrm{ln}\mathrm{ln}\left|q^2\right|`$ should be large
$$\mathrm{ln}(1/x)>>\mathrm{ln}\mathrm{ln}\left|q^2/\lambda ^2\right|.$$
(2.2)
But, on other hand ,
$$\mathrm{ln}(1/x)<<\mathrm{ln}\left|q^2/\lambda ^2\right|.$$
(2.3)
The leading, able to compensate smallness of $`\alpha _s(\lambda )<<1`$, contributions give integration over wide range $`\lambda ^2<<k_i^2<<q^2`$, where $`k_i^2>0`$ is the ‘mass’ of real, i.e. time-like, gluon. If the time needed to capture the parton into the hadron is $`(1/\mathrm{\Lambda })`$ then the gluon should decay if $`k_i^2>>\lambda ^2`$. This leads to creation of (mini)jets. The mean multiplicity $`\overline{n}_j`$ in the QCD jets is high if the gluon ‘mass’ $`|k|`$ is high: $`\mathrm{ln}\overline{n}_j\sqrt{\mathrm{ln}(k^2/\lambda ^2)}`$.
Rising multiplicity may (i) rise number of (mini)jets $`\nu `$ and/or (ii) rise the mean value mass of (mini)jets $`<|k_i|>`$. We will see that the mechanism (ii) would be favorable. This is consequence of the Markovian character of considered process.
But rising mean value of gluon masses $`|k_i|`$ decrease the range of integrability over $`k_i`$, i.e. violate the condition (2.2) for fixed $`x`$. One can remain the LLA taking $`x0`$. But this may contradict to (2.3), i.e. in any case the LLA becomes invalid in the VHM domain and the next to leading order corrections should be taken into account.
Noting that the LLA gives main contribution, that the rising multiplicity leads to the infrared domain, where the soft gluons creation becomes dominant.
## 3 Preliminary notes
First of all, neglecting the vacuum effects, we introduce definite uncertainty to the formalism. It is reasonable to define the level of strictness of our computations. Let us introduce for this purpose the generating function $`T_{ab}(x,q^2;z)`$:
$$D_{ab}(x,q^2;n)=\frac{1}{2\pi i}\frac{dz}{z^{n+1}}T_{ab}(x,q^2;z).$$
(3.1)
Having large $`n`$ the integral may be calculated by the saddle point method. The smallest solution $`z_c`$ of the equation
$$n=z\frac{}{z}\mathrm{ln}T_{ab}(x,q^2;z)$$
(3.2)
defines the asymptotic over $`n`$ behavior:
$$D_{ab}(x,q^2;n)e^{n\mathrm{ln}z_c(x,q^2;n)}.$$
(3.3)
Using the statistical interpretation of $`z_c`$ as the fugacity it is natural to write:
$$\mathrm{ln}z_c(x,q^2;n)=\frac{C_{ab}(x,q^2;n)}{\overline{n}_{ab}(x,q^2)}.$$
(3.4)
Notice that the solution of eq.(3.2) $`z_c(x,q^2;n)`$ should be the increasing function of $`n`$. At first glance this follows from positivity of all $`D_{ab}(x,q^2;n)`$. But actually this assumes that $`T_{ab}(x,q^2;z)`$ is a regular function of z at $`z=1`$. This is a natural assumption considering just the pQCD predictions.
Therefore,
$$D_{ab}(x,q^2;n)e^{\frac{n}{\overline{n}_{ab}(x,q^2)}C_{ab}(x,q^2;n)}.$$
(3.5)
This form of $`D_{ab}(x,q^2;n)`$ is useful since usually $`C_{ab}(x,q^2;n)`$ is the slowly varying function of $`n`$. So, for Poisson distribution $`C_{ab}(x,q^2;n)\mathrm{ln}n`$. For KNO scaling we have $`C_{ab}(x,q^2;n)=const.`$ over $`n`$.
We would like to note, that neglecting effects of vacuum polarization we introduce into the exponent so high uncertainty assuming $`nn_p`$ that it is reasonable perform the calculations with the exponential accuracy. So, we would calculate
$$\overline{\mu }_{ab}(x,q^2;n)=\mathrm{ln}\frac{D_{ab}(x,q^2;n)}{D_{ab}(x,q^2)}=\frac{n}{\overline{n}_{ab}(x,q^2)}C_{ab}(x,q^2;n)(1+O(1/n))$$
(3.6)
The $`n`$ dependence of $`C_{ab}(x,q^2;n)`$ defines the asymptotic behavior of $`\overline{\mu }_{ab}(x,q^2;n)`$ and calculation of its explicit form would be our aim.
## 4 Correlation functions
Considering particles creation in the DIS processes one should distinguish correlation of particles in the (mini)jets and the correlations between (mini)jets. We will describe the jet correlations. One should introduce the $`\nu `$ jets creation inclusive cross section $`\mathrm{\Phi }_\nu ^{(r)_\nu }(k_1,k_2,\mathrm{}.,k_\nu ;q^2,x)`$, where $`k_i`$, $`i=1,2,\mathrm{},n`$ are the jets 4-momentum. Having $`\mathrm{\Phi }_\nu `$ we can find the correlations function $`N_\nu ^{(r)}(k_1,k_2,\mathrm{}.,k_\nu ;q^2,x)`$, where $`(r)=r_1,\mathrm{},r_\nu `$ and $`r_i=(q,g)`$ defines the sort of created color particle. It is useful introduce the generating functional
$$F^{ab}(q^2,x;w)=\underset{n}{}d\mathrm{\Omega }_n(k)\underset{i=1}{\overset{n}{}}w^{r_i}(k_i)|a_n^{ab}(k_1,k_2,\mathrm{}.,k_n;q^2,x)|^2,$$
(4.1)
where $`a_n^{ab}`$ is the amplitude, $`d\mathrm{\Omega }_n(k)`$ is the phase space volume and $`w^{r_i}(k_i)`$ are the arbitrary functions. It is evident,
$$F^{ab}(q^2,x;w)|_{w=1}=D^{ab}(q^2,x).$$
(4.2)
The inclusive cross sections
$$\mathrm{\Phi }_\nu ^{(r)}(k_1,k_2,\mathrm{}.,k_\nu ;q^2,x)=\underset{1=1}{\overset{\nu }{}}\frac{\delta }{\delta w^{r_i}(k_i)}F^{ab}(q^2,x;w)|_{w=1}.$$
(4.3)
The correlation function
$$N_\nu ^{(r)}(k_1,k_2,\mathrm{}.,k_\nu ;q^2,x)=\underset{1=1}{\overset{\nu }{}}\frac{\delta }{\delta w^{r_i}(k_i)}\mathrm{ln}F^{ab}(q^2,x;w)|_{w=1}.$$
(4.4)
We will use this definitions to find the partial structure functions $`D^{ab}(q^2,x;n)`$.
It will be useful to introduce the Laplace transform over variable $`\mathrm{ln}(1/x)`$:
$$F^{ab}(q^2,x;w)=_{\mathrm{Rej}<0}\frac{dj}{2\pi i}(\frac{1}{x})^jf^{ab}(q^2,j;w)$$
(4.5)
The expansion parameter of our problem $`\alpha _s\mathrm{ln}(q^2/\lambda ^2)1`$. By this reason one should take into account all possible cuttings of the ladder diagrams. So, calculating $`D^{ab}(q^2,x)`$in the LLA all possible cuttings of sceleton ladder diagrams is defined by the factor :
$$\frac{1}{\pi }\left\{\mathrm{\Gamma }_r^{ab}G_r\mathrm{\Gamma }_r^{ab}\right\},$$
(4.6)
i.e. the cutting line may get not only through the exact Green function $`G_r(k_i^2)`$ but through the exact vertex functions $`\mathrm{\Gamma }_r^{ab}(q_i,q_{i+1},k_i)`$ also ($`q_i^2,q_{i+1}^2`$ are negative). We have in the LLA that
$$\lambda ^2<<q_i^2<<q_{i+1}^2<<q^2$$
and
$$xx_{i+1}x_i1.$$
Following to our approximation, see previous section, we would not distinguish the way as cut line go through the Born amplitude
$$a_r^{ab})=\{(\mathrm{\Gamma }_r^{ab})^2G_r\}.$$
We will simply associate $`w^r\mathrm{Ima}_\mathrm{r}^{\mathrm{ab}})`$ to each rung of the ladder.
Considering the asymptotics over $`n`$ the time-like partons virtuality $`k_iq_i^2/y_i`$ should be maximal. Here $`y_i`$ is the fraction of the longitudinal momentum of the jet. Then, slightly limiting the jets phase space,
$$\mathrm{ln}k_i^2=\mathrm{ln}\left|q_{i+1}\right|^2(1+O(\mathrm{ln}(1/x)/|q_{i+1}|^2)).$$
(4.7)
In result, introducing useful in the LLA variable $`\tau _i=\mathrm{ln}(q_i^2/\mathrm{\Lambda }^2)`$, where $`\alpha _s(q^2)=1/\beta \tau `$, $`\beta =(11N/3)(2n_f/3)`$, we can find following set of equations:
$$\tau \frac{}{\tau }f_{ab}(q^2,j;w)=\underset{c,r}{}\phi _{ac}^r(j)w^r(\tau )f_{ab}(q^2,x;w),$$
(4.8)
where
$$\phi _{ac}^r(j)=\phi _{ac}(j)=_0^1\frac{dx}{x}x^jP_{ac}(x)$$
(4.9)
and $`P_{ac}(x)`$ is the regular kernel of the Bethe-Salpeter equation . At $`w=1`$ this equation is the ordinary one for $`D^{ab}(q^2,x)`$.
We will search the correlation functions from eq.(4.8) in terms of Laplace transform $`n_{ab}^{(r)_\nu }(k_1,k_2,\mathrm{}.,k_\nu ;q^2,j)`$. Let us write:
$$f_{ab}(q^2,j;w)=d_{ab}(q^2,j)\mathrm{exp}\{\underset{\nu }{}\frac{1}{\nu !}\underset{i=1}{\overset{\nu }{}}\left(\frac{d\tau _i}{\tau _i}(w^{r_i}(\tau _i)1)\right)n_{ab}^{(r)_\nu }(k_1,k_2,\mathrm{}.,k_\nu ;q^2,j)\}$$
(4.10)
Inserting (4.10) into (4.8) and expanding over $`(w1)`$ we find the sequence of coupled equation.
Omitting the cumbersome calculations, we write in the LLA that
$$\varphi _{ab}^{(r)_\nu }(\tau _1,\tau _2,\mathrm{}.,\tau _\nu ;q^2,j)=d_{ac_1}(j,\tau _1)\phi _{c_1c_2}^{r_1}(j)d_{c_2c_3}(j,\tau _2)\mathrm{}\phi _{c_\nu c_{\nu +1}}^{r_\nu }(j)d_{c_{\nu +1}b}(j,\tau _{\nu +1}).$$
(4.11)
One should take into account the conservation laws:
$$\tau _1\tau _2\mathrm{}\tau _{\nu +1}=\tau ,\tau _1<\tau _2<\mathrm{}<\tau _{\nu +1}<\tau .$$
(4.12)
Computing the Laplace transform of this expression we find $`\mathrm{\Phi }_{ab}^{(r)_\nu }(\tau _1,\tau _2,\mathrm{}.,\tau _\nu ;q^2,x)`$.
The kernel $`d_{ab}(j,\tau )`$ was introduced in (4.11). Let us write it in the form:
$$d_{ab}(j,\tau )=\underset{\sigma =\pm }{}\sigma \frac{d_{ab}(j)}{\nu _+\nu _{}}\tau ^{\nu _\sigma (j)},$$
(4.13)
where
$$d_{qq}^\sigma =\nu _s\phi _{gg},d_{qg}^\sigma =\nu _s\phi _{qq},d_{qg}^\sigma =\phi _{gq},d_{gq}^\sigma =\phi _{qg}$$
(4.14)
and
$$\nu _\sigma =\frac{1}{2}\left\{\phi _{qq}+\phi _{gg}+\sigma \left[(\phi _{qq}\phi _{gg})^24n_f\phi _{qg}\phi _{gq}\right]^{1/2}\right\}.$$
(4.15)
If $`x<<1`$, then $`(j1)<<1`$ are essential. In this case ,
$$\phi _{gg}\phi _{gq}>>\phi _{qg}\phi _{qq}=O(1).$$
(4.16)
This means the gluon jets dominance and
$$n_{gg}^g=\phi _{gg}+O(1).$$
(4.17)
One can find following estimation of the two-jet correlation function:
$$n_{ab}^{r_1r_2}(\tau _1,\tau _2;ȷ,\tau )=O\left(\mathrm{max}\{(\tau _1/\tau )^{\phi _{gg}},(\tau _2/\tau )^{\phi _{gg}},(\tau _1/\tau _2)^{\phi _{gg}}\}\right\}.$$
(4.18)
This correlation function is small since in the LLA $`\tau _1<\tau _2<\tau `$. This means that the jet correlations becomes high iff the mass of correlated jets are comparable. But this condition shrinks the range of integration over $`\tau `$ and by this reason one may neglect the ‘short-range’ correlations among jets. Therefore, as follows from (4.10),
$$f_{ab}(q^2,j;w)=d_{gg}(\tau ,j)\mathrm{exp}\left\{\phi _{gg}_{\tau _0}^\tau \frac{d\tau ^{}}{\tau ^{}}w^g(\tau ^{})\right\}$$
(4.19)
We will use this expression to find the multiplicity distribution in the DIS domain.
## 5 Generating function in the DIS kinematics
To describe particles production one should replace $`w^r\mathrm{Ima}_\mathrm{r}^{\mathrm{ab}})`$ on $`w_n^r\mathrm{Ima}_\mathrm{r}^{\mathrm{ab}})`$, where $`w_n^r`$ is the $`probability`$ of $`n`$ particles production,
$$\underset{n}{}w_n^r=1.$$
(5.1)
Having $`\nu `$ jets one should take into account the conservation condition $`n_1+n_2+\mathrm{}+n_\nu =n`$. By this reason the generating functions formalism is useful. In result one can find that if we take (4.19)
$$w^g=w^g(\tau ,z),w^g(\tau ,z)|_{z=1}=1,$$
(5.2)
then $`f_{ab}(q^2,j;w)`$ defined by (4.19) is the generating functional of the multiplicity distribution in the ‘$`j`$ representation’. In this expression $`w^g(\tau ,z)`$ is the generating function of the multiplicity distribution in the jet of mass $`|k|=\lambda e^{\tau /2}`$.
In result, see (4.5),
$$F^{ab}(q^2,x;w)_{\mathrm{Rej}<0}\frac{dj}{2\pi i}(1/x)^je^{\phi _{gg}\omega (\tau ,z)}$$
(5.3)
where
$$\omega (\tau ,z)=_{\tau _0}^\tau \frac{d\tau ^{}}{\tau ^{}}w^g(\tau ^{},z).$$
(5.4)
Noting the normalization condition (5.2),
$$\omega (\tau ,z=1)=\mathrm{ln}\tau .$$
(5.5)
The integral (5.3) may be calculated by steepest descent method. It is not hard to see that
$$jj_c=1+\left\{4N\omega (\tau ,z)/\mathrm{ln}(1/x)\right\}^{1/2}$$
(5.6)
is essential. Notice that $`j1<<1`$ should be essential we find, instead of the constraint (2.2), that
$$\omega (\tau ,z)<<\mathrm{ln}(1/x).$$
(5.7)
In the frame of this constraint,
$$F^{ab}(q^2,x;w)\mathrm{exp}\left\{4\sqrt{N\omega (\tau ,z)\mathrm{ln}(1/x)}\right\}.$$
(5.8)
Generally speaking, exist such values of $`z`$ that $`j_c11`$. This is possible if $`\omega (\tau ,z)`$ is a regular function of $`z`$ at $`z=1`$. Then $`z_c`$ should be the increasing function of $`n`$ and consequently $`\omega (\tau ,z_c)`$ would be the increasing function of $`n`$. Therefore, one may expect that in the VHM domain $`j_c11`$.
Then $`j1+\omega (\tau ,z)/\mathrm{ln}(1/x)`$ would be essential in the integral (5.3). This leads to following estimation:
$$F^{ab}(q^2,x;w)e^{\omega (\tau ,z)}.$$
But this is impossible since $`F^{ab}(q^2,x;w)`$ should be the increasing function of $`z`$. This shows that the estimation (5.8) has a finite range of validity.
Solution of this problem with unitarity is evident. One should take into account correlations among jets considering the expansion (4.10). Indeed, smallness of $`n_{ab}^{(r)_\nu }`$ may be compensated by large values of $`_i^\nu w^{r_i}(\tau _i,z)`$ in the VHM domain.
## 6 Conclusion
We can conclude that our LLA is applicable in the VHM domain till
$$\omega (\tau ,z)<<\mathrm{ln}(1/x)<<\tau =\mathrm{ln}(q^2/\lambda ).$$
(6.1)
The mean multiplicity of gluons created in the DIS kinematics
$$\overline{n}_g(\tau ,x)=\frac{}{z}\mathrm{ln}F^{ab}(q^2,x;w)|_{z=1}=\omega _1(\tau )\sqrt{4N\mathrm{ln}(1/x)/\mathrm{ln}\tau }>>\omega _1(\tau ),$$
(6.2)
where
$$\omega _1(\tau )=_{\tau _0}^\tau \frac{d\tau _1}{\tau _1}\overline{n}_j(\tau )$$
(6.3)
and the mean gluon multiplicity in the jet $`\overline{n}_j(\tau )`$ has following estimation :
$$\mathrm{ln}\overline{n}_j(\tau )\sqrt{\tau }$$
(6.4)
Inserting (6.4) into (6.3),
$$\omega _1(\tau )=\overline{n}_j(\tau )/\sqrt{\tau }.$$
Therefore, noting (2.3),
$$\overline{n}_g(\tau ,x)\overline{n}_j(\tau )\sqrt{4N\mathrm{ln}(1/x)/\tau \mathrm{ln}\tau }<<\overline{n}_j(\tau ).$$
(6.5)
This means that the considered ‘t-channel’ ladder is important in the narrow domain of multiplicities
$$n\overline{n}_g<<\overline{n}_j.$$
(6.6)
So, in the VHM domain $`n>>\overline{n}_g`$ one should consider
(i) The ladder diagrams with small number of rungs;
(ii) To take into account the malty-jet correlations assuming that increasing multiplicity leads to the increasing number of rungs in the ladder diagram.
To choose one of this possibilities one should consider the structure of $`\omega (\tau ,z)`$ much more carefully. This will be done in the consequent paper.
We can conclude that the VHM domain multiplicities production unavoidably destroy the ladder LLA. To conserve this leading approximation one should choose $`x0`$ and, in result, to get to the multi-ladder diagrams, since in this case $`\alpha _s\mathrm{ln}(q^2/\lambda ^2)1`$ and $`\alpha _s\mathrm{ln}(1/x)1`$. Such theory was considered in .
Acknowledgments
We are grateful to V.G.Kadyshevski for interest to discussed in the paper questions. We would like to note with gratitude that the conversations with E.Kuraev was always interesting and important. |
warning/0003/quant-ph0003056.html | ar5iv | text | # 1 Introduction
## 1 Introduction
In a recent paper, we have introduced a new treatment of systems of compounded angular momentum which leads to very generalized formulas for the states and operators for such systems . In the paper, we worked these quantities out explicitly for the cases of spin $`0`$ and spin $`1`$ resulting from the addition of two spins of $`1/2`$ each. We first obtained the generalized probability amplitudes describing results of measurements on such systems, then used these to derive the matrix treatment of the systems. The forms that we obtained for the vector states and operators proved to be entirely different from the standard forms. However, as might be expected, the results of calculations of measurable quantities are the same. Nevertheless, the question how the vectors and operators belonging to the standard treatment of spin addition are related to the new forms requires an answer. Indeed, since we consider the treatment of spin addition by means of probability amplitudes as being the foundation of any matrix treatment, it is necessary to derive the standard quantities by the new approach.
In this paper, we demonstrate that the standard matrix treatment of compounded spin is indeed derivable by our method. The use of this approach not only yields the standard treatment, but also produces results more generalized than any in the literature. However, these results reduce to the standard forms in an appropriate limit. The systems on which the theory developed is tested are the triplet and singlet states resulting from the addition of the spins of two spin-$`1/2`$ systems. It is clear from the application of the theory to these cases how the extension to arbitrary systems of compounded spin is achieved.
The organization of the paper is as follows. After the introduction in Section $`1`$, we give in Section $`2`$ a brief description of those features of the Landé approach to quantum mechanics that we shall use to develop our treatment. In Section $`3`$, we turn our attention to a review of the work we have so far done on systems of compounded angular momentum. We remind ourselves of the expressions for the probability amplitudes for the addition of general angular momentum in Section $`3.1`$, and of the probability amplitudes for spin addition in Section $`3.2`$.
In Section $`4`$, we look at the way the transformation from wave or probability- amplitude mechanics to matrix mechanics is achieved. We sketch in Section $`4.1`$ the derivation of matrix mechanics from probability-amplitudes mechanics for simple systems. In Section $`4.2`$, we derive the standard form of matrix mechanics for systems of compounded spin \- however, the new results are more generalized than the standard ones.
The results of Section $`4`$ are employed on actual systems in Section $`5`$. The test systems are the singlet and triplet states arising from the addition of the spins of two spin -$`1/2`$ systems. The matrix operator is common to both cases, and is calculated in Section $`5.1`$. The vectors states are obtained in Section $`5.2.`$
The results obtained in Section $`5`$ are more generalized than the standard forms found in the literature. In Section $`6`$ we demonstrate how to reduce these results to the standard forms.
We end the paper with a Discussion and Conclusion in Section $`7`$.
## 2 Basic Theory
### 2.1 The Landé Approach to Quantum Mechanics
The basic theory underlying our work derives from the interpretation of quantum mechanics due to Landé \[2-5\]. Among many features of the Landé approach is the assumption that wave functions and eigenfunctions in quantum mechanics are probability amplitudes. Any probability amplitude connects two states - one state pertaining to the situation that obtains before a measurement is made, and the other to the state that results from the measurement. Thus, an energy eigenfunction $`\varphi _E(𝐫)`$ resulting from solution of the time-independent Schrödinger equation is a probability amplitude connecting two states: the state defined by the eigenvalue $`E`$ is the initial state, while the final state is characterized by the position eigenvalue $`𝐫`$. Thus $`\left|\varphi _E(𝐫)\right|^2d𝐫`$ is the probability that if the system is initially in the state corresponding to $`E`$, a measurement of the position gives $`𝐫`$ in the volume element $`d𝐫`$.
The Landé interpretation of quantum mechanics is based on the principle that nature is ultimately indeterministic and should be described by a theory that is fundamentally probabilistic. Consequently, the description of measurements can only be given in probabilistic terms through probability amplitudes. For a particular system, the different sets of probability amplitudes connecting different measurable quantities are inter-related in the following way.
Let a quantum system have the observables $`A`$, $`B`$ and $`C`$ which have the respective eigenvalue spectra $`A_1`$, $`A_2,`$…,$`A_N,`$ $`B_1`$, $`B_2,`$…,$`B_N`$ and $`C_1`$, $`C_2,..,C_N.`$ If the system is initially in the state corresponding to the eigenvalue $`A_i`$, a measurement of $`B`$ yields any of the eigenvalues $`B_j`$ with probabilities determined by the probability amplitudes $`\eta (A_i;B_j).`$ A measurement of $`C`$ results in one of the eigenvalues $`C_j`$ with probabilities determined by the probability amplitudes $`\psi (A_i;C_j).`$ If the system is initially in the state corresponding to the eigenvalue $`B_i`$, a measurement of $`C`$ gives any of the eigenvalues $`C_j`$ with probabilities determined by the probability amplitudes $`\xi (B_i;C_j).`$ The probability amplitudes display a two-way symmetry contained in the Hermiticity condition
$$\psi (C_j;A_i)=\psi ^{}(A_i;C_j).$$
(1)
These probability amplitudes are orthogonal:
$$\underset{j=1}{\overset{N}{}}\psi ^{}(A_i;C_j)\psi (A_k;C_j)=\delta _{ik}.$$
(2)
The law that connects the three sets of probability amplitudes is
$$\psi (A_i;C_n)=\underset{j=1}{\overset{N}{}}\eta (A_i;B_j)\xi (B_j;C_n).$$
(3)
Though the features of the Landé approach highlighted above refer to probability amplitudes that correspond to a discrete final eigenvalue spectrum, there is no essential difference if this spectrum is instead continuous. In fact, if the observable $`C`$ is the position $`𝐫`$, and if, as is customary, we ignore the initial state $`A_i`$ in the labelling, Eq. (3) becomes
$$\psi (𝐫)=\underset{j=1}{\overset{N}{}}\eta _j\xi _j(𝐫),$$
(4)
where we have set $`\eta _j=\eta (B_j)`$ and $`\xi _j=\xi (B_j;C_n).`$ We recognize this equation as the law of interference of probabilities. In the Landé formalism, this important relation is derived, not assumed.
The relation Eq. (4) is, of course, the basis for the transformation of representation from wave to matrix mechanics for the case where the eigenfunctions $`\xi _j(𝐫)`$ are known from solution of some eigenvalue equation. By the same token, its parent relation Eq. (3) is the basis for the transformation of representation from probability-amplitude mechanics to matrix mechanics in all cases, irrespective of whether or not a differential eigenvalue equation exists for the probability amplitudes $`\xi (B_j;C_n).`$ Indeed, this is the relation on which we have based the derivation of the matrix theory of spin from probability amplitudes \[1,7-11\].
## 3 Review of Previous Results on Angular Momentum Addition
### 3.1 General Theory
In a previous paper , we derived a matrix treatment of spin addition which resulted in new forms for the vectors and the operators, apart from throwing light on the theory of angular momentum addition. In this section, we review the new treatment of spin addition in order to present those results which will be needed in the development of the present work.
We consider first the case of general angular momentum addition. Let a system have the total angular momentum $`𝐉`$ resulting from adding the angular momenta $`𝐉_1`$ and $`𝐉_2`$ of subsystems $`1`$ and $`2`$. Thus,
$$𝐉=𝐉_1+𝐉_2.$$
(5)
The quantum numbers of the angular momenta of the subsystems are $`j_1`$ and $`j_2`$, while that of the angular momentum of the total system is $`j`$. The $`z`$ components of these respective angular momenta are characterized by the quantum numbers $`m_1`$, $`m_2`$ and $`M`$. For the time being, we shall assume that $`𝐉,`$ $`𝐉_1`$ and $`𝐉_2`$ are orbital angular momenta. The subsystems $`1`$ and $`2`$ are characterized by the angular variables $`(\theta _1,\phi _1)`$ and $`(\theta _2,\phi _2)`$ respectively. The standard expression for the wave function of the coupled system is
$$\mathrm{\Psi }_{j_1j_2jM}(\theta _1,\phi _1,\theta _2,\phi _2)=\underset{m_1}{}C(j_1j_2j;m_1m_2M)\varphi _{j_1m_1}^{(1)}(\theta _1,\phi _1)\varphi _{j_2m_2}^{(2)}(\theta _2,\phi _2),$$
(6)
where we have used the notation in Rose for the Clebsch-Gordan coefficients $`C(j_1j_2j;m_1m_2M)`$. If we are dealing with orbital angular momentum, the $`\varphi _{j_1m_1}^{(1)}(\theta _1,\phi _1)`$ and the $`\varphi _{j_2m_2}^{(2)}(\theta _2,\phi _2)`$ are spherical harmonics.
In the Landé interpretation, the function $`\mathrm{\Psi }_{j_1j_2jM}(\theta _1,\phi _1,\theta _2,\phi _2)`$ is a probability amplitude. Its expression in terms of an expansion must be of the general structure of Eq. (3). Therefore, we rewrite Eq. (6) in the following way:
$`\mathrm{\Psi }(j_1,j_2,j,M;\theta _1,\phi _1,\theta _2,\phi _2)={\displaystyle \underset{m_1}{}}\chi (j_1,j_2,j,M;j_1,m_1,j_2,m_2)`$
$`\times \mathrm{\Phi }(j_1,m_1,j_2,m_2;\theta _1,\phi _1,\theta _2,\phi _2),`$ (7)
where
$$\chi (j_1,j_2,j,M;j_1,m_1,j_2,m_2)=C(j_1j_2j;m_1m_2M)$$
(8)
and
$$\mathrm{\Phi }(j_1,m_1,j_2,m_2;\theta _1,\phi _1,\theta _2,\phi _2)=\varphi _{j_1m_1}^{(1)}(\theta _1,\phi _1)\varphi _{j_2m_2}^{(2)}(\theta _2,\phi _2).$$
(9)
Then the various quantities have the following interpretations:
The function $`\mathrm{\Psi }(j_1,j_2,j,M;\theta _1,\phi _1,\theta _2,\phi _2)`$ is a probability amplitude characterized by an initial state corresponding to the quantum numbers $`(j_1,j_2,j,M)`$, and a final state corresponding to the eigenvalues $`(\theta _1,\phi _1,\theta _2,\phi _2).`$ In the initial state, $`j`$ is the quantum number for the total angular momentum, $`M\mathrm{}`$ is the projection of the total angular momentum along the $`z`$ direction, $`j_1`$ is the quantum number of subsystem $`1`$ and $`j_2`$ is the quantum number of subsystem $`2`$. In the state which results from the measurement, $`(\theta _1,\phi _1)`$ is the angular position of subsystem 1, while $`(\theta _2,\phi _2)`$ is the angular position of subsystem $`2`$. Thus, this probability amplitude gives the probability for obtaining specified angular positions of systems $`1`$ and $`2`$ upon measurement if the initial state of the compound system is defined by the quantum numbers $`(j_1,j_2,j,M)`$.
The Clebsch-Gordan coefficient $`\chi (j_1,j_2,j,M;j_1,m_1,j_2,m_2)`$ is a probability amplitude characterized by an initial state corresponding to $`(j_1,j_2,j,M)`$ and a final state defined by $`(j_1,m_1,j_2,m_2)`$. In the state resulting from the measurement, the angular momentum quantum number of subsystem $`1`$ is $`j_1`$, while its component in the $`z`$ direction is $`m_1\mathrm{}`$, and the angular momentum quantum number of subsystem $`2`$ is $`j_2`$, while its $`z`$ projection is $`m_2\mathrm{}`$. This probability amplitude thus gives the probability of obtaining specified projections of the angular momenta of the subsystems along the $`z`$ axis starting from a state of the compound system defined by the quantum numbers $`(j_1,j_2,j,M).`$
The function $`\mathrm{\Phi }(j_1,m_1,j_2,m_2;\theta _1,\phi _1,\theta _2,\phi _2)`$ is a probability amplitude with an initial state defined by $`(j_1,m_1,j_2,m_2)`$ and a final state defined by the eigenvalues $`(\theta _1,\phi _1,\theta _2,\phi _2).`$ This probability amplitude thus gives the probability of obtaining specified angular positions of the subsystems starting from a state characterized by specified projections of these subsystems along the $`z`$ direction.
In labelling the various probability amplitudes, we may reduce on the clutter by suppressing those quantum numbers which do not change at all during the measurement. Thus, we omit $`j_1`$ and $`j_2`$. However, we retain the subscript $`j`$ because for given $`j_1`$ and $`j_2`$, several values of $`j`$ are possible within the limits
$$j_1+j_2j\left|j_1j_2\right|.$$
(10)
With these changes, Eq. (7) becomes
$$\mathrm{\Psi }(j,M;\theta _1,\phi _1,\theta _2,\phi _2)=\underset{m_1}{}\chi (j,M;m_1,m_2)\mathrm{\Phi }(m_1,m_2;\theta _1,\phi _1,\theta _2,\phi _2).$$
(11)
We have elsewhere interpreted the probability amplitude $`\mathrm{\Psi }(j,M;\theta _1,\phi _1,\theta _2,\phi _2)`$ as a special form of the probability amplitude $`\mathrm{\Psi }(j(\theta ,\phi ),M;\theta _1,\phi _1,\theta _2,\phi _2)`$. The former quantity is specialized because it pertains to a situation where projections of the total angular momentum are measured with respect to the $`z`$ direction (for which $`\theta =\phi =0`$), while the latter corresponds to these projections being measured with respect to the arbitrary vector $`\widehat{𝐚}`$ whose polar angles are $`(\theta ,\phi )`$. To define the generalized probability amplitude corresponding to the latter case, we add the superscript $`\widehat{𝐚}`$ to $`M`$ . The expansion for the generalized probability amplitude is thus
$$\mathrm{\Psi }(j,M^{(\widehat{𝐚})};\theta _1,\phi _1,\theta _2,\phi _2)=\underset{m_1,m_2}{}\chi (j,M^{(\widehat{𝐚})};m_1,m_2)\mathrm{\Phi }(m_1,m_2;\theta _1,\phi _1,\theta _2,\phi _2).$$
(12)
In this expansion, the projections of the angular momenta of the subsystems are not necessarily measured with respect to the direction $`\widehat{𝐚}`$. In fact the projection of subsystem $`1`$ need not be measured with respect to the same vector as the projection of subsystem $`2`$. In general, the projection of the angular momentum of subsystem $`1`$ is measured relative to the direction $`\widehat{𝐠}_1`$ and that of subsystem $`2`$ with respect to the direction $`\widehat{𝐠}_1`$. For this reason, the condition that $`\chi (j,M^{(\widehat{𝐚})};m_1,m_2)`$ vanishes unless $`m_1+m_2=M`$ is generally not satisfied. Furthermore, the summation is generally a double summation, running over indices corresponding to the spin projections of subsystems $`1`$ and $`2`$ respectively. If the projections of subsystem $`1`$ with respect to $`\widehat{𝐠}_1`$ are $`(m_1)_\alpha ^{(\widehat{𝐠}_1)}`$, while those of subsystem $`2`$ with respect to $`\widehat{𝐠}_2`$ are $`(m_2)_\alpha ^{^{}}^{(\widehat{𝐠}_2)}`$, Eq. (12) is modified to
$`\mathrm{\Psi }(j,M^{(\widehat{𝐚})};\theta _1,\phi _1,\theta _2,\phi _2)={\displaystyle \underset{\alpha ,\alpha ^{}}{}}\chi (j,M^{(\widehat{𝐚})};(m_1)_\alpha ^{(\widehat{𝐠}_1)},(m_2)_\alpha ^{^{}}^{(\widehat{𝐠}_2)})`$
$`\times \mathrm{\Phi }((m_1)_\alpha ^{(\widehat{𝐠}_1)},(m_2)_\alpha ^{^{}}^{(\widehat{𝐠}_2)};\theta _1,\phi _1,\theta _2,\phi _2).`$ (13)
But in the special event that $`\widehat{𝐠}_1=\widehat{𝐠}_2=\widehat{𝐚}`$, then the functions $`\chi (j,M^{(\widehat{𝐚})};(m_1)_\alpha ^{(\widehat{𝐚})},(m_2)_\alpha ^{^{}}^{(\widehat{𝐚})})`$ are generalized Clebsch-Gordan coefficients and it is once more true that
$$(m_1)_\alpha ^{(\widehat{𝐚})}+(m_2)_\alpha ^{^{}}^{(\widehat{𝐚})})=M^{(\widehat{𝐚})}.$$
(14)
### 3.2 Theory for Spin
We now consider the case of spin. For a measurement on a simple system, the initial state is defined by a spin projection with respect to a given initial direction, while the final state is defined by a spin projection with respect to a new final direction. Thus, in the theory outlined in the previous section, probability amplitudes corresponding to spin projection measurements replace the spherical harmonics.
For a system of compounded spin, the total spin is
$$𝐒=𝐒_1+𝐒_2.$$
(15)
The quantum numbers of the spins are $`s`$, $`s_1`$ and $`s_2`$ for the total system, subsystem $`1`$ and subsystem $`2`$, respectively.
Suppose that the projections of the combined spin are initially known with respect to the direction of the vector $`\widehat{𝐚}`$, whose polar angles are $`(\theta ,\phi )`$. Let the projection of the total spin in that direction be $`M_i^{(\widehat{𝐚})}\mathrm{}`$. We proceed to measure the projection of the spin of subsystem $`1`$ with respect to the direction $`\widehat{𝐜}_1`$ (defined by the angles $`(\theta _1,\phi _1)`$) and the projection of the spin of subsystem $`2`$ with respect to the direction $`\widehat{𝐜}_2`$ (polar angles $`(\theta _2,\phi _2)`$). The projections that result from the measurement are identified by their corresponding quantum numbers and the vectors with respect to which they are measured. Thus, the probability amplitude for this measurement is $`\mathrm{\Psi }(s,M_i^{(\widehat{𝐚})};(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)})`$, where $`(u,v=1,2,\mathrm{}).`$ The generalized probability amplitude Eq. (13) is expressed as
$`\mathrm{\Psi }(j,M_i^{(\widehat{𝐚})};(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)})={\displaystyle \underset{\alpha ,\alpha ^{^{}}}{}}\chi (s,M_i^{(\widehat{𝐚})};(m_1)_\alpha ^{(\widehat{𝐠}_1)},(m_2)_\alpha ^{^{}}^{(\widehat{𝐠}_2)})`$
$`\times \mathrm{\Phi }((m_1)_\alpha ^{(\widehat{𝐠}_1)},(m_2)_\alpha ^{^{}}^{(\widehat{𝐠}_2)};(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)}).`$ (16)
Since the direction vectors $`\widehat{𝐠}_1`$ and $`\widehat{𝐠}_2`$ are arbitrary, they may be chosen for best convenience. The obvious choice is $`\widehat{𝐠}_1=\widehat{𝐠}_2=\widehat{𝐤}`$. We observe that if this is the case, and in addition $`\widehat{𝐚}=\widehat{𝐤}`$, then the $`\chi `$’s become Clebsch-Gordan coefficients. Since the Clebsch-Gordan coefficients vanish unless $`(m_1)^{(\widehat{𝐤})}+(m_2)^{(\widehat{𝐤})}=M^{(\widehat{𝐤})}`$, the double summation effectively becomes a single summation. In fact, if $`\widehat{𝐠}_1`$ and $`\widehat{𝐠}_2`$ are arbitrary, then the probability amplitudes $`\chi (s,M_i^{(\widehat{𝐚})};(m_1)_\alpha ^{(\widehat{𝐠}_1)},(m_2)_\alpha ^{^{}}^{(\widehat{𝐠}_2)})`$ and $`\mathrm{\Psi }(s,M_i^{(\widehat{𝐚})};(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)})`$ are essentially identical, since they differ only in the choice of arbitrary vectors along which the spin projections of subsystems $`1`$ and $`2`$ are measured. In practice, it is essential to make the choice $`\widehat{𝐠}_1=\widehat{𝐠}_2=\widehat{𝐤}`$, in order to convert $`\chi (s,M_i^{(\widehat{𝐚})};(m_1)_\alpha ^{(\widehat{𝐠}_1)},(m_2)_\alpha ^{^{}}^{(\widehat{𝐠}_2)})`$ to $`\chi (s,M_i^{(\widehat{𝐚})};(m_1)_\alpha ^{(\widehat{𝐤})},(m_2)_\alpha ^{}^{(\widehat{𝐤})})`$ which can be expressed in terms of Clebsch-Gordan coefficients . By this means it is possible to find an expression for $`\mathrm{\Psi }(s,M_i^{(\widehat{𝐚})};(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)})`$ (or equivalently $`\chi (s,M_i^{(\widehat{𝐚})};(m_1)_\alpha ^{(\widehat{𝐠}_1)},(m_2)_\alpha ^{^{}}^{(\widehat{𝐠}_2)})`$).
The actual form of the generalized spin probability amplitudes has been obtained in Ref. . It is
$`\mathrm{\Psi }(s,M_i^{(\widehat{𝐚})};(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)})=`$ (17)
$`{\displaystyle \underset{\alpha ,\alpha ^{}}{}}{\displaystyle \underset{l}{}}\zeta (s,M_i^{(\widehat{𝐚})};s,M_l^{(\widehat{𝐤})})\vartheta (s,M_l^{(\widehat{𝐤})};(m_1)_\alpha ^{(\widehat{𝐤})},(m_2)_\alpha ^{}^{(\widehat{𝐤})})`$
$`\times \mathrm{\Phi }((m_1)_\alpha ^{(\widehat{𝐤})},(m_2)_\alpha ^{}^{(\widehat{𝐤})};(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)}).`$
The quantities in expression (17) are defined as follows. The quantity $`\zeta (s,M_i^{(\widehat{𝐚})};s,M_l^{(\widehat{𝐤})})`$ is the probability amplitude that if the total spin is $`s`$ and its projection along the vector $`\widehat{𝐚}`$ is $`M_i^{(\widehat{𝐚})}\mathrm{}`$, a measurement of its projection along the $`z`$ axis gives $`M_l^{(\widehat{𝐤})}\mathrm{}.`$ The quantity $`\vartheta (s,M_l^{(\widehat{𝐤})};(m_1)_\alpha ^{(\widehat{𝐤})},(m_2)_\alpha ^{}^{(\widehat{𝐤})})`$ is actually the standard Clebsch-Gordan coefficient for the case at hand. It is obtained from $`\chi (s,M_i^{(\widehat{𝐚})};(m_1)_\alpha ^{(\widehat{𝐠}_1)},(m_2)_\alpha ^{^{}}^{(\widehat{𝐠}_2)})`$ by setting $`\widehat{𝐚}=\widehat{𝐠}_1=\widehat{𝐠}_2=\widehat{𝐤}`$.
For future convenience we rewrite Eq. (17) as
$`\mathrm{\Psi }(s,M_i^{(\widehat{𝐚})};(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)})={\displaystyle \underset{\alpha ,\alpha ^{}}{}}\chi (s,M_i^{(\widehat{𝐚})};(m_1)_\alpha ^{(\widehat{𝐤})},(m_2)_\alpha ^{}^{(\widehat{𝐤})})`$
$`\times \mathrm{\Phi }((m_1)_\alpha ^{(\widehat{𝐤})},(m_2)_\alpha ^{}^{(\widehat{𝐤})};(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)}),`$ (18)
where
$`\chi (s,M_i^{(\widehat{𝐚})};(m_1)_\alpha ^{(\widehat{𝐤})},(m_2)_\alpha ^{}^{(\widehat{𝐤})})={\displaystyle \underset{l}{}}\zeta (s,M_i^{(\widehat{𝐚})};s,M_l^{(\widehat{𝐤})})`$
$`\times \vartheta (s,M_l^{(\widehat{𝐤})};(m_1)_\alpha ^{(\widehat{𝐤})},(m_2)_\alpha ^{}^{(\widehat{𝐤})}).`$ (19)
If we compare Eq. (18) with the fundamental expansion Eq. (3), we see that the intermediate observable which we are using to achieve the expansion is the combination of spin projections of systems $`1`$ and $`2`$ with respect to the $`z`$ axis. The notation is simpler if we use the symbol $`B`$ for this observable. However, because of the double summation, the symbol has two subscripts. If the subsystems $`1`$ and $`2`$ are both spin-$`1/2`$ systems, the values of $`B`$ are
$$B_{11}=((m_1)_1^{(\widehat{𝐤})},(m_2)_1^{(𝐤)})=((+\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(𝐤)}),$$
(20)
$$B_{12}=((m_1)_1^{(\widehat{𝐤})},(m_2)_2^{(𝐤)})=((+\frac{1}{2})^{(\widehat{𝐤})},(\frac{1}{2})^{(𝐤)}),$$
(21)
$$B_{21}=((m_1)_2^{(\widehat{𝐤})},(m_2)_1^{(𝐤)})=((\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(𝐤)})$$
(22)
and
$$B_{22}=((m_1)_2^{(\widehat{𝐤})},(m_2)_2^{(𝐤)})=((\frac{1}{2})^{(\widehat{𝐤})},(\frac{1}{2})^{(𝐤)}).$$
(23)
This makes it easier to denote values pertaining to just one subsystem. For a particular value $`B_{\alpha \alpha ^{}}`$, we shall denote the value corresponding to the subsystem $`w`$ by $`(B_{\alpha \alpha ^{}})_w`$, where $`w=1,2`$. This means that for subsystem $`1`$,
$$(B_{11})_1=(B_{12})_1=(+\frac{1}{2})^{(\widehat{𝐤})}$$
(24)
and
$$(B_{21})_1=(B_{22})_1=(\frac{1}{2})^{(\widehat{𝐤})},$$
(25)
while for subsystem $`2`$,
$$(B_{11})_2=(B_{21})_2=(+\frac{1}{2})^{(\widehat{𝐤})}$$
(26)
and
$$(B_{12})_2=(B_{22})_2=(\frac{1}{2})^{(\widehat{𝐤})}.$$
(27)
For the sake of convenience, we set $`A_i=(s,M_i^{(\widehat{𝐚})})`$. Then the probability amplitude Eq. (18) becomes
$$\mathrm{\Psi }(A_i;(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)})=\underset{\alpha ,\alpha ^{}}{}\chi (A_i;B_{\alpha \alpha ^{}})\mathrm{\Phi }(B_{\alpha \alpha ^{}};(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)}).$$
(28)
We remind ourselves that according to Eq. (9),
$`\mathrm{\Phi }(B_{\alpha \alpha ^{}};(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)})=\mathrm{\Phi }((m_1)_\alpha ^{(\widehat{𝐤})},(m_2)_\alpha ^{}^{(\widehat{𝐤})};(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)})`$
$`=`$ $`\varphi _1((m_1)_\alpha ^{(\widehat{𝐤})};(m_1)_u^{(\widehat{𝐜}_1)})\varphi _2((m_2)_\alpha ^{}^{(𝐤)}:(m_2)_v^{(\widehat{𝐜}_2)}).`$ (29)
We remark that, purely for convenience, we have altered the notation slightly. Thus, $`\varphi _i=\varphi ^{(i)}`$ $`(i=1,2)`$.
## 4 From Probability-Amplitude Mechanics to Matrix Mechanics
### 4.1 Theory For Simple Systems
In order to move from wave to matrix mechanics, an expansion of the eigenfunction or wave function in terms of some basis set is necessary. To obtain the matrix theory of orbital angular momentum we use the spherical harmonics as the basis set. This kind of procedure was thought impossible for spin, because spin is not describable by eigenfunctions resulting from an eigenvalue equation. But in our work\[1, 7-11\], we have shown how, by using the probability amplitudes for measurements on spin systems, we can derive the matrix treatment of spin in the same way as the matrix treatment of orbital angular momentum is obtained.
The relation Eq. (3) is the basis of the transformation from amplitude to matrix mechanics. We now review how we use it to obtain the matrix treatment of a simple quantum system. This review is needed because it is the foundation of the more involved derivation of the standard matrix treatment of compounded spin from generalized probability amplitudes.
Let us suppose that have a quantum system possessing the observables $`A`$, $`B`$ and $`C`$. We assume that as we are measuring values of $`C`$, we are measuring values of a quantity $`T(C)`$ which is a function of $`C`$. Let $`T(C)`$ take upon measurement the values $`T_n`$ determined by the values $`C_n`$ of $`C`$. Thus $`T_n=T(C_n)`$. The expectation value of $`T`$ is
$$T(C)=\underset{n=1}{\overset{N}{}}\left|\psi (A_i;C_n)\right|^2T_n.$$
(30)
where $`N`$ is the total number of eigenvalues of $`C`$.
If we use the expansions
$$\psi ^{}(A_i;C_n)=\underset{j=1}{\overset{N}{}}\eta ^{}(A_i;B_j)\xi ^{}(B_j;C_n)$$
(31)
and
$$\psi (A_i;C_n)=\underset{j^{}=1}{\overset{N}{}}\eta (A_i;B_j^{})\xi (B_j^{};C_n),$$
(32)
we find
$$T(C)=\underset{j=1}{\overset{N}{}}\underset{j^{}=1}{\overset{N}{}}\eta ^{}(A_i;B_j)T_{jj^{}}\eta (A_i;B_j^{}),$$
(33)
where
$$T_{jj^{}}=\underset{n=1}{\overset{N}{}}\xi ^{}(B_j;C_n)T_n\xi (B_j^{};C_n).$$
(34)
Hence
$$T(C)=[\eta (A_i)]^{}[T][\eta (A_i)],$$
(35)
where the state is
$$[\eta (A_i)]=\left(\begin{array}{c}\eta (A_i;B_1)\\ \eta (A_i;B_2)\\ ..\\ \eta (A_i;B_N)\end{array}\right),$$
(36)
and the operator is
$$[T]=\left(\begin{array}{cccc}T_{11}& T_{12}& \mathrm{}& T_{1N}\\ T_{21}& T_{22}& \mathrm{}& T_{2N}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ T_{N1}& T_{N2}& \mathrm{}& T_{NN}\end{array}\right).$$
(37)
Thus, by means of the probability amplitudes for a quantum system, we can derive its matrix treatment. We note the convention of enclosing a quantity in brackets in order to denote its matrix representation.
### 4.2 Theory for Systems of Compounded Spin
The theory in the previous section will now be extended so as to yield the derivation of the matrix treatment of systems of compounded spin. As usual, we go through the expectation value in order to obtain the matrix form of the probability amplitudes for compounded spin, and of the operators for quantities that may be measured on the systems.
Suppose that a compounded spin is obtained by adding the spins $`s_1`$ and $`s_2`$. Suppose that initially the total spin is $`s`$ and its projection with respect to the vector $`\widehat{𝐚}`$ is $`M\mathrm{}`$. Subsequently, the spin projection of the spin of system $`1`$ is measured along the vector $`\widehat{𝐜}_1`$ and the spin projection of the spin of system $`2`$ is measured along the vector $`\widehat{𝐜}_2`$. At the same time, the quantity $`R((m_1)^{(\widehat{𝐜}_1)},(m_2)^{(\widehat{𝐜}_2)})`$ is measured. This quantity is measured on the separate systems $`1`$ and $`2`$. It is constructed from the quantity $`r^{(1)}((m_1)^{(\widehat{𝐜}_1)})`$, which is measured on system $`1`$ and the quantity $`r^{(2)}((m_2)^{(\widehat{𝐜}_2)})`$ measured on system $`2`$. Therefore, we write
$$R=R(r^{(1)}((m_1)^{(\widehat{𝐜}_1)}),r^{(2)}((m_2)^{(\widehat{𝐜}_2)}).$$
(38)
The values of $`r^{(1)}((m_1)^{(\widehat{𝐜}_1)})`$ are independent of the values of $`r^{(2)}((m_2)^{(\widehat{𝐜}_2)})`$: any value $`r^{(1)}((m_1)_u^{(\widehat{𝐜}_1)})`$ and any value $`r^{(2)}((m_2)_v^{(\widehat{𝐜}_2)})`$ can result together from the measurements. Therefore, the expectation value of $`R`$ is
$`R={\displaystyle \underset{u}{}}{\displaystyle \underset{v}{}}\mathrm{\Psi }^{}(A_i;(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)})\mathrm{\Psi }(A_i;(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)})`$
$`\times R(r^{(1)}((m_1)_u^{(\widehat{𝐜}_1)}),r^{(2)}((m_2)_v^{(\widehat{𝐜}_2)}),`$ (39)
where $`R(r^{(1)}((m_1)_u^{(\widehat{𝐜}_1)}),r^{(2)}((m_2)_v^{(\widehat{𝐜}_2)})`$ is an actual value of $`R((m_1)^{(\widehat{𝐜}_1)},(m_2)^{(\widehat{𝐜}_2)}).`$
The probability amplitude is given by Eq.(28). Using Eq. (29), we obtain the expansions
$`\mathrm{\Psi }^{}(A_i;(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)})={\displaystyle \underset{\alpha ,\alpha ^{}}{}}\chi ^{}(A_i;B_{\alpha \alpha ^{}})`$
$`\times \varphi _1^{}((B_{\alpha \alpha ^{}})_1;(m_1)_u^{(\widehat{𝐜}_1)})\varphi _2^{}((B_{\alpha \alpha ^{}})_2;(m_2)_v^{(\widehat{𝐜}_2)})`$ (40)
and
$`\mathrm{\Psi }(A_i;(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)})={\displaystyle \underset{\beta ,\beta ^{}}{}}\chi (A_i;B_{\beta \beta ^{}})`$
$`\times \varphi _1((B_{\beta \beta ^{}})_1;(m_1)_u^{(\widehat{𝐜}_1)})\varphi _2((B_{\beta \beta ^{}})_2;(m_2)_v^{(\widehat{𝐜}_2)}).`$ (41)
The expectation value becomes
$`R`$ $`=`$ $`{\displaystyle \underset{\alpha ,\alpha ^{}}{}}{\displaystyle \underset{\beta ,\beta ^{}}{}}\chi ^{}(A_i;B_{\alpha \alpha ^{}})\chi (A_i;B_{\beta \beta ^{}})`$
$`\times {\displaystyle \underset{u}{}}{\displaystyle \underset{v}{}}\{\varphi _1^{}((B_{\alpha \alpha ^{}})_1;(m_1)_u^{(\widehat{𝐜}_1)})\varphi _1((B_{\beta \beta ^{}})_1;(m_1)_u^{(\widehat{𝐜}_1)})`$
$`\times \varphi _2^{}((B_{\alpha \alpha ^{}})_2;(m_2)_v^{(\widehat{𝐜}_2)})\varphi _2((B_{\beta \beta ^{}})_2;(m_2)_v^{(\widehat{𝐜}_2)})`$ (42)
$`\times R(r^{(1)}((m_1)_u^{(\widehat{𝐜}_1)})r^{(2)}((m_2)_v^{(\widehat{𝐜}_2)})\}.`$
When $`R`$ is factorizable, so that
$$R(r^{(1)}((m_1)^{(\widehat{𝐜}_1)}),r^{(2)}((m_2)^{(\widehat{𝐜}_2)})=r^{(1)}((m_1)^{(\widehat{𝐜}_1)})r^{(2)}((m_2)^{(\widehat{𝐜}_2)}),$$
(43)
we can write Eq. (42) as
$`R={\displaystyle \underset{\alpha }{}}{\displaystyle \underset{\alpha ^{^{}}}{}}\chi ^{}(A_i;B_{\alpha \alpha ^{}})\chi (A_i;B_{\beta \beta ^{}}){\displaystyle \underset{u}{}}\varphi _1^{}((B_{\alpha \alpha ^{}})_1;(m_1)_u^{(\widehat{𝐜}_1)})`$
$`\times \varphi _1((B_{\beta \beta ^{}})_1;(m_1)_u^{(\widehat{𝐜}_1)})r^{(1)}((m_1)_u^{(\widehat{𝐜}_1)})`$
$`\times {\displaystyle \underset{v}{}}\varphi _2^{}((B_{\alpha \alpha ^{}})_2;(m_2)_v^{(\widehat{𝐜}_2)})\varphi _2((B_{\beta \beta ^{}})_2;(m_2)_v^{(\widehat{𝐜}_2)})r^{(2)}((m_2)_v^{(\widehat{𝐜}_2)}).`$ (44)
Thus
$$R=\underset{\alpha ,\alpha ^{}}{}\underset{\beta ,\beta ^{}}{}\chi ^{}(A_i;B_{\alpha \alpha ^{}})I_{\alpha \alpha ^{}\beta \beta ^{}}^{(1)}I_{\alpha \alpha ^{}\beta \beta ^{}}^{(2)}\chi (A_i;B_{\beta \beta ^{}}),$$
(45)
where
$`I_{\alpha \alpha ^{}\beta \beta ^{}}^{(1)}`$ $`=`$ $`{\displaystyle \underset{u}{}}\varphi _1^{}((B_{\alpha \alpha ^{}})_1;(m_1)_u^{(\widehat{𝐜}_1)})r^{(1)}((m_1)_u^{(\widehat{𝐜}_1)})`$ (46)
$`\times \varphi _1((B_{\beta \beta ^{}})_1;(m_1)_u^{(\widehat{𝐜}_1)})`$
and
$`I_{\alpha \alpha ^{}\beta \beta ^{}}^{(2)}`$ $`=`$ $`{\displaystyle \underset{v}{}}\varphi _2^{}((B_{\alpha \alpha ^{}})_2;(m_2)_v^{(\widehat{𝐜}_2)})r^{(2)}((m_2)_v^{(\widehat{𝐜}_2)})`$ (47)
$`\times \varphi _2((B_{\beta \beta ^{}})_2;(m_2)_v^{(\widehat{𝐜}_2)}).`$
In order to treat $`I_{\alpha \alpha ^{}\beta \beta ^{}}^{(1)}`$, we introduce the observable $`D`$ which corresponds to spin projections of subsystem $`1`$ with respect to the direction $`\widehat{𝐝\text{,}}`$ whose polar angles are $`(\theta _d,\phi _d)`$. We use this observable to expand $`\varphi _1`$ by means of formula (3). We note that the values $`D_p`$ are
$$D_p=(m_1)_p^{(\widehat{𝐝})}\mathrm{}.$$
(48)
The expansions of $`\varphi _1`$ and $`\varphi _1^{}`$ are
$$\varphi _1^{}((B_{\alpha \alpha ^{}})_1;(m_1)_u^{(\widehat{𝐜}_1)})=\underset{p}{}\eta _1^{}((B_{\alpha \alpha ^{}})_1;D_p)\xi _1^{}(D_p;(m_1)_u^{(\widehat{𝐜}_1)})$$
(49)
and
$$\varphi _1((B_{\beta \beta ^{}})_1;(m_1)_u^{(\widehat{𝐜}_1)})=\underset{p^{^{}}}{}\eta _1((B_{\beta \beta ^{}})_1;D_p^{^{}})\xi _1(D_p^{^{}};(m_1)_u^{(\widehat{𝐜}_1)}).$$
(50)
Applying the theory outlined in Section $`4.1`$, we obtain
$$I_{\alpha \alpha ^{}\beta \beta ^{}}^{(1)}=\underset{p}{}\underset{p^{^{}}}{}\eta _1^{}((B_{\alpha \alpha ^{}})_1;D_p)r_{pp^{^{}}}^{(1)}\eta _1((B_{\beta \beta ^{}})_1;D_p^{^{}}),$$
(51)
where
$$r_{pp^{^{}}}^{(1)}=\underset{u}{}\xi _1^{}(D_p;(m_1)_u^{(\widehat{𝐜}_1)})r^{(1)}((m_1)_u^{(\widehat{𝐜}_1)})\xi _1(D_p^{^{}};(m_1)_u^{(\widehat{𝐜}_1)}).$$
(52)
Hence
$$I_{\alpha \alpha ^{}\beta \beta ^{}}^{(1)}=[\eta _1((B_{\alpha \alpha ^{}})_1)]^{}[r^{(1)}][\eta _1((B_{\beta \beta ^{}})_1)],$$
(53)
where
$$[\eta _1(B_{\alpha \alpha ^{}})]=\left(\begin{array}{c}\eta _1((B_{\alpha \alpha ^{}})_1;D_1)\\ \eta _1((B_{\alpha \alpha ^{}})_1;D_2)\\ :\\ \eta _1((B_{\alpha \alpha ^{}})_1;D_N)\end{array}\right)$$
(54)
and
$$[r^{(1)}]=\left(\begin{array}{cccc}r_{11}^{(1)}& r_{12}^{(1)}& ..& r_{1N}^{(1)}\\ r_{21}^{(1)}& r_{22}^{(1)}& ..& r_{2N}^{(1)}\\ ..& ..& ..& ..\\ r_{N1}^{(1)}& r_{N2}^{(1)}& ..& r_{NN}^{(1)}\end{array}\right).$$
(55)
Here $`N`$ is the total number of states of the observable $`D`$. In the case where subsystem $`1`$ is a spin-$`1/2`$ system, $`N=2`$.
To deal with $`I_{\alpha \alpha ^{}\beta \beta ^{}}^{(2)}`$, we introduce the observable $`F`$ defined by the unit vector $`\widehat{𝐟}`$ whose polar angles are $`(\theta _f,\phi _f)`$. The eigenvalues of $`F`$ are
$$F_q=(m_2)_q^{(\widehat{𝐟})}\mathrm{}.$$
(56)
Expanding $`\varphi _2`$ over the states of $`F`$, we get
$$\varphi _2^{}(B_{\alpha \alpha ^{}};(m_2)_v^{(\widehat{𝐜}_2)})=\underset{q}{}\eta _2^{}((B_{\alpha \alpha ^{}})_2;F_q)\xi _2^{}(F_q;(m_2)_v^{(\widehat{𝐜}_2)})$$
(57)
and
$$\varphi _2(B_{\beta \beta ^{}};(m_2)_v^{(\widehat{𝐜}_2)})=\underset{q^{^{}}}{}\eta _2((B_{\beta \beta ^{}})_2;F_q^{^{}})\xi _2(F_q^{^{}};(m_2)_v^{(\widehat{𝐜}_2)}).$$
(58)
This means that
$$I_{\alpha \alpha ^{}\beta \beta ^{}}^{(2)}=\underset{q}{}\underset{q^{^{}}}{}\eta _2^{}((B_{\alpha \alpha ^{}})_2;F_q)r_{qq^{^{}}}^{(2)}\eta _2((B_{\beta \beta ^{}})_2;F_q^{^{}}),$$
(59)
where
$$r_{qq^{^{}}}^{(2)}=\underset{v}{}\xi _2^{}(F_q;(m_2)_v^{(\widehat{𝐜}_2)})r^{(2)}((m_2)_v^{(\widehat{𝐜}_2)})\xi _2(F_q^{^{}};(m_2)_v^{(\widehat{𝐜}_2)}).$$
(60)
In view of Eq. (59), the expression for $`I_{\alpha \alpha ^{}\beta \beta ^{}}^{(2)}`$ is
$$I_{\alpha \alpha ^{}\beta \beta ^{}}^{(2)}=[\eta _2((B_{\alpha \alpha ^{}})_2)]^{}[r^{(2)}][\eta _2((B_{\beta \beta ^{}})_2)],$$
(61)
where
$$[\eta _2((B_{\alpha \alpha ^{}})_2)]=\left(\begin{array}{c}\eta _2((B_{\alpha \alpha ^{}})_2;F_1)\\ \eta _2((B_{\alpha \alpha ^{}})_2;F_2)\\ :\\ \eta _2((B_{\alpha \alpha ^{}})_2;F_M)\end{array}\right)$$
(62)
and
$$[r^{(2)}]=\left(\begin{array}{cccc}r_{11}^{(2)}& r_{12}^{(2)}& ..& r_{1M}^{(2)}\\ r_{21}^{(2)}& r_{22}^{(2)}& ..& r_{2M}^{(2)}\\ ..& ..& ..& ..\\ r_{M1}^{(2)}& r_{M2}^{(2)}& ..& r_{MM}^{(2)}\end{array}\right).$$
(63)
Here $`M`$ is the number of states of $`F`$.
Collecting all the results together, we find that
$`R={\displaystyle \underset{\alpha ,\alpha ^{}}{}}{\displaystyle \underset{\beta ,\beta ^{}}{}}\chi ^{}(A_i;B_{\alpha \alpha ^{}})[\eta _1((B_{\alpha \alpha ^{}})_1)]^{}[r^{(1)}][\eta _1((B_{\beta \beta ^{}})_1)]`$
$`\times [\eta _2((B_{\alpha \alpha ^{}})_2)]^{}[r^{(2)}][\eta _2((B_{\beta \beta ^{}})_2)]\chi (A_i;B_{\beta \beta ^{}})`$
$`=`$ $`\left({\displaystyle \underset{\alpha ,\alpha ^{}}{}}\chi ^{}(A_i;B_{\alpha \alpha ^{}})[\eta _1((B_{\alpha \alpha ^{}})_1)]^{}[\eta _2((B_{\alpha \alpha ^{}})_2)]^{}\right)[r^{(1)}][r^{(2)}]`$
$`\times \left({\displaystyle \underset{\beta ,\beta ^{}}{}}\chi (A_i;B_{\beta \beta ^{}})[\eta _1((B_{\beta \beta ^{}})_1)][\eta _2((B_{\beta \beta ^{}})_2)]\right)`$
$`=`$ $`[\mathrm{\Psi }(A_i;(m_1)^{(\widehat{𝐜}_1)},(m_2)^{(\widehat{𝐜}_2)})]^{}[R(r^{(1)}((m_1)^{(\widehat{𝐜}_1)}),r^{(2)}((m_2)^{(\widehat{𝐜}_2)}))]`$ (64)
$`\times [\mathrm{\Psi }(A_i;(m_1)^{(\widehat{𝐜}_1)},(m_2)^{(\widehat{𝐜}_2)})],`$
where
$$[\mathrm{\Psi }(A_i;(m_1)^{(\widehat{𝐜}_1)},(m_2)^{(\widehat{𝐜}_2)})]=\left(\underset{\alpha ,\alpha ^{}}{}\chi (A_i;B_{\alpha \alpha ^{}})[\eta _1((B_{\alpha \alpha ^{}})_1)][\eta _2((B_{\alpha \alpha ^{}})_2)]\right)$$
(65)
and
$$[R(r^{(1)}((m_1)^{\widehat{𝐜}_1}),r^{(2)}((m_2)^{(\widehat{𝐜}_2)}))]=[r^{(1)}][r^{(2)}].$$
(66)
We see that in Eq. (65) we have obtained the generalized vector state corresponding to the probability amplitude $`\mathrm{\Psi }(A_i;(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)}).`$ In addition, in Eqs. (52), 55, (60) and (63), we have obtained the generalized operator for any observable of the system which is a function of the spin projections of the subsystems.
We note that by definition, the vectors relating to subsystem $`1`$ act only on one another and on the operator corresponding to this subsystem. The same holds for the quantities corresponding to subsystem $`2`$. To emphasize this fact, we introduce the labels $`1`$ and $`2`$ to distinguish the corresponding quantities. Thus, we get
$$[\mathrm{\Psi }(A_i;(m_1)^{(\widehat{𝐜}_1)},(m_2)^{(\widehat{𝐜}_2)}]=\left(\underset{\alpha ,\alpha ^{}}{}\chi (A_i;B_{\alpha \alpha ^{}})[\eta _1((B_{\alpha \alpha ^{}})_1)]_1[\eta _2((B_{\alpha \alpha ^{}})_2)]_2\right)$$
(67)
and
$$[R(r^{(1)}((m_1)^{(\widehat{𝐜}_1)}),r^{(2)}((m_2)^{(\widehat{𝐜}_2)})]=[r^{(1)}((m_1)^{(\widehat{𝐜}_1)})]_1[r^{(2)}((m_2)^{(\widehat{𝐜}_2)})]_2.$$
(68)
These results are dependent on the condition that $`R`$ is factorizable; if this is not the case, it is not so easy to transform Eq. (42) to matrix form. But in that case, we can use the alternative approach of Ref. ; we then end up with $`3`$\- or $`4`$\- dimensional matrix representations which appply whether $`R`$ is factorizable or not.
## 5 Application to Actual Systems
### 5.1 The Matrix Operator
The results derived in the last section will now be used to obtain specific operators and vectors. For a particular case, the operator is calculated by explicitly working out the matrix elements $`r_{pp^{^{}}}^{(1)}`$ and $`r_{qq^{^{}}}^{(2)}`$ . The cases at hand are of a system of total spin $`0`$, and of a system of total spin $`1`$ (with three possible values of the magnetic quantum number) obtained by adding two spins of $`1/2`$ each. The form of the operator is independent of whether the total spin is $`0`$ or $`1`$. We therefore derive this quantity first.
We start with the matrix element $`r_{pp^{^{}}}^{(1)}.`$ In order to obtain this quantity, we require the forms of the probability amplitudes $`\xi `$. Since both systems $`1`$ and $`2`$ are spin-$`1/2`$ systems, these are obtained from the generalized spin-$`1/2`$ amplitudes, whose explicit forms we have already worked out .
We first recall the details of the probability amplitudes for spin $`1/2`$. We consider system $`1`$. Let the spin projection be initially known with respect to $`𝐝`$; it is subsequently measured with respect to $`\widehat{𝐜}_1`$. The probability amplitude that it will be found upon measurement to be up with respect to $`\widehat{𝐜}_1`$ is $`\xi _1((+\frac{1}{2})^{(\widehat{𝐝})};(+\frac{1}{2})^{(\widehat{𝐜}_1)})`$. The other three probability amplitudes are therefore $`\xi _1((+\frac{1}{2})^{(\widehat{𝐝})};(\frac{1}{2})^{(\widehat{𝐜}_1)})`$, $`\xi _1((\frac{1}{2})^{(\widehat{𝐝})};(+\frac{1}{2})^{(\widehat{𝐜}_1)})`$ and $`\xi _1((\frac{1}{2})^{(\widehat{𝐝})};(\frac{1}{2})^{(\widehat{𝐜}_1)}).`$ These probability amplitudes come in a variety of forms, depending on the phase choice made when they are being derived . One form is the following:
$$\xi _1((+\frac{1}{2})^{(\widehat{𝐝})};(+\frac{1}{2})^{(\widehat{𝐜}_1)})=\mathrm{cos}\theta _d/2\mathrm{cos}\theta _1/2+e^{i(\phi _d\phi _1)}\mathrm{sin}\theta _d/2\mathrm{sin}\theta _1/2,$$
(69)
$$\xi _1((+\frac{1}{2})^{(\widehat{𝐝})};(\frac{1}{2})^{(\widehat{𝐜}_1)})=\mathrm{cos}\theta _d/2\mathrm{sin}\theta _1/2+e^{i(\phi _d\phi _1)}\mathrm{sin}\theta _d/2\mathrm{cos}\theta _1/2,$$
(70)
$$\xi _1((\frac{1}{2})^{(\widehat{𝐝})};(+\frac{1}{2})^{(\widehat{𝐜}_1)})=\mathrm{sin}\theta _d/2\mathrm{cos}\theta _1/2+e^{i(\phi _d\phi _1)}\mathrm{cos}\theta _d/2\mathrm{sin}\theta _1/2$$
(71)
and
$$\xi _1((\frac{1}{2})^{(\widehat{𝐝})};(\frac{1}{2})^{(\widehat{𝐜}_1)})=\mathrm{sin}\theta _d/2\mathrm{sin}\theta _1/2+e^{i(\phi _d\phi _1)}\mathrm{cos}\theta _d/2\mathrm{cos}\theta _1/2.$$
(72)
Since in this case system $`2`$ is also a spin-$`1/2`$ system, the probability amplitudes corresponding to it are identical in form to Eqs. (69) - (72). To obtain them, we merely make the following change to the labels: $`2`$ replaces subscript $`1`$; $`\widehat{𝐟}`$ replaces $`\widehat{𝐝}`$, so that $`f`$ replaces $`d`$; and $`\widehat{𝐜}_2`$ replaces $`\widehat{𝐜}_1`$.
For the case of spin $`1/2`$, the summation over $`u`$ which appears in the expression for $`r_{pp^{^{}}}^{(1)}`$ contains only two terms. $`u=1`$ corresponds to the outcome $`(+\frac{1}{2})^{(\widehat{𝐜}_1)}`$ while $`u=2`$ corresponds to $`(\frac{1}{2})^{(\widehat{𝐜}_1)}.`$ Thus
$`r_{11}^{(1)}`$ $`=`$ $`\left|\xi _1((+{\displaystyle \frac{1}{2}})^{(\widehat{𝐝})};(+{\displaystyle \frac{1}{2}})^{(\widehat{𝐜}_1)})\right|^2r^{(1)}((+{\displaystyle \frac{1}{2}})^{(\widehat{𝐜}_1)})`$ (73)
$`+\left|\xi _1((+{\displaystyle \frac{1}{2}})^{(\widehat{𝐝})};({\displaystyle \frac{1}{2}})^{(\widehat{𝐜}_1)})\right|^2r^{(1)}(({\displaystyle \frac{1}{2}})^{(\widehat{𝐜}_1)}).`$
The values of the summation indices $`p`$ and $`p^{^{}}`$are such that $`p,p^{^{}}=+1`$ corresponds to $`(+\frac{1}{2})^{(\widehat{𝐝})}`$, while $`p,p^{^{}}=2`$ corresponds to $`(\frac{1}{2})^{(\widehat{𝐝})}`$. Hence
$`r_{11}^{(1)}=[\mathrm{cos}^2(\theta _d\theta _1)/2\mathrm{sin}\theta _d\mathrm{sin}\theta _1\mathrm{sin}^2(\phi _d\phi _1)/2]r^{(1)}((+{\displaystyle \frac{1}{2}})^{(\widehat{𝐜}_1)})`$ (74)
$`+[\mathrm{sin}^2(\theta _d\theta _1)/2+\mathrm{sin}\theta _d\mathrm{sin}\theta _1\mathrm{sin}^2(\phi _d\phi _1)/2]r^{(1)}(({\displaystyle \frac{1}{2}})^{(\widehat{𝐜}_1)}).`$
Similarly,
$`r_{12}^{(1)}=\xi _1^{}((+{\displaystyle \frac{1}{2}})^{(\widehat{𝐝})};(+{\displaystyle \frac{1}{2}})^{(\widehat{𝐜}_1)})\xi _1(({\displaystyle \frac{1}{2}})^{(\widehat{𝐝})};(+{\displaystyle \frac{1}{2}})^{(\widehat{𝐜}_1)})r^{(1)}((+{\displaystyle \frac{1}{2}})^{(\widehat{𝐜}_1)})`$
$`+\xi _1^{}((+{\displaystyle \frac{1}{2}})^{(\widehat{𝐝})};({\displaystyle \frac{1}{2}})^{(\widehat{𝐜}_1)})\xi _1(({\displaystyle \frac{1}{2}})^{(\widehat{𝐝})};({\displaystyle \frac{1}{2}})^{(\widehat{𝐜}_1)})r^{(1)}(({\displaystyle \frac{1}{2}})^{(\widehat{𝐜}_1)})`$
$`=`$ $`[{\displaystyle \frac{1}{2}}\mathrm{sin}\theta _d\mathrm{cos}\theta _1+{\displaystyle \frac{1}{2}}\mathrm{sin}\theta _1\mathrm{cos}\theta _d\mathrm{cos}(\phi _d\phi _1)`$ (75)
$`+{\displaystyle \frac{i}{2}}\mathrm{sin}\theta _1\mathrm{sin}(\phi _d\phi _1)]r^{(1)}((+{\displaystyle \frac{1}{2}})^{(\widehat{𝐜}_1)})`$
$`+[{\displaystyle \frac{1}{2}}\mathrm{sin}\theta _d\mathrm{cos}\theta _1{\displaystyle \frac{1}{2}}\mathrm{sin}\theta _1\mathrm{cos}\theta _d\mathrm{cos}(\phi _d\phi _1)`$
$`{\displaystyle \frac{i}{2}}\mathrm{sin}\theta _1\mathrm{sin}(\phi _d\phi _1)]r^{(1)}(({\displaystyle \frac{1}{2}})^{(\widehat{𝐜}_1)}),`$
$$r_{21}^{(1)}=r_{12}^{(1)}$$
(76)
and
$`r_{22}^{(1)}=[\mathrm{sin}^2(\theta _d\theta _1)/2+\mathrm{sin}\theta _d\mathrm{sin}\theta _1\mathrm{sin}^2(\phi _d\phi _1)/2]r^{(1)}((+{\displaystyle \frac{1}{2}})^{(\widehat{𝐜}_1)})`$
$`+[\mathrm{cos}^2(\theta _d\theta _1)/2\mathrm{sin}\theta _d\mathrm{sin}\theta _1\mathrm{sin}^2(\phi _d\phi _1)/2]r^{(1)}(({\displaystyle \frac{1}{2}})^{(\widehat{𝐜}_1)}).`$ (77)
The elements of $`[r^{(2)}]_2`$ are identical in form to those of $`[r^{(1)}]_1.`$ The difference is that in system $`2`$ the vector $`\widehat{𝐟}`$ plays the role that the vector $`\widehat{𝐝}`$ plays in system $`1`$. Thus, wherever $`d`$ appears, it is replaced by $`f`$. Wherever the label $`1`$ appears, it is replaced by $`2`$. Wherever $`(\pm \frac{1}{2})^{(\widehat{𝐜}_1)}`$ appears, it is replaced by $`(\pm \frac{1}{2})^{(\widehat{𝐜}_2)}.`$ Finally wherever $`r^{(1)}`$ appears, it is replaced by $`r^{(2)}`$. Thus, the elements of $`[r^{(2)}]_2`$ are
$`r_{11}^{(2)}=[\mathrm{cos}^2(\theta _f\theta _2)/2\mathrm{sin}\theta _d\mathrm{sin}\theta _1\mathrm{sin}^2(\phi _f\phi _2)/2]r^{(2)}((+{\displaystyle \frac{1}{2}})^{(\widehat{𝐜}_2)})`$
$`+[\mathrm{sin}^2(\theta _f\theta _2)/2+\mathrm{sin}\theta _d\mathrm{sin}\theta _1\mathrm{sin}^2(\phi _f\phi _2)/2]r^{(2)}(({\displaystyle \frac{1}{2}})^{(\widehat{𝐜}_2)}),`$ (78)
$`r_{12}^{(2)}=[{\displaystyle \frac{1}{2}}\mathrm{sin}\theta _f\mathrm{cos}\theta _2+{\displaystyle \frac{1}{2}}\mathrm{sin}\theta _2\mathrm{cos}\theta _f\mathrm{cos}(\phi _f\phi _2)`$ (79)
$`+{\displaystyle \frac{i}{2}}\mathrm{sin}\theta _2\mathrm{sin}(\phi _f\phi _2)]r^{(2)}((+{\displaystyle \frac{1}{2}})^{(\widehat{𝐜}_2)})`$
$`+[{\displaystyle \frac{1}{2}}\mathrm{sin}\theta _f\mathrm{cos}\theta _2{\displaystyle \frac{1}{2}}\mathrm{sin}\theta _2\mathrm{cos}\theta _f\mathrm{cos}(\phi _f\phi _2)`$
$`{\displaystyle \frac{i}{2}}\mathrm{sin}\theta _2\mathrm{sin}(\phi _f\phi _2)]r^{(2)}(({\displaystyle \frac{1}{2}})^{(\widehat{𝐜}_2)}),`$
$$r_{21}^{(2)}=r_{12}^{(2)}$$
(80)
and
$`r_{22}^{(2)}=[\mathrm{sin}^2(\theta _f\theta _2)/2+\mathrm{sin}\theta _f\mathrm{sin}\theta _2\mathrm{sin}^2(\phi _f\phi _2)/2]r^{(2)}((+{\displaystyle \frac{1}{2}})^{(\widehat{𝐜}_2)})`$
$`+[\mathrm{cos}^2(\theta _f\theta _2)/2\mathrm{sin}\theta _f\mathrm{sin}\theta _2\mathrm{sin}^2(\phi _f\phi _2)/2]r^{(2)}(({\displaystyle \frac{1}{2}})^{(\widehat{𝐜}_2)}).`$ (81)
### 5.2 The Vector States
#### 5.2.1 The Triplet State
We start our calculation of the states by looking at the probability amplitudes corresponding to the triplet state, defined by the quantum numbers $`s=1,`$ $`M^{(\widehat{𝐚})}=0,\pm 1.`$ We first consider the case $`M^{(\widehat{𝐚})}=1.`$
##### The $`M^{(\widehat{𝐚})}=1`$ State
For this case, we write,
$$\mathrm{\Psi }(1,1^{(\widehat{𝐚})};(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)})=\mathrm{\Psi }(s=1,M=1^{(\widehat{𝐚})};(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)}),$$
(82)
and the generalized probability amplitude is
$`\mathrm{\Psi }(1,1^{(\widehat{𝐚})};(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)})=`$
$`{\displaystyle \underset{\alpha ,\alpha ^{}}{}}\left[{\displaystyle \underset{l}{}}\zeta (1,1^{(\widehat{𝐚})};1,M_l^{(\widehat{𝐤})})\vartheta (1,M_l^{(\widehat{𝐤})};(m_1)_\alpha ^{(\widehat{𝐤})},(m_2)_\alpha ^{}^{(\widehat{𝐤})})\right]`$
$`\times \mathrm{\Phi }((m_1)_\alpha ^{(\widehat{𝐤})},(m_2)_\alpha ^{}^{(\widehat{𝐤})});(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)}).`$ (83)
Thus,
$`\mathrm{\Psi }(1,1^{(\widehat{𝐚})};(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)})={\displaystyle \underset{\alpha ,\alpha ^{}}{}}\chi (1,1^{(\widehat{𝐚})};B_{\alpha \alpha ^{}})`$
$`\times \mathrm{\Phi }(B_{\alpha \alpha ^{}};(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)}),`$ (84)
where
$$\chi (1,1^{(\widehat{𝐚})};B_{\alpha \alpha ^{}})=\underset{l}{}\zeta (1,1^{(\widehat{𝐚})};1,M_l^{(\widehat{𝐤})})\vartheta (1,M_l^{(\widehat{𝐤})};B_{\alpha \alpha ^{}}).$$
(85)
With the probability amplitude expressed in the form Eq. (84), which is identical in form to Eq. (28), the transformation to matrix form is straightforwardly achieved. We find that the matrix form of the probability amplitude is
$$[\mathrm{\Psi }(1,1^{(\widehat{𝐚})};(m_1)^{(\widehat{𝐜}_1)},(m_2)^{(\widehat{𝐜}_2)})]=\underset{\alpha ,\alpha ^{}}{}\chi (1,1^{(\widehat{𝐚})};B_{\alpha \alpha ^{}})[\eta _1((B_{\alpha \alpha ^{}})_1)]_1[\eta _2((B_{\alpha \alpha ^{}})_2)]_2.$$
(86)
As $`B_{\alpha \alpha ^{}}`$ takes the values Eqs. (20) - (23), we have
$`[\mathrm{\Psi }(1,1^{(\widehat{𝐚})};(m_1)^{(\widehat{𝐜}_1)},(m_2)^{(\widehat{𝐜}_2)})]=`$ $`\chi (1,1^{(\widehat{𝐚})};(+\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})})`$
$`\times [\eta _1((+\frac{1}{2})^{(\widehat{𝐤})})]_1[\eta _2((+\frac{1}{2})^{(\widehat{𝐤})})]_2`$
$`+\chi (1,1^{(\widehat{𝐚})};(+\frac{1}{2})^{(\widehat{𝐤})},(\frac{1}{2})^{(\widehat{𝐤})})[\eta _1((+\frac{1}{2})^{(\widehat{𝐤})})]_1[\eta _2((\frac{1}{2})^{(\widehat{𝐤})})]_2`$
$`+\chi (1,1^{(\widehat{𝐚})};(\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})})[\eta _1((\frac{1}{2})^{(\widehat{𝐤})})]_1[\eta _2((+\frac{1}{2})^{(\widehat{𝐤})})]_2`$
$`+\chi (1,1^{(\widehat{𝐚})};(\frac{1}{2})^{(\widehat{𝐤})},(\frac{1}{2})^{(\widehat{𝐤})})[\eta _1((\frac{1}{2})^{(\widehat{𝐤})})]_1[\eta _2((\frac{1}{2})^{(\widehat{𝐤})})]_2,`$
where
$$[\eta _1((\pm \frac{1}{2})^{(\widehat{𝐤})})]_1=\left(\begin{array}{c}\eta _1((\pm \frac{1}{2})^{(\widehat{𝐤})};(+\frac{1}{2})^{(\widehat{𝐝})})\\ \eta _1((\pm \frac{1}{2})^{(\widehat{𝐤})};(\frac{1}{2})^{(\widehat{𝐝})})\end{array}\right)_1$$
(88)
and
$$[\eta _2((\pm \frac{1}{2})^{(\widehat{𝐤})})]_2=\left(\begin{array}{c}\eta _2((\pm \frac{1}{2})^{(\widehat{𝐤})});(+\frac{1}{2})^{(\widehat{𝐟})})\\ \eta _2((\pm \frac{1}{2})^{(\widehat{𝐤})});(\frac{1}{2})^{(\widehat{𝐟})})\end{array}\right)_2.$$
(89)
The $`\eta _1`$’s and $`\eta _2`$’s are known. They are just the spin-$`1/2`$ probability amplitudes and are essentially identical to the $`\xi `$’s, Eqs. (69) - (72). The only difference is in the direction vectors. In the labelling of the arguments for the $`\eta ^{}`$s, the initial direction corresponds to the $`z`$ axis, so that its direction vector is $`\widehat{𝐤}`$; the final directions are defined by $`\widehat{𝐝}`$ and $`\widehat{𝐟}`$ for system $`1`$ and system $`2`$ respectively. From Eqs. (69) -(72), with the arguments appropriately changed, we get
$$[\eta _1((+\frac{1}{2})^{(\widehat{𝐤})}]_1=\left(\begin{array}{c}\mathrm{cos}\theta _d/2\\ \mathrm{sin}\theta _d/2\end{array}\right)_1\text{}$$
(90)
$$\text{ }[\eta _1(\frac{1}{2})^{(\widehat{\text{k}})}]_1=\left(\begin{array}{c}\mathrm{sin}\theta _d/2e^{i\phi _d}\\ \mathrm{cos}\theta _d/2e^{i\phi _d}\end{array}\right)_1,$$
(91)
$$[\eta _2(+\frac{1}{2})^{(\widehat{𝐤})}]_2=\left(\begin{array}{c}\mathrm{cos}\theta _f/2\\ \mathrm{sin}\theta _f/2\end{array}\right)_2$$
(92)
and
$$\text{ }[\eta _2((\frac{1}{2})^{(\widehat{\text{k}})})]_2=\left(\begin{array}{c}\mathrm{sin}\theta _f/2e^{i\phi _f}\\ \mathrm{cos}\theta _f/2e^{i\phi _f}\end{array}\right)_2.$$
(93)
It only remains to compute the $`\chi `$’s. According to Eq. (85)
$`\chi (1,1^{(\widehat{𝐚})};(+\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})})=\chi (1,1^{(\widehat{𝐚})};(+\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})})`$ (94)
$`=`$ $`\zeta (1,1^{(\widehat{𝐚})};1,1^{(\widehat{𝐤})})\vartheta (1,1^{(\widehat{𝐤})};(+\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})})`$
$`+\zeta (1,1^{(\widehat{𝐚})};1,0^{(\widehat{𝐤})})\vartheta (1,0^{(\widehat{𝐤})};(+\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})})`$
$`+\zeta (1,1^{(\widehat{𝐚})};1,(1)^{(\widehat{𝐤})})\vartheta (1,(1)^{(\widehat{𝐤})};(+\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})}).`$
The angles defining $`\widehat{𝐚}`$ are $`(\theta ,\phi ).`$ As we have shown ,
$$\zeta (1,1^{(\widehat{𝐚})};1,1^{(\widehat{𝐤})})=\mathrm{cos}^2\theta /2e^{i\phi }$$
(95)
$$\zeta (1,1^{(\widehat{𝐚})};1,0^{(\widehat{𝐤})})=\frac{1}{\sqrt{2}}\mathrm{sin}\theta $$
(96)
and
$$\zeta (1,1^{(\widehat{𝐚})};1,(1)^{(\widehat{𝐤})})=\mathrm{sin}^2\theta /2e^{i\phi }.$$
(97)
The $`\vartheta `$’s are Clebsch-Gordan coefficients. As a result
$$\vartheta (1,1^{(\widehat{𝐤})};(+\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})})=C(\frac{1}{2}\frac{1}{2}1;\frac{1}{2}\frac{1}{2}1)=1,$$
(98)
$$\vartheta (1,0^{(\widehat{𝐤})};(+\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})})=C(\frac{1}{2}\frac{1}{2}1;\frac{1}{2}\frac{1}{2}0)=0$$
(99)
and
$$\vartheta (1,(1)^{(\widehat{𝐤})};(+\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})})=C(\frac{1}{2}\frac{1}{2}1;\frac{1}{2}\frac{1}{2}1)=0.$$
(100)
Hence,
$$\chi (1,1^{(\widehat{𝐚})};(+\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})})=\frac{1}{2}\mathrm{sin}\theta .$$
(101)
Similarly,
$`\chi (1,1^{(\widehat{𝐚})};(+\frac{1}{2})^{(\widehat{𝐤})},(\frac{1}{2})^{(\widehat{𝐤})})=\zeta (1,1^{(\widehat{𝐚})};1,1^{(\widehat{𝐤})})\vartheta (1,1^{(\widehat{𝐤})};(+\frac{1}{2})^{(\widehat{𝐤})},(\frac{1}{2})^{(\widehat{𝐤})})`$
$`+\zeta (1,1^{(\widehat{𝐚})};1,0^{(\widehat{𝐤})})\vartheta (1,0^{(\widehat{𝐤})};(+\frac{1}{2})^{(\widehat{𝐤})},(\frac{1}{2})^{(\widehat{𝐤})})`$
$`+\zeta (1,1^{(\widehat{𝐚})};1,(1)^{(\widehat{𝐤})})\vartheta (1,(1)^{(\widehat{𝐤})};(+\frac{1}{2})^{(\widehat{𝐤})},(\frac{1}{2})^{(\widehat{𝐤})}).`$ (102)
This is the same as the expression for $`\chi (1,1^{(\widehat{𝐚})};(+\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})})`$, except for the change in the $`\vartheta `$’s. Since
$$\vartheta (1,+1^{(\widehat{𝐤})};(+\frac{1}{2})^{(\widehat{𝐤})},(\frac{1}{2})^{(\widehat{𝐤})})=C(\frac{1}{2}\frac{1}{2}1;\frac{1}{2},\frac{1}{2}1)=0,$$
(103)
$$\vartheta (1,(1)^{(\widehat{𝐤})};(+\frac{1}{2})^{(\widehat{𝐤})},(\frac{1}{2})^{(\widehat{𝐤})})=C(\frac{1}{2}\frac{1}{2}1;\frac{1}{2},\frac{1}{2},1)=0$$
(104)
and
$$\vartheta (1,0^{(\widehat{𝐤})};(+\frac{1}{2})^{(\widehat{𝐤})},(\frac{1}{2})^{(\widehat{𝐤})})=C(\frac{1}{2}\frac{1}{2}1;\frac{1}{2},\frac{1}{2}0)=\frac{1}{\sqrt{2}},$$
(105)
we find that
$$\chi (1,1^{(\widehat{𝐚})};(+\frac{1}{2})^{(\widehat{𝐤})},(\frac{1}{2})^{(\widehat{𝐤})})=\frac{1}{2}\mathrm{sin}\theta .$$
(106)
In the same way
$`\chi (1,+1^{(\widehat{𝐚})};(\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})})=\zeta (1,+1^{(\widehat{𝐚})};1,+1^{(\widehat{𝐤})})\vartheta (1,+1^{(\widehat{𝐤})};(\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})})`$
$`+\zeta (1,+1^{(\widehat{𝐚})};1,0^{(\widehat{𝐤})})\vartheta (1,0^{(\widehat{𝐤})};(\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})})`$
$`+\zeta (1,+1^{(\widehat{𝐚})};1,(1)^{(\widehat{𝐤})})\vartheta (1,(1)^{(\widehat{𝐤})};(\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})}).`$ (107)
In this case, we have
$$\vartheta (1,+1^{(\widehat{𝐚})};(\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})})=C(\frac{1}{2}\frac{1}{2}1;\frac{1}{2}\frac{1}{2}1)=0,$$
(108)
$$\vartheta (1,(1)^{(\widehat{𝐚})};(\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})})=C(\frac{1}{2}\frac{1}{2}1;\frac{1}{2}\frac{1}{2},1)=0$$
(109)
and
$$\vartheta (1,0^{(\widehat{𝐚})};(\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})})=C(\frac{1}{2}\frac{1}{2}1;\frac{1}{2}\frac{1}{2}0)=\frac{1}{\sqrt{2}}.$$
(110)
Thus,
$$\chi (1,+1^{(\widehat{𝐚})};(\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})})=\frac{1}{2}\mathrm{sin}\theta .$$
(111)
Finally
$`\chi (1,1^{(\widehat{𝐚})};(\frac{1}{2})^{(\widehat{𝐤})},(\frac{1}{2})^{(\widehat{𝐤})})=`$
$`\zeta (1,1^{(\widehat{𝐚})};1,1^{(\widehat{𝐤})})\vartheta (1,+1^{(\widehat{𝐤})};(\frac{1}{2})^{(\widehat{𝐤})},(\frac{1}{2})^{(\widehat{𝐤})})`$
$`+\zeta (1,1^{(\widehat{𝐚})};1,0^{(\widehat{𝐤})})\vartheta (1,0^{(\widehat{𝐤})};(\frac{1}{2})^{(\widehat{𝐤})},(\frac{1}{2})^{(\widehat{𝐤})})`$
$`+\zeta (1,1^{(\widehat{𝐚})};1,(1)^{(\widehat{𝐤})})\vartheta (1,(1)^{(\widehat{𝐤})};(\frac{1}{2})^{(\widehat{𝐤})},(\frac{1}{2})^{(\widehat{𝐤})}).`$ (112)
With
$$\vartheta (1,1^{(\widehat{𝐤})};(\frac{1}{2})^{(\widehat{𝐤})},(\frac{1}{2})^{(\widehat{𝐤})})=C(\frac{1}{2}\frac{1}{2}1;\frac{1}{2},\frac{1}{2}1)=0,$$
(113)
$$\vartheta (1,(1)^{(\widehat{𝐤})};(\frac{1}{2})^{(\widehat{𝐤})},(\frac{1}{2})^{(\widehat{𝐤})})=C(\frac{1}{2}\frac{1}{2}1;\frac{1}{2},\frac{1}{2},1)=1$$
(114)
and
$$\vartheta (1,0^{(\widehat{𝐤})};(\frac{1}{2})^{(\widehat{𝐤})},(\frac{1}{2})^{(\widehat{𝐤})})=C(\frac{1}{2}\frac{1}{2}1;\frac{1}{2},\frac{1}{2}0)=0.$$
(115)
we obtain
$$\chi (1,1^{(\widehat{𝐚})};(\frac{1}{2})^{(\widehat{𝐤})},(\frac{1}{2})^{(\widehat{𝐤})})=\mathrm{sin}^2\theta /2e^{i\phi }.$$
(116)
Combining all these results together, we find that the matrix state for $`s=1`$, $`M^{(\widehat{𝐚})}=1`$ is
$`[\mathrm{\Psi }(1,1^{(\widehat{𝐚})};(m_1)^{\widehat{𝐜}_1},(m_2)^{\widehat{𝐜}_2})]=`$ $`\mathrm{cos}^2\theta /2e^{i\phi }[\eta _1((+\frac{1}{2})^{(\widehat{𝐤})})]_1[\eta _2((+\frac{1}{2})^{(\widehat{𝐤})})]_2`$
$`+{\displaystyle \frac{1}{2}}\mathrm{sin}\theta [\eta _1((+\frac{1}{2})^{(\widehat{𝐤})})]_1[\eta _2((\frac{1}{2})^{(\widehat{𝐤})})]_2`$
$`+{\displaystyle \frac{1}{2}}\mathrm{sin}\theta [\eta _1((\frac{1}{2})^{(\widehat{𝐤})})]_1[\eta _2((+\frac{1}{2})^{(\widehat{𝐤})})]_2`$
$`+\mathrm{sin}^2\theta /2e^{i\phi }[\eta _1((\frac{1}{2})^{(\widehat{𝐤})})]_1[\eta _2((\frac{1}{2})^{(\widehat{𝐤})})]_2.`$
Thus the generalised form of the triplet state for $`M^{(\widehat{𝐚})}=1`$ is,
$`[\mathrm{\Psi }(1,1^{(\widehat{𝐚})};(m_1)^{(\widehat{𝐜}_1)},(m_2)^{(\widehat{𝐜}_2)})]=`$ $`\mathrm{cos}^2\theta /2e^{i\phi }\left(\begin{array}{c}\mathrm{cos}\theta _d/2\\ \mathrm{sin}\theta _d/2\end{array}\right)_1\left(\begin{array}{c}\mathrm{cos}\theta _f/2\\ \mathrm{sin}\theta _f/2\end{array}\right)_2`$ (122)
$`+{\displaystyle \frac{1}{2}}\mathrm{sin}\theta \left(\begin{array}{c}\mathrm{cos}\theta _d/2\\ \mathrm{sin}\theta _d/2\end{array}\right)_1\left(\begin{array}{c}\mathrm{sin}\theta _f/2e^{i\phi _f}\\ \mathrm{cos}\theta _f/2e^{i\phi _f}\end{array}\right)_2`$ (132)
$`+{\displaystyle \frac{1}{2}}\mathrm{sin}\theta \left(\begin{array}{c}\mathrm{sin}\theta _d/2e^{i\phi _d}\\ \mathrm{cos}\theta _d/2e^{i\phi _d}\end{array}\right)_1\left(\begin{array}{c}\mathrm{cos}\theta _f/2\\ \mathrm{sin}\theta _f/2\end{array}\right)_2`$
$`+\mathrm{sin}^2\theta /2e^{i\phi }\left(\begin{array}{c}\mathrm{sin}\theta _d/2e^{i\phi _d}\\ \mathrm{cos}\theta _d/2e^{i\phi _d}\end{array}\right)_1\left(\begin{array}{c}\mathrm{sin}\theta _f/2e^{i\phi _f}\\ \mathrm{cos}\theta _f/2e^{i\phi _f}\end{array}\right)_2.`$ (137)
##### The $`M^{(\widehat{𝐚})}=0`$ State
For this case, the probability amplitude is
$`\mathrm{\Psi }(1,0^{(\widehat{𝐚})};(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)})={\displaystyle \underset{\alpha ,\alpha ^{}}{}}\{{\displaystyle \underset{l}{}}\zeta (1,0^{(\widehat{𝐚})};1,M_l^{(\widehat{𝐤})})`$
$`\times \vartheta (1,M_l^{(\widehat{𝐤})};(m_1)_\alpha ^{(\widehat{𝐤})},(m_2)_\alpha ^{}^{(\widehat{𝐤})})\}`$
$`\times \mathrm{\Phi }((m_1)_\alpha ^{(\widehat{𝐤})},(m_2)_\alpha ^{}^{(\widehat{𝐤})};(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)})`$
$`=`$ $`{\displaystyle \underset{j}{}}\chi (1,0^{(\widehat{𝐚})};B_{\alpha \alpha ^{}})\mathrm{\Phi }(B_{\alpha \alpha ^{}};(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)}),`$ (139)
where
$$\chi (1,0^{(\widehat{𝐚})};B_{\alpha \alpha ^{}})=\underset{l}{}\zeta (1,0^{(\widehat{𝐚})};1,M_l^{(\widehat{𝐤})})\vartheta (1,M_l^{(\widehat{𝐤})};B_{\alpha \alpha ^{}}).$$
(140)
The $`\zeta `$’s change because they are functions of the initial-state quantum numbers. Thus, we have
$$\zeta (1,0^{(\widehat{𝐚})};1,1^{(\widehat{𝐤})})=\frac{1}{\sqrt{2}}\mathrm{sin}\theta e^{i\phi }$$
(141)
$$\zeta (1,0^{(\widehat{𝐚})};1,0^{(\widehat{𝐤})})=\mathrm{cos}\theta $$
(142)
and
$$\zeta (1,0^{(\widehat{𝐚})};1,(1)^{(\widehat{𝐤})})=\frac{1}{\sqrt{2}}\mathrm{sin}\theta e^{i\phi }.$$
(143)
However, the $`\vartheta `$’s do not change. Thus, using Eqs. (98) - (100) for the $`\vartheta `$’s, we find that
$`\chi (1,0^{(\widehat{𝐚})};(+\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})})=`$ (144)
$`\zeta (1,0^{(\widehat{𝐚})};1,1^{(\widehat{𝐤})})\vartheta (1,1^{(\widehat{𝐤})};(+\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})})`$
$`+\zeta (1,0^{(\widehat{𝐚})};1,0^{(\widehat{𝐤})})\vartheta (1,0^{(\widehat{𝐤})};(+\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})})`$
$`+\zeta (1,0^{(\widehat{𝐚})};1,(1)^{(\widehat{𝐤})})\vartheta (1,(1)^{(\widehat{𝐤})};(+\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})})`$
$`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\mathrm{sin}\theta e^{i\phi }.`$
The other $`\chi `$’s are found to be
$`\chi (1,0^{(\widehat{𝐚})};(+\frac{1}{2})^{(\widehat{𝐤})},(\frac{1}{2})^{(\widehat{𝐤})})`$ $`=`$ $`\chi (1,0^{(\widehat{𝐚})};(\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})})`$ (145)
$`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\mathrm{cos}\theta `$
and
$$\chi (1,0^{(\widehat{𝐚})};(\frac{1}{2})^{(\widehat{𝐤})},(\frac{1}{2})^{(\widehat{𝐤})})=\frac{1}{\sqrt{2}}\mathrm{sin}\theta e^{i\phi }.$$
(146)
Thus
$`[\mathrm{\Psi }(1,0^{(\widehat{𝐚})};(m_1)^{(\widehat{𝐜}_1)},(m_2)^{(\widehat{𝐜}_2)}]=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\mathrm{sin}\theta e^{i\phi }\left(\begin{array}{c}\mathrm{cos}\theta _d/2\\ \mathrm{sin}\theta _d/2\end{array}\right)_1\left(\begin{array}{c}\mathrm{cos}\theta _f/2\\ \mathrm{sin}\theta _f/2\end{array}\right)_2`$ (151)
$`+{\displaystyle \frac{1}{\sqrt{2}}}\mathrm{cos}\theta \left(\begin{array}{c}\mathrm{cos}\theta _d/2\\ \mathrm{sin}\theta _d/2\end{array}\right)_1\left(\begin{array}{c}\mathrm{sin}\theta _f/2e^{i\phi _f}\\ \mathrm{cos}\theta _f/2e^{i\phi _f}\end{array}\right)_2`$ (161)
$`+{\displaystyle \frac{1}{\sqrt{2}}}\mathrm{cos}\theta \left(\begin{array}{c}\mathrm{sin}\theta _d/2e^{i\phi _d}\\ \mathrm{cos}\theta _d/2e^{i\phi _d}\end{array}\right)_1\left(\begin{array}{c}\mathrm{cos}\theta _f/2\\ \mathrm{sin}\theta _f/2\end{array}\right)_2`$
$`+{\displaystyle \frac{1}{\sqrt{2}}}\mathrm{sin}\theta e^{i\phi }\left(\begin{array}{c}\mathrm{sin}\theta _d/2e^{i\phi _d}\\ \mathrm{cos}\theta _d/2e^{i\phi _d}\end{array}\right)_1\left(\begin{array}{c}\mathrm{sin}\theta _f/2e^{i\phi _f}\\ \mathrm{cos}\theta _f/2e^{i\phi _f}\end{array}\right)_2.`$ (166)
##### The $`M^{(\widehat{𝐚})}=1`$ State
For this case, the probability amplitude is
$`\mathrm{\Psi }(1,(1)^{(\widehat{𝐚})};(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)})={\displaystyle \underset{\alpha ,\alpha ^{}}{}}\{{\displaystyle \underset{l}{}}\zeta (1,(1)^{(\widehat{𝐚})};1,M_l^{(\widehat{𝐤})})`$
$`\times \vartheta (1,M_l^{(\widehat{𝐤})};(m_1)_\alpha ^{(\widehat{𝐤})},(m_2)_\alpha ^{}^{(\widehat{𝐤})})\}`$
$`\times \mathrm{\Phi }((m_1)_\alpha ^{(\widehat{𝐤})},(m_2)_\alpha ^{}^{(\widehat{𝐤})};(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)})`$
$`=`$ $`{\displaystyle \underset{\alpha ,\alpha ^{}}{}}\chi (1,(1)^{(\widehat{𝐚})};B_{\alpha \alpha ^{}})\mathrm{\Phi }(B_{\alpha \alpha ^{}};(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)}),`$ (168)
where
$$\chi (1,(1)^{(\widehat{𝐚})};B_{\alpha \alpha ^{}})=\chi (1,(1)^{(\widehat{𝐚})};B_{\alpha \alpha ^{}}):=\underset{l}{}\zeta (1,(1)^{(\widehat{𝐚})};1,M_l^{(\widehat{𝐤})})\vartheta (1,M_l^{(\widehat{𝐤})};B_{\alpha \alpha ^{}})$$
(169)
The $`\zeta `$’s for this case are
$$\zeta (1,(1)^{(\widehat{𝐚})};1,1^{(\widehat{𝐤})})=\mathrm{sin}^2\theta /2e^{i\phi },$$
(170)
$$\zeta (1,(1)^{(\widehat{𝐚})};1,0^{(\widehat{𝐤})})=\frac{1}{\sqrt{2}}\mathrm{sin}\theta $$
(171)
and
$$\zeta (1,(1)^{(\widehat{𝐚})};1,(1)^{(\widehat{𝐤})})=\mathrm{cos}^2\frac{\theta }{2}e^{i\phi }.$$
(172)
Since the $`\vartheta `$’s remain the same, it follows that
$$\chi (1,(1)^{(\widehat{𝐚})};(+\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})}))=\mathrm{sin}^2\theta /2e^{i\phi },$$
(173)
$`\chi (1,(1)^{(\widehat{𝐚})};(+\frac{1}{2})^{(\widehat{𝐤})},(\frac{1}{2})^{(\widehat{𝐤})}))=\chi (1,(1)^{(\widehat{𝐚})};(\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})}))`$ (174)
$`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{sin}\theta ,`$
and
$$\chi (1,(1)^{(\widehat{𝐚})};(\frac{1}{2})^{(\widehat{𝐤})},(\frac{1}{2})^{(\widehat{𝐤})}))=\mathrm{cos}^2\theta /2e^{i\phi }.$$
(175)
As a result, we obtain for the matrix state
$`[\mathrm{\Psi }(1,(1)^{(\widehat{𝐚})};(m_1)^{(\widehat{𝐜}_1)},(m_2)^{(\widehat{𝐜}_2)})]=`$ $`\mathrm{sin}^2\theta /2e^{i\phi }\left(\begin{array}{c}\mathrm{cos}\theta _d/2\\ \mathrm{sin}\theta _d/2\end{array}\right)_1\left(\begin{array}{c}\mathrm{cos}\theta _f/2\\ \mathrm{sin}\theta _f/2\end{array}\right)_2`$ (180)
$`+{\displaystyle \frac{1}{2}}\mathrm{sin}\theta \left(\begin{array}{c}\mathrm{cos}\theta _d/2\\ \mathrm{sin}\theta _d/2\end{array}\right)_1\left(\begin{array}{c}\mathrm{sin}\theta _f/2e^{i\phi _f}\\ \mathrm{cos}\theta _f/2e^{i\phi _f}\end{array}\right)_2`$ (190)
$`+{\displaystyle \frac{1}{2}}\mathrm{sin}\theta \left(\begin{array}{c}\mathrm{sin}\theta _d/2e^{i\phi _d}\\ \mathrm{cos}\theta _d/2e^{i\phi _d}\end{array}\right)_1\left(\begin{array}{c}\mathrm{cos}\theta _f/2\\ \mathrm{sin}\theta _f/2\end{array}\right)_2`$
$`\mathrm{cos}^2\theta /2e^{i\phi }\left(\begin{array}{c}\mathrm{sin}\theta _d/2e^{i\phi _d}\\ \mathrm{cos}\theta _d/2e^{i\phi _d}\end{array}\right)_1\left(\begin{array}{c}\mathrm{sin}\theta _f/2e^{i\phi _f}\\ \mathrm{cos}\theta _f/2e^{i\phi _f}\end{array}\right)_2.`$ (195)
#### 5.2.2 The Singlet State
Having obtained the triplet states, we now seek the singlet state. The general formula is Eq. (67). The generalized probability amplitude for the singlet state is
$`\mathrm{\Psi }(s=0,M=0;(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)})=\mathrm{\Psi }(0,0;(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)})`$
$`=`$ $`{\displaystyle \underset{\alpha ,\alpha ^{}}{}}\chi (0,0;(m_1)_\alpha ^{(\widehat{𝐤})},(m_2)_\alpha ^{}^{(\widehat{𝐤})})\mathrm{\Phi }((m_1)_\alpha ^{(\widehat{𝐤})},(m_2)_\alpha ^{}^{(\widehat{𝐤})};(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)}).`$ (197)
The $`\chi `$’s are now directly Clebsch-Gordan coefficients for the case of total spin $`0`$ and subspins $`1/2`$ and $`1/2`$. Thus,
$$\chi (0,0;(+\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})})=C(\frac{1}{2}\frac{1}{2}0,\frac{1}{2}\frac{1}{2}0)=0,$$
(198)
$$\chi (0,0;(+\frac{1}{2})^{(\widehat{𝐤})},(\frac{1}{2})^{(\widehat{𝐤})})=C(\frac{1}{2}\frac{1}{2}0,\frac{1}{2},\frac{1}{2}0)=\frac{1}{\sqrt{2}},$$
(199)
$$\chi (0,0;(\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})})=C(\frac{1}{2}\frac{1}{2}0,\frac{1}{2}\frac{1}{2}0)=\frac{1}{\sqrt{2}}$$
(200)
and
$$\chi (0,0;(\frac{1}{2})^{(\widehat{𝐤})},(\frac{1}{2})^{(\widehat{𝐤})})=C(\frac{1}{2}\frac{1}{2}0,\frac{1}{2},\frac{1}{2}0)=0.$$
(201)
This means that the generalized probability amplitude is
$`\mathrm{\Psi }(0,0;(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)})=\frac{1}{\sqrt{2}}\mathrm{\Phi }((+\frac{1}{2})^{(\widehat{𝐤})},(\frac{1}{2})^{(\widehat{𝐤})};(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)})`$
$`\frac{1}{\sqrt{2}}\mathrm{\Phi }((\frac{1}{2})^{(\widehat{𝐤})},(+\frac{1}{2})^{(\widehat{𝐤})};(m_1)_u^{(\widehat{𝐜}_1)},(m_2)_v^{(\widehat{𝐜}_2)}).`$ (202)
Hence, the matrix state is
$`[\mathrm{\Psi }(0,0;(m_1)^{(\widehat{𝐜}_1)},(m_2)^{(\widehat{𝐜}_2)})]=`$ $`\frac{1}{\sqrt{2}}[\eta _1((+\frac{1}{2})^{(\widehat{𝐤})})]_1[\eta _2((\frac{1}{2})^{(\widehat{𝐤})})]_2`$
$`\frac{1}{\sqrt{2}}[\eta _1((\frac{1}{2})^{(\widehat{𝐤})})]_1[\eta _2((+\frac{1}{2})^{(\widehat{𝐤})})]_2.`$
Therefore, the generalized matrix form for the singlet state is
$`[\mathrm{\Psi }(0,0;(m_1)^{(\widehat{𝐜}_1)},(m_2)^{(\widehat{𝐜}_2)})]=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{c}\mathrm{cos}\theta _d/2\\ \mathrm{sin}\theta _d/2\end{array}\right)_1\left(\begin{array}{c}\mathrm{sin}\theta _f/2e^{i\phi _f}\\ \mathrm{cos}\theta _f/2e^{i\phi _f}\end{array}\right)_2`$ (213)
$`{\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{c}\mathrm{sin}\theta _d/2e^{i\phi _d}\\ \mathrm{cos}\theta _d/2e^{i\phi _d}\end{array}\right)_1\left(\begin{array}{c}\mathrm{cos}\theta _f/2\\ \mathrm{sin}\theta _f/2\end{array}\right)_2.`$
## 6 Recovery of Standard Results
It is now easy to see how the standard results come about from the current ones. First of all, if $`\widehat{𝐚}`$ is along the $`z`$ axis, so that $`\theta =\phi =0`$, the triplet states become.
$$[\mathrm{\Psi }(1,1^{(\widehat{𝐚})};(m_1)^{(\widehat{𝐜}_1)},(m_2)^{(\widehat{𝐜}_2)}]=\left(\begin{array}{c}\mathrm{cos}\theta _d/2\\ \mathrm{sin}\theta _d/2\end{array}\right)_1\left(\begin{array}{c}\mathrm{cos}\theta _f/2\\ \mathrm{sin}\theta _f/2\end{array}\right)_2,$$
(215)
$`[\mathrm{\Psi }(1,0^{(\widehat{𝐚})};(m_1)^{(\widehat{𝐜}_1)},(m_2)^{(\widehat{𝐜}_2)})]=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{c}\mathrm{cos}\theta _d/2\\ \mathrm{sin}\theta _d/2\end{array}\right)_1\left(\begin{array}{c}\mathrm{sin}\theta _f/2e^{i\phi _f}\\ \mathrm{cos}\theta _f/2e^{i\phi _f}\end{array}\right)_2`$ (225)
$`+{\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{c}\mathrm{sin}\theta _d/2e^{i\phi _d}\\ \mathrm{cos}\theta _d/2e^{i\phi _d}\end{array}\right)_1\left(\begin{array}{c}\mathrm{cos}\theta _f/2\\ \mathrm{sin}\theta _f/2\end{array}\right)_2`$
and
$$[\mathrm{\Psi }(1,(1)^{(\widehat{𝐚})};(m_1)^{(\widehat{𝐜}_1)},(m_2)^{(\widehat{𝐜}_2)})]=\left(\begin{array}{c}\mathrm{sin}\theta _d/2e^{i\phi _d}\\ \mathrm{cos}\theta _d/2e^{i\phi _d}\end{array}\right)_1\left(\begin{array}{c}\mathrm{sin}\theta _f/2e^{i\phi _f}\\ \mathrm{cos}\theta _f/2e^{i\phi _f}\end{array}\right)_2,$$
(227)
while the singlet state remains
$`[\mathrm{\Psi }(0,0;(m_1)^{(\widehat{𝐜}_1)},(m_2)^{(\widehat{𝐜}_2)})]=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{c}\mathrm{cos}\theta _d/2\\ \mathrm{sin}\theta _d/2\end{array}\right)_1\left(\begin{array}{c}\mathrm{sin}\theta _f/2e^{i\phi _f}\\ \mathrm{cos}\theta _f/2e^{i\phi _f}\end{array}\right)_2`$ (237)
$`{\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{c}\mathrm{sin}\theta _d/2e^{i\phi _d}\\ \mathrm{cos}\theta _d/2e^{i\phi _d}\end{array}\right)_1\left(\begin{array}{c}\mathrm{cos}\theta _f/2\\ \mathrm{sin}\theta _f/2\end{array}\right)_2.`$
The operator, Eqs. (73) - (81), remains unchanged.
We recover the standard formulas if in addition, $`\widehat{𝐝}=\widehat{𝐟}=\widehat{𝐤}`$ :
$$[\mathrm{\Psi }(1,1^{(\widehat{𝐚})};(m_1)^{(\widehat{𝐜}_1)},(m_2)^{(\widehat{𝐜}_2)})]=\left(\begin{array}{c}1\\ 0\end{array}\right)_1\left(\begin{array}{c}1\\ 0\end{array}\right)_2,$$
(239)
$`[\mathrm{\Psi }(1,0^{(\widehat{𝐚})};(m_1)^{(\widehat{𝐜}_1)},(m_2)^{(\widehat{𝐜}_2)})]`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{c}1\\ 0\end{array}\right)_1\left(\begin{array}{c}0\\ 1\end{array}\right)_2`$ (249)
$`+{\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{c}0\\ 1\end{array}\right)_1\left(\begin{array}{c}1\\ 0\end{array}\right)_2,`$
$$[\mathrm{\Psi }(1,(1)^{(\widehat{𝐚})};(m_1)^{(\widehat{𝐜}_1)},(m_2)^{(\widehat{𝐜}_2)})]=\left(\begin{array}{c}0\\ 1\end{array}\right)_1\left(\begin{array}{c}0\\ 1\end{array}\right)_2$$
(250)
and
$`[\mathrm{\Psi }(0,0;(m_1)^{(\widehat{𝐜}_1)},(m_2)^{(\widehat{𝐜}_2)})]`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{c}1\\ 0\end{array}\right)_1\left(\begin{array}{c}0\\ 1\end{array}\right)_2`$ (260)
$`{\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{c}0\\ 1\end{array}\right)_1\left(\begin{array}{c}1\\ 0\end{array}\right)_2.`$
In this limit, the elements of $`[r^{(1)}]_1`$ are
$$r_{11}^{(1)}=\mathrm{cos}^2\theta _1/2r^{(1)}((+\frac{1}{2})^{(\widehat{𝐜}_1)})\mathrm{sin}^2\theta _1/2r^{(1)}((\frac{1}{2})^{(\widehat{𝐜}_1)}),$$
(262)
$$r_{12}^{(1)}=\frac{1}{2}\mathrm{sin}\theta _1e^{i\phi _1}(r^{(1)}((+\frac{1}{2})^{(\widehat{𝐜}_1)})r^{(1)}((\frac{1}{2})^{(\widehat{𝐜}_1)})),$$
(263)
$$r_{21}^{(1)}=r_{12}^{(1)}$$
(264)
and
$$r_{22}^{(1)}=r_{11}^{(1)}.$$
(265)
The elements of $`[r^{(2)}]_2`$ are
$$r_{11}^{(2)}=\mathrm{cos}^2\theta _2/2r^{(2)}((+\frac{1}{2})^{(\widehat{𝐜}_2)})\mathrm{sin}^2\theta _2/2r^{(2)}((\frac{1}{2})^{(\widehat{𝐜}_2)})$$
(266)
$$r_{12}^{(2)}=\frac{1}{2}\mathrm{sin}\theta _2e^{i\phi _2}(r^{(2)}((+\frac{1}{2})^{(\widehat{𝐜}_2)})r^{(2)}((\frac{1}{2})^{(\widehat{𝐜}_2)}))$$
(267)
$$r_{21}^{(2)}=r_{12}^{(2)}$$
(268)
$$r_{22}^{(2)}=r_{11}^{(2)}$$
(269)
In the event that the quantities $`r_1`$ and $`r_2`$ are spin projections, we may assign the values $`+1`$ if the projections are up with respect to the respective unit vectors $`\widehat{𝐜}_1`$ and $`\widehat{𝐜}_2`$ and $`1`$ if they are down with respect to these vectors. Thus, $`r^{(1)}((\pm \frac{1}{2})^{(\widehat{𝐜}_1)})=\pm 1.`$ In that case, the generalized operator $`[r^{(1)}]_1`$ has the elements
$$r_{11}^{(1)}=\mathrm{cos}(\theta _d\theta _1)2\mathrm{sin}\theta _d\mathrm{sin}\theta _1\mathrm{sin}^2((\phi _d\phi _1)/2),$$
(270)
$$r_{12}^{(1)}=\mathrm{sin}\theta _d\mathrm{cos}\theta _1+\mathrm{sin}\theta _1\mathrm{cos}\theta _d\mathrm{cos}(\phi _d\phi _1)+i\mathrm{sin}\theta _1\mathrm{sin}(\phi _d\phi _1),$$
(271)
$$r_{21}^{(1)}=r_{12}^{(1)}$$
(272)
and
$$r_{22}^{(1)}=r_{11}^{(1)}.$$
(273)
Exactly the same expressions hold for the operator $`[r^{(2)}]_2`$, except that the subscript $`d`$ is replaced by the subscript $`f`$, and where the numeral $`1`$ does not give the row or column of a matrix element, it is replaced by $`2`$.
In the limit $`\widehat{𝐝}=\widehat{𝐟}=\widehat{𝐤}`$, the operators become
$$[r^{(1)}]_1=\left(\begin{array}{cc}\mathrm{cos}\theta _1& \mathrm{sin}\theta _1e^{i\phi _1}\\ \mathrm{sin}\theta _1e^{i\phi _1}& \mathrm{cos}\theta _1\end{array}\right)_1$$
(274)
and
$$[r^{(2)}]_2=\left(\begin{array}{cc}\mathrm{cos}\theta _2& \mathrm{sin}\theta _2e^{i\phi _2}\\ \mathrm{sin}\theta _2e^{i\phi _2}& \mathrm{cos}\theta _2\end{array}\right)_2,$$
(275)
the well-known standard forms. Thus, we see that the standard results are obtained easily and logically from this approach. These states are all normalized to unity, as is easily proved. Also, they are mutually orthogonal.
## 7 Discussion and Conclusion
In this paper, we have derived the standard matrix treatment of spin addition from probability amplitudes. This confirms the fact, first brought out in Ref. , that spin theory can be based on probability amplitudes. It also confirms the correctness of the Landé approach to quantum mechanics.
A very important observation arising from this paper is that the standard results for spin matrix mechanics are only a special case of more generalized ones. Despite that calculations can be successfully performed with the standard quantities even in ignorance of this fact, our understanding of spin theory is incomplete until we take this fact on board. Although the results in this paper relate to spin addition, the general observation that the standard theory of angular momentum addition is not generalized enough is true. Therefore, more generalized results await the application of the current approach to the addition of spins other than those corresponding to spin-$`1/2`$ systems. By the same token, the addition of spin and orbital angular momentum will lead to more generalized results. One can extend this observation to the case of the addition of three or more spins. The elucidation of angular momentum theory cannot be regarded as complete until the task of obtaining the generalized results is finished.
From the generalized probability amplitudes derived in Ref. , we presented two different matrix treatments for the triplets states, and one matrix treatment for the singlet state. We found that we could express the singlet state or the triplet states in terms of $`4\times 4`$ operators and vectors with four rows each. In addition, we could express the triplet states by means of $`3\times 3`$ operators and vectors with three rows. In both cases, the total space was not decomposed into two spaces corresponding to the constituent subsystems $`1`$ and $`2`$. But in the standard treatment, the operator is the product of an operator in the space of subsystem $`1`$ and of an operator in the space of subsystem $`2`$. The state consists of terms which are products of vectors in the subspaces of systems $`1`$ and $`2`$. There is thus this difference between the standard treatment and the new treatments in Ref. . This difference appears to be far from trivial. In the present generalized standard treatment, we could only succeed in deriving results by assuming that the operator of the arbitrary observable $`R`$ was factorizable into factors depending on the spaces of subsystem $`1`$ and of subsystem $`2`$. In the new treatments of Ref. , this was not necessary. Thus, when we need to deal with ”entangled” observables $`R`$, we need to resort to the new treatments, or to forgo matrix mechanics and use probability-amplitude mechanics.
Our work highlights the power of the Landé interpretation of quantum mechanics. This approach continues to surprise, and it is all but certain that it has new results to yield when applied to areas of quantum mechanics other than spin theory.
## 8 References
1. Mweene H. V., ”New Treatment of Systems of Compounded Angular Momentum”, quant-ph/9907082
2. Landé A., ”From Dualism To Unity in Quantum Physics”, Cambridge University Press, 1960.
3. Landé A., ”New Foundations of Quantum Mechanics”, Cambridge University Press, 1965.
4. Landé A., ”Foundations of Quantum Theory,” Yale University Press, 1955.
5. Landé A., ”Quantum Mechanics in a New Key,” Exposition Press, 1973.
6. Rose M. E., ”Elementary Theory of Angular Momentum”, John Wiley and Sons, Inc. (New York), 1957
7. Mweene H. V., ”Derivation of Spin Vectors and Operators From First Principles”, quant-ph/9905012
8. Mweene H. V., ”Generalized Spin-1/2 Operators and Their Eigenvectors”, quant-ph/9906002
9. Mweene H. V., ”Vectors and Operators for Spin 1 Derived From First Principles”, quant-ph/9906043
10. Mweene H. V., ”Alternative Forms of Generalized Vectors and Operators for Spin 1/2”, quant-ph/9907031
11. Mweene H. V., ”Spin Description and Calculations in the Landé Interpretation of Quantum Mechanics”, quant-ph/9907033 |
warning/0003/hep-th0003171.html | ar5iv | text | # HD-THEP-00-17 HU-EP-00/20 hep-th/0003171 Two-loop Feynman Diagrams in Yang-Mills Theory from Bosonic String Amplitudes
## 1 Introduction
Achievements developed by Bern and Kosower allow to deduce the contributions of a large number of one-loop Feynman diagrams in Yang-Mills gauge theory from a single string scattering diagram . This was first noticed by analyzing amplitudes of some heterotic string model, but later on it was realized that only the bosonic degrees of freedom were relevant in the appropriate field theoretical limit. In this sense one can say that the bosonic string amplitude adds up implicitly all particle diagrams of a given loop order. To extract these contributions one has to get rid of the massive modes of the string as well as its tachyonic excitation. If one were able to compute the corresponding zero Regge slope limit of the entire amplitude also in higher loop orders, a tremendous simplification of the computation of loop corrections in gauge theories might follow, which could be of great impact in perturbative techniques. It was further noticed that the string theoretical input to this method can be reduced to the quantum mechanics of particles moving on some world line which is reminiscent of the string world sheet at infinite string tension. Thus the World Line Formalism was established and its relation to string theory explored . During this investigation it was also realized that the tachyonic mode of the string could be employed to calculate scattering amplitudes of scalar field theory despite its unphysical mass. As well a conclusive first quantized treatment of QED could be defined without regarding the string theoretical origin of the formalism any more. On the other hand, the foremost challenge of the entire enterprise still remains unsolved, which is the generalization of the Bern Kosower rules of Yang-Mills theory to higher loop orders. In fact the World Line formalism did allow higher order computations in QED and scalar field theory . It was also possible to extract such information directly from string diagrams . Any attempt to master two-loop Yang-Mills theory starting from string diagrams remained inconclusive up to now. In a worldine approach to two-loop Yang-Mills theory was proposed, whose connection to string theory has not yet been explored.
In this article which is an extension of relevant parts of we try to define such a generalization in the spirit of the mentioned earlier works. Thus we first review parts of these and deduce a general procedure how to project the string amplitude onto the massless vector boson mode. Our method attempts to unify the formerly employed ones, reproducing all their results and making the two-loop extension straightforward at the first glance. We are able to present an exact computation of all contributions of the two-loop two gluon string amplitude, that are relevant to obtain the two-loop coefficient of the Yang-Mills $`\beta `$ function. But in the end we shall uncover two severe difficulties which we are unable to deal with, preventing us from identifying two-loop Bern-Kosower rules. One of these is the unanswered question of how the gauge choice of field theory enters into the string amplitude. We demonstrate that a comparison of Feynman diagrams of a particular “topology” to their string counterpart does not allow any of the so-called covariant background gauges. At the one-loop level the diagrams were reproduced in the particular Feynman background gauge and this was therefore supposed to be the preferred gauge in which the string amplitudes naturally appear in the field theory limit. The second open question is, how the string diagrams deal with renormalization. While string theory contains a scale at which the massive string modes cut off the “divergencies” of local field theories, these require renormalization. This, for instance, calls for the presence of counter term insertions contributing at the two-loop level of perturbation theory, a procedure which in general will depend on the renormalization scheme chosen. In which manner these contributions are included in string diagrams is rather obscure. Using background techniques, like we do in this paper, one also has to inspect IR regularization, before one can compare with the conventional $`\beta `$ function .
A further technical obstruction is our present inability to give a consistent treatment of four-gluon vertices. These do not generically occurr in string diagrams which are built up by sewing together three-point string vertices. In the limit when the geometry of the string world sheet tends to a diagram involving four-point vertices some of the moduli are frozen and their integrations have to be removed, which leaves one with divergencies that cannot be explained or regularized in an obvious manner. Thus our work unfortunately remains inconclusive in these respects but still allows some insight into the problems that prevent a further progress so far.
The work on these topics is still in progress and we would like to express our appreciation for the discussion on this paper with P. Di Vecchia, A. Lerda, R. Marotta and R. Russo in fall 1998 in Copenhagen.
## 2 The field theoretical limit of string amplitudes
The observation that field theoretical amplitudes can be recovered from string theory traces back to Scherk , who found that expressions obtained from the dual operator formalism coincide precisely with field theortical result for tree-level and one-loop diagrams of $`\mathrm{\Phi }^3`$ theory in the limit of vanishing Regge slope. One only has to introduce a relation between the string and field theoretical couplings, the whole kinematics of Feynman diagrams follows automatically.
From the string theoretical point of view one expects that in this limit the massless modes of the string particle spectrum form an effective low energy field theory. In the simple bosonic model these particle-like states are the scalar tachyonic excitation and the massless vector boson of the spectrum of the bosonic string. When the string tension goes to infinity, i.e. $`\alpha ^{}0`$, the former will be identified as a divergent $`1/\alpha ^{}`$ contribution of the amplitude while the latter is given by its constant term.
We shall identify these divergent, respectively constant contributions by introducing proper time variables, playing the same role as Schwinger proper times (SPT) in field theory. These are introduced when rewriting the momentum integrals over the internal momenta as
$`{\displaystyle \frac{d^dp}{(2\pi )^d}\frac{p_\mu p_\nu \mathrm{}}{p^2+m^2}}={\displaystyle _0^{\mathrm{}}}𝑑t{\displaystyle \frac{d^dp}{(2\pi )^d}p_\mu p_\nu \mathrm{}e^{t(p^2+m^2)}}.`$ (1)
One can now perform the Gaussian momentum integrations to obtain a SPT integral. This gives the type of integral which is naturally derived from string diagrams, by the conservation of difficulties it is technically not easier to solve than the Feynman momentum integral itself. The integration region is singled out by reducing the moduli space which has to be integrated over in the amplitude to that small region where the SPT variables stay finite . The SPT variables are defined in string theory by identifying
$`\delta \tau (L_0a)=\delta \tau \alpha ^{}\left(p^2+m^2\right)\delta t\left(p^2+m^2\right),`$ (2)
or
$`t=\alpha ^{}\tau +t_0=\alpha ^{}\mathrm{ln}|z|+t_0.`$ (3)
Logarithms of moduli correspond to proper time variables in the field theory. This reveals a truly geometrical interpretation of Feynman diagrams, which in this picture represent particles whose propagation is parametrized by proper time variables. The length of a particular propagator is defined by the fact that $`x^{L_0}`$, $`x`$ ranging from $`0`$ to $`1`$, is the operator that propagates the external states in the string diagrams along the boundary of the open world sheet. It leads to
$`t=\alpha ^{}\mathrm{ln}(x)\text{and}x=e^{t/\alpha ^{}}.`$ (4)
We thus have found means to extract the tachyonic and vector contribution from the string amplitude in terms of different powers in the moduli of the world sheet. The tachyon part of the amplitude is exponentially divergent in $`\alpha ^{}`$, or proportional to $`1/x`$, the gluon part is constant when $`x0`$ and all massive states are exponentially suppressed. The full integration region in the moduli space is defined by proper times $`t[0,\mathrm{}]`$, which translates to $`x[0,1]`$. The insertion points of the external states are being integrated over all connected components of the world sheet boundary. On the other hand, from the point of view of field theory we could already be satisfied with $`x[0,ϵ]`$.
We shall now set up a systematic three step procedure to extract the field theoretical contribution from string amplitudes following the sewing procedure in the form of . We first replace moduli by proper time variables according to (4) and then, secondly, eliminate all terms proportional to higher powers of moduli, keeping only those which are constant for gluon diagrams or the divergent parts, proportional to inverse powers of moduli, for scalar field theory. This method does to the present stage not enable to extract any mixed diagram, which involves couplings of two different types of particles. Finally we determine the integration region of the moduli. This program will be demonstrated to work on scalar theory and one-loop Yang-Mills diagrams, before we come to our main topic, the computation of two-loop Yang-Mills diagrams and the extraction of renormalization constants.
## 3 Scalar theory
The limit leading to the tachyonic mode of the string spectrum can be used to reproduce results for single Feynman diagrams as well as formulas for complete $`n`$-point functions, which are known from the World Line approach to field theory . All field theoretical results quoted will refer to the $`\mathrm{\Phi }^3`$ theory defined by the Lagrangian
$`={\displaystyle \frac{1}{2}}_\mu \mathrm{\Phi }_\mu \mathrm{\Phi }{\displaystyle \frac{m^2}{2}}\mathrm{\Phi }^2+{\displaystyle \frac{\lambda }{3!}}\mathrm{\Phi }^3.`$ (5)
We first briefly discuss the World Line Formalism, show how the Green’s function of the bosonic string reduces to the World Line Green’s function , and afterwards proceed to Feynman diagrams .
### 3.1 The bosonic Green’s function
The Green’s function of the bosonic string can be written
$`𝒢^{(h)}(z_i,z_j)=\mathrm{ln}\left({\displaystyle \frac{E^{(h)}(z_i,z_j)}{\sqrt{V_i^{}(0)V_j^{}(0)}}}\right){\displaystyle \frac{1}{2}}{\displaystyle _{z_i}^{z_j}}\omega ^\mu (2\pi \mathrm{}\left(\tau _{\mu \nu }\right))^1{\displaystyle _{z_i}^{z_j}}\omega ^\nu .`$ (6)
It is discussed in greater detail in the appendix A which we refer to for the geometrical description of world sheets. While the full string amplitude does not depend on any change of coordinate, the tiny part which is extracted from it in the chosen limit does, and one has to take a choice which set of local coordinates $`V_i(z)`$ one wants to use. The tree-level Green’s function reduces to the inverse of the Laplacian on the flat plane:
$`𝒢^{(0)}(z_1,z_2)=\mathrm{ln}\left|{\displaystyle \frac{z_1z_2}{\sqrt{V_1^{}(0)V_2^{}(0)}}}\right|.`$ (7)
The local coordinates in general are specified in a way that the insertion points of the external states are being integrated over only one component of the boundary, which corresponds in the field theoretical picture to selecting a particular loop for any external particle. If one then neglects the radii of the isometrical circles, the boundary component of a single loop becomes the interval between two fixed points on the real axis, both, circles and fixed points, referring to the Schottky map of this loop. We then use the explicit form
$`V_i^{}(0)=\left|{\displaystyle \frac{(z_i\eta _j)(z_i\xi _j)}{(\xi _j\eta _j)}}\right|`$ (8)
for the local coordinates. They specify the loop $`j`$ that is generated by the Schottky maps whose fixed points are $`\xi _j`$ and $`\eta _j`$ and to which the external state $`i`$ is attached. In fact, they are just the inverse of the first abelian differentials on the world aheet
$`V_i^{}(0)=\left({\displaystyle \frac{\omega ^j(z_i)}{dz_i}}\right)^1.`$ (9)
On the tree-level world sheet one then defines SPT variables by $`\mathrm{ln}(z_{1,2})=t_{1,2}/\alpha ^{}`$ and find
$`𝒢^{(0)}(z_1,z_2)\mathrm{ln}\left|\sqrt{{\displaystyle \frac{z_1}{z_2}}}\right|=\pm {\displaystyle \frac{(t_1t_2)}{2\alpha ^{}}},`$ (10)
depending on which $`z_i`$ is the larger one. This matches precisely to the tree-level World Line Green’s function.
Proceeding to loop-level needs an interpretation of the further moduli of the world sheet, which are the fixed points and multipliers of the generators of the Schottky group. The latter are connected to the length of that loop by
$`T=\alpha ^{}\mathrm{ln}|k|,`$ (11)
while the fixed points shall be treated later. We now have to expand the prime form and the period matrix of the world sheet into a power series in the multipliers of the Schottky group and neglect all terms which are suppressed when $`\alpha ^{}0`$. The expansion is discussed thoroughly in appendix A.2. Using the results one finds:
$`E(z_1,z_2)`$ $`=`$ $`(z_1z_2)+o(k_i),`$ (12)
$`\left(2\pi \mathrm{}\left(\tau _{\mu \nu }\right)\right)`$ $`=`$ $`\delta _{\mu \nu }\mathrm{ln}(k_\mu )(1\delta _{\mu \nu })\mathrm{ln}\left|{\displaystyle \frac{(\eta _\mu \eta _\nu )(\xi _\mu \xi _\nu )}{(\eta _\mu \xi _\nu )(\xi _\mu \eta _\nu )}}\right|+o(k_i),`$
$`{\displaystyle _{z_1}^{z_2}}\omega ^\mu `$ $`=`$ $`\mathrm{ln}\left|{\displaystyle \frac{(z_1\eta _\mu )(z_2\xi _\mu )}{(z_1\xi _\mu )(z_2\eta _\mu )}}\right|+o(k_i).`$
This simplifies for the one-loop result by choosing $`\eta =0`$ and $`\xi =\mathrm{}`$, so that we obtain:
$`E^{(1)}(z_1,z_2)`$ $``$ $`(z_1z_2),`$ (13)
$`\left(2\pi \mathrm{}\left(\tau _{11}^{(1)}\right)\right)^1`$ $``$ $`{\displaystyle \frac{1}{\alpha ^{}T}},`$
$`\left({\displaystyle _{z_1}^{z_2}}\omega ^1\right)^{(1)}`$ $``$ $`{\displaystyle \frac{|t_1t_2|}{\alpha ^{}}}.`$
This has to be substituted into (6) and we recover the World Line one-loop Green’s function. A more complete discussion of the whole procedure is given in and also the case of higher loop orders is verified. This establishes the relation
$`𝒢^{(g)}(z_1,z_2){\displaystyle \frac{1}{2\alpha ^{}}}G_B^{(g)}(t_1,t_2)`$ (14)
between the Green’s function of the bosonic string and that of the scalar World Line formalism of field theory to any order in perturbation theory.
We now proceed to discuss the two-loop case. We always divide the integration measure of the integrals over moduli by the volume of the modular group of the Riemann surface by fixing three of the four fixed points to the values
$`\eta _2=0,\xi _2=\mathrm{},\xi _1=1`$ (15)
on the complex sphere, calling $`\eta _1=\eta `$. This standard choice is illustrated in figure 1.
If we for instance intend to compute the diagram of figure 2 we can cut the world sheet as in figure 3. The conformal invariance allows to cut in a way that $`z_1,z_2,\eta <1`$.
We then read off
$`0=\eta _2<z_2<\eta _1=\eta <z_1<\xi _1=1<\xi _2=\mathrm{}`$ (16)
and satisfy this relation by defining new moduli $`A_i[0,1]`$ for $`i=1,2,3`$, which we substitute for the former ones:
$`z_1=A_1,\eta =A_1A_2,z_2=A_1A_2A_3.`$ (17)
The $`A_i`$ are interpreted as sewing parameter which parametrize the length of the according propagator. Their logarithms are translated into the proper time variables $`t_i`$ of the field theoretical diagrams. The loops support two more such SPT variables by their Schottky multipliers which are replaced by $`t_4`$ and $`t_5`$ so that the prescription altogether reads
$`t_i`$ $`=`$ $`\alpha ^{}\mathrm{ln}(A_i),\text{for}i=1,2,3,`$ (18)
$`t_5`$ $`=`$ $`\alpha ^{}\mathrm{ln}(k_1)t_1t_2,`$
$`t_4`$ $`=`$ $`\alpha ^{}\mathrm{ln}(k_2)t_1t_2t_3.`$
This parametrization is depicted in figure 4.
Any other logarithm occurring in the formula for the Green’s function will be translated analogously:
$`\alpha ^{}\mathrm{ln}(z_1z_2)`$ $`=`$ $`t_1+\alpha ^{}\mathrm{ln}(1A_2A_3),`$ (19)
$`\alpha ^{}\mathrm{ln}(z_1z_2)`$ $`=`$ $`2t_1t_2t_3,`$
$`\alpha ^{}\mathrm{ln}\left({\displaystyle \frac{z_1}{z_2}}\right)`$ $`=`$ $`t_2+t_3,`$
$`\alpha ^{}\mathrm{ln}\left({\displaystyle \frac{(z_2\eta )(z_11)}{(z_21)(z_1\eta )}}\right)`$ $`=`$ $`t_2+\alpha ^{}\mathrm{ln}\left({\displaystyle \frac{(1A_3)(1A_1)}{(1A_1A_2A_3)(1A_2)}}\right)`$
and we finally find
$`𝒢^{(2)}(z_1,z_2)`$ $`=`$ $`\mathrm{ln}(z_1z_2){\displaystyle \frac{1}{2}}\mathrm{ln}(z_1z_2){\displaystyle \frac{1}{\mathrm{ln}(k_1)\mathrm{ln}(k_2)\mathrm{ln}^2(\eta )}}`$
$`\times (\mathrm{ln}\left({\displaystyle \frac{(z_1\eta )(z_21)}{(z_11)(z_2\eta )}}\right)\mathrm{ln}\left({\displaystyle \frac{z_1}{z_2}}\right)\mathrm{ln}(\eta )\mathrm{ln}\left({\displaystyle \frac{z_1}{z_2}}\right)\mathrm{ln}(k_1)`$
$`\mathrm{ln}\left({\displaystyle \frac{(z_1\eta )(z_21)}{(z_11)(z_2\eta )}}\right)\mathrm{ln}(k_2))+o(\alpha ^0)`$
$`=`$ $`{\displaystyle \frac{1}{2\alpha ^{}}}{\displaystyle \frac{t_1t_2(t_3+t_4+t_5)+t_1t_3(t_4+t_5)+t_2t_4(t_3+t_5)+t_3t_4t_5}{(t_1+t_2)(t_3+t_4+t_5)+(t_3+t_4)t_5}}+o(\alpha ^0).`$
This is precisely the result of the Gaussian integration which has to be performed in field theory for the particular combination of propagators which belongs to the diagram drawn in figure 2:
$`{\displaystyle \frac{d^dpd^dk}{(2\pi )^{2d}}\frac{1}{p^2k^2(qk)^2(pq)^2(kp)^2}}={\displaystyle _0^{\mathrm{}}}\left({\displaystyle \underset{i=1}{\overset{5}{}}}dt_i\right)\left((t_1+t_2)(t_3+t_4+t_5)+(t_3+t_4)t_5\right)^{d/2}`$
$`\times \mathrm{exp}\left(q^2{\displaystyle \frac{t_1t_2(t_3+t_4+t_5)+t_1t_3(t_4+t_5)+t_2t_4(t_3+t_5)+t_3t_4t_5}{(t_1+t_2)(t_3+t_4+t_5)+(t_3+t_4)t_5}}\right).`$
### 3.2 Scalar Feynman diagrams
The relation between the string coupling $`g_S`$ and the field theoretical coupling $`\lambda `$ of the $`\mathrm{\Phi }^3`$ theory from (5) has to be chosen
$`\lambda =16(2\alpha ^{})^{(d6)/4}g_S.`$ (21)
Furthermore a treatment of the integration measure and the normalization constants must be defined. The two-loop two-tachyon amplitude is constructed by inserting two scalar vertex operators into the path integral and it follows
$`𝒜_2^{(2)}(p,p)=C_2\left(2g_S(2\alpha ^{})^{(d2)/4}\right)^2{\displaystyle _\mathrm{\Gamma }}\left[dm\right]_2^2\mathrm{exp}\left(2\alpha ^{}p^2𝒢^{(2)}(z_1,z_2)\right).`$ (22)
The nature of the normalization factor
$`C_2\left((2\pi )^2\sqrt{2\alpha ^{}}\right)^dg_S^2`$ (23)
will also be discussed later on more thoroughly. The integration measure is given by:
$`\left[dm\right]_2^2`$ $`=`$ $`{\displaystyle \frac{dk_1dk_2d\eta }{k_1^2k_2^2(\eta 1)^2}}{\displaystyle \frac{dz_1dz_2}{V_1^{}(0)V_2^{}(0)}}(1k_1)^2(1k_2)^2`$
$`\times \text{det}^{d/2}\left(i\tau _{\mu \nu }\right){\displaystyle \underset{\beta }{}}\left({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left(1k_\beta ^n\right)^{2d}\left(1k_\beta \right)^2\right).`$
The volume of the projective group has already been divided out by introducing the standard choice (15) of coordinates. For the local coordinates we use the usual $`V_i^{}(0)=z_i`$. Next we substitute the sewing parameters from (17) and find
$`{\displaystyle _\mathrm{\Gamma }}{\displaystyle \frac{dk_1dk_2d\eta dz_1dz_2}{k_1k_2\eta z_1z_2}}`$ $`=`$ $`{\displaystyle _0^1}{\displaystyle \underset{i=1}{\overset{3}{}}}\left({\displaystyle \frac{dA_i}{A_i}}\right){\displaystyle \frac{dk_1dk_2}{k_1k_2}}=\left({\displaystyle \frac{1}{\alpha ^{}}}\right)^5{\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{i=0}{\overset{5}{}}}dt_i,`$ (25)
$`(2\pi )^d\text{det}^{d/2}\left(i\tau _{\mu \nu }\right)`$ $`=`$ $`\left(\mathrm{ln}(k_1)\mathrm{ln}(k_2)\mathrm{ln}^2(\eta )\right)^{d/2}+o(k_1,k_2)`$
$`=`$ $`\left((t_1+t_2)(t_3+t_3+t_5)+(t_3+t_4)t_5\right)^{d/2}+o(k_1,k_2).`$
As we intend to project exclusively onto the tachyonic excitations we can employ their mass-shell condition
$`\alpha ^{}m^2=1`$ (26)
to define a particle mass for the scalar field. We then discard all terms proportional to powers of any moduli and retain a field theory diagram of a massive scalar field, while all other contributions are suppressed exponentially or, in the case of the gluon, stay finite when $`\alpha ^{}0`$. In regard of the prefactor we can omitt terms proportional to $`k_1,k_2`$ in the integrand
$`{\displaystyle \underset{\beta }{}}\left({\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}\left(1k_\beta ^i\right)^{2d}\left(1k_\beta \right)^2\right)(1k_1)^2(1k_2)^2=1+o(k_1,k_2)`$ (27)
and following (17) and (18) we get:
$`{\displaystyle \frac{\eta }{k_1k_2}}=\mathrm{exp}\left(m^2{\displaystyle \underset{i=1}{\overset{5}{}}}t_i\right).`$ (28)
The final result for the two-loop two-tachyon amplitude reads:
$`𝒜_2^{(2)}(p^2)`$ $`=`$ $`{\displaystyle \frac{\lambda ^4}{2^9(4\pi )^d}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{i=1}{\overset{5}{}}}\left(dt_ie^{m^2t_i}\right)\left((t_1+t_2)(t_3+t_3+t_5)+(t_3+t_4)t_5\right)^{d/2}`$
$`\times \mathrm{exp}\left(p^2{\displaystyle \frac{t_1t_2(t_3+t_4+t_5)+t_1t_3(t_4+t_5)+t_2t_4(t_3+t_5)+t_3t_4t_5}{(t_1+t_2)(t_3+t_4+t_5)+(t_3+t_4)t_5}}\right).`$
It reproduces the field theoretical result for the Feynman diagrams from figure 2 including any combinatorial factors. In this method is also applied to reducible diagrams and the same positive result was confirmed.
This procedure not only works for $`\mathrm{\Phi }^3`$ theory but can also be extended to scalar $`\mathrm{\Phi }^4`$ theory by a different matching of the coupling constants. A problem arises with the moduli or SPT variables of those propagators which shrink to zero length. Their integrations need to be regularized by hand, a method which is conceptually not very satisfactory. Still it can be performed to get contributions consistent with results for scalar field theories, while we will not be able to extrapolate the method to gauge theories. Another open question which one should be able to tackle already in two-loop scalar diagrams concerns the presence, or absence, of counterterm diagrams. It has not been investigated, how their contributions to the two-loop scattering are somehow implicitly contained in the string diagram. If not, the missing part could depend on the renormalization scheme.
## 4 Pure Yang-Mills gauge theory
After having analyzed the scalar field theory in terms of Green’s functions and Feynman diagrams we shall now treat the most interesting case of Yang-Mills gauge theory. The challenging aim of our investigation is the generalization of the Bern-Kosower rules which allow a greatly simplified computation of one-loop $`n`$-point functions in pure gauge theories . We shall present the means to calculate exactly all contributions of the two-point two-loop string amplitude that survive in the field theory limit. In this context we shall particularly have to point out the difficulties arising from the unknown relation between different choices for the local coordinate on the world sheet and different gauge choices in the field theory. Furthermore the problem of constructing four-gluon vertices in string theory will be left a riddle.
Starting point is the $`n`$-gluon $`h`$-loop amplitude of the bosonic string which is given by the expectation value of $`n`$ gluon vertex operators computed in the background of a bosonic string theory on a Genus $`h`$ Riemann surface. The contractions of the world sheet coordinates are again performed using the Green’s function from (6). We briefly collect the constituents of this amplitude
$`𝒜_n^{(h)}(p_1,\mathrm{},p_n)`$ $`=`$ $`N^h\text{tr}\left(\lambda ^{a_1}\mathrm{}\lambda ^{a_n}\right)C_h𝒩_0^n{\displaystyle _\mathrm{\Gamma }}\left[dm\right]_h^n{\displaystyle \underset{i<j}{\overset{n}{}}}\mathrm{exp}\left(2\alpha ^{}p_ip_j𝒢^{(h)}(z_i,z_j)\right)`$
$`\times \mathrm{exp}({\displaystyle \underset{ij}{\overset{n}{}}}\sqrt{2\alpha ^{}}(p_jϵ_i)_{z_i}𝒢^{(h)}(z_i,z_j)+{\displaystyle \frac{1}{2}}{\displaystyle \underset{ij}{\overset{n}{}}}(ϵ_iϵ_j)_{z_i}_{z_j}𝒢^{(h)}(z_i,z_j))_{\text{m.l.}}.`$
The prefactor $`N^h\text{tr}\left(\lambda ^{a_1}\mathrm{}\lambda ^{a_n}\right)`$ is the Chan Paton factor of the diagram with external gauge charges $`\lambda ^{a_i}`$ of some $`SU(N)`$ gauge group, whose adjoint representation is normalized as follows:
$`\text{tr}\left(\lambda ^{a_i}\lambda ^{a_j}\right)={\displaystyle \frac{1}{2}}\delta _{a_ia_j}.`$ (31)
The $`ϵ_i`$ are the polarization vectors of the external gluons and from the expansion of the exponential one only has to keep the terms that are multilinear in the $`ϵ_i`$. The general integration measure can be written in terms of the Schottky parameters that span the moduli space $`\mathrm{\Gamma }`$:
$`\left[dm\right]_h^n`$ $``$ $`{\displaystyle \underset{m=1}{\overset{h}{}}}\left({\displaystyle \frac{dk_md\xi _md\eta _m}{k_m^2(\xi _m\eta _m)^2}}(1k_m)^2\right){\displaystyle \underset{m=1}{\overset{n}{}}}(dz_m){\displaystyle \frac{1}{dV_{abc}}}`$
$`\times \left(det(i\tau _{\mu \nu })\right)^{d/2}{\displaystyle \underset{\beta }{}}\left({\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}(1k_\beta ^m)^d{\displaystyle \underset{m=2}{\overset{\mathrm{}}{}}}(1k_\beta ^m)^2\right).`$
It contains the normalization constant involving the product over primary classes and the volume of the modular group $`dV_{abc}`$ in the denominator. The bosonic Green’s function is again be given by
$`𝒢^{(h)}(z_i,z_j)=\mathrm{ln}\left({\displaystyle \frac{E^{(h)}(z_i,z_j)}{\sqrt{V_i^{}(0)V_j^{}(0)}}}\right){\displaystyle \frac{1}{2}}{\displaystyle _{z_i}^{z_j}}\omega ^\mu (2\pi \mathrm{}\left(\tau _{\mu \nu }\right))^1{\displaystyle _{z_i}^{z_j}}\omega ^\nu `$ (33)
Note the difference in the depence on the local coordinates compared to , which has to be obeyed to get correct results in higher than one-loop order . We are discussing off-shell amplitudes and therefore do not demand the mass-shell and transversality relations
$`p^2=0,pϵ=0`$ (34)
to hold, although for the sake of brevity we often restrict ourselves to the term proportional to $`(p_jϵ_i)(ϵ_jp_i)`$. The derivation of the normalization constant is found in . It uses tree-level three- and four-point amplitudes and factorization arguments to fix the dependence of the prefactor on the scale $`\alpha ^{}`$ and the dimensionless string coupling constant
$`g_S{\displaystyle \frac{g}{2}}(2\alpha ^{})^{1d/4}.`$ (35)
It is then found:
$`𝒩_0`$ $`=`$ $`2g_S(2\alpha ^{})^{(d6)/4},`$ (36)
$`C_h`$ $`=`$ $`\left((2\pi )^h\sqrt{2\alpha ^{}}\right)^dg_S^{2h2}.`$
The first factor reproduces the normalization of the Fourier transformations which come with loop calculations and also assigns the correct physical dimension to the amplitude, while the second carries the necessary power in the string coupling that is dictated by the Euler characteristic $`\chi ()=22h`$ of the Riemann surface.
### 4.1 One-loop diagrams
Following the rules we have established when adressing scalar field theory we shall now proceed to gauge fields . We expand the integrand of the amplitude in the moduli, keep only the finite and non-vanishing part of the expansion when $`\alpha ^{}0`$ and translate everything into the language of SPT variables. The results allow to confirm the Ward identities of the Feynman background gauge and to read off the coefficient of the Yang-Mills $`\beta `$ function from the wave function renormalization constant of the gauge field, which at one-loop order receives contributions only from a single Feynman diagram.
We first specify the formulas for the measure (4), the Green’s function (33) and the constants (36):
$`\left[dm\right]_1^n`$ $`=`$ $`{\displaystyle \underset{m=1}{\overset{n}{}}}\left(dz_m\right){\displaystyle \frac{1}{dV_{abc}}}{\displaystyle \frac{dkd\eta d\xi }{k^2(\eta \xi )^2}}\left({\displaystyle \frac{\mathrm{ln}(k)}{2\pi }}\right)^{d/2}{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}\left(1k^m\right)^{2d}`$
$`=`$ $`{\displaystyle \underset{m=2}{\overset{n}{}}}\left(dz_m\right){\displaystyle \frac{dk}{k^2}}\left({\displaystyle \frac{\mathrm{ln}(k)}{2\pi }}\right)^{d/2}{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}\left(1k^m\right)^{2d},`$
$`𝒢^{(1)}(z_i,z_j)`$ $`=`$ $`\mathrm{ln}(z_iz_j){\displaystyle \frac{1}{2}}\mathrm{ln}(z_iz_j)+{\displaystyle \frac{1}{2\mathrm{ln}(k)}}\mathrm{ln}^2\left({\displaystyle \frac{z_i}{z_j}}\right)`$
$`+\mathrm{ln}\left({\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\left(1k^m\frac{z_i}{z_j}\right)\left(1k^m\frac{z_j}{z_i}\right)}{\left(1k^m\right)^2}}\right),`$
$`N_0^nC_h`$ $`=`$ $`\left(2\pi \sqrt{2\alpha ^{}}\right)^d(2\alpha ^{}g^2)^{n/2}.`$
We have used the usual fixed point choice $`\eta =0`$ and $`\xi =\mathrm{}`$ as well as the Lovelace type local coordinates $`V_i^{}(0)=z_i`$. Further $`z_2=1`$ has been fixed by a projective transformation.
The only diagram can be parametrized as in figure 5 which leads to the integration region $`\mathrm{\Gamma }=\{z_1|0<z_1<1\}`$. Substituting the sewing parameter $`z_1=A`$ into (4.1) and defining the proper times
$`T=\alpha ^{}\mathrm{ln}(k),t=\alpha ^{}\mathrm{ln}(A),[0,1][\mathrm{},0],`$ (38)
we find
$`\left[dm\right]_1^2`$ $`=`$ $`{\displaystyle \frac{dkdz_1}{k^2}}\left({\displaystyle \frac{\mathrm{ln}(k)}{2\pi }}\right)^{d/2}{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}\left(1k^m\right)^{2d}`$
$`=`$ $`\left({\displaystyle \frac{1}{\alpha ^{}}}\right)^2(2\pi )^{d/2}{\displaystyle \frac{AT^{d/2}dTdt}{k}}\left(1+(d2)k+o\left(k^2\right)\right),`$
$`𝒢^{(1)}(z_1,z_2)`$ $`=`$ $`\mathrm{ln}(1A){\displaystyle \frac{1}{2}}\mathrm{ln}(A)+{\displaystyle \frac{1}{2\mathrm{ln}(k)}}\mathrm{ln}^2\left(A\right)+\mathrm{ln}\left({\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\left(1k^m/A\right)\left(1Ak^m\right)}{\left(1k^m\right)^2}}\right)`$
$`=`$ $`\mathrm{ln}(1A)+{\displaystyle \frac{t}{2\alpha ^{}}}{\displaystyle \frac{t^2}{2\alpha ^{}T}}+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left(\mathrm{ln}\left(1Ak^n\right)+\mathrm{ln}\left(1k^n/A\right)2\mathrm{ln}\left(1k^n\right)\right).`$
In the same manner one can take the derivatives of the Green’s function. We now only have to keep the term in the integrand that is proportional to $`A^0,k^0,\alpha ^0`$, and thus finally find:
$`𝒜_2^{(2)}(p^2)`$ $`=`$ $`N\text{tr}\left(\lambda ^{a_1}\lambda ^{a_2}\right){\displaystyle \frac{4g^2}{(4\pi )^{d/2}}}\left((ϵ_2p_1)(ϵ_1p_2)ϵ^2p^2\right){\displaystyle _0^{\mathrm{}}}𝑑TT^{d/2}{\displaystyle _0^{\mathrm{}}}𝑑t`$
$`\times \left(\left({\displaystyle \frac{1}{2}}{\displaystyle \frac{t}{T}}\right)^2(d2)2+o(k,A)\right)\mathrm{exp}\left(p^2\left(t{\displaystyle \frac{t^2}{T}}+o(\alpha ^{},k,A)\right)\right).`$
Rescaling the integration over $`t`$ by a factor $`T`$ leads to the result of . All integrations can easily be performed using some formulas given in appendix D and the precise value that comes out of the computation of the sum of the two Feynman diagrams of Yang-Mills theory in the Feynman background gauge , drawn in 6, is recovered.
In particular from the wave function one-loop renormalization constant
$`Z_A=1+{\displaystyle \frac{g^2N}{(4\pi )^2}}{\displaystyle \frac{11}{3}}{\displaystyle \frac{1}{ϵ}}`$ (41)
the correct coefficient of the Yang-Mills $`\beta `$ function is deduced. It is interesting to note that these notions obviously depend on the choice of local coordinates $`V_i(z)`$.
We now only briefly adress the one-loop three-gluon diagrams. This is the first and most simple example how to construct four-gluon vertices in string theory, which will be defined by extracting those regions in the moduli space where a propagator between two three-gluon vertices is short in terms of field theoretical proper time. We use the parametrization of figure 7.
The integration region for the sewing parameters from (17) will be as usually $`[0,1]^3`$. They are defined by
$`z_1=A_1A_2,z_2=A_2`$ (42)
and lead to proper time variables
$`t_i`$ $`=`$ $`\alpha ^{}\mathrm{ln}(A_i),\text{for}i=1,2,`$ (43)
$`t_3`$ $`=`$ $`\alpha ^{}\mathrm{ln}(k)t_1t_2.`$
We then get the sum of the two Feynman diagrams without four-gluon vertices by expanding the integrand and taking the appropriate limit. On the other hand, if we let $`\alpha ^{}0`$ such that e.g. $`A_1`$ stays finite, we have $`t_10`$, although the insertion points $`z_1`$ and $`z_2`$ remain widely separated on the world sheet. The problem of this naive definition of four-gluon vertices in string amplitudes, the so-called pinching, is that it remains completely unclear what should happen to the free modulus $`A_1`$. In the term in the integrand, that stays finite when $`A_11`$, is chosen to be relevant and the rest is being ignored. In fact, this method did allow to verify the relation between the three and four gluon diagrams deriving from the Ward identities of the Feyman background gauge for Yang-Mills theories and might thus be called heuristically successfull.
### 4.2 Two-loop vacuum diagrams
The most simple two-loop Feynman diagrams one can think of are the vacuum diagrams drawn in figure 8.
They were discussed in . The generalization will not be straightforward as we shall have to deal with the announced difficulties concerning the diagram (b) with a four-gluon vertex. Another more fundamental question arises in the context of gauge invariance. The diagrams (a) and (b) clearly have to correspond to different regions in the moduli space of the string world sheet, simply because the number of field theoretical propagators is not the same, and therefore a different number of proper time variables has to be introduced. On the other hand their respective contributions in the field theoretical calculation will depend on the gauge one has chosen. One might expect that the statement which was true at one-loop level, that string contribution just come out in the Feynman background gauge, is still valid. The least one would like to require would be that added up they give a gauge invariant result that can be compared to field theory. Actually the question cannot be answered in the case of the vacuum diagrams, as their contributions vanish identically in dimensional regularization. This follows the principle that vacuum diagrams only give an irrelevant phase factor to the scattering matrix. So in the end there is no result to be extracted from this treatment, which we still undertake to point out the ealier mentioned problems that shall come up again whith the two-loop two-point function in the next chapter.
We now specialize all the relevant expressions in the amplitude, dropping the colour factor which is empty. To then work out the amlitude
$`𝒜_0^{(2)}=C_2{\displaystyle _\mathrm{\Gamma }}\left[dm\right]_2^0={\displaystyle \frac{g^2}{4}}\left(2\pi \sqrt{2\alpha ^{}}\right)^{2d}(2\alpha ^{})^2{\displaystyle _\mathrm{\Gamma }}\left[dm\right]_2^0`$ (44)
we only need the measure of the integration over the moduli of which some are fixed as in (15). It involves the period matrix
$`2\pi \mathrm{}\left(\tau _{11}\right)`$ $`=`$ $`{\displaystyle \frac{\alpha ^{}\mathrm{ln}(k_1)}{\alpha ^{}}}{\displaystyle \frac{2k_2(\eta 1)^2}{\eta }}+o(k_1^2,k_2^2),`$ (45)
$`2\pi \mathrm{}\left(\tau _{12}\right)`$ $`=`$ $`2\pi \mathrm{}\left(\tau _{21}\right)={\displaystyle \frac{\alpha ^{}\mathrm{ln}(\eta )}{\alpha ^{}}}{\displaystyle \frac{2k_1k_2(\eta +1)(\eta 1)^3}{\eta ^2}}+o(k_1^2,k_2^2),`$
$`2\pi \mathrm{}\left(\tau _{22}\right)`$ $`=`$ $`{\displaystyle \frac{\alpha ^{}\mathrm{ln}(k_2)}{\alpha ^{}}}{\displaystyle \frac{2k_1(\eta 1)^2}{\eta }}+o(k_1^2,k_2^2),`$
and its determinant
$`\text{det}^{d/2}(i\tau _{\mu \nu })`$ $`=`$ $`(2\pi )^d\left({\displaystyle \frac{\alpha ^{}\mathrm{ln}(k_1)\alpha ^{}\mathrm{ln}(k_2)\alpha ^2\mathrm{ln}^2(\eta )}{\alpha ^2}}\right)^{d/2}`$
$`\times (1\alpha ^{}d({\displaystyle \frac{(\eta 1)^2(\alpha ^{}\mathrm{ln}(k_1)k_1+\alpha ^{}\mathrm{ln}(k_2)k_2)}{\eta (\alpha ^{}\mathrm{ln}(k_1)\alpha ^{}\mathrm{ln}(k_2)\alpha ^2\mathrm{ln}^2(\eta ))}}`$
$`{\displaystyle \frac{2k_1k_2\alpha ^{}\mathrm{ln}(\eta )(\eta +1)(\eta 1)^3}{\eta ^2(\alpha ^{}\mathrm{ln}(k_1)\alpha ^{}\mathrm{ln}(k_2)\alpha ^2\mathrm{ln}^2(\eta ))}})+o(\alpha ^2,k_1^2,k_2^2)).`$
Together we find:
$`\left[dm\right]_2^0`$ $`=`$ $`{\displaystyle \frac{dk_1dk_2d\eta }{k_1^2k_2^2(1\eta )^2}}\text{det}^{d/2}(i\tau _{\mu \nu })`$
$`\times {\displaystyle \underset{\beta }{}}\left({\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}(1k_\beta ^m)^d{\displaystyle \underset{m=2}{\overset{\mathrm{}}{}}}(1k_\beta ^m)^2\right)(1k_1)^2(1k_2)^2`$
$`=`$ $`{\displaystyle \frac{dk_1dk_2d\eta }{k_1^2k_2^2}}\text{det}^{d/2}(i\tau _{\mu \nu })(1+(d2)(k_1+k_2)`$
$`+((d2)^2+{\displaystyle \frac{d(1\eta )^2(1+\eta ^2)}{\eta ^2}})k_1k_2+o(k_1^2,k_2^2)).`$
The derivation of these expressions is summarized in appendix A. We have always added powers of $`\alpha ^{}`$ to get finite SPT variables in the end. After substituting for the integration over moduli
$`{\displaystyle _\mathrm{\Gamma }}{\displaystyle \frac{dk_1dk_2d\eta }{k_1k_2\eta }}=\left({\displaystyle \frac{1}{\alpha ^{}}}\right)^3{\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{i=1}{\overset{3}{}}}(dt_i)`$ (48)
and regarding all relevant powers of the string tension, we are left with an over all factor $`1/\alpha ^{}`$.
The parametrization of figure 9 using $`\eta =A`$ translates into
$`t_1=\alpha ^{}\mathrm{ln}(A),t_2=\alpha ^{}\mathrm{ln}(k_1)t_1,t_3=\alpha ^{}\mathrm{ln}(k_2)t_1.`$ (49)
Now we have to extract the term proportional to $`A^0,k^0,\alpha ^{}`$ from the remaining integrand, which should belong to that part of the amplitude, i.e. that region in the moduli space, where all proper times stay finite. This corresponds to the sum of the two diagrams in (a) of figure 8. We get
$`𝒜_0^{(2)}|_{\text{(a)}}={\displaystyle \frac{g^2}{(4\pi )^d}}d(d2){\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{i=1}{\overset{3}{}}}(dt_i){\displaystyle \frac{2t_1+t_2+t_3}{\left(t_1t_2+t_1t_3+t_2t_3\right)^{1+d/2}}},`$ (50)
which differs from the result in by the symmetrized integration region. The integral vanishes anyway in dimensional regularization.
Investigating the possible pinching contributions to diagram (b) of figure 8 we find the curious situation that three different limits of the string amplitude contribute to a single Feynman diagram. The regions where $`t_2`$ or $`t_3`$ vanish, are defined by finite values for the moduli $`k_1/A`$ and $`k_2/A`$. We set $`k_iAk_i`$, expand in $`A`$ and the other multiplier, keeping the finite part. The expansion in $`\alpha ^{}`$ has to be performed without any prefactor in $`\alpha ^{}`$ as the missing proper time integration leads to the missing of a factor $`1/\alpha ^{}`$. The integrand we get is finite when $`k_i1`$ in both cases and the $`k`$-integration can be split off from the rest, giving the unique result
$`𝒜_0^{(2)}|_{\text{(b)}}^{t_{2,3}}={\displaystyle \frac{g^2}{(4\pi )^d}}(d2){\displaystyle \frac{dk}{k^2}_0^{\mathrm{}}\underset{i=1}{\overset{2}{}}(dt_i)\left(t_1t_2\right)^{d/2}}.`$ (51)
The question what should happen to the integration over $`k`$ was answered in in the sense that one has to take an infinitesimal region around $`k=1`$ in the integrand and drop the integration. Even more problematic is the treatment of the integral over the free modulus in the case of the third possible pinching. This we get by having $`t_10`$, i.e. $`A`$ finite. We do the same expansion as before:
$`𝒜_0^{(2)}|_{\text{(b)}}^{t_1}`$ $`=`$ $`{\displaystyle \frac{g^2}{(4\pi )^d}}{\displaystyle \frac{dA}{A^2(1A)^2}\left(d(1+A^4)2d(1+A^2)A+(d^22d+4)A^2\right)}`$
$`\times {\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}(dt_i)\left(t_1t_2\right)^{d/2}.`$
The integral over $`A`$ is divergent at $`0`$ and $`1`$. In this was cured by a rather arbitrary zeta function regularization after expanding the integrand around $`A=1`$ and omitting the divergency at $`A=0`$ completely:
$`{\displaystyle _{1ϵ}^1}{\displaystyle \frac{dA}{(1A)^2}}\zeta (0)={\displaystyle \frac{1}{2}}.`$ (53)
Similar divergent integrals have also been encountered in scalar $`\mathrm{\Phi }^4`$ theory and similar methods, involving “world sheet cut-offs”, have been used for regularization . The result is
$`𝒜_0^{(2)}|_{\text{(b)}}^{t_1}={\displaystyle \frac{g^2}{2(4\pi )^d}}(d2)^2{\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}(dt_i)\left(t_1t_2\right)^{d/2}.`$ (54)
As we are unable to propose any reasonable alternatives, we do not intend to criticise these methods in detail. As mentioned a comparison of the contributions found from the string amplitude to field theory is impossible anyway, both are vanishing by definition. But we have to conclude, that certain divergent integrals over free moduli are appearing, if one tries to follow the strategies of the naive pinching procedures. These divergencies appear exactly in two different types which cover all the cases we investigated. If the proper time variable which is associated with the multiplier $`k`$ of a loop becomes small, we find an integrand finite when $`k1`$ and the integration divergent of the kind $`dk/k^2`$, while any other sewing parameter $`\eta `$ leads to integrands that diverge like $`d\eta /(\eta ^2(1\eta )^2)`$. This shall be verified in the next chapter to be a generic feature of the pinching. The nature of the diverging integrals is very similar to those types of integrals that have to be performed when doing the trace over the Hilbert space of the intermediate string states when sewing together the world sheets of two strings, thus building up world sheets of higher Genus. In this sense we believe them to be related to the tachyon exchange in the shrinking propagator.
### 4.3 Problems concerning the overlapping of isometric circles
In this subsection we will explain a problem which arises when performing the $`k0`$ limit of the string worldsheet. Up to now, and as well in the following, we completely ignore the requirement, that the circles cut out off the complex plane around the fixed points in the Schottky construction of the Riemann surface are not overlapping. Strictly speaking, this induces further restrictions on the moduli space, the allowed values of Schottky parameters, which in principle can be relevant also in the low energy limit. For all contributions to scalar field theory this apprantly has not been the case, because the naive sewing has been successful. But as we are unable in general to resolve the same task for Yang-Mills diagrams, this cannot be ruled out. In fact, in the prescription derived in was used, which implied a further restriction on the sewing parameters, as compared to our results above. But, luckily or not, the difference can be tracked down to an overall factor of 2 from symmetrizing the integrands in , which does not appear to be very significant in deciding whether the modification is necessary or not. Still it is worthwhile to be investigated in more detail.
In the Schottky parametrization was employed with circles different from the isometric ones. The latter are located at infinity, which makes the same treatment more difficult. Starting from the parametrization of diagram 9 one adds the isometric circle around $`\eta _1=\eta `$ and $`\xi _1=1`$ and a circle of radius $`\sqrt{k_2}`$ around $`\eta _2=0`$ and its image with infinite radius but finite position at $`1/\sqrt{k_2}`$ around $`\xi _2=\mathrm{}`$. Now $`\eta `$ can no longer vary freely between 0 and 1, but the two circles are further required not to intersect. This translates into the inequalities
$`\sqrt{k_2}{\displaystyle \frac{\eta \sqrt{k_1}}{1\sqrt{k_1}}},0{\displaystyle \frac{1\sqrt{k_1}}{1+\sqrt{k_1}}}\left(1\eta \right).`$ (55)
The second inequality reduces to $`\eta 1`$ for small $`k_1`$, whereas the first one is still a complicated relation between fixed points and multipliers, and several regions of their values will contribute. In the single inequality
$`0\sqrt{k_2}\sqrt{k_1}\eta 1`$ (56)
was used. It is not sufficient to solve the above requirements, but was enough to produce results consistent with scalar field theory expectations. The same inequality, translated to a different region of integration, was also used in for the vacuum diagrams of the Yang-Mills theory, but again we cannot judge their results to be decisive. Anyway, in more complicated diagrams, a systematic evaluation of all regions in the moduli space which respect the generalized version of the above inequalities may be necessary. One can easily convince oneself that for diagrams with external legs, the inequalities get much more complicated and a variety of distinctions arises. For some scalar $`\mathrm{\Phi }^4`$ diagrams this task has been performed in . Also the methods we shall be employing to solve the unrestricted integrals do not allow a straightforward extension to this case, which leaves us with very little hope to be able to evaluate the analogous integrals.
### 4.4 Two-loop diagrams involving external gluons
We now come to our main topic, two-loop string diagrams with two external gluon vertex operators inserted. We shall again proceed along the lines we have established in the chapter on scalar theory and also we shall have to deal with the problems we already encountered in the previous two sections about Yang-Mills diagrams. This method should in principle allow to compute the two-loop cofficient of the Yang-Mills $`\beta `$ function
$`\beta =g\mu {\displaystyle \frac{Z_g}{\mu }}=g\left(\beta _0\left({\displaystyle \frac{g}{4\pi }}\right)^2+\beta _1\left({\displaystyle \frac{g}{4\pi }}\right)^4+o\left(g^6\right)\right).`$ (57)
It can be extracted from the two point function, if the field theoretical gauge is any background type gauge. We shall admit an arbitrary covariant gauge of the gauge field fluctuations, which is still allowed after the background gauge is chosen. In this kind of gauge the dependence of the charge renormalization constant on the subtraction scale can be deduced from the gauge field wave function renormalization constant alone. This in turn is given by the two point function:
$`Z_A=1+{\displaystyle \frac{\beta _0}{ϵ}}\left({\displaystyle \frac{g}{4\pi }}\right)^2+{\displaystyle \frac{\beta _1}{2ϵ}}\left({\displaystyle \frac{g}{4\pi }}\right)^4+o\left(g^6\right).`$ (58)
The coefficients take the values
$`\beta _0={\displaystyle \frac{11}{3}},\beta _1={\displaystyle \frac{34}{3}}.`$ (59)
They are gauge invariant and therefore pose a good criterion to test the validity of the methods employed to higher loop order. All purely gluonic diagrams that contribute in field theory are drawn in figure 10.
The further diagrams involving ghost fields are obtained by substituting gluon loops by ghost loops. All such are listed in reference , where also their contributions in the Feynman background gauge are displayed. By the help of M. Peter from the Heidelberg University we are also able to discuss the contributions of these diagrams in an arbitrary covariant background gauge. In table 1 they are summarized for $`d=42ϵ`$, using the same letters for the diagrams as in reference , which are not to be confused with our notation of figure 10. We see that to the leading order in the $`1/ϵ`$ expansion all diagrams are transverse, while this order necessarily adds up to zero. The next to leading order is not transverse diagram by diagram but the sum of all contributions of course is. It remains unclear what we should expect to happen to the counter term diagrams, a problem that could not be investigated at one-loop order and has not been in scalar field theory. We should be able to notice their absence by a deviation of our result from the correct $`\beta `$ function coefficient by just their contribution. A very natural assumption seems, that we are dealing with a “naked” theory with string derived diagrams, where no counter terms are present. It is unclear, how the string diagrams, which orginally have a natural cut-off, implement renormalization. The explicit form of counter terms conventionally depends on the renormalization scheme employed, which would pose another puzzle. In other words: If the effects of the one-loop renormalization are included automatically, the result is universal, while if not, it will depend on the scheme. As we are unable to resolve the problems concerning four-gluon vertices, we shall not be in the position to unravel this question anyway.
We now take (4) for $`h=2`$ and $`n=2`$ and use results from appendix A employing the usual coordinate fixing (15) to get
$`𝒜_2^{(2)}(p,p)`$ $`=`$ $`N^2\text{tr}\left(\lambda ^{a_1}\lambda ^{a_2}\right)C_2𝒩_0^2{\displaystyle _\mathrm{\Gamma }}{\displaystyle \frac{dk_1dk_2d\eta dz_1dz_2}{k_1^2k_2^2(1\eta )^2}}(1k_1)^2(1k_2)^2`$
$`\times \text{det}^{d/2}(i\tau _{\mu \nu }){\displaystyle \underset{\beta }{}}\left({\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}(1k_\beta ^m)^d{\displaystyle \underset{m=2}{\overset{\mathrm{}}{}}}(1k_\beta ^m)^2\right)\mathrm{exp}\left(2\alpha ^{}p^2𝒢^{(2)}(z_1,z_2)\right)`$
$`\times \left((2\alpha ^{})(ϵ_1p_2)(p_2ϵ_1)_{z_1}𝒢^{(2)}(z_1,z_2)_{z_2}𝒢^{(2)}(z_1,z_2)+(ϵ_1ϵ_2)_{z_1}_{z_2}𝒢^{(2)}(z_1,z_2)\right),`$
where the expansion in the multipliers and $`\alpha ^{}`$ is partly done already. The determinant of the period matrix can be read off from (4.2). The prime form is
$`E^{(2)}(z_1,z_2)`$ $`=`$ $`{\displaystyle \frac{\alpha ^{}\mathrm{ln}(z_1z_2)}{\alpha ^{}}}(z_1z_2)^2({\displaystyle \frac{k_1(\eta 1)^2}{(z_1\eta )(z_2\eta )(z_11)(z_21)}}+{\displaystyle \frac{k_2}{z_1z_2}}`$
$`+{\displaystyle \frac{k_1k_2(\eta 1)^2}{\eta ^2}}({\displaystyle \frac{\eta ^2+z_1z_2}{z_1z_2(z_11)(z_21)}}+{\displaystyle \frac{\eta ^2(1+z_1z_2)}{z_1z_2(z_1\eta )(z_2\eta )}}))+o(k_1^2,k_2^2)`$
and the abelian integrals follow
$`\left({\displaystyle _{z_1}^{z_2}}\omega ^1\right)^{(2)}`$ $`=`$ $`{\displaystyle \frac{\alpha ^{}}{\alpha ^{}}}\mathrm{ln}\left({\displaystyle \frac{(z_1\eta )(z_21)}{(z_11)(z_2\eta )}}\right){\displaystyle \frac{k_2(\eta 1)(z_2z_1)(z_1z_2+\eta )}{\eta z_1z_2}}`$
$`+{\displaystyle \frac{k_1k_2(\eta 1)^3(z_1z_2)}{\eta ^2}}\left({\displaystyle \frac{\eta +1}{(z_11)(z_21)}}+{\displaystyle \frac{\eta (\eta +1)}{(z_1\eta )(z_2\eta )}}\right)+o(k_1^2,k_2^2),`$
$`\left({\displaystyle _{z_1}^{z_2}}\omega ^2\right)^{(2)}`$ $`=`$ $`{\displaystyle \frac{\alpha ^{}}{\alpha ^{}}}\mathrm{ln}\left({\displaystyle \frac{z_1}{z_2}}\right)+{\displaystyle \frac{k_1(\eta 1)^2(z_1z_2)}{\eta }}\left({\displaystyle \frac{1}{(z_11)(z_21)}}+{\displaystyle \frac{\eta }{(z_1\eta )(z_2\eta )}}\right)`$
$`+{\displaystyle \frac{k_1k_2(\eta 1)^2(z_1z_2)}{\eta ^2}}\left(\eta +1+{\displaystyle \frac{\eta (\eta +1)}{z_1z_2}}\right)+o(k_1^2,k_2^2),`$
finally the inverse of the period matrix:
$`\left(2\pi \mathrm{}\left(\tau _{\mu \nu }\right)\right)^1`$ $`=`$ $`\text{det}^1\left(2\pi \mathrm{}\left(\tau _{\mu \nu }\right)\right)\left(\begin{array}{ccc}\tau _{22}& \tau _{12}& \\ \tau _{21}& \tau _{11}& \end{array}\right),`$ (65)
$`\text{det}^1\left(2\pi \mathrm{}\left(\tau _{\mu \nu }\right)\right)`$ $`=`$ $`\left({\displaystyle \frac{\alpha ^{}\mathrm{ln}(k_1)\alpha ^{}\mathrm{ln}(k_2)\alpha ^2\mathrm{ln}^2(\eta )}{\alpha ^2}}\right)^1`$
$`\times (1\alpha ^{}({\displaystyle \frac{2(\eta 1)^2(\alpha ^{}\mathrm{ln}(k_1)k_1+\alpha ^{}\mathrm{ln}(k_2)k_2)}{\eta (\alpha ^{}\mathrm{ln}(k_1)\alpha ^{}\mathrm{ln}(k_2)\alpha ^2\mathrm{ln}^2(\eta ))}}`$
$`{\displaystyle \frac{4k_1k_2\alpha ^{}\mathrm{ln}(\eta )(\eta +1)(\eta 1)^3}{\eta ^2(\alpha ^{}\mathrm{ln}(k_1)\alpha ^{}\mathrm{ln}(k_2)\alpha ^2\mathrm{ln}^2(\eta ))}})+o(\alpha ^2,k_1^2,k_2^2)).`$
For the product over the primary classes we use the expansion
$`{\displaystyle \underset{\beta }{}}\left({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1k_\beta ^n)^d{\displaystyle \underset{n=2}{\overset{\mathrm{}}{}}}(1k_\beta ^n)^2\right)(1k_1)^2(1k_2)^2`$
$`=`$ $`1+(d2)(k_1+k_2)+\left((d2)^2+{\displaystyle \frac{d(1\eta )^2(1+\eta ^2)}{\eta ^2}}\right)k_1k_2+o(k_1^2,k_2^2)`$
and the prefactor is
$`C_2𝒩_0^2={\displaystyle \frac{g^4}{4}}\left((2\pi )\sqrt{2\alpha ^{}}\right)^{2d}(2\alpha ^{})^3.`$ (67)
All expansions that have to be performed in the following are clearly impossible to be done by hand, so we regularly used the algebraic computer program MAPLE. We now treat all the diagrams of figure 10 one by one and keep track not to loose any contributions from the string amplitude. We start with (a), parametrize according to figure 3, substitute sewing parameters as in (17) and finally define proper time variables by (18). The integration is transformed as
$`{\displaystyle _\mathrm{\Gamma }}{\displaystyle \frac{dk_1dk_2d\eta dz_1dz_2}{k_1k_2\eta z_1z_2}}=\left({\displaystyle \frac{1}{\alpha ^{}}}\right)^5{\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{i=1}{\overset{5}{}}}(dt_i).`$ (68)
This corresponds to a sewing procedure that sewed the diagram together as depicted in figure 4. We next choose the insertion points of the external states to lie at the “large” loop, i.e. $`\eta _2=0`$ and $`\xi _2=\mathrm{}`$. The possibly astonishing fact, that the fixing of coordinates on the string world sheet leads to a distiction of different loops when looking for contributions to particular field theoretical Feynman diagrams has already been discussed in . It has been found how the local coordinates have to be chosen in order to get correct Green’s functions, which in our case requires us to take $`V_i^{}(0)=z_i`$, as expected. The only way to have the external states on different loops is then given by the parametrization of figure 3, except up to interchanging the external states, which does not give any different result, as has been checked explicitely in all cases. Using further (19) and substitute everything into (4.4) we find
$`𝒜_2^{(2)}\left(p^2\right)|_{\text{(a)}}`$ $`=`$ $`N^2\delta _{a_1a_2}{\displaystyle \frac{g^4}{(4\pi )^d}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{i=1}{\overset{5}{}}}(dt_i)`$
$`\times {\displaystyle \frac{\left((ϵ_1p_2)(ϵ_2p_1)P_5^{(a)}(t_1,\mathrm{},t_5,d)+ϵ^2P_4^{}^{(a)}(t_1,\mathrm{},t_5,d)\right)}{\left((t_1+t_2)(t_3+t_4+t_5)+(t_3+t_4)t_5\right)^{d/2+3}}}`$
$`\times \mathrm{exp}\left(p^2{\displaystyle \frac{t_1t_2(t_3+t_4+t_5)+t_1t_3(t_4+t_5)+t_2t_4(t_3+t_5)+t_3t_4t_5}{(t_1+t_2)(t_3+t_4+t_5)+(t_3+t_4)t_5}}\right)`$
$`=`$ $`N^2\delta _{a_1a_2}\left({\displaystyle \frac{g}{4\pi }}\right)^4\left({\displaystyle \frac{4\pi }{\sqrt{p^2}}}\right)^ϵ\left((ϵ_1p_2)(ϵ_2p_1)I_5^{(a)}(ϵ)+(ϵ^2p^2)I_4^{}^{(a)}(ϵ)\right).`$
The integrals $`I^{(a)}`$ must be computed in dimensional regularization. Each of the polynomials $`P^{(a)}`$ contains a couple of hundred terms and is therefore not written explicitely, instead we only display its degree as an index. In the second line we have extracted the momentum dependence in $`d=4+ϵ`$ dimensions by rescaling the integration variables.
Power counting reveals that we have to deal at least with some logarithmic UV divergencies proportional to $`1/ϵ`$ for small values of the SPT variables, whereas for $`t_i\mathrm{}`$ all integrals are IR finite. In a first step one can substitute
$`t_1`$ $`=`$ $`tx_1x_2x_3x_4,`$ (70)
$`t_2`$ $`=`$ $`t(1x_1)x_2x_3x_4,`$
$`t_3`$ $`=`$ $`t(1x_2)x_3x_4,`$
$`t_4`$ $`=`$ $`t(1x_3)x_4,`$
$`t_5`$ $`=`$ $`t(1x_4),`$
$`{\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{i=1}{\overset{5}{}}}(dt_i)`$ $`=`$ $`{\displaystyle _0^1}{\displaystyle \underset{i=1}{\overset{4}{}}}(dx_i)x_2x_3^2x_4^3{\displaystyle _0^{\mathrm{}}}𝑑tt^4`$
to do the trivial integration over the sum of all $`t_i`$, which gives
$`I_5^{(a)}(ϵ)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑tt^{1ϵ}e^t{\displaystyle _0^1}{\displaystyle \underset{i=1}{\overset{4}{}}}(dx_i){\displaystyle \frac{x_2x_3^2x_4^3\overline{P}_5^{(a)}(x_1,x_2,x_3,x_4,ϵ)}{\left(x_4(1x_4(1x_2x_3(1x_2x_3)))\right)^{5+ϵ/2}}}`$
$`\times \left({\displaystyle \frac{\overline{P}_3(x_1,x_2,x_3,x_4)}{x_4(1x_4(1x_2x_3(1x_2x_3)))}}\right)^ϵ,`$
where the $`\overline{P}^{(a)}`$ are the numerators after the substitution. Investigating possible poles and finding them harmless we can expand the integrand and replace
$`\left({\displaystyle \frac{\overline{P}_3(x_1,x_2,x_3,x_4)}{x_4(1x_4(1x_2x_3(1x_2x_3)))}}\right)^ϵ=1+o(ϵ).`$ (72)
The integration over $`t`$ can be done using (D.4) and yields
$`\mathrm{\Gamma }(ϵ)={\displaystyle \frac{1}{ϵ}}\gamma +o(ϵ).`$ (73)
We shall later find that the first term of the integrand in (4.4) after substituting according to (72) has only $`1/ϵ`$ divergencies. This allows to extract the $`1/ϵ^2`$ divergencies of the integrals, which originate from the poles of this first term, without regarding corrections proportional to $`ϵ`$ from the second factor. This calculation is therefore able to reveal the leading order divergencies. It is completed in appendix B. There we also present an algorithm who enables an exact computation of all the integrals considered. Both methods lead to consistent results that are independently obtained. Hence we are very confident about the correctness of our calculations.
Regarding figure 4 we notice that diagram (a) can have pinching limits conbtributing to diagram (c) if $`t_i0`$ for $`i\{1,\mathrm{},4\}`$, while it reduces to diagram (e) when $`t_50`$. Further we have to perform the limit in a way that $`t_1`$ and $`t_30`$ or $`t_2`$ and $`t_40`$ to end up with diagram (d). The choices $`t_1,t_2,t_30`$ are obtained by assigning finite values to the sewing parameters $`A_i`$ and after permuting the integration variables they are all of the form
$`I_4^{(c)}(ϵ)`$ $`=`$ $`{\displaystyle \frac{dA}{A^2(1A)^2}_0^{\mathrm{}}\underset{i=1}{\overset{4}{}}(dt_i)\frac{P_4^{(c)}(t_1,\mathrm{},t_4,A,ϵ)}{\left((t_1+t_2)(t_3+t_4)+t_3t_4\right)^{4+ϵ/2}}}`$
$`\times \mathrm{exp}\left({\displaystyle \frac{P_3^{(c)}(t_1,\mathrm{},t_4)}{(t_1+t_2)(t_3+t_4)+t_3t_4}}\right).`$
The exponent and the denominator are simply got by setting the appropriate proper time variables in (4.4) to zero, whereas the numerator has to be computed separately by adding a power of $`\alpha ^{}`$ into the expansion of the integrand. The integration of those sewing variables that are left finite displays precisely the divergencies already encountered in the previous chapter when we discussed two-loop vacuum diagrams. We do not decide what to do with them at this stage but expand the integrand around $`A=1`$ and keep the finite as well as the divergent term. We then compute the proper time integration and mark the respective terms of the expansion in $`A`$ by their possibly divergent prefactors
$`O_2{\displaystyle \frac{dA}{(1A)^2}},O_0{\displaystyle 𝑑A}.`$ (75)
We have for some cases even computed the integrations over proper times without doing any expansion of the integrand in the free moduli finding a more complicated structure of five instead of two terms but no deeper insight into the uncertainties of the pinching procedure. Therefore we restrict ourselves at this point to state the two types of terms we mentioned. The expression resulting from letting $`t_40`$, which contributes to diagram (c), differs only in the respect, that keeping $`k_2/(A_1A_2A_3)`$ finite we end up with the undefined integral $`dk_2/k_2^2`$, which we, too, already found in the two-loop vacuum case. Then, of course, also $`P_4^{(c)}`$ depends on $`k_2`$ but is finite in the limit $`k_21`$. We expand around $`k_2=1`$, mark the finite term by the factor
$`O_1{\displaystyle 𝑑k_2}`$ (76)
and omit terms of the order $`o(k_21)`$.
We are now only left with contributions to (d) and (e). The case $`t_50`$ is completely analogous to the one mentioned previously with the only exception that the integration looks somewhat easier:
$`I_4^{(e)}(ϵ)`$ $`=`$ $`{\displaystyle \frac{dk_1}{k_1^2}_0^{\mathrm{}}\underset{i=1}{\overset{4}{}}(dt_i)\frac{P_4^{(e)}(t_1,\mathrm{},t_4,k_1,ϵ)}{\left((t_1+t_2)(t_3+t_4)\right)^{4+ϵ/2}}\mathrm{exp}\left(\frac{P_3^{(e)}(t_1,\mathrm{},t_4)}{(t_1+t_2)(t_3+t_4)}\right)}.`$ (77)
We shall explore this by factorizing the integrals. In the last case of diagram (d) we do not find any contributions at all, simply because there are two proper time variables that are absent. This leads to two factors of $`\alpha ^{}`$ in front of the integral, which then vanishes proportional to the inverse string tension.
Having completed the study of contributions to diagram (a) and any diagrams including four-gluon vertices that arise from it, we now proceed to diagram (b) and its possible pinchings. We have to distinguish, which internal propagator the external states are sitting at, by regarding both of the two parametrizations drawn in figure 11 and 13.
Firgures 12 and 14 demonstrate how the left and right loop are to be distinguished. The sewing parameters of the configuration (i) are given by
$`\eta =A_1,z_1=A_1A_2,z_2=A_1A_2A_3,`$ (78)
those of (ii) by
$`z_1=A_1,z_2=A_1A_2,\eta =A_1A_2A_3.`$ (79)
For both cases all but one proper time variables can be defined in a unique manner
$`t_i`$ $`=`$ $`\alpha ^{}\mathrm{ln}(A_i)\text{for}i=1,2,3,`$ (80)
$`t_4`$ $`=`$ $`\alpha ^{}\mathrm{ln}(k_2)t_1t_2t_3,`$
but the one of the smaller loop is different:
$`t_5`$ $`=`$ $`\alpha ^{}\mathrm{ln}(k_1)t_1\text{for (i)},`$ (81)
$`t_5`$ $`=`$ $`\alpha ^{}\mathrm{ln}(k_1)t_1t_2t_3\text{for (ii)}.`$
This makes a distinction in the translation of further logarithms of moduli appearing in the integrand of the amplitude necessary. Configuration (i) has to be used with
$`\alpha ^{}\mathrm{ln}(z_1z_2)`$ $`=`$ $`t_1t_2+\alpha ^{}\mathrm{ln}(1A_3),`$ (82)
$`\alpha ^{}\mathrm{ln}(z_1z_2)`$ $`=`$ $`2t_12t_2t_3,`$
$`\alpha ^{}\mathrm{ln}\left({\displaystyle \frac{z_1}{z_2}}\right)`$ $`=`$ $`t_3,`$
$`\alpha ^{}\mathrm{ln}\left({\displaystyle \frac{(z_1\eta )(z_21)}{(z_11)(z_2\eta )}}\right)`$ $`=`$ $`\alpha ^{}\mathrm{ln}\left({\displaystyle \frac{(1A_2)(1A_1A_2A_3)}{(1A_1A_2)(1A_2A_3)}}\right),`$
whereas (ii) with:
$`\alpha ^{}\mathrm{ln}(z_1z_2)`$ $`=`$ $`t_1+\alpha ^{}\mathrm{ln}(1A_2),`$ (83)
$`\alpha ^{}\mathrm{ln}(z_1z_2)`$ $`=`$ $`2t_1t_2,`$
$`\alpha ^{}\mathrm{ln}\left({\displaystyle \frac{z_1}{z_2}}\right)`$ $`=`$ $`t_2,`$
$`\alpha ^{}\mathrm{ln}\left({\displaystyle \frac{(z_1\eta )(z_21)}{(z_11)(z_2\eta )}}\right)`$ $`=`$ $`t_2+\alpha ^{}\mathrm{ln}\left({\displaystyle \frac{(1A_2A_3)(1A_1A_2)}{(1A_1)(1A_3)}}\right).`$
Despite the variables $`t_2`$ and $`t_4`$ featuring independently in the case (i), the integrand will only depend on their sum $`t_2+t_4`$, which is also true for $`t_1`$ and $`t_3`$ in the case of (ii). This reflects that the two propagators parametrized by these two proper times carry the same field theoretical momentum. By doing the usual expansion of the integrand we then obtain the amplitude
$`𝒜_2^{(2)}\left(p^2\right)|_{\text{(b)}}`$ $`=`$ $`N^2\delta _{a_1a_2}{\displaystyle \frac{g^4}{(4\pi )^d}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{i=1}{\overset{5}{}}}(dt_i)`$
$`\times {\displaystyle \frac{\left((ϵ_1p_2)(ϵ_2p_1)P_5^{(b)}(t_1,\mathrm{},t_5,d)+ϵ^2P_4^{}^{(b)}(t_1,\mathrm{},t_5,d)\right)}{\left((t_1+t_5)(t_2+t_3+t_4)+t_1t_5\right)^{d/2+3}}}`$
$`\times \mathrm{exp}\left(p^2{\displaystyle \frac{t_3((t_1+t_5)(t_2+t_4)+t_1t_5)}{(t_1+t_5)(t_2+t_3+t_4)+t_1t_5}}\right)`$
$`=`$ $`N^2\delta _{a_1a_2}\left({\displaystyle \frac{g}{4\pi }}\right)^4\left({\displaystyle \frac{4\pi }{\sqrt{p^2}}}\right)^ϵ\left((ϵ_1p_2)(ϵ_2p_1)I_5^{(b)}(ϵ)+(ϵ^2p^2)I_4^{}^{(b)}(ϵ)\right).`$
The sewing configurations (i) und (ii) lead to different explicit polynomials in the numerator but all are of identical structure, so that we do not have to treat them separately in this general discussion. We now follow the same strategies as earlier and cut our arguments short accordingly. From the figures 12 and 14 one immediately notices all relevant pinching contributions. Having $`t_20`$ or $`t_40`$ in region (i) we get diagram (c), if both vanish (d). The same is true in region (ii) for $`t_1`$ and $`t_3`$. The moduli divergencies we find are identical to those found in the pinching contributions derived from diagram (a). We use again for vanishing $`t_1`$, $`t_2`$ or $`t_3`$ the prefactors $`O_0`$ and $`O_2`$ for the divergency and if $`t_4`$ is small we have $`O_1`$. A contribution to diagram (e) does not exist and contributions to (d) go to zero proportional to $`\alpha ^{}`$ again.
We have finally got all possible contributions to field theoretical Feynman diagrams that can be extracted from the appropriate string amplitude given a particular choice of local coordinates and fixed points. The results of the calculation of the integrals are summarized in table 2 and 3, while the details of the integration are postponed to appendix B.
All terms proportional to the Euler constant or $`\mathrm{ln}\left(4\pi /p^2\right)`$ are dropped as they come out correctly automatically and can be restored easily. We have also changed our conventions in the favour of using $`d=42ϵ`$ in order to compare results to the field theory. The given numbers still leave a lot of space to possible interpretations. Even only regarding the contributions to diagram (a) from figure 10 reveal that there is no choice of the field theoretical gauge parameter, which lets the results coincide with the sum of the diagrams (i) + (j) + (k) of the same topology from table 1, if one also demands the sum of the two parametrizations (i) and (ii) of diagram (b) to coincide with (a) and (b) from table 1. An idenfication diagram by diagram seems to be ruled out therefore. The next would be to try to compare the sum of all diagrams. Unambiguously are
$`\left({\displaystyle \frac{35}{12}}+{\displaystyle \frac{11}{6}}\right)ϵ^2+\left({\displaystyle \frac{335}{24}}+{\displaystyle \frac{127}{18}}{\displaystyle \frac{14}{9}}\right)ϵ^1={\displaystyle \frac{19}{4}}ϵ^2+{\displaystyle \frac{467}{24}}ϵ^1`$ (85)
from diagram (a) and (b) without any four-gluon vertices, while the diagrams including such give
$`\left(3\left(O_2{\displaystyle \frac{9}{24}}+O_0{\displaystyle \frac{12}{16}}\right)+O_1{\displaystyle \frac{39}{48}}\right)ϵ^2+\left(3\left(O_2{\displaystyle \frac{87}{48}}+O_0{\displaystyle \frac{92}{32}}\right)+O_1{\displaystyle \frac{315}{96}}\right)ϵ^1.`$ (86)
They depend on the undefined diverging integrals over the free moduli, which should be replaced by some finite number in the manner of a regularization prescription. The most simple criterion for consistency could be seen to be the vanishing of the leading order of the expansion in $`1/ϵ`$. As we deal with two unknown parameters this cannot give a unique answer and the general simplicity of the occurring numbers, which on the other hand might seem encouraging, allows more speculation about such a regularization prescription than we dare to present in this spot. For instance, using only the quadratic divergency $`O_2`$ and dropping $`O_0`$ completely, then demanding the $`1/ϵ^2`$ term to vanish, one finds $`(7/2)ϵ^1`$. A surprisingly simple number, but it is simply wrong, neither is such a kind of discusion sufficiently rigorous in any manner.
There are a couple of possibly sensible modifications of the general expansion method we used, that may have passed unnoticed in scalar theory and the one-loop Yang-Mills computation. One can for instance think of the ambiguity concerning which loop an external particle sits at and demand to add up all possibilities to attach a given number of external states to a diagram, by distinguishing all the different loops, even if the topolgy of the diagram created is identical. Following this would mean that one had to choose different local coordinates in the vicinity if the external states
$`V_i^{}(0)=\left|{\displaystyle \frac{(z_i\eta _1)(z_i\xi _1)}{\eta _1\xi _1}}\right|=\left|{\displaystyle \frac{(z_i\eta )(z_i1)}{\eta 1}}\right|.`$ (87)
If we omit reducible diagrams, whose identification has been demonstrated in , the existing possibilities are listed diagrammatically in appendix C, where we also give the appropriate sewing translation and the results that are computed by solving the integrals which are found after expanding in all the relevant parameters. We have not displayed the pinching contributions which can be derived from these diagrams. The whole situation seems far too involved to allow any conclusion concerning the related field theoretical gauge choice and its relation to the coordinate fixing on the world sheet of the string, as the divergencies that are apearing when performing the pinching of these kind of diagrams are of the form
$`{\displaystyle 𝑑A\frac{1}{A^3(1A)}}\text{and}{\displaystyle 𝑑A\frac{1}{A(1A)^3}},`$ (88)
which is different from that already encountered. We can neither draw any new conclusion with respect to the values of the ambiguous integrals.
At the moment we are therefore unable to give any definite result.
Acknowledgements
We would like to thank Stefan Heusler for collaboration in earlier stages of this work and Markus Peter for substantial support in evaluating various integrations. M. G. Schmidt thanks Paolo Di Vecchia for exciting lectures and discussions 1996 in Heidelberg and kind hospitality in Copenhagen. We further thank him, A. Lerda, R. Marotta and R. Russo for important communications, as well as J. Fuchs and C. Schweigert for helpful explanations.
## Appendix A The Schottky representation of Riemann surfaces
The most important issues about the Schottky group and the representation of Riemann surfaces by its means are given in , while the method used to expand geometrical quantities in the Schottky multipliers is not published in the literature. We follow in most respects. Some more general useful facts about abelian integrals and the prime form are detailed in .
### A.1 The Schottky group
The convenient parametrization of automorphisms $`t(z)`$ of the compactified complex plane is written
$`t(z)={\displaystyle \frac{az+b}{cz+d}}\text{with}a,b,c,d\text{and}adbc=1.`$ (A.1)
The matrices
$`\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\text{with}\text{det}\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)=1`$ (A.6)
are in $`SL(2,)`$, which is even isomorphic to the automorphism group. The Schottky parameters are then defined by
$`{\displaystyle \frac{t(z)\xi }{t(z)\eta }}=k{\displaystyle \frac{z\xi }{z\eta }},`$ (A.7)
where $`\xi `$ and $`\eta `$ are the fixed points of the mapping $`t(z)`$, $`k`$ its multiplier. The invariance under $`\eta \xi `$ and $`kk^1`$ allows to choose $`|k|1`$. The relation to the former parametrization is found to be
$`t(z)={\displaystyle \frac{\eta (z\xi )k\xi (z\eta )}{(z\xi )k(z\eta )}}.`$ (A.8)
And we read off:
$`a`$ $`=`$ $`{\displaystyle \frac{\eta k\xi }{\sqrt{|k|}|\eta \xi |}},b={\displaystyle \frac{\eta \xi (1k)}{\sqrt{|k|}|\eta \xi |}},`$
$`c`$ $`=`$ $`{\displaystyle \frac{1k}{\sqrt{|k|}|\eta \xi |}},d={\displaystyle \frac{k\eta \xi }{\sqrt{|k|}|\eta \xi |}}.`$
The Schottky group $`G_S^g`$ is defined by composition and inversion of a given set of generators $`\{t_1,\mathrm{},t_g\}`$. It therefore consists of all mappings of the form $`t_{i_1}^{j_1}t_{i_2}^{j_2}\mathrm{}t_{i_m}^{j_m}`$, having indices in $`\{1,\mathrm{},g\}`$ and integer exponents.
Projective mappings transform circles into circles and for each there is a particular pair of isometric circles which have identical radii. The interior of the original circle is mapped onto the exterior of its image and the exterior of the origin onto the interior of its image. This induces a one to one identification of points and a fundamental region can be chosen to be the exterior of both circles which are identified. This construction of equivalence classes is a manifold that has the topology of a Riemann surface $`\mathrm{\Sigma }_g`$ of Genus $`g`$, if this is the number of identified, disjoint circles. For technical reasons one excludes the set of points $`𝒟`$, where the orbits of the Schottky group become dense:
$`\{C\{\mathrm{}\}\}/G_S𝒟\mathrm{\Sigma }_g`$ (A.9)
The attractive and repulsive fixed points then lie inside the original, respectively the image isometric circle. Their radii $`r`$ and $`\overline{r}`$ and the coordinates of their centres are:
$`r`$ $`=`$ $`\overline{r}=\sqrt{|k|}{\displaystyle \frac{|\eta \xi |}{|1k|}},`$ (A.10)
$`c`$ $`=`$ $`{\displaystyle \frac{\xi k\eta }{1k}},\overline{c}={\displaystyle \frac{\eta k\xi }{1k}}.`$
The dual of the homology basis of non contractable cycles is the cohomology basis in the sense of de Rham cohomology. It consists of the first abelian differentials
$`\omega _\mu {\displaystyle \underset{G_S(0,\mu )}{}}\left({\displaystyle \frac{1}{zt(\eta _\mu )}}{\displaystyle \frac{1}{zt(\xi _\mu )}}\right)dz.`$ (A.11)
The sum is performed over all Schottky mappings that do not carry a power of $`t_\mu `$ on their right. The $`\omega _\mu `$ are analytic on all of $`\mathrm{\Sigma }_g`$. One further has to demand the normalization of the integrals
$`{\displaystyle _{a_\nu }}\omega ^\mu =2\pi i\delta _{\mu \nu }`$ (A.12)
along the paths $`a_\nu `$, that constitute one half of the homology basis to render the $`\omega ^\mu `$ unique. The other half of the integrals
$`{\displaystyle _{b_\nu }}\omega _\mu 2\pi i\tau _{\mu \nu }`$ (A.13)
are the entries of the period matrix $`\tau _{\mu \nu }`$ of the Riemann surface. Another geometrical function we shall need is the prime form which is characterized by its local behaviour
$`E(z_i,z_j)(z_iz_j),\text{for}z_iz_j`$ (A.14)
and its transformation property under conformal changes of coordinates
$`E(V_i(z_i),V_j(z_j))={\displaystyle \frac{E(z_i,z_j)}{\sqrt{V_i^{}|_{V_i^1(z_i)}V_j^{}|_{V_j^1(z_j)}}}}.`$ (A.15)
The defintion is unique, if also the invariance under transport around $`a_\nu `$ cycles and a phase factor for $`b_\nu `$ cycles is demanded. One can always choose coordinates so that all the fixed points and the centre coordinates of isometrical circles are lying on the real axis, which greatly simplifies all figures we display. The coordinate invariance of string theory further allows to fix three of the moduli of the world sheet by global diffeomorphisms, which is conveniently employed to fix one at $`0`$, another at $`\mathrm{}`$ and finally one at $`1`$.
The two generalizations of the representation of Riemann surfaces one eventually needs in string theory are those to surfaces with boundaries and to surfaces that are not orientable. The former type of world sheets for open string theories is obtained by cutting holes into the complex plane, the closed string worldsheet. If one exclusively has to deal with open strings it is convenient to take only half of the complex plane and choose all the fixed points of the Schottky group that generates the desired loops as well as the insertion points of the external states to lie on its boundary, taken to be the real axis. Passing from points to equivalence classes then includes identifying $`2g`$ semi-circles whose centre coordinates are real.
Non-orientable manifolds are constructed by including inversions into the generators of the Schottky group, which results in inserting projective planes $`^2`$ into the world sheet.
### A.2 Expanding in Schottky multipliers
We now demonstrate the techniques which allow to expand the geometrical quantities involved in the two-loop string amplitude into power series in the two Schottky multipliers $`k_1`$ and $`k_2`$, which then enables us to extract the relevant part of the amplitude. It is necessary to identify all terms up to the order $`k_1^1k_2^1`$, to cancel the prefactors coming from the integration measure of the amplitude.
We exploit the isomorphism from the projective group onto $`SL(2,)`$ and use the basis invariance of trace and determinant. One can read off the matrix $`T(z)`$ corresponding to a Schottky generator $`t(z)`$ from (A.8)
$`T(z)={\displaystyle \frac{1}{\sqrt{|k|}|\eta \xi |}}\left(\begin{array}{cc}\xi k\eta & \eta \xi (1k)\\ 1k& \eta +k\xi \end{array}\right)`$ (A.18)
and gets:
$`{\displaystyle \frac{\text{det}\left(T(z)\right)}{(\text{tr}\left(T(z)\right))^2}}={\displaystyle \frac{k(\eta \xi )^2}{((1+k)(\eta \xi ))^2}}={\displaystyle \frac{k}{(1+k)^2}}=k+o(k^2).`$ (A.19)
This gives a first order expression for the multiplier $`k(t_1(z)t_2(z))`$ of $`t(z)=t_1(z)t_2(z)`$:
$`k(t_1(z)t_2(z))`$ $`=`$ $`{\displaystyle \frac{\text{det}(T_1(z)T_2(z))}{(\text{tr}(T_1(z)T_2(z)))^2}}+o(k^2(t_1(z)),k^2(t_2(z)))`$
$`=`$ $`{\displaystyle \frac{\text{det}(T_1(z))\text{det}(T_2(z))}{(\text{tr}(T_1(z)T_2(z)))^2}}+o(k^2(t_1(z)),k^2(t_2(z)))`$
$`=`$ $`k(t_1(z))k(t_2(z)){\displaystyle \frac{(\eta _1\xi _1)^2(\eta _2\xi _2)^2}{(\xi _1\eta _2)^2(\xi _2\eta _1)^2}}+o(k^2(t_1(z)),k^2(t_2(z))),`$
where we had to use
$`\text{tr}(T_1(z)T_2(z))={\displaystyle \frac{1}{\sqrt{|k_1||k_2|}|\eta _1\xi _1||\eta _2\xi _2|}}\left((\xi _1\eta _2)(\xi _2\eta _1)+o(k(t_1(z)),k(t_2(z)))\right).`$ (A.21)
We simplify our notation and write $`k_ik(t_i(z))`$, dropping the argument $`z`$. Multipliers of mappings of the kind $`t_1t_2^1`$ are obtained by replacing $`\eta _2\xi _2`$ in (A.2):
$`k(t_1t_2^1)=k_1k_2{\displaystyle \frac{(\eta _1\xi _1)^2(\eta _2\xi _2)^2}{(\xi _1\xi _2)^2(\eta _1\eta _2)^2}}+o(k_1^2,k_2^2).`$ (A.22)
In fact, the two cases we have computed are already covering all types of multipliers needed. The symmetries of (A.2) and (A.22) show that
$`k(t_1t_2)`$ $`=`$ $`k(t_2t_1)=k(t_1^1t_2^1)=k(t_2^1t_1^1),`$ (A.23)
$`k(t_1^1t_2)`$ $`=`$ $`k(t_1t_2^1)=k(t_2^1t_1)=k(t_2t_1^1)`$
holds. One easily deduces the complete set of Schottky mappings that might have multipliers which are only of first order in both the multipliers of the two generators to be the subset
$`G_S^2G_S^{(2)}\{id,t_1,t_1^1,t_2,t_2^1,t_1t_2,t_2t_1,t_1^1t_2,t_2t_1^1,t_1t_2^1,t_2^1t_1.t_1^1t_2^1,t_2^1t_1^1\}`$
A useful asymptotic expression for the mapping itself is
$`t(z)=\xi +k(\eta \xi ){\displaystyle \frac{z\xi }{z\eta }}+o(k^2),`$ (A.24)
which is obtained from (A.8). Because of $`|k|1`$ this explaines the asymptotic properties of a Schottky map:
$`\underset{n\mathrm{}}{lim}t^n(z)=\xi ,\underset{n\mathrm{}}{lim}t^n(z)=\eta .`$ (A.25)
We shall also need
$`{\displaystyle \frac{(z_1t(w_1))(z_2t(w_2))}{(z_1t(w_2))(z_2t(w_1))}}=1+k{\displaystyle \frac{(\xi \eta )^2(z_1z_2)(w_1w_2)}{(z_1\xi )(z_2\xi )(w_1\eta )(w_2\eta )}}+o(k^2).`$
Now we can expand the period matrix, the abelian integrals, the normalization constant $`𝒩`$ of the partition function and the prime form to first order. An explicit expression for the prime form can easily be found from the definition given in the previous section:
$`E(z_1,z_2)(z_1z_2){\displaystyle \underset{tG_S\{id\}}{}}\sqrt{{\displaystyle \frac{(z_1t(z_2))(z_2t(z_1))}{(z_1t(z_1))(z_2t(z_2))}}},`$ (A.26)
which has the expansion
$`E^{(2)}(z_1,z_2)`$ $`=`$ $`(z_1z_2){\displaystyle \underset{tG_S^{(2)}\{id\}}{}}\sqrt{{\displaystyle \frac{(z_1t(z_2))(z_2t(z_1))}{(z_1t(z_1))(z_2t(z_2))}}}+o(k_1^2,k_2^2)`$
$`=`$ $`(z_1z_2)\left(1k_1{\displaystyle \frac{(\eta _1\xi _1)^2(z_1z_2)^2}{(z_1\eta _1)(z_2\eta _1)(z_1\xi _1)(z_2\xi _1)}}\right)`$
$`\times \left(1k_2{\displaystyle \frac{(\eta _2\xi _2)^2(z_1z_2)^2}{(z_1\eta _2)(z_2\eta _2)(z_1\xi _2)(z_2\xi _2)}}\right)`$
$`\times \left(1k(t_1t_2){\displaystyle \frac{(\eta _2\xi _1)^2(z_1z_2)^2}{(z_1\eta _2)(z_2\eta _2)(z_1\xi _1)(z_2\xi _1)}}\right)`$
$`\times \left(1k(t_1t_2){\displaystyle \frac{(\eta _1\xi _2)^2(z_1z_2)^2}{(z_1\eta _1)(z_2\eta _1)(z_1\xi _2)(z_2\xi _2)}}\right)`$
$`\times \left(1k(t_1t_2^1){\displaystyle \frac{(\xi _2\xi _1)^2(z_1z_2)^2}{(z_1\xi _2)(z_2\xi _2)(z_1\xi _1)(z_2\xi _1)}}\right)`$
$`\times \left(1k(t_1t_2^1){\displaystyle \frac{(\eta _2\eta _1)^2(z_1z_2)^2}{(z_1\eta _2)(z_2\eta _2)(z_1\eta _1)(z_2\eta _1)}}\right)+o(k_1^2,k_2^2).`$
The normalization constant of the partition function is being combined with two factors from the integration measure
$`𝒩(1k_1)^2(1k_2)^2{\displaystyle \underset{\beta }{}}\left({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1k_\beta ^n)^d{\displaystyle \underset{n=2}{\overset{\mathrm{}}{}}}(1k_\beta ^n)^2\right)(1k_1)^2(1k_2)^2,`$ (A.28)
where the product over $`\beta `$ extends over the primary classes of the Schottky group, which are equivalence classes under conjugation. Those which contribute to the required order may be represented by $`G_S^{\text{pr}}\{t_1,t_2,t_1t_2,t_1^1t_2,t_1t_2^1\}`$. The expansion reads:
$`𝒩^{(2)}(1k_1)^2(1k_2)^2`$ $`=`$ $`{\displaystyle \underset{\beta G_S^{\text{pr}}}{}}\left({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1k_\beta ^n)^d{\displaystyle \underset{n=2}{\overset{\mathrm{}}{}}}(1k_\beta ^n)^2\right)(1k_1)^2(1k_2)^2+o(k_1^2,k_2^2)`$ (A.29)
$`=`$ $`1+(d2)(k_1+k_2)+((d2)^2+d({\displaystyle \frac{(\eta _1\xi _1)^2(\eta _2\xi _2)^2}{(\eta _1\eta _2)^2(\xi _1\xi _2)^2}}`$
$`+{\displaystyle \frac{(\xi _1\eta _1)^2(\xi _2\eta _2)^2}{(\xi _1\eta _2)^2(\xi _2\eta _1)^2}}))k_1k_2+o(k_1^2,k_2^2).`$
The integrals over the abelian differentials can be explicitely computed elementarily. They are then defined by sums over logarithms
$`{\displaystyle _{z_1}^{z_2}}\omega ^\mu {\displaystyle \underset{tG_S(0,\mu )}{}}\mathrm{ln}\left({\displaystyle \frac{(z_2t(\xi _\mu ))(z_1t(\eta _\mu ))}{(z_1t(\xi _\mu ))(z_2t(\eta _\mu ))}}\right),`$ (A.30)
where the sum is performed over elements of $`G_S^{(2)}`$ which do not have a power of $`t_\mu `$ to the right, e.g. for $`\mu =1`$ just over $`\{id,t_2,t_2^1,t_1t_2,t_1^1t_2,t_1t_2^1,t_1^1t_2^1\}`$. The expansion is:
$`\mathrm{exp}\left({\displaystyle _{z_1}^{z_2}}\omega ^1\right)^{(2)}`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \underset{tG_S^{(2)}(0,1)}{}}\mathrm{ln}\left({\displaystyle \frac{(z_2t(\xi _1))(z_1t(\eta _1))}{(z_1t(\xi _1))(z_2t(\eta _1))}}\right)\right)+o(k_1^2,k_2^2)`$
$`=`$ $`{\displaystyle \frac{(z_2\xi _1)(z_1\eta _1)}{(z_2\eta _1)(z_1\xi _1)}}\left(1k_2{\displaystyle \frac{(\eta _2\xi _2)^2(\xi _1\eta _1)(z_1z_2)}{(\eta _2\xi _1)(\eta _2\eta _1)(z_1\xi _2)(z_2\xi _2)}}\right)`$
$`\times \left(1k_2{\displaystyle \frac{(\eta _2\xi _2)^2(\xi _1\eta _1)(z_1z_2)}{(\xi _2\xi _1)(\xi _2\eta _1)(z_1\eta _2)(z_2\eta _2)}}\right)`$
$`\times \left(1k(t_1t_2){\displaystyle \frac{(\eta _2\xi _1)^2(\xi _1\eta _1)(z_1z_2)}{(\eta _2\xi _1)(\eta _2\eta _1)(z_1\xi _1)(z_2\xi _1)}}\right)`$
$`\times \left(1k(t_1t_2){\displaystyle \frac{(\xi _2\eta _1)^2(\xi _1\eta _1)(z_1z_2)}{(\xi _2\xi _1)(\xi _2\eta _1)(z_1\eta _1)(z_2\eta _1)}}\right)`$
$`\times \left(1k(t_1t_2^1){\displaystyle \frac{(\xi _2\xi _1)^2(\xi _1\eta _1)(z_1z_2)}{(\xi _2\xi _1)(\xi _2\eta _1)(z_1\xi _1)(z_2\xi _1)}}\right)`$
$`\times \left(1k(t_1t_2^1){\displaystyle \frac{(\eta _2\eta _1)^2(\xi _1\eta _1)(z_1z_2)}{(\eta _2\xi _1)(\eta _2\eta _1)(z_1\eta _1)(z_2\eta _1)}}\right)+o(k_1^2,k_2^2),`$
$`\mathrm{exp}\left({\displaystyle _{z_1}^{z_2}}\omega ^2\right)^{(2)}`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \underset{tG_S^{(2)}(0,2)}{}}\mathrm{ln}\left({\displaystyle \frac{(z_2t(\xi _2))(z_1t(\eta _2))}{(z_1t(\xi _2))(z_2t(\eta _2))}}\right)\right)+o(k_1^2,k_2^2)`$
$`=`$ $`(12).`$
The period matrix looks explicitly
$`2\pi i\tau _{\mu \nu }\delta _{\mu \nu }\mathrm{ln}(k_\mu )+{\displaystyle \underset{tG_S(\mu ,\nu )}{}}\mathrm{ln}\left({\displaystyle \frac{(\eta _\mu t(\eta _\nu ))(\xi _\mu t(\xi _\nu ))}{(\eta _\mu t(\xi _\nu ))(\xi _\mu t(\eta _\nu ))}}\right)`$ (A.33)
with summation over elements of $`G_S^{(2)}`$ that neither have a power of $`t_\mu `$ to the left nor one of $`t_\nu `$ to the right, for diagonal elements of $`\tau _{\mu \nu }`$ also $`id`$ is excluded. For instance for $`\mu =1,\nu =1`$ we only have to sum over $`\{t_2,t_2^1\}`$ while $`\mu =1,\nu =2`$ just allows $`\{id,t_2t_1,t_2t_1^1,t_2^1t_1,t_2^1t_1^1\}`$. The expansion is obtained to be
$`2\pi i\tau _{11}^{(2)}`$ $`=`$ $`\mathrm{ln}(k_1)+{\displaystyle \underset{tG_S^{(2)}(1,1)}{}}\mathrm{ln}\left({\displaystyle \frac{(\eta _1t(\eta _1))(\xi _1t(\xi _1))}{(\eta _1t(\xi _1))(\xi _1t(\eta _1))}}\right)+o(k_1^2,k_2^2)`$
$`=`$ $`\mathrm{ln}(k_1)+2{\displaystyle \frac{(\xi _1\eta _1)^2(\xi _2\eta _2)^2}{(\xi _1\xi _2)(\eta _1\eta _2)(\eta _1\xi _2)(\xi _1\eta _2)}}+o(k_1^2,k_2^2),`$
$`2\pi i\tau _{22}^{(2)}`$ $`=`$ $`(12),`$
$`2\pi i\tau _{12}^{(2)}`$ $`=`$ $`2\pi i\tau _{21}^{(2)}={\displaystyle \underset{tG_S^{(2)}(1,2)}{}}\mathrm{ln}\left({\displaystyle \frac{(\eta _1t(\eta _2))(\xi _1t(\xi _2))}{(\eta _1t(\xi _2))(\xi _1t(\eta _2))}}\right)+o(k_1^2,k_2^2)`$ (A.34)
$`=`$ $`\mathrm{ln}\left({\displaystyle \frac{(\eta _1\eta _2)(\xi _1\xi _2)}{(\eta _1\xi _2)(\xi _1\eta _2)}}\right)+2k(t_1t_2){\displaystyle \frac{(\eta _1\xi _1)(\xi _2\eta _2)}{(\eta _1\eta _2)(\xi _1\xi _2)}}`$
$`+2k(t_1t_2^1){\displaystyle \frac{(\eta _1\xi _1)(\xi _2\eta _2)}{(\xi _1\eta _2)(\eta _1\xi _2)}}+o(k_1^2,k_2^2).`$
In all the given expressions we still have got the freedom to fix three of the four fixed points to arbitrary values and will conveniently choose these to be $`0,1,\mathrm{}`$.
## Appendix B Computation of SPT integrals
In this appendix we explicitely compute the integrals over proper time variables we have obtained in chapter 4.4 by performing the low energy limit of the two-loop string diagram extracting the gluon contribution. We use the conventional dimensional regularization scheme in $`d=4+ϵ`$ dimensions. It is usually not necessary to distinguish between IR and UV by using different $`ϵ`$ parameter as long as one can be sure that the overall amplitude is IR finite, which is obvious from (4.4). For computing the integrals we therefore freely use the formulas (D), (D) and (D) by analytic extension outside the region, where the integrals on the left hand sides of these equations converge. We then finally expand the result in powers of $`ϵ`$ and extract the minimal subtraction terms. This procedure is nothing but the usual one.
### B.1 Computation of the leading divergence
We have found the term proportional to $`(ϵp)(pϵ)`$ in (4.4) up to prefactors of the kind
$`{\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{i=1}{\overset{5v}{}}}(dt_i){\displaystyle \frac{\text{Polynomial of (5-}v\text{)th degree in}t_i}{(\text{Polynomial of second degree in}t_i)^{5v+ϵ/2}}}`$ (B.1)
$`\times \mathrm{exp}\left({\displaystyle \frac{\text{Polynomial of third degree in}t_i}{\text{Polynomial of second degree}t_i}}\right),`$
if $`v=0,1`$ is the number of four gluon vertices we consider. The term proportional to $`ϵ^2p^2`$ can be computed in completely the same manner and for the sake of brevity we concentrate on the one given above. If we are only interested in the leading $`1/ϵ^2`$ divergence, we can meanwhile omitt the exponential from the integrand of (B.2) and only after performing all but the last integration use it as an IR regulator. We have already done the substitution that lead to (4.4) and justified to replace the second term deriving from the exponent as in (72) by $`1+o(ϵ)`$. Then reverting the first substitution we get (4.4) back up to the simplified exponent:
$`{\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{i=1}{\overset{5v}{}}}(dt_i){\displaystyle \frac{\text{Polynomial of (5-}v\text{)th degree in}t_i}{(\text{Polynom of second degree in}t_i)^{5v+ϵ/2}}}\mathrm{exp}\left({\displaystyle \underset{i=1}{\overset{5}{}}}t_i\right).`$ (B.2)
We actually need this integral only in an infinitesimal vicinity of the origin, where we are allowed to expand the integrand because of coninuity. If we take care of not getting any divergencies from large proper times we can finally extend the integration of the modified integrand to the full interval from zero to infinity again. We had to check that all the conditions mentioned are satisfied holding in all the cases we shall be regarding in the following. As the integrals obtained by these means can then be computed, we get a result for the $`1/ϵ^2`$ divergent term that can be compared to results we shall get later on by more sophisticated methods that are less transparent on the other hand. We write the integrals thus obtained in the following manner
$`J_5(ϵ)={\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{i=1}{\overset{5}{}}}(dt_i){\displaystyle \frac{P_5(t_1,\mathrm{},t_5,ϵ)}{(P_2(t_1,t_2,t_3,t_4,t_5))^{5+ϵ/2}}}.`$ (B.3)
In this subsection we use the notation $`J_{5v}^{(x)}|_i\mathrm{}(ϵ)`$ for these integrals without the exponential factor in the integrand, using the lower index for the number $`5v`$ of integrations and the additional ones $`i\mathrm{}`$ for the possibly missing proper time variables, as well as an additional letter $`(x)`$ for the according diagram of figure 10, sticking strictly to the notation introduced in (18) and (80). In the end we shall reintroduce the exponential as cut-off without change of notation for the integral. The polynomials are carrying indices signaling their degree.
The diagram (a) has the contribution given in (4.4) which has a polynomial
$`P_2^{(a)}(t_1,\mathrm{},t_5)`$ $`=`$ $`\alpha ^2\left(\mathrm{ln}(k_1)\mathrm{ln}(k_2)\mathrm{ln}^2(\eta )\right)`$
$`=`$ $`(t_1+t_2)(t_3+t_4+t_5)+(t_3+t_4)t_5`$
$``$ $`a(t_1+t_2)+b`$
in the denominator. For the polynomial $`P_5^{(a)}`$ in the numarator we use
$`P_5^{(a)}(t_1,\mathrm{},t_5,ϵ)`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{3}{}}}c_it_1^i={\displaystyle \underset{i=0}{\overset{3}{}}}{\displaystyle \underset{j=0}{\overset{3i}{}}}c_{ij}t_1^it_2^j={\displaystyle \underset{i=0}{\overset{3}{}}}{\displaystyle \underset{j=0}{\overset{3i}{}}}{\displaystyle \underset{k=0}{\overset{3,5ij}{}}}c_{ijk}t_1^it_2^jt_3^k`$
$`=`$ $`{\displaystyle \underset{i=0}{\overset{3}{}}}{\displaystyle \underset{j=0}{\overset{3i}{}}}{\displaystyle \underset{k=0}{\overset{3,5ij}{}}}{\displaystyle \underset{l=0,2ijk}{\overset{3k}{}}}c_{ijkl}t_1^it_2^jt_3^kt_4^l.`$
The summation limits are to be understood in the sense that $`0,2ijk`$ stands for the larger value of $`0`$ or $`2ijk`$, alternatively, if upper limits are double valued, they stand for the lower one. We now use for the integrations over $`t_1`$ and $`t_2`$ the formula (D), for that over $`t_3`$ we have (D) and finally for $`t_4`$ it is (D):
$`J_5^{(a)}(ϵ)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{i=1}{\overset{5}{}}}(dt_i){\displaystyle \frac{_ic_it_1^i}{(a(t_1+t_2)+b)^{5+ϵ/2}}}`$
$`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{i=2}{\overset{5}{}}}(dt_i){\displaystyle \underset{ij}{}}c_{ij}t_2^j(at_2+b)^{i4ϵ/2}a^{i1}B(i+1,4i+ϵ/2)`$
$`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑t_3𝑑t_4𝑑t_5{\displaystyle \underset{ijk}{}}c_{ijk}t_3^kb^{i+j3ϵ/2}a^{ij2}`$
$`\times B(i+1,4i+ϵ/2)B(j+1,3ij+ϵ/2)`$
$`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑t_4𝑑t_5{\displaystyle \underset{ijkl}{}}c_{ijkl}t_4^{i+j+k+l2ϵ/2}t_5^{i+j3ϵ/2}(t_4+t_5)^{ij2}`$
$`\times B(i+1,4i+ϵ/2)B(j+1,3ij+ϵ/2)B(k+1,4k+ϵ/2)`$
$`\times {}_{2}{}^{}F_{1}^{}(i+j+2,k+1;5+ϵ/2;t_5/(t_4+t_5))`$
$`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑t_5{\displaystyle \underset{ijkl}{}}c_{ijkl}t_5^{i+j+k+l6ϵ}B(i+1,4i+ϵ/2)`$
$`\times B(j+1,3ij+ϵ/2)B(k+1,4k+ϵ/2)`$
$`\times B(3kl+ϵ/2,i+j+k+l1ϵ/2)`$
$`\times {}_{3}{}^{}F_{2}^{}(i+j+2,k+1,3kl+ϵ/2;5+ϵ/2,i+j+2;1).`$
For doing the $`t_4`$ integration we had to substitute
$`t_4`$ $`=`$ $`{\displaystyle \frac{(1x)t_5}{x}},dt_4={\displaystyle \frac{t_5dx}{x^2}},[0,\mathrm{}][1,0].`$ (B.7)
The cut-off integral is
$`{\displaystyle _0^{\mathrm{}}}𝑑t_5t_5^{1ϵ}e^{t_5}`$ $`=`$ $`\mathrm{\Gamma }(ϵ)={\displaystyle \frac{1}{ϵ}}\gamma +o(ϵ).`$ (B.8)
To compute the integrals derived from (a) by pinching is even easier because of their polynomials being of lower degree. For $`t_10`$ we get a contribution to diagram (c), having the polynomial
$`P_2^{(c)}(t_2,t_3,t_4,t_5)`$ $`=`$ $`t_2(t_3+t_4+t_5)+(t_3+t_4)t_5`$
$``$ $`at_2+b`$
in the denominator and
$`P_4^{(c)}(t_2,t_3,t_4,t_5,ϵ)`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{2}{}}}c_it_2^i={\displaystyle \underset{i=0}{\overset{2}{}}}{\displaystyle \underset{j=0}{\overset{1}{}}}c_{ij}t_2^it_3^j={\displaystyle \underset{i=0}{\overset{2}{}}}{\displaystyle \underset{j=0}{\overset{1}{}}}{\displaystyle \underset{k=0,2ij}{\overset{2,4ij}{}}}c_{ijk}t_2^it_3^jt_4^k`$ (B.10)
in the numerator, where all coefficients whose indices do not satisfy $`j+k2`$ vanish. The integrals are solved by (D),(D) and (D) using again the substitution (B.7).
$`J_4^{(c)}|_1(ϵ)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{i=2}{\overset{5}{}}}(dt_i){\displaystyle \frac{_ic_it_2^i}{(at_2+b)^{4+ϵ/2}}}`$
$`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑t_3𝑑t_4𝑑t_5{\displaystyle \underset{ij}{}}c_{ij}t_3^jb^{i3ϵ/2}a^{i1}B(i+1,3i+ϵ/2)`$
$`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑t_4𝑑t_5{\displaystyle \underset{ijk}{}}c_{ijk}t_4^{i+j+k2ϵ/2}t_5^{i3ϵ/2}(t_4+t_5)^{i1}`$
$`\times B(i+1,3i+ϵ/2)B(j+1,3j+ϵ/2)`$
$`\times {}_{2}{}^{}F_{1}^{}(i+1,j+1;4+ϵ/2;t_5/(t_4+t_5))`$
$`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑t_5{\displaystyle \underset{ijk}{}}c_{ijk}t_5^{i+j+k5ϵ}B(i+1,3i+ϵ/2)B(j+1,3j+ϵ/2)`$
$`\times B(i+j+k1ϵ/2,2jk+ϵ/2)`$
$`\times {}_{3}{}^{}F_{2}^{}(i+1,j+1,2jk+ϵ/2;4+ϵ/2,i+1;1).`$
The $`t_5`$ integration again has to be treated by the cut-off prescription as in (B.8). The integral deduced from $`t_20`$ has identical structure as this one with $`t_10`$, while $`t_30`$ can be handled even more simply. This then is equivalent to $`t_40`$ again. We get in both cases
$`P_2^{(c)}(t_1,t_2,t_4,t_5)`$ $`=`$ $`(t_1+t_2)(t_4+t_5)+t_4t_5`$
$``$ $`a(t_1+t_2)+b`$
in the denominator and
$`P_4^{(c)}(t_1,t_2,t_4,t_5,ϵ)`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{1}{}}}c_it_1^i={\displaystyle \underset{i=0}{\overset{1}{}}}{\displaystyle \underset{j=0}{\overset{2i}{}}}c_{ij}t_1^it_2^j={\displaystyle \underset{i=0}{\overset{1}{}}}{\displaystyle \underset{j=0}{\overset{2i}{}}}{\displaystyle \underset{k=2ij}{\overset{2,4ij}{}}}c_{ijk}t_1^it_2^jt_4^k`$ (B.13)
in the numerator. The first integrations can all be done applying (D). Afterwards the cut-off is employed and the result reads
$`J_4^{(c)}|_3(ϵ)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑t_1𝑑t_2𝑑t_4𝑑t_5{\displaystyle \frac{_ic_it_1^i}{(a(t_1+t_2)+b)^{4+ϵ/2}}}`$
$`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑t_2𝑑t_4𝑑t_5{\displaystyle \underset{ij}{}}c_{ij}t_2^ja^{i1}(at_2+b)^{i3ϵ/2}B(i+1,3i+ϵ/2)`$
$`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑t_4𝑑t_5{\displaystyle \underset{ijk}{}}c_{ijk}t_4^{i+j+k2ϵ/2}t_5^{i+j2ϵ/2}(t_4+t_5)^{ij2}`$
$`\times B(i+1,3i+ϵ/2)B(j+1,2ij+ϵ/2)`$
$`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑t_5{\displaystyle \underset{ijk}{}}c_{ijk}t_5^{i+j+k5ϵ}B(i+1,3i+ϵ/2)`$
$`\times B(j+1,2ij+ϵ/2)B(i+j+k1ϵ/2,3k+ϵ/2).`$
We shall not present the computation of the contributions to diagram (e) as they are vanishing, which makes it a little boring. This is in fact a consequence of numerous cancellations, whereas the vanishing of contributions to diagram (d) is due to the overall prefactor $`\alpha ^{}`$ as already mentioned.
We now proceed to diagram (b) of figure 10, where we get
$`P_2^{(a)}(t_1,\mathrm{},t_5)`$ $`=`$ $`\alpha ^2\left(\mathrm{ln}(k_1)\mathrm{ln}(k_2)\mathrm{ln}^2(\eta )\right)`$
$`=`$ $`(t_1+t_5)(t_2+t_3+t_4)+t_1t_5a(t_2+t_3+t_4)+b`$
in the denominator which suggests to integrate by the order $`t_2t_3t_4t_1`$. The polynomial in the numerator then is
$`P_5^{(b)}(t_1,\mathrm{},t_5,ϵ)`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{3}{}}}c_it_2^i={\displaystyle \underset{i=0}{\overset{3}{}}}{\displaystyle \underset{j=0}{\overset{3i}{}}}c_{ij}t_2^it_3^j`$
$`=`$ $`{\displaystyle \underset{i=0}{\overset{3}{}}}{\displaystyle \underset{j=0}{\overset{3i}{}}}{\displaystyle \underset{k=0}{\overset{3ij}{}}}c_{ijk}t_2^it_3^jt_4^k={\displaystyle \underset{i=0}{\overset{3}{}}}{\displaystyle \underset{j=0}{\overset{3i}{}}}{\displaystyle \underset{k=0}{\overset{3ij}{}}}{\displaystyle \underset{l=0,2ijk}{\overset{3,5ijk}{}}}c_{ijkl}t_2^it_3^jt_4^kt_1^l.`$
The integrations can be performed be using exclusively (D) and reveal
$`\overline{J}_5^{(b)}(ϵ)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑t_1𝑑t_2𝑑t_3𝑑t_4𝑑t_5{\displaystyle \frac{_ic_it_2^i}{(a(t_2+t_3+t_4)+b)^{5+ϵ/2}}}`$
$`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑t_1𝑑t_3𝑑t_4𝑑t_5{\displaystyle \underset{ij}{}}c_{ij}t_3^ja^{i1}(a(t_3+t_4)+b)^{i4ϵ/2}`$
$`\times B(i+1,4i+ϵ/2)`$
$`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑t_1𝑑t_4𝑑t_5{\displaystyle \underset{ijk}{}}c_{ijk}t_4^ka^{ij2}(at_4+b)^{ijk3}`$
$`\times B(i+1,4i+ϵ/2)B(j+1,3ij+ϵ/2)`$
$`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑t_1𝑑t_5{\displaystyle \underset{ijkl}{}}c_{ijkl}t_1^{i+j+k+l2ϵ/2}t_5^{i+j+k2ϵ/2}(t_1+t_5)^{ijk3}`$
$`\times B(i+1,4i+ϵ/2)B(j+1,3ij+ϵ/2)`$
$`\times B(k+1,2ijk+ϵ/2)`$
$`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑t_5{\displaystyle \underset{ijkl}{}}c_{ijkl}t_5^{i+j+k+l6ϵ}B(i+1,4i+ϵ/2)`$
$`\times B(j+1,3ij+ϵ/2)B(k+1,2ijk+ϵ/2)`$
$`\times B(4l+ϵ/2,i+j+k+l1ϵ/2),`$
where again the final $`t_5`$ integration has to be done by (B.8). All contributions to diagrams including four gluon vertices can be reduced to cases we have already covered when discussing those derived from (a). The necessary replacements are quite obvious.
### B.2 Exact computation of SPT integrals
For the integrals from (4.4) we can define an algorithm that even allows an exact computation. This will reveal further information to us, which hints, how the rest of the integrals corresponding to all the other diagrams can be solved by explicit calculation, using only the formulas displayed in appendix D. We have to thank M. Peter from Heidelberg for his advice and active help on this subject. Our notation is adopted from the previous chapter, only replacing $`J`$ by $`I`$.
The most important observation is that the proper time integrals can formally be written as ordinary momentum integrals of a scalar field theory, which are known from solving the integrals emerging from Feynman rules for such a theory. To do this we have to treat each term of the polynomials in the numerators separately and exploit the identity
$`I(q^2,n_1,\mathrm{},n_5)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d^Dpd^Dk}{(2\pi )^{2D}}}{\displaystyle \frac{1}{\left(p^2\right)^{n_1}\left((pq)^2\right)^{n_2}\left(k^2\right)^{n_3}\left((kq)^2\right)^{n_4}\left((pk)^2\right)^{n_5}}}`$
$`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{i=1}{\overset{5}{}}}\left({\displaystyle \frac{dt_it_i^{n_i}}{\mathrm{\Gamma }(n_i)}}\right)\left((t_1+t_2)(t_3+t_4+t_5)+(t_3+t_4)t_5\right)^{D/2}`$
$`\times \mathrm{exp}\left(q^2{\displaystyle \frac{t_1t_2(t_3+t_4+t_5)+t_1t_3(t_4+t_5)+t_2t_4(t_3+t_5)+t_3t_4t_5}{(t_1+t_2)(t_3+t_4+t_5)+(t_3+t_4)t_5}}\right).`$
The SPT integral of the second line with $`D=10+ϵ`$ exactly resembles the integrals we are trying to solve when dealing with Yang-Mills theory. We use the variable $`t_i`$ for the propagator with exponent $`n_i`$ as in (1) and then do the Gaussian integrations by (D.8). Such momentum integrals one knows how to handle. The recursion relation
$`0`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d^Dpd^Dk}{(2\pi )^{2D}}}{\displaystyle \frac{}{p_\mu }}\left({\displaystyle \frac{p_\mu k_\mu }{\left(p^2\right)^{n_1}\left((pq)^2\right)^{n_2}\left(k^2\right)^{n_3}\left((kq)^2\right)^{n_4}\left((pk)^2\right)^{n_5}}}\right)`$
$`=`$ $`(Dn_1n_22n_5)I(q^2,n_1,\mathrm{},n_5)`$
$`n_1\left(I(q^2,n_1+1,n_2,n_3,n_4,n_51)I(q^2,n_1+1,n_2,n_31,n_4,n_5)\right)`$
$`n_2\left(I(q^2,n_1,n_2+1,n_3,n_4,n_51)I(q^2,n_1,n_2+1,n_3,n_41,n_5)\right)`$
relates any given integral successively to a number of integrals that has $`n_5=0`$, $`n_3=0`$ or $`n_4=0`$. In the first case both integrals factorize as one propagator vanishes, in the latter case we get two one-loop integrals that can be easily computed, which will be done explicitly in the following. This recursion thus allows to reduce all terms to a small number of simple types, but it creates a couple of thousand terms of such. Therefore we have used the algebraic computer programs MAPLE and FORM to do the task of performing the recursion and the expansion of its results in $`1/ϵ`$. The results are summarized in several tables chapter 4.4.
We next demonstrate how the pinching contributions of (a) can be calculated and shall find that we have to explicitely compute only the two intgrals just cited. To find the most appropriate order of substitutions for the proper time variables it is essential to understand the translation into momentum integrals. Take the case $`t_10`$ whose contribution is
$`I_4^{(c)}|_1(ϵ)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{i=2}{\overset{5}{}}}\left(dt_i\right){\displaystyle \frac{P_4^{(c)}(t_2,\mathrm{},t_5,ϵ)}{\left(t_2(t_3+t_4+t_5)+(t_3+t_4)t_5\right)^{4+ϵ/2}}}`$
$`\times \mathrm{exp}\left({\displaystyle \frac{t_2t_4(t_3+t_5)+t_3t_4t_5}{t_2(t_3+t_4+t_5)+(t_3+t_4)t_5}}\right)`$
and use
$`t_2=xt,t_5=(1x)t,t_3tt_3,t_4tt_4,`$ (B.21)
rescaling $`t_2`$ and $`t_5`$ to a unit square. On the contrary, a simultaneous scaling of all the integration variables would have not been successfull. The integration over $`t`$ is trivial and gives
$`I_4^{(c)}|_1(ϵ)=\mathrm{\Gamma }(ϵ){\displaystyle _0^1}𝑑x{\displaystyle _0^{\mathrm{}}}𝑑t_3𝑑t_4\left({\displaystyle \frac{t_3t_4+x(1x)t_4}{t_3+t_4+x(1x)}}\right)^ϵ{\displaystyle \frac{\overline{P}_4^{(c)}(t_3,t_4,x,ϵ)}{\left(t_3+t_4+x(1x)\right)^{4+ϵ/2}}}.`$ (B.22)
We next define
$`\overline{P}_4^{(c)}(t_3,t_4,x,ϵ){\displaystyle \underset{i}{}}c_i(t_3,x,ϵ)t_4^i{\displaystyle \underset{ij}{}}c_{ij}(x,ϵ)t_4^it_3^j{\displaystyle \underset{ijk}{}}c_{ijk}(ϵ)t_4^it_3^jx^k`$ (B.23)
and integrate over $`t_3`$ and $`t_4`$ by using (D):
$`I_4^{(c)}|_1(ϵ)`$ $`=`$ $`\mathrm{\Gamma }(ϵ){\displaystyle _0^1}𝑑x{\displaystyle \underset{ijk}{}}c_{ijk}(ϵ)x^{i+j+k2+ϵ/2}(1x)^{i+j2+ϵ/2}`$
$`\times B(i+1+ϵ,3i+ϵ/2)B(j+1,2ijϵ/2)`$
$`=`$ $`\mathrm{\Gamma }(ϵ){\displaystyle \underset{ijk}{}}c_{ijk}(ϵ)B(i+1+ϵ,3i+ϵ/2)B(j+1,2ijϵ/2)`$
$`\times B(i+j1+ϵ/2,i+j+k1+ϵ/2).`$
The cases $`t_2,t_3,t_40`$ can be treated in precisely the same manner by a simple permutation of indices, which can be read of from (4.4) by inspecting (B.2). We only remain with computing the contribution coming from $`t_50`$ to the diagram (e):
$`I_4^{(e)}|_5(ϵ)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{i=1}{\overset{4}{}}}\left(dt_i\right){\displaystyle \frac{P_4^{(e)}(t_1,\mathrm{},t_4,ϵ)}{\left((t_1+t_2)(t_3+t_4)\right)^{4+ϵ/2}}}`$
$`\times \mathrm{exp}\left({\displaystyle \frac{t_1(t_2t_3+t_2t_4+t_3t_4)+t_2t_3t_4}{(t_1+t_2)(t_3+t_4)}}\right)`$
$`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{i=1}{\overset{4}{}}}\left(dt_i\right){\displaystyle \frac{P_4^{(e)}(t_1,\mathrm{},t_4,ϵ)}{\left((t_1+t_2)(t_3+t_4)\right)^{4+ϵ/2}}}\mathrm{exp}\left({\displaystyle \frac{t_1t_2}{t_1+t_2}}+{\displaystyle \frac{t_3t_4}{t_3+t_4}}\right).`$
We rescale $`t_1`$ and $`t_2`$
$`t_1=xt,t_2=(1x)t,`$ (B.26)
and as the polynomial $`P_4^{(e)}`$ in the numerator is always quadratic in these two variables, the integration over $`t`$ can be split off from the rest:
$`I_4^{(e)}|_5(ϵ)`$ $`=`$ $`\mathrm{\Gamma }\left(ϵ/2\right){\displaystyle _0^1}𝑑x{\displaystyle _0^{\mathrm{}}}𝑑t_3𝑑t_4\left(x(1x)\right)^{ϵ/2}`$
$`\times {\displaystyle \frac{\overline{P}_4^{(e)}(t_3,t_4,x,ϵ)}{(t_3+t_4)^{4+ϵ/2}}}\mathrm{exp}\left({\displaystyle \frac{t_3t_4}{t_3+t_4}}\right).`$
After another substitution
$`t_3=y\rho ,t_4=(1y)\rho ,`$ (B.28)
and defining
$`\overline{P}_4^{(e)}(t_3,t_4,x,ϵ){\displaystyle \underset{i}{}}c_i(y,ϵ)x^i{\displaystyle \underset{ij}{}}c_{ij}(ϵ)x^iy^j`$ (B.29)
we finally obtain
$`I_4^{(e)}|_5(ϵ)`$ $`=`$ $`\mathrm{\Gamma }\left(ϵ/2\right)^2{\displaystyle _0^1}𝑑x𝑑y{\displaystyle \underset{ij}{}}c_{ij}(ϵ)\left(x(1x)\right)^{ϵ/2}\left(y(1y)\right)^{ϵ/2}x^iy^j`$
$`=`$ $`\mathrm{\Gamma }\left(ϵ/2\right)^2{\displaystyle \underset{ij}{}}c_{ij}(ϵ)B(i+1+ϵ/2,1+ϵ/2)B(j+1+ϵ/2,1+ϵ/2).`$
We have such got all the pinching contributions derived from (a).
The only case we have not covered yet is the diagram (b) itself. All its pinching limits can again be reduced to the former $`t_10`$ integral type. From (4.4) we get, following the notation of figure 12,
$`I_5^{(b)}(ϵ)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{i=1}{\overset{5}{}}}(dt_i){\displaystyle \frac{P_5^{(b)}(t_1,\mathrm{},t_5,ϵ)}{\left((t_1+t_5)(t_2+t_3+t_4)+t_1t_5\right)^{5+ϵ/2}}}`$
$`\times \mathrm{exp}\left({\displaystyle \frac{t_3((t_1+t_5)(t_2+t_4)+t_1t_5)}{(t_1+t_5)(t_2+t_3+t_4)+t_1t_5}}\right).`$
In fact, by explicit inspection the polynomial $`P_5^{(b)}`$ turns out to depend only on the sum of $`t_2`$ and $`t_4`$. This corresponds to the fact that from the point of view of the field theory, one of the two variables is superfluous, as both propagators carry the same momentum, are thus identical. We then use their sum as a new variable
$`t_2=xt,t_4=(1x)t,t_1tt_1,t_3tt_3,t_5tt_5.`$ (B.32)
The integrand does not depend on $`x`$ and the integration over $`t`$ is trivial, so that we get
$`I_5^{(b)}(ϵ)`$ $`=`$ $`\mathrm{\Gamma }(ϵ){\displaystyle _0^{\mathrm{}}}𝑑t_1𝑑t_3𝑑t_5{\displaystyle \frac{\overline{P}_5^{(b)}(t_1,t_3,t_5,ϵ)}{\left((t_1+t_5)(t_3+1)+t_1t_5\right)^{5+ϵ/2}}}`$
$`\times \left({\displaystyle \frac{t_3(t_1+t_5+t_1t_5)}{(t_1+t_5)(t_3+1)+t_1t_5}}\right)^ϵ.`$
We then proceed as usual
$`\overline{P}_5^{(b)}(t_1,t_3,t_5,ϵ){\displaystyle \underset{i}{}}c_i(t_1,t_5,ϵ)t_3^i{\displaystyle \underset{ij}{}}c_{ij}(t_5,ϵ)t_3^it_1^j{\displaystyle \underset{ijk}{}}c_{ijk}(ϵ)t_3^it_1^jt_5^k,`$ (B.34)
integrate over $`t_3`$ by (D), next over $`t_1`$ using (D) and finally over $`t_5`$ substituting
$`t_5={\displaystyle \frac{y}{1y}},dt_5={\displaystyle \frac{dy}{(1y)^2}},[0,\mathrm{}][0,1]`$ (B.35)
and by (D). The result of all this is
$`I_5^{(b)}(ϵ)`$ $`=`$ $`\mathrm{\Gamma }(ϵ){\displaystyle \underset{ijk}{}}c_{ijk}(ϵ)B(i+1+ϵ,4i+ϵ/2)B(j+1,4j+ϵ/2)`$
$`\times B(j+k3ϵ/2,4k+ϵ/2)`$
$`\times {}_{3}{}^{}F_{2}^{}(i+1+ϵ,j+1,j+k3ϵ/2;5+ϵ/2,j+1;1).`$
We have been able to exactly solve all the SPT integrals.
## Appendix C Two-loop diagrams with external states attached to the small loop
This is an extension to chapter 4.4, where we treat diagrams with one or both external states sitting at the interior, “small” loop. The procedure is very much equivalent to the usual and differs mainly in the choice of local coordinates one uses in the vicinity of the external legs of the diagrams, as well as in the parametrization of the sewing variables that is derived fom this. We shall therefore be very brief and present all possible combinations of contributions in terms of the world sheet diagrams we have introduced earlier, then display in table 4 the appropriate sewing parameters and finally simply cite the results one obtains from solving the integrals in table 5. For the local coordinates of those external states that are supposed to be attached to the interior loop we now take
$`V_i^{}(0)=\left|{\displaystyle \frac{(z_i1)(z_i\eta )}{1\eta }}\right|,`$ (C.1)
while for those which are sitting at the large loop we have $`V_i^{}(0)=z_i`$ as before. The parametrization we define in table 4 is given by the methods of , which necessarily leads to the correct Green’s function.
## Appendix D Useful formulas and functions
Formulas used to compute SPT integrals :
$`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{x^{\mu 1}}{(\beta x+1)^\nu }}𝑑x`$ $`=`$ $`\beta ^\mu B(\mu ,\nu \mu ),`$
$`\left[\mathrm{}(\nu )>\mathrm{}(\mu )>0,|\mathrm{arg}(\beta )|<\pi \right],`$
$`{\displaystyle _0^{\mathrm{}}}x^{\nu 1}(\beta +x)^\mu (\gamma +x)^\rho 𝑑x`$ $`=`$ $`\beta ^\mu \gamma ^{\nu \rho }B(\nu ,\mu \nu +\rho )`$
$`\times {}_{2}{}^{}F_{1}^{}(\mu ,\nu ;\mu +\rho ;1\gamma /\beta ),`$
$`[\mathrm{}(\nu )>0,\mathrm{}(\mu )>\mathrm{}(\nu \rho ),`$
$`|\mathrm{arg}(\beta )|<\pi ,|\mathrm{arg}(\gamma )|<\pi ],`$
$`{\displaystyle _0^1}x^{\rho 1}(1x)^{\sigma 1}{}_{2}{}^{}F_{1}^{}(\alpha ,\beta ;\gamma ;x)𝑑x`$ $`=`$ $`B(\rho ,\sigma ){}_{3}{}^{}F_{2}^{}(\alpha ,\beta ,\rho ;\gamma ,\rho +\sigma ;1),`$
$`\left[\mathrm{}(\rho )>0,\mathrm{}(\sigma )>0,\mathrm{}(\gamma +\sigma \alpha \beta )>0\right].`$
The Gamma function:
$`\mathrm{\Gamma }(z)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑te^tt^{z1},\left[\mathrm{}(z)>0\right].`$ (D.4)
The Euler Beta function:
$`B(\mu ,\nu )`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(\mu )\mathrm{\Gamma }(\nu )}{\mathrm{\Gamma }(\mu +\nu )}}`$
$`=`$ $`{\displaystyle _0^1}x^{\nu 1}(1x)^{\mu 1}𝑑x,\left[\mathrm{}(\nu ),\mathrm{}(\mu )>0\right].`$
The generalized hypergeometric series:
$`{}_{p}{}^{}F_{q}^{}(\alpha _1,\mathrm{},\alpha _p;\beta _1,\mathrm{},\beta _q;z)={\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(\alpha _1)_k\mathrm{}(\alpha _p)_k}{(\beta _1)_k\mathrm{}(\beta _q)_k}}{\displaystyle \frac{z^k}{k!}},`$ (D.6)
where $`(\alpha )_k\alpha (\alpha +1)\mathrm{}(\alpha +k1)`$ and $`(\alpha )_01`$. For $`{}_{2}{}^{}F_{1}^{}`$ there exists the integral representation :
$`{}_{2}{}^{}F_{1}^{}(\alpha ,\beta ;\gamma ;z)`$ $`=`$ $`{\displaystyle \frac{1}{B(\beta ,\gamma \beta )}}{\displaystyle _0^1}t^{\beta 1}(1t)^{\gamma \beta 1}(1tz)^\alpha 𝑑t,`$
$`\left[\mathrm{}(\gamma )>\mathrm{}(\beta )>0\right].`$
For the integration of Gaussian integrals after applying the Schwinger trick to Feynman propagators
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d^dpd^dk}{(2\pi )^{2d}}}\mathrm{exp}\left(\alpha p^2\beta k^2\delta q^2+2\gamma pk+2xpq+2ykq\right)=`$ (D.8)
$`\left(\alpha \beta \gamma ^2\right)^{d/2}\mathrm{exp}\left(q^2{\displaystyle \frac{\alpha \delta \beta \delta \gamma ^22xy\gamma x^2\beta y^2\alpha }{\alpha \beta \gamma ^2}}\right).`$
The Riemann Zeta function:
$`\zeta (z)={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}n^z,\left[\mathrm{}(z)>1\right].`$ (D.9) |
warning/0003/nucl-th0003034.html | ar5iv | text | # Double giant dipole resonances in time-dependent density-matrix theory
## Abstract
The strength functions of the double giant dipole resonances (DGDR) in <sup>16</sup>O and <sup>40</sup>Ca are calculated with the use of an extended version of the time-dependent Hartree-Fock theory known as the time-dependent density-matrix theory. The calculations are done in a self-consistent manner, in which the same Skyrme force as that used for a mean-field potential is used as an effective interaction for a two-body correlation function. It is found that the DGDR in <sup>16</sup>O has a large width due to the Landau damping, although the centroid energy of the strength distribution is close to twice the energy of the giant dipole resonance (GDR) calculated in RPA. The DGDR in <sup>40</sup>Ca is found more harmonic than that in <sup>16</sup>O: the strength function of the DGDR in <sup>40</sup>Ca is similar to what is predicted from the strength function of the GDR in RPA.
PACS numbers: 21.60.Jz, 24.10.Cn, 24.30.Cz
Keywords: giant dipole resonance, double phonon state, extended time-dependent Hartree-Fock theory
The double phonon states of giant resonances have become the subject of a number of recent experimental and theoretical investigations . Microscopic calculations of the strength functions of double giant dipole resonances (DGDR) have also been done based on the shell model and quasiparticle-phonon models . However, few microscopic studies have been reported, in which a single-particle basis and a residual interaction are treated in such a self-consistent manner as used in random-phase-approximation (RPA) calculations for giant resonances . We have recently proposed a self-consistent approach based on an extended version of the time-dependent Hartree-Fock theory (TDHF) known as the time-dependent density-matrix theory (TDDM) , in which the same Skyrme force as that used for the calculation of a mean-field potential is used as a residual interaction for a two-body correlation function. We applied the model to the double giant quadrupole resonances (DGQR) in <sup>16</sup>O and <sup>40</sup>Ca and showed that the DGQR’s in these nuclei have strong harmonic properties. The aim of this paper is to report the results of the application of the TDDM approach to the DGDR’s in <sup>16</sup>O and <sup>40</sup>Ca.
The formulation of TDDM is based on the truncation of the hierarchy of reduced density matrices, in which genuine correlated parts in a three-body density matrix and higher reduced density matrices are neglected . The TDDM equations thus determine the time evolution of a one-body density matrix $`\rho `$ and a two-body correlation function $`C_2`$ defined by $`C_2=\rho _2A[\rho \rho ]`$, where $`A[\rho \rho ]`$ is the antisymmetrized product of the one-body density matrices and $`\rho _2`$ is a two-body density matrix. In TDDM, further truncation is made by expanding $`\rho `$ and $`C_2`$ with a finite number of single-particle states $`\{\psi _\alpha \}`$ as
$`\rho (11^{},t)`$ $`=`$ $`{\displaystyle \underset{\alpha \alpha ^{}}{}}n_{\alpha \alpha ^{}}(t)\psi _\alpha (1,t)\psi _\alpha ^{}^{}(1^{},t),`$ (1)
$`C_2(121^{}2^{},t)={\displaystyle \underset{\alpha \beta \alpha ^{}\beta ^{}}{}}C_{\alpha \beta \alpha ^{}\beta ^{}}(t)\psi _\alpha (1,t)\psi _\beta (2,t)\psi _\alpha ^{}^{}(1^{},t)\psi _\beta ^{}^{}(2^{},t),`$ (2)
where the numbers denote space, spin and isospin coordinates. The time evolution of $`\rho `$ and $`C_2`$ is determined by the following three coupled equations :
$`i\mathrm{}{\displaystyle \frac{}{t}}\psi _\alpha (1,t)=h(1,t)\psi _\alpha (1,t),`$ (3)
$`i\mathrm{}\dot{n}_{\alpha \alpha ^{}}={\displaystyle \underset{\beta \gamma \delta }{}}[\alpha \beta |v|\gamma \delta C_{\gamma \delta \alpha ^{}\beta }C_{\alpha \beta \gamma \delta }\gamma \delta |v|\alpha ^{}\beta ],`$ (4)
$`i\mathrm{}\dot{C}_{\alpha \beta \alpha ^{}\beta ^{}}=B_{\alpha \beta \alpha ^{}\beta ^{}}+P_{\alpha \beta \alpha ^{}\beta ^{}}+H_{\alpha \beta \alpha ^{}\beta ^{}},`$ (5)
where $`h(1,t)`$ is the mean-field hamiltonian and $`v`$ the residual interaction. The term $`B_{\alpha \beta \alpha ^{}\beta ^{}}`$ on the right-hand side of Eq.(5) represents the Born terms (the first-order terms of $`v`$). The terms $`P_{\alpha \beta \alpha ^{}\beta ^{}}`$ and $`H_{\alpha \beta \alpha ^{}\beta ^{}}`$ in Eq.(5) contain $`C_{\alpha \beta \alpha ^{}\beta ^{}}`$ and represent higher-order particle-particle (and hole-hole) and particle-hole type correlations, respectively. Thus full two-body correlations including those induced by the Pauli exclusion principle are taken into account in the equation of motion for $`C_{\alpha \beta \alpha ^{}\beta ^{}}`$. The explicit expressions for $`B_{\alpha \beta \alpha ^{}\beta ^{}}`$, $`P_{\alpha \beta \alpha ^{}\beta ^{}}`$ and $`H_{\alpha \beta \alpha ^{}\beta ^{}}`$ are given in Ref.. The small amplitude limit of TDDM was investigated and it was shown that TDDM can be reduced to the second RPA - in such a limit. The number of two-body matrices treated in TDDM grows very rapidly with increasing mass number, restricting the application of TDDM to light nuclei for the present. To solve the coupled equations, we use the Skyrme interaction of the form
$`v(𝐫𝐫^{})`$ $`=`$ $`t_0(1+x_0P_\sigma )\delta ^3(𝐫𝐫^{})+{\displaystyle \frac{1}{2}}t_1\{k^2\delta ^3(𝐫𝐫^{})+\delta ^3(𝐫𝐫^{})k^2\}`$ (6)
$`+`$ $`t_2𝐤^{}\delta ^3(𝐫𝐫^{})𝐤+{\displaystyle \frac{1}{2}}t_3\rho \left({\displaystyle \frac{𝐫+𝐫^{}}{2}}\right)\delta ^3(𝐫𝐫^{}),`$ (7)
where $`𝐤=(_𝐫_𝐫^{})/2i`$ acts to the right and $`𝐤^{}=(_𝐫^{}_𝐫)/2i`$ acts to the left. The factor 1/2 on the density dependent term contains the contribution of a rearrangement effect . We use the parameter set of the Skyrme III force (SKIII) . The spin-orbit force is neglected. We assume that the motion of the DGDR is generated by a two-body operator $`\widehat{D}^2`$:
$`|\mathrm{\Psi }(t=0)=e^{ik\widehat{D}^2}|\mathrm{\Phi }_0,`$ (8)
where $`\widehat{D}`$ is a one-body dipole operator and $`|\mathrm{\Phi }_0`$ the ground-state wave function. The initial conditions for solving the coupled equations Eqs.(3)-(5) are determined with the use of the above wave function. We evaluate the initial values of $`C_{\alpha \beta \alpha ^{}\beta ^{}}`$
$`C_{\alpha \beta \alpha ^{}\beta ^{}}(t=0)=\mathrm{\Psi }(t=0)|a_\alpha ^{}^+a_\beta ^{}^+a_\beta a_\alpha |\mathrm{\Psi }(t=0),`$ (9)
assuming that $`|\mathrm{\Phi }_0`$ is the Hartree-Fock (HF) ground-state wave function. At first order of $`k`$, the initial condition for $`C_{\alpha \beta \alpha ^{}\beta ^{}}`$ becomes
$`C_{\mu \nu \rho \sigma }`$ $`=`$ $`\mathrm{\Psi }|a_\rho ^+a_\sigma ^+a_\nu a_\mu |\mathrm{\Psi }`$ (10)
$`=`$ $`2ik\{\mu |D|\rho \nu |D|\sigma \mu |D|\sigma \nu |D|\rho \}`$ (11)
$`C_{\rho \sigma \mu \nu }`$ $`=`$ $`\mathrm{\Psi }|a_\mu ^+a_\nu ^+a_\sigma a_\rho |\mathrm{\Psi }`$ (12)
$`=`$ $`2ik\{\rho |D|\mu \sigma |D|\nu \rho |D|\nu \sigma |D|\mu \},`$ (13)
where $`\rho `$ and $`\sigma `$ refer to unoccupied single-particle states, and $`\mu `$ and $`\nu `$ refer to occupied ones. We choose the dipole operator in the above equations such as $`D=\tau _zz`$. Other elements of the initial $`C_{\alpha \beta \alpha ^{}\beta ^{}}`$ vanish at first order of $`k`$. Similarly, non-varnishing initial values of $`n_{\alpha \alpha ^{}}`$ become
$`n_{\mu \rho }`$ $`=`$ $`\mathrm{\Psi }|a_\rho ^+a_\mu |\mathrm{\Psi }`$ (14)
$`=`$ $`2ik{\displaystyle \underset{\nu }{}}\mu |D|\nu \nu |D|\rho `$ (15)
$`n_{\rho \mu }`$ $`=`$ $`\mathrm{\Psi }|a_\mu ^+a_\rho |\mathrm{\Psi }`$ (16)
$`=`$ $`2ik{\displaystyle \underset{\nu }{}}\rho |D|\nu \nu |D|\mu .`$ (17)
For the initial $`\psi _\alpha `$’s we use the HF single-particle wave functions. The strength function of the DGDR, defined by
$`S_2(E)={\displaystyle \underset{n}{}}|\mathrm{\Phi }_n|\widehat{D}^2|\mathrm{\Phi }_0|^2\delta (EE_n),`$ (18)
is given by the Fourier transform of the time-dependent two-body dipole moment $`D_2(t)`$ as
$`S_2(E)={\displaystyle \frac{1}{\pi k\mathrm{}}}{\displaystyle _0^{\mathrm{}}}D_2(t)\mathrm{sin}{\displaystyle \frac{Et}{\mathrm{}}}dt,`$ (19)
where $`D_2`$ is given by
$`D_2(t)`$ $`=`$ $`\mathrm{\Psi }(t)|\widehat{D}^2|\mathrm{\Psi }(t)`$ (20)
$`=`$ $`{\displaystyle \underset{\alpha \alpha ^{}}{}}\alpha |D^2|\alpha ^{}n_{\alpha ^{}\alpha }+{\displaystyle \underset{\alpha \beta \alpha ^{}\beta ^{}}{}}\alpha |D|\alpha ^{}\beta |D|\beta ^{}\{A[n_{\alpha ^{}\alpha }n_{\beta ^{}\beta }]+C_{\alpha ^{}\beta ^{}\alpha \beta }\}.`$ (21)
The terms without $`C_{\alpha \beta \alpha ^{}\beta ^{}}`$ in the above equation have negligible contribution to the Fourier transformation in Eq.(14). The $`k`$ dependence of $`S_2(E)`$ thus obtained is negligible as long as $`k`$ is sufficiently small. The energy-weighted sum rule (EWSR) for the DGDR is given as
$`{\displaystyle _0^{\mathrm{}}}ES_2(E)𝑑E`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{\Phi }_0|[\widehat{D}^2,[H,\widehat{D}^2]|\mathrm{\Phi }_0`$ (22)
$`=`$ $`{\displaystyle \frac{2\mathrm{}^2}{m}}\mathrm{\Phi }_0|\widehat{D}^2|\mathrm{\Phi }_0+4(t_1+t_2)\mathrm{\Phi }_0|\widehat{R}\widehat{D}^2|\mathrm{\Phi }_0,`$ (23)
where $`H`$ is the total hamiltonian and $`m`$ the nucleon mass. The second term on the right-hand side of Eq.(16) is due to the momentum dependence of the Skyrme force and $`\widehat{R}`$ is the following two-body operator
$`\widehat{R}={\displaystyle \underset{ip,jn}{}}\delta ^3(𝐫_i𝐫_j).`$ (24)
The EWSR value is evaluated with the use of the HF wave function for $`|\mathrm{\Phi }_0`$. The second term on the right-hand side of Eq.(16) has a contribution of about 30% to the total EWSR value.
To solve the coupled equations Eqs.(3)-(5), we use a minimum number of single-particle states: the $`1s,1p,2s`$ and $`1d`$ single-particle orbits for <sup>16</sup>O and the $`1s,1p,2s,1d,2p`$ and $`1f`$ orbits for <sup>40</sup>Ca. To check the validity of such truncation of the single-particle space, we performed RPA calculations for the GDR strength functions in <sup>16</sup>O and <sup>40</sup>Ca using the time-dependent RPA equations in the same truncated space. The obtained results were compared with those of the TDHF calculations which correspond to continuum RPA calculations . The fractions of the EWSR values depleted in the energy interval $`040`$MeV were turned out to be 90% in <sup>16</sup>O and 93% in <sup>40</sup>Ca, respectively. These values are sufficiently large and comparable with the TDHF values which are close to 100% . However, the excitation energies of the GDR’s in RPA were slightly larger than those in TDHF. To adjust the excitation energies of the GDR’s to the TDHF values, we reduced the parameter $`x_0`$ of the spin-dependent term of the Skyrme force Eq.(6). The obtained value of $`x_0`$ was 0.3 instead of the original value of 0.45. We use this reduced value of $`x_0`$ in the following calculations. The spin-dependent term of the Skyrme force has a negligible contribution to the mean-field potential in spin-isospin symmetric nuclei like <sup>16</sup>O and <sup>40</sup>Ca considered here and, therefore, the single-particle wave functions are not affected by the reduction of the strength of the spin-dependent term. The integration in Eq.(14) is performed for a finite time interval of $`1.52\times 10^{21}`$s. As a result $`S_2(E)`$ has small fluctuations. To reduce the fluctuations in $`S_2(E)`$, we multiply $`D_2(t)`$ by a damping factor $`e^{\mathrm{\Gamma }t/2}`$ before performing the time integration. This corresponds to smoothing the strength function with a width $`\mathrm{\Gamma }`$. We use $`\mathrm{\Gamma }=1`$MeV. Other calculational details are explained in our previous publications .
The strength distribution of the DGQR in <sup>16</sup>O calculated in TDDM is shown in Fig.1 (thick solid line). The bump seen around $`E=45`$MeV corresponds to the DGDR. The strength function $`S_1(E)`$ of the GDR obtained from the time-dependent RPA calculation is also shown in Fig.1 (dotted line). The width of the GDR is small and nearly equal to the width due to the smoothing and the finite time integration as explained above. The thin vertical bar at $`E=45.8`$MeV indicates the location of the DGDR strength predicted from the GDR shown in Fig.1. The fraction of the EWSR value of the DGDR depleted in the energy interval 10–60MeV is 82%. The centroid energy of the DGDR strength distribution between 35MeV and 55MeV is 44.6MeV. The energy difference $`\mathrm{\Delta }E=45.8`$MeV$`44.6`$MeV$`=1.2`$MeV is small but slightly larger than the value $`\mathrm{\Delta }E0.8`$MeV obtained from the work done by de Souza Cruz and Weiss using the generator coordinate method and the value $`\mathrm{\Delta }E<0.7`$MeV obtained from the formula given by Bertsch and Feldmeire . Though the centroid energy of the DGDR is close to twice the GDR energy, the DGDR in <sup>16</sup>O has a large width due to the Landau damping.
The strength distribution of the DGQR in <sup>40</sup>Ca calculated in TDDM is shown in Fig.2 (solid line). The bump seen around $`E=40`$MeV corresponds to the DGDR. The strength function $`S_1(E)`$ of the GDR in <sup>40</sup>Ca calculated in the time-dependent RPA is also shown in Fig.2 (dotted line). The GDR strength is split into two peaks in the case of <sup>40</sup>Ca. A similar split is seen in the TDHF calculation for <sup>40</sup>Ca . The fraction of the EWSR value of the DGDR depleted in the energy interval 10–60MeV is 88%. In Fig.3 the strength function of the DGDR (thick solid line) is compared with what is expected from that of the GDR shown in Fig.2: The thin vertical bars in Fig.3 indicate the locations and relative strengths of the DGDR predicted from the GDR in RPA, assuming that the GDR consists of the two discrete components. The centroid of the DGDR strength distribution in the energy range 30–50MeV is 39.6MeV, while that of the GDR prediction is 39.0MeV. The energy difference $`\mathrm{\Delta }E=0.6`$MeV may be larger than 0.2MeV obtained from Refs. but is smaller than 1.0MeV given by the shell model calculation for <sup>40</sup>Ca . Though the main peak in TDDM located at 40MeV has more strengths than that predicted from the RPA calculation, the DGDR seems to have weaker Landau damping in <sup>40</sup>Ca than in <sup>16</sup>O and the shape of the strength distribution is similar to what is expected from the GDR in RPA. This means that the dipole mode becomes fairly harmonic in the mass region of <sup>40</sup>Ca. A similar conclusion was obtained by de Souza Cruz and Weiss comparing the DGDR in <sup>16</sup>O with that in <sup>40</sup>Ca. The small anharmonicities of the DGDR in Ca have also been pointed out by Catara et al. using a boson expansion approach and the anharmonicities of the DGDR’s in heavier nuclei have recently been studied by several groups .
In summary, the strength functions of the DGDR’s in <sup>16</sup>O and <sup>40</sup>Ca were calculated in TDDM in a self-consistent manner, in which the same Skyrme force as that used for the calculation of the mean-field potential was used for the two-body correlation function. It was pointed out that the strength function of the DGDR is obtained from the Fourier transform of the time-dependent two-body dipole moment. It was found that in both nuclei the excitation energies of the DGDR’s are very close to twice those of the GDR’s. It was also found that the DGDR in <sup>16</sup>O has a large width due to the Landau damping, indicating the anharmonicity of the dipole mode in <sup>16</sup>O. |
warning/0003/astro-ph0003042.html | ar5iv | text | # Low Frequency Radio Emission of Pulsar PSR J1907+0919 Associated with the Magnetar SGR 1900+14
## 1. Introduction
Among known soft gamma- ray repeaters to date, only SGR 1900+14 and SGR 1806-20 are the objects for which a secular spin-down of the pulse periods (5.16 s and 7.47 s accordingly) with $`\dot{P}`$ of order $`10^{10}`$s/s was detected and thereby was established that these SGRs are neutron stars with a superstrong magnetic field of order $`10^{15}`$G (Kouveliotou et al. 1998; Kouveliotou et al. 1999), called as a ”magnetars” (Duncan & Thompson 1992).
Since the end of 1998 we carried out the observations of the SGR 1900+14 at low frequency (111 MHz) and have detected the periodic pulsed radio emission from this magnetar (Shitov 1999). In this paper we report the results of our observations obtained till August 1999.
## 2. Observations
The observations of the magnetar SGR 1900+14 were started since 1998 December using the BSA radio telescope operating now at the new frequency of 111 MHz. (In 1998 October BSA - wavelength dipole array with dimensions of 187 x 384 m, transit time of $`200s/Cos(\delta )`$ \- was reconstructed to shift the previous operating frequency of 102.5 MHz to the new one of 111 MHz (Kutuzov et al. 1999) ) We used 128 x 20 kHz filterbank receiver and multichannel recording system to record the individual pulses during the BSA transit time. Different sampling intervals of 22.6816 and 20.1728 ms were used. All further processing procedures: cleaning of the noise signal in a each frequency channel, dedispersion, spectral analysis, the folding of dedispersed data for integrated pulse profile and the timing analysis were made ”off line”.
## 3. The detection of radio pulses from SGR 1900+14
For the first time we have detected pulsed with 5.161-s period radio emission from SGR 1900+14 in 1998 December 12, when as result of dedispersion procedures the integrated pulse profile with a good enough signal/noise ratio appeared with a dispersion measure ( DM ) of about $`280pccm^3`$. It was surprising that the revealed profile was quite narrow, with the width of about $`7.^{}5`$. The next radio detection of SGR 1900+14 like the first one was only in 1999 January 6 (note that between this dates the antenna array was often covered by hoar-frost). Fig. 1 demonstrates these two detection of the SGR signals at 111 MHz.
During subsequent regular observations the pulses from SGR 1900+14 with the same DM values and similar pulse width of about 100 ms were detected repeatedly. The estimated for the best records mean flux density at 111 MHz is approximately 50 mJy. The measured value of $`DM=281.4(9)pccm^3`$. So far as the period derivative $`\dot{P}`$ was known for the magnetar with the limited accuracy the pulses were recorded in a unpredictable phases (Fig.2, left) till the middle of January 1999 when $`P`$ and $`\dot{P}`$ were improved as a result of the timing analysis.
Fig.2 (right) demonstrates some pulse profiles of SGR 1900+14, the phase of which aligned in an accordance with the timing solution.
There were a number of records in which we have detected interpulses with amplitude (and, of course, dispersion measure) like the main pulse, and which were located in different phases. Two examples of these interpulses are seen in Fig. 2 (09 01 99 and 17 06 99).
## 4. Timing analysis
For timing analysis of our data as the initial parameters of SGR 1900+14 we used the values of P = 5.160199(2) s, $`\dot{P}=1.14(23)10^{10}`$ s/s, MJD = 51056.0, obtained by Kouveliotou et al. (1999) in 1998 August 28. The exact position of the magnetar: $`R.A.(J2000)=19h07m14.^\mathrm{s}33`$, $`Dec.(J2000)=09{}_{}{}^{}19{}_{}{}^{}20.^{\prime \prime }1`$ , obtained with the VLA observations of an outburst of relativistic particles from SGR 1900+14 (Frail, Kulkarni & Bloom 1999) we used in our analysis. The first timing solution for $`P`$ and $`\dot{P}`$ was found for the first records, obtained in the date interval Dec. 12 1998 - Jan. 14 1999, after that in the subsequent observations with improved parameters the pulses were detected in the precalculated phases. In IAUC 7110 we have reported the solution: P = 5.161297854 (83)s, $`\dot{P}`$ = 1.23228 (34)s/s, which was found for the interval Dec. 12 1998 - Feb. 4 1999.
In this paper we present the results, obtained during Dec. 12 1998 \- July 30 1999. We have selected 16 best records and found for this dates interval the following timing solution: P = 5.161297899 (67) s, $`\dot{P}=1.23197(15)10^{10}`$ s/s, $`\ddot{P}=+0.53(14)10^{20}`$ s/s/s, MJD = 51159.4605 with the RMS of phase residual of 28 ms. The inclusion of the second period derivative essentially improved the timing in this interval of dates (see Fig. 3). For the epoch MJD = 51056.0, found parameters give P = 5.1601969 s, that is in a good agreement with Kouveliotou’s et al. (1999) measurements.
## 5. Conclusions
New radio pulsar PSR J1907+0919 associated with the soft gamma repeater SGR 1900+14 is representative of a new class of pulsars with a superstrong magnetic field, slow down value of which for this pulsar is $`8.110^{14}`$ G. Presented results confirm that this object is a magnetar (Duncan & Thompson 1992; Kouveliotou et al. 1999). The pulsar distance of about 5.8 kpc determined from dispersion measure $`DM=281pccm^3`$ supports the suggested earlier in a number of papers genetic connection of SGR 1900+14 with supernova remnants SNR G42.8+0.6. As timing analysis have shown there is no evidence for binary orbital motion of this pulsar, at least with $`P_{orb}<250`$ days and with $`a\mathrm{sin}i>60`$ ms.
### Acknowledgments.
This work was supported partly by grant INTAS 96-0154.
## References
Duncan, R. C., & Thompson, C. 1992, ApJ, 392, L9
Frail, D.A., Kulkarni, S.R. & Bloom, J.S. 1999, Nature, 398, 127
Hurley, K., Kouveliotou, C., Murakami, T., & Strohmayer, T. 1998, IAU Circ. 7001
Kouveliotou, C., et al. 1998, Nature, 393, 235
Kouveliotou, C., Strohmayer, T., Hurley, K., van Paradijs, J., Finger, M. H., Dieters, S., Woods, P., Thompson, C. & Duncan, R. C. 1999, ApJ, 510, L115
Kutuzov, S. M., et al. 1999 in Lebedev Physical Institute Proc., submitted
Shitov, Yu.P. 1999, IAU Circ. 7110 |
warning/0003/hep-ph0003298.html | ar5iv | text | # Extending the Velocity-dependent One-scale String Evolution Model
## I Introduction
The velocity-dependent one-scale (VOS) model provides the most convenient and reliable method by which to calculate the large-scale quantitative properties of a string network in cosmological and other contexts. It is widely used for making quantitative predict ions of the potential observational implications of cosmic strings. Given its simplicity, it is remarkable how well the VOS model performs when tested against high resolution numerical simulations of string networks. It is well-known that string evolution is a complex physical process with a build-up of small-scale structure on the strings, which is very computationally demanding to model accurately. Analytic approaches like the VOS model, abandon the possibility of describing the statistical physics of the string network accurately and concentrate instead on its thermodynamics. In other words, a small number of macroscopic quantities are selected and the microscopic string equations of motion are used to derive evolution equations for these averaged quantities. The price to be paid in this approach is that the averaging process introduces phenomenological parameters whose values are not specified by the model itself. Instead, one must still fix these parameters by direct comparison with numerical simulations.
The VOS model is a generalization of the ‘one-scale’ model pioneered by Kibble (see also ref.) which describes string motion in terms of a single correlation length $`L`$. By incorporating a variable rms velocity $`v`$, the VOS model extends its validity into early regimes with frictional damping and across the important matter-radiation transition, thus giving a quantitative picture of the complete history of a cosmic string network. Other analytic approaches to string evolution have attempted to incorporate the additional small-scale structure seen in numerical simulations. This includes a ‘kink-counting’ model, a functional approach , a ‘three-scale’ model and a ‘wiggly’ model . While these are important for characterising detailed network features, they introduce a significant number of further phenomenological parameters which must be fixed by simulations (and which remain rather uncertain). Nevertheless, for describing the large-scale properties of a long-string network, the VOS model has proved to be sufficient for a good quantitative fit using only a single parameter, the loop chopping efficiency $`\stackrel{~}{c}`$.
The purpose of the present paper is, first, to provide a concise exposition of the VOS string evolution model. We summarise how it can be applied to describe cosmological string evolution, including late times with a cosmological constant, and we present the very different histories of both GUT- and electroweak-scale strings. Secondly, we propose an improvement of the VOS model by presenting a new ansatz for the momentum parameter $`k`$, which we justify both analytically and numerically. Thirdly, we present a further extension incorporating radiation backreaction, which provides small corrections to the cosmological scaling laws and which also compares favourably with published results of global string simulations. Finally, we review generalizations in a curved FRW spacetime, giving some further asymptotic scaling solutions. We report on detailed comparisons between numerical string simulations and the VOS model elsewhere .
## II The VOS string network model
The velocity-dependent one-scale model has been described in considerable detail elsewhere , so here we limit ourselves to a brief summary which highlights the features that will be important for what follows. Also for simplicity, we will only discuss the evolution of the long string network, even though this formalism is also applicable to the loop population, mutatis mutandis. We will discuss this case in detail elsewhere .
### A The averaged evolution equations
Averaged quantities which we could use to describe the string network are its energy $`E`$ and RMS velocity $`v`$ defined by
$$E=\mu a(\tau )ϵ𝑑\sigma ,v^2=\frac{\dot{𝐱}^2ϵ𝑑\sigma }{ϵ𝑑\sigma },$$
(1)
where the string trajectory $`𝐱(\sigma ,t)`$ is parametrised by the worldsheet coordinates $`\sigma `$ and $`t`$ and the ‘energy density’ $`ϵ(\sigma ,t)`$ gives the string length per unit $`\sigma `$ along the string.
Any string network divides fairly neatly into two distinct populations, long (or ‘infinite’) strings and small closed loops with corresponding quantities denoted by a subscript $`\mathrm{}`$ and $`\mathrm{}`$ respectively. The long string network is a Brownian random walk on large scales and can be characterised by a correlation length $`L`$. This can be used to replace the energy $`E_{\mathrm{}}=\rho _{\mathrm{}}V`$ in long strings in our averaged description, that is,
$$\rho _{\mathrm{}}\frac{\mu }{L^2}.$$
A phenomenological term must then be included to account for the loss of energy from long strings by the production of loops, which are much smaller than $`L`$. A ‘loop chopping efficiency’ parameter $`\stackrel{~}{c}`$ is introduced to characterise this loop production as
$$\left(\frac{d\rho _{\mathrm{}}}{dt}\right)_{\mathrm{to}\mathrm{loops}}=\stackrel{~}{c}v_{\mathrm{}}\frac{\rho _{\mathrm{}}}{L}.$$
(2)
In this approximation, we would expect the loop parameter $`\stackrel{~}{c}`$ to remain constant irrespective of the cosmic regime, because it is multiplied by factors which determine the string network self-interaction rate.
From the microscopic string equations of motion, one can then average to derive the evolution equation for the correlation length $`L`$,
$$2\frac{dL}{dt}=2HL(1+v_{\mathrm{}}^2)+\frac{L}{\mathrm{}_\mathrm{f}}v_{\mathrm{}}^2+\stackrel{~}{c}v_{\mathrm{}},$$
(3)
where $`H`$ is the Hubble parameter and $`\mathrm{}_\mathrm{f}`$ is a friction damping length scale. The first term in (3) is due to the stretching of the network by the Hubble expansion which is modulated by the redshifting of the string velocity. The second term is due to frictional interactions by a high density of background particles scattering off the strings. The friction length scale $`\mathrm{}_\mathrm{f}`$ (defined in ref. ) typically depends on the background temperature as $`\mathrm{}_\mathrm{f}(1)\mu T^3`$, so that it grows with the scale factor as $`a^3`$. It usually becomes irrelevant after a time $`t_{}(G\mu )^1t_c`$ (with $`t_c`$ being the epoch at which the network was formed) which is a very short time for GUT-scale strings but can be as late as $`t_0`$ for electroweak strings. Note also that equation (3) is valid for an arbitrary flat FRW model with $`H`$ given by the Friedmann equation,
$$H^2\left(\frac{\dot{a}}{a}\right)^2=H_0^2\left(\mathrm{\Omega }_{\mathrm{R0}}a^4+\mathrm{\Omega }_{\mathrm{M0}}a^3\right)+\frac{1}{3}\mathrm{\Lambda }.$$
(4)
where $`\mathrm{\Omega }_{\mathrm{R0}}`$ and $`\mathrm{\Omega }_{\mathrm{M0}}`$ are the fractional radiation and matter densities today at $`t_0`$, $`\mathrm{\Lambda }`$ is the cosmological constant and we take $`a(t_0)=1`$.
One can also derive an evolution equation for the long string velocity with only a little more than Newton’s second law
$$\frac{dv_{\mathrm{}}}{dt}=\left(1v_{\mathrm{}}^2\right)\left[\frac{k}{L}\left(2H+\frac{1}{\mathrm{}_\mathrm{f}}\right)v_{\mathrm{}}\right],$$
(5)
where $`k`$ is called the ‘momentum parameter’. The first term is the acceleration due to the curvature of the strings and the second damping term is from both the expansion and background friction. The parameter $`k`$ is defined by
$$k\frac{(1\dot{𝐱}^2)(\dot{𝐱}𝐮)}{v(1v^2)},$$
(6)
with $`\dot{𝐱}`$ the microscopic string velocity and $`𝐮`$ a unit vector parallel to the curvature radius vector. In previous work , we left $`k`$ as a second phenomenological parameter, while pointing out that it is related to small-scale structure and also demonstrating specified asymptotic dependencies on the velocity. In the next section, however, we justify an accurate ansatz for $`k`$ which removes this additional freedom. For most relativistic regimes relevant to cosmic strings it is sufficient to define it as follows:
$$k_{\mathrm{rel}}(v)=\frac{2\sqrt{2}}{\pi }\frac{18v^6}{1+8v^6}.$$
(7)
In the extreme friction-dominated case ($`v0`$), we have the nonrelativistic limit $`k_{\mathrm{nr}}=2\sqrt{2}/\pi `$ with a more complicated ansatz than (7) interpolating between these limits for intermediate regimes.
Finally, we end this summary by noting that the VOS model has been extensively compared with the results of numerical simulations and shown to provide a good fit to the large-scale properties of a string network. In particular, it matches well the evolution between asymptotic regimes as a network passes through the matter–radiation transition. Comparisons with numerical simulations confirm the constancy of the only free parameter, the loop chopping efficiency $`\stackrel{~}{c}`$, and fix its value to be
$$\stackrel{~}{c}=0.23\pm 0.04.$$
(8)
The VOS model for any flat FRW cosmology, then, consists of the evolution equations (3) and (5) with the parameters $`c`$ specified in (8) and $`k`$ given by (7) (or a more accurate general expression given below) and the scale factor $`a`$ satisfying the Friedmann equation (4).
### B Scale-invariant solutions
We now start to use the VOS model to provide a general overview of the evolution of string networks in various cosmological scenarios. First we analyse some basic late-time properties of cosmic string networks, neglecting the effect of friction due to particle scattering. A crucial question is whether or not they can reach a ‘scale invariant’ attractor solution which is required, among other things, for a Harrison-Zel’dovich spectrum of primordial density fluctuations to be generated. This can be discussed by analysing the VOS equations (3) and (5).
Scale-invariant solutions of the form $`Lt`$, $`LH^1`$ or $`Ld_H`$, together with $`v_{\mathrm{}}=const.`$, only appear to exist when the scale factor is a power law of the form
$$a(t)t^\beta ,\beta =const.,0<\beta <1.$$
(9)
This condition implies that
$$LtH^1d_H,$$
(10)
with the proportionality factors dependent on $`\beta `$. It is useful to introduce the following useful parameters to describe the relative correlation length and densities, defining them respectively as
$$L=\gamma t,\zeta \gamma ^2=\rho _{\mathrm{}}t^2/\mu .$$
(11)
By looking for stable fixed points in the VOS equations, we can express the actual scaling solutions in the following implicit form:
$$\gamma ^2=\frac{k(k+\stackrel{~}{c})}{4\beta (1\beta )},v^2=\frac{k(1\beta )}{\beta (k+\stackrel{~}{c})},$$
(12)
where $`k`$ is the constant value of $`k(v)`$ given by solving the second (implicit) equation for the velocity. Although it may not be obvious by inspection, it is easy to verify numerically that this solution is well-behaved and stable for all realistic parameter values.
If the scale factor is not a power law, then simple scale-invariant solutions like (12) do not exist. Physically this happens because the network dynamics are unable to adapt rapidly enough to the changes in the background cosmology. A prime example of this is, of course, the transition between the radiation and matter-dominated eras. The evolution of a GUT-scale network shown in fig. 1 illustrates asymptotic regimes in which string evolution is scale-invariant, as well as the matter–radiation transition where it is not. Note that for realistic cosmological parameter choices, a string network today is still only slowly approaching its asymptotic matter density. Since, the cosmological importance of the changes in the network properties during the matter–radiation transition cannot be over-emphasised, it is important to calculate these accurately either with direct numerical simulations or with the VOS or similar analytic model. The same is also true for late time curvature or cosmological constant domination.
### C Friction-dominated scaling solutions
During friction-dominated epochs one has different ‘scaling’ solutions. However, these are no longer ‘scale-invariant’, since in this case the network retains a memory of its initial conditions, and in particular of the epoch of formation. This can be conveniently expressed by a parameter
$$\theta \left(\frac{t_c}{t_{Pl}}\right)^{1/2},$$
(13)
which is essentially the value of the ratio of the damping terms due to friction and Hubble damping, measured at the epoch of string formation.
Thus in realistic cosmological contexts one can have two different regimes . The first is a ‘stretching’ regime,
$$\frac{L}{L_c}=\left(\frac{t}{t_c}\right)^{1/2},v=\frac{t}{\theta L_c},$$
(14)
which is an early-time, transient period which will occur when the initial string density and velocity are sufficiently low—for example, as a result of a slow first-order phase transition. In this case the network starts out with a correlation length significantly larger than the damping length and so is ‘frozen’, and is conformally stretched. However the damping length is growing as $`\mathrm{}_ft^{3/2}`$, so it quickly catches up with it, ending this regime. However, this can last for many orders of magnitude in time for electroweak-scale networks. Although this is not cosmologically relevant except for extremely light strings, the analogous regime in the matter-dominated case would be
$$Lt^{2/3},vt^{4/3}.$$
(15)
The true attractor solution for a friction-dominated epoch, which follows the stretching regime (if this exists) is the Kibble regime, which in the radiation era has the form
$$\frac{L}{L_c}=\left[\frac{2k_{nr}(\stackrel{~}{c}+k_{nr})}{3\theta }\right]^{1/2}\left(\frac{t}{t_c}\right)^{5/4},v=\left[\frac{3k_{nr}}{2\theta (\stackrel{~}{c}+k_{nr})}\right]^{1/2}\left(\frac{t}{t_c}\right)^{1/4},$$
(16)
where $`k_{nr}`$ is the value of the momentum parameter in the nonrelativistic limit given above. In this case the correlation length stays halfway between the damping length and the horizon length. Again there is a matter era analogue, which as the form
$$Lt^{3/2},vt^{1/2},$$
(17)
but this is rarely relevant cosmologically. In the extreme case of an electroweak scale network, friction domination ends after radiation-matter equality, but not far enough away from it for this regime to be reached.
These points are illustrated in fig. 2 where we see the evolution of an electroweak-scale string network. It is interesting to note the scaling behaviour for two cases with very different initial conditions from extreme first- and second-order transitions. The high density strings from a second-order transition quickly approach the Kibble scaling regime (16) discussed above. However, the low density strings from a first-order transition begin in a distinct stretching regime (14) and persist in it with their density falling slowly until it matches that for the attractor Kibble regime. In this case, the string network retains a ‘memory’ of its initial density for about ten orders of magnitude in cosmic time. However, even during the friction-dominated regime the network is able to erase this memory once the Kibble regime is reached. On the other hand, the memory of the epoch of formation is not erased—one could in principle recover it by measuring the parameter $`\theta `$. This can also be seen for GUT strings in fig. 1, although the network does not have time to relax into a definite scaling regime before friction-domination ends. For electroweak strings, there are also interesting departures from scaling behaviour at the matter–radiation transition and the network remains friction-dominated until about three orders of magnitude in time afterwards. Also note that in the $`\mathrm{\Omega }_m=1`$ case the strings only reach the relativistic regime at about the present time, and in the observationally preferred case $`\mathrm{\Omega }_m=0.2,\mathrm{\Omega }_\mathrm{\Lambda }=0.8`$ they are always non-relativistic. This point is crucial, among other things, for a quantitative analysis of the evolution of superconducting strings and vortons .
### D A cosmological constant
We can also use the VOS model in a flat background to discuss the domination at late times by a cosmological constant $`\mathrm{\Lambda }`$ , a model for which there appears to be growing observational evidence. In the extreme asymptotic case when the universe is inflating we have $`a\mathrm{exp}(Ht)`$ with $`H=\sqrt{\mathrm{\Lambda }/3}`$. The network will ‘freeze out’ and will simply be conformally stretched, that is,
$$La,v_{\mathrm{}}a^1,$$
(18)
where, as soon as the strings become nonrelativistic $`k_{\mathrm{nr}}=2\sqrt{2}/\pi `$, their product satisfies
$$Lv_{\mathrm{}}=\frac{2\sqrt{2}}{\pi }H.$$
(19)
Of course, this solution will only apply at early times actually during inflation. At the present time we will only be slowly approaching a new stretching regime, so we have to solve the VOS model explicitly. This is shown for a model in which $`\mathrm{\Omega }_\mathrm{\Lambda }=0.8`$, as a dashed line at late times in figs. 12, as well as in detail for GUT-scale strings in fig. 3. It is clear that there is a significant fall in the string density and velocity, an effect which, for example, would affect the large-angle anisotropies in the cosmic microwave sky. The evolution of the string network is clearly not scale-invariant during any period after $`t_{\mathrm{eq}}`$.
## III The momentum parameter
Having introduced the VOS model and some of its key cosmological implications, we now turn to a more detailed discussion of the so-called ‘momentum parameter’ $`k`$ defined in (6) and which is important for solving (5). String velocities are determined by the current acceleration to which they are subjected from the local string curvature, as well as their ‘bulk’ momentum left-over from previous accelerations. Heuristically, we can imagine separating the velocity into these the curvature ‘c’ and bulk ‘p’ contributions as $`\dot{𝐱}=\dot{𝐱}_\mathrm{c}+\dot{𝐱}_\mathrm{p}`$. In the extreme friction-dominated limit, the velocity is entirely due to the curvature and reaches a limiting average velocity set by the friction length scale, $`v_\mathrm{c}\dot{𝐱}_\mathrm{c}^2^{1/2}=\mathrm{}_f/L`$. However, as the velocity increases towards relativistic values we can expect the momentum contribution to become larger. Let us suppose that their relative contribution is proportional to some power law of the total velocity $`v`$, that is, $`v_\mathrm{p}/v_\mathrm{c}v^\alpha `$ where $`\alpha `$ is clearly greater than unity. In this case, one can find after some straightforward (although tedious) manipulations that the following approximate relation holds for the momentum parameter $`k`$
$$k\frac{12^\alpha v^{2\alpha }}{1+2^\alpha v^{2\alpha }}.$$
(20)
In this we have also used the fact that flat spacetime analytic calculations have shown that
$$k(1/\sqrt{2})=0,$$
(21)
which holds exactly. Note that the above expression means that (20) can also be approximately written as
$$k\frac{1(v_p/v_c)^2}{1+(v_p/v_c)^2}.$$
(22)
We can determine $`\alpha `$ by studying the well-known helicoidal string solution in flat space but perturbed by a frictional force. This solution is
$$𝐱=(A\mathrm{sin}\sigma (\mathrm{cos}t+\eta ),A\mathrm{cos}\sigma (\mathrm{cos}t+\eta ),\sqrt{1A^2}\sigma ),$$
(23)
where $`0A1`$ and $`\eta `$ is a small perturbation, which vanishes if there is no friction. Here, $`A=1`$ corresponds to a circular loop, while $`A=0`$ is a static straight string. The evolution equation for the perturbation $`\eta `$ has the form
$$\ddot{\eta }+2\dot{\eta }\frac{\mathrm{sin}t\mathrm{cos}t}{1A^2\mathrm{sin}^2t}+\eta \frac{(1A^2)\mathrm{sin}^2t\mathrm{cos}^2t}{1A^2\mathrm{sin}^2t}=\frac{\mathrm{sin}t}{\mathrm{}_f}(1A^2\mathrm{sin}^2t),$$
(24)
where $`\mathrm{}_f`$ is the friction length scale. This can then be solved numerically, and from this solution one can calculate $`k`$. By changing parameters $`A`$ and $`\mathrm{}_f`$ one can do this for a wide range of velocities, and hence obtain a plot of $`k=k(v)`$. This is plotted in Fig. 4 and compared with cases $`\alpha =2`$ and $`\alpha =3`$ of our ansatz (20). It can be seen that $`\alpha =3`$ provides an excellent approximation in this regime.
There is, however, one problem with this simple ansatz, namely that it would give $`k(0)=1`$. Even though one might naively expect this to be the correct limit, it is not so. This can be seen easily as follows. Assume that the velocity of say a loop is determined only by curvature, that is, neglect the momentum contribution. (This should be valid in the non-relativistic case.) Then $`k`$ will be approximately given by
$$k\frac{<|\dot{𝐱}^2|(1\dot{𝐱}^2)>}{v(1v^2)}.$$
(25)
This quantity can be easily calculated for the analytic helicoidal solution, yielding
$$k=\frac{2\sqrt{2}}{\pi }\frac{12A^2/3}{1A^2/2}.$$
(26)
Hence we find in the small amplitude limit $`A0`$ that
$$k_{\mathrm{nr}}=\frac{2\sqrt{2}}{\pi }0.9.$$
(27)
We also note in passing that in the relativistic limit $`A1`$ this same calculation would give $`k4\sqrt{2}(3/\pi )0.6`$ instead of the true value $`k=0`$, which clearly demonstrates the importance of momentum in the relativistic case.
The final issue to be considered is the transition between the two regimes. The only reliable way of studying this issue is through direct measurement in a string network simulations with ultra-high resolution. We shall report on the details of this elsewhere . Here we simply point out that we do confirm the value $`k_{\mathrm{nr}}`$ as the non-relativistic limit. It is then easy to find a fitting function for the transition between the regimes which has the correct asymptotic limits described previously, that is,
$$k(v)=\frac{2\sqrt{2}}{\pi }(1v^2)(1+2\sqrt{2}v^3)\frac{18v^6}{1+8v^6}.$$
(28)
The additional factors are required to reproduce both the relativistic and non-relativistic limits accurately. Note that if one is only interested in the relativistic regime (say for GUT-scale cosmic strings, as in the present paper) then the simpler expression (7), that is
$$k_{\mathrm{rel}}(v)=\frac{2\sqrt{2}}{\pi }\frac{18v^6}{1+8v^6},$$
(29)
should be sufficiently accurate to provide reliable results. On the other hand, a reliable approximation for small non-relativistic velocities in the friction-dominated limit extends (27) as
$$k_{\mathrm{nr}}(v)=\frac{2\sqrt{2}}{\pi }(1v^2).$$
(30)
We plot these three ansatze in fig. 5, and also confirm the validity of $`k_{\mathrm{rel}}`$ and $`k_{\mathrm{nr}}`$ in the appropriate limits.
Before we end this section, however, it is wise to discuss the interpretation of parameter $`k`$ in this model. It should be kept in mind that this is, ab initio, a phenomenological parameter, which accounts for a number of non-trivial effects related to the presence of small-scale structures on the strings. By construction, our model does not explicitly account for these small-scale effects, and hence they end up somehow encoded in $`k`$. One should not therefore infer too much from the aesthetic qualities of the function $`k(v)`$ we find—it is simply a phenomenological parameter that does a good job. Presumably this parameter will have a much clearer interpretation in the context of a proper wiggly string evolution model .
## IV The effect of radiation back-reaction
We now turn to some further extensions of the VOS model. The effect of gravitational back-reaction on the long-string network can be included in the evolution equation for the correlation length (3) in the same way as previously achieved for the evolution of the length of a string loop. For gravitational radiation the following term can be added to the right-hand side of (3)
$$2\left(\frac{dL}{dt}\right)_{\mathrm{gr}}8\mathrm{\Sigma }_{\mathrm{gr}}v_{\mathrm{}}^6=8\stackrel{~}{\mathrm{\Gamma }}G\mu v_{\mathrm{}}^6.$$
(31)
Here, $`\stackrel{~}{\mathrm{\Gamma }}`$ is a constant which is a long-string analogue of the $`\mathrm{\Gamma }65`$ found for the radiative decay of strings. Of course, $`\stackrel{~}{\mathrm{\Gamma }}`$ will be affected by a number of physical factors such as the presence of small-scale structure, but we can expect it to satisfy $`\stackrel{~}{\mathrm{\Gamma }}\text{ }<\mathrm{\Gamma }`$ and it would be surprising if it were very much smaller. Clearly, due to the high velocity power ($`v^6`$) involved in radiative backreaction, this term will not be important in any regime where string motion is strongly friction-dominated (and hence non-relativistic). We note also that there is an interesting coincidence in the ansatz (7) for the momentum parameter $`k`$ which too has a $`v^6`$ power, but we are unsure as yet whether this has any deeper significance.
For global string radiation into Goldstone bosons or axions, the corresponding radiative decay term at a time $`t`$ will be
$$2\left(\frac{dL}{dt}\right)_{\mathrm{ax}}8\mathrm{\Sigma }_{\mathrm{ax}}v_{\mathrm{}}^6=\frac{8\stackrel{~}{\mathrm{\Gamma }}v_{\mathrm{}}^6}{2\pi \mathrm{ln}(t/\delta )},$$
(32)
where the logarithmic term arises because of the long-range fields of the global string and $`\delta `$ is the string width. For cosmological GUT-scale strings, the backreaction term for local strings is $`\mathrm{\Gamma }G\mu 10^4`$ whereas for global strings it is about three orders of magnitude larger.
Note that the velocity equation has no correction at this order due to the gravitational back-reaction effects. Such effects are already included through the string curvature, which acts as a source for the velocity equation (i.e., the $`1/L`$ term), which will be different in this case.
Remarkably, the inclusion of the back-reaction term does not affect the existence of a scale-invariant attractor solution. However, it does of course influence the quantitative values of the scaling parameters, as well as the timescale necessary for this solution to be reached. For example, the inclusion of back-reaction can make the approach to scaling much faster.
If the gravitational back-reaction is non-zero, one can distinguish two asymptotic cases. Firstly, if $`\mathrm{\Sigma }`$ is small (of order unity at most) then the effect of back-reaction on the scaling solution will also be small. This will be the case, for example, for most local or global string networks in a cosmological context. We can express this as
$$\gamma ^2\gamma _0^2\left(1+\mathrm{\Delta }\right),v^2v_0^2\left(1\mathrm{\Delta }\right),$$
(33)
where $`\gamma _0`$ and $`v_0`$ are the “unperturbed” scaling values, given by eqns. (12), and the back-reaction correction has the form
$$\mathrm{\Delta }=8\beta v_0^5\mathrm{\Sigma }=8\beta \left[\frac{k(1\beta )}{\beta (k+\stackrel{~}{c})}\right]^{5/2}\mathrm{\Sigma }.$$
(34)
On the other hand, for large enough values of $`\mathrm{\Sigma }`$, the back-reaction term will dominate the evolution equation for the string length scale $`L`$, and the attractor scale-invariant solution has a different form altogether. It is not possible to write this solution in closed form, even expressing $`k`$ implicitly as above. However, it is possible to write it as a series. The dominant term and the first correction take the form
$$\gamma =\frac{k}{2\beta }\left[\frac{8\beta \mathrm{\Sigma }}{k(1\beta )}\right]^{1/7}\left(1+\mathrm{\Delta }_1+\mathrm{}\right),$$
(35)
$$v=\left[\frac{k(1\beta )}{8\beta \mathrm{\Sigma }}\right]^{1/7}\left(1\mathrm{\Delta }_1+\mathrm{}\right),$$
(36)
with
$$\mathrm{\Delta }_1=\frac{1}{2^{6/7}7}(k+\stackrel{~}{c})\left[\frac{\beta }{k(1\beta )}\right]^{5/7}\mathrm{\Sigma }^{2/7}.$$
(37)
In fig. 6 we plot the approach to scaling of some relevant string networks in the radiation and matter eras. The different timescales for convergence towards the attractor solution are clearly noticeable. Here, we have chosen initial conditions that would correspond to somewhat extreme first and second order phase transitions. Also note that for the radiation era we have neglected the effect of friction due to particle scattering, in order to reproduce the initial conditions often used in numerical simulations of string networks. For each of the cases above, three curves are plotted, corresponding to the values $`\mathrm{\Sigma }=0`$ (no back-reaction), $`\mathrm{\Sigma }=1.25`$ (close to the maximum value that can be accurately described by the scaling solution (3334)) and $`\mathrm{\Sigma }=5.5`$ (beyond which the scaling solution (3537) becomes accurate).
For large $`\mathrm{\Sigma }`$, the effects of back-reaction seen in fig. 6 on could be quite dramatic for the string network density, but note that they are much less drastic for the string velocities. In particular, we emphasize that gravitational back-reaction alone does not slow down a string network to non-relativistic speeds—only a friction-dominated regime can achieve this.
Interestingly, there has been recent work on numerical simulations of global string networks which explore the strong backreaction regime described by (3537). These authors report a surprisingly low string density relative to the gauged case. For their expanding universe simulations in the realistic case with periodic boundary conditions, they find the following radiation and matter era densities respectively,
$$\zeta _{\mathrm{rad}}=0.9\pm 0.1,\zeta _{\mathrm{mat}}=0.5\pm 0.1.$$
(38)
These results are perfectly consistent (within the estimated error bars) with our extended VOS model if we adopt a back-reaction parameter
$$\mathrm{\Sigma }_{\mathrm{ax}\text{-}\mathrm{sim}}3.$$
(39)
Indeed, this corresponds to the approximate average value for $`\mathrm{\Sigma }_{\mathrm{ax}}`$ that one would estimate for simulations of this resolution. Present limitations on numerical dynamic range give the upper bound $`\mathrm{ln}(t/\delta )5`$ (at the end of the simulation), implying $`\mathrm{\Sigma }_{\mathrm{ax}\text{-}\mathrm{sim}}2`$ throughout.
It is important to note, however, that the immediate extrapolation of these results (38) to a cosmological context would be erroneous. For example, for GUT-scale global strings, at the present time we can expect $`\mathrm{ln}(t/\delta )100`$, which implies the appropriate backreaction parameter in this case will be $`\mathrm{\Sigma }_{\mathrm{ax}}0.1`$. Such cosmic global strings are firmly in the regime (3334) where backreaction effects are small and, in this case, these would reduce the local string density by less than 10%. We make a more detailed comparison with elsewhere .
## V String networks in general FRW spacetimes
In this section we discuss the behaviour of our model in more general FRW universes, and in particular in open universes .
The evolution equations will obviously be affected by the different behaviour of the scale factor, as given by the Friedmann equation
$$H^2\left(\frac{\dot{a}}{a}\right)=H_0^2\left(\mathrm{\Omega }_{\mathrm{R0}}a^4+\mathrm{\Omega }_{\mathrm{M0}}a^3+\mathrm{\Omega }_{\mathrm{Q0}}a^m\right)+\frac{1}{3}\mathrm{\Lambda }Ka^2.$$
(40)
Note that we are allowing for curvature, and also for an extra fluid whose energy density decays as $`a^m`$. It should be kept in mind that the Friedman equation should, in general, contain a contribution for the string density, since it is possible that this becomes cosmologically important.
However, apart from these effects, one must also include an additional correction due to the curvature . One should note that this is essentially the curvature radius of the strings, $`L`$, divided by the radius of spatial curvature of the universe,
$$=\frac{H^1}{1\mathrm{\Omega }^{1/2}}.$$
(41)
Indeed, after a certain amount of algebra, one finds correction terms that are of the form
$$w=1(1\mathrm{\Omega })(HL)^2.$$
(42)
Note that $`\mathrm{\Omega }`$ denotes the total density of the universe. For a universe with a critical density, $`\mathrm{\Omega }=1`$, we have $`w=1`$.
The evolution equation for the correlation length $`L`$ now takes the form
$$2\frac{dL}{dt}=2HL+\frac{L}{\mathrm{}_d}\frac{v_{\mathrm{}}^2}{w^2}+\stackrel{~}{c}v_{\mathrm{}}.$$
(43)
For simplicity wa have also defined a damping length, including both the effects of Hubble damping and friction,
$$\frac{1}{\mathrm{}_d}=2H+\frac{1}{\mathrm{}_f},$$
(44)
with the friction length scale $`\mathrm{}_f`$ being defined in .
Similarly, the velocity equation becomes
$$\frac{dv_{\mathrm{}}}{dt}=\left(1\frac{v_{\mathrm{}}^2}{w^2}\right)\left(w^2\frac{k}{L}\frac{v_{\mathrm{}}}{\mathrm{}_d}\right).$$
(45)
Note that these are valid for any cosmological scenario<sup>*</sup><sup>*</sup>*There are some additional subtleties involved when discussing the mechanism of loop production in the case of Minkowski space string networks, which make it quite different from any cosmological scenario. We shall discuss this important point elsewhere . We do expect the loop chopping efficiency $`\stackrel{~}{c}`$ to be a constant, regardless of the cosmological model, since it is supposed to be reflecting a rather deep and fundamental property of the evolution of a network. Indeed, we think that whether or not one finds a constant chopping efficiency can in some sense be seen as a measure of how accurately the analytic modelling is reproducing the true dynamics of the network.
We can now re-examine the question of the existence of ‘scale invariant’ attractor solutions. Again, scaling solutions of the form $`Lt`$, $`LH^1`$ or $`Ld_H`$, together with $`v_{\mathrm{}}=const.`$ will only exist provided one has
$$a(t)t^\beta ,\beta =const.,0<\beta <1,$$
(46)
but now we also require
$$\mathrm{\Omega }=const.$$
(47)
The simplest example of the second condition is of course a flat, $`\mathrm{\Omega }_{\mathrm{M0}}=1`$ universe, but there are examples of cosmological models which have attractors other than $`\mathrm{\Omega }=1`$ . In any case, note that there can be additional relations between the values of $`\beta `$ and $`\mathrm{\Omega }`$ for specific models. Writing $`L=\gamma t`$ as before, the scaling solution is now given in the implicit form
$$\gamma ^2=w^2\frac{k(k+\stackrel{~}{c})}{4\beta (1\beta )},v^2=w^2\frac{k(1\beta )}{\beta (k+\stackrel{~}{c})},$$
(48)
where $`k`$ is (implicitly) the constant value of $`k(v)`$ for the appropriate value of velocity, and
$$w=\frac{2(1\beta )}{(1\mathrm{\Omega })\beta k(k+\stackrel{~}{c})}\left[\left(1+\frac{(1\mathrm{\Omega })\beta k(k+\stackrel{~}{c})}{(1\beta )}\right)^{1/2}1\right].$$
(49)
Again, although it may not be immediately obvious, it can be checked numerically that this solution is well-behaved for all sensible values of the parameters. If the two conditions above do not hold, then a scaling solution will not exist.
We should also mention another cosmologically important solution: in an open universe with $`\mathrm{\Omega }0`$, $`at`$, the asymptotic solution is
$$L=At\left(\mathrm{ln}t\right)^{1/2}A=\left[\frac{k_{nr}\stackrel{~}{c}}{2(1k_{nr})}\right]^{1/2}2.13\stackrel{~}{c}^{1/2},$$
(50)
$$v_{\mathrm{}}=B\left(\mathrm{ln}t\right)^{1/2},B=\left[\frac{k_{nr}(1k_{nr})}{2\stackrel{~}{c}}\right]^{1/2}0.21\stackrel{~}{c}^{1/2},$$
(51)
with $`k_{nr}`$ given by(27). Note that this is not a scale-invariant solution, since $`H^1=t`$ and $`d_H=t\mathrm{ln}t`$. In other words, by looking at the network one would be able to determine when the curvature-dominated period had started.
## VI Discussion and conclusions
In this paper we have presented a modified version of the velocity-dependent one-scale (VOS) model which depends on a single free parameter, the loop chopping efficiency $`\stackrel{~}{c}`$. We have tested it against the largest and most accurate numerical simulations to date , and we find that it provides a good fit to the large-scale scaling properties of the string network in both the radiation and matter epochs, as well as in the transition between the two eras—we will describe these tests elsewhere . These facts and its intrinsic simplicity make this model particularly suited for any analytic or semi-analytic study of cosmic strings where one is only interested in the large-scale features of the network.
We have re-analysed some simple evolutionary properties of cosmic string networks in the light of the VOS model and corresponding numerical simulations. An important conclusion to note is that any realistic cosmic string network is not scaling at any time from just before the epoch of equal matter and radiation through to the present day. This is something that must be properly taken into account particularly when discussing string-seeded structure formation scenarios with GUT-scale strings. The extended VOS model is also valid when deviations from scaling are even larger at late times in a universe which becomes dominated by curvature or a cosmological constant.
Finally, we considered the effects of radiation back-reaction on the scaling properties of the long string network, and we have shown that although the existence (or otherwise) of a scale-invariant attractor solution will not be affected, the quantitative scaling properties can be. In some cases, the suppression of string density can be quite dramatic (as we saw for small-scale global string simulations), although the string velocities always remain relativistic. For the most part, however, the density of a cosmic string network, whether local or global is only affected slightly by radiation backreaction effects.
Despite the many virtues of the VOS model, we are aware, of course, that the small number of available degrees of freedom means that this model is unable to provide a proper description of the small-scale properties of the network; these are important in a number of cosmological scenarios (and sometimes even crucial). Nevertheless, we believe that the phenomenological parameter $`k`$ does encode some important small-scale structure effects, though clearly a more detailed analytic and numerical study is still required. A number of possible approaches to the problem of string small-scale structure have been suggested in the literature , and our own analysis using Carter’s elastic string model will be discussed in a forthcoming publication .
###### Acknowledgements.
We would like to thank Pedro Avelino, Brandon Carter, Jonathan Moore, Levon Pogosian, Tanmay Vachaspati and Proty Wu for useful conversations. C.M. also acknowledges discussions with a number of participants in the EC Summer School ‘Multi-fractals—Mathematics and Applications’, held at the Isaac Newton Institute. C.M. is funded by FCT (Portugal) under ‘Programa PRAXIS XXI’ (grant no. PRAXIS XXI/BPD/11769/97). This work was performed on COSMOS, the Origin2000 owned by the UK Computational Cosmology Consortium, supported by Silicon Graphics/Cray Research, HEFCE and PPARC. |
warning/0003/math0003220.html | ar5iv | text | # Special Lagrangian submanifolds and Algebraic Complexity one Torus Actions
## 1 Introduction
In our previous paper we have shown how to use torus actions on non-compact Calabi-Yau manifolds to construct SLag submanifolds. In all the ensuing discussions by SLag submanifolds we mean the following: Let $`(M^{2n},\omega ,\phi )`$ be a complex manifold with a Kahler form $`\omega `$ and a non-vanishing holomorphic $`(n,0)`$-form $`\phi `$. Then a submanifold $`L^n`$ of $`M`$ is SLag if it satisfies $`\omega |_L=0,Im\phi |_L=0`$. We refer the reader to and for a discussion of SLag submanifolds. In the first part of this paper we develop one general setup, which generalizes all examples in :
Consider a compact algebraic manifold $`M^{2n}`$ with an algebraic $`(n1)`$-torus action (this is an algebraic complexity one space). We will show (see Theorem 1) that there is a meromorphic section $`\sigma `$ of the canonical bundle of $`M`$, which is invariant under the $`T`$-action. Any such $`\sigma `$ defines a divisor $`D`$ and on the complement $`M^{}=MD`$ $`\sigma `$ gives a trivialization of the canonical bundle. Let $`H^1(M^{},)=0`$. Take any $`T`$-invariant Kahler form $`\omega `$ on $`M^{}`$. From we deduce that $`M^{}`$ has a SLag fibration with respect to $`\omega `$. We will investigate how the fibers compactify in $`M`$. We treat 2 cases : If $`\omega `$ is the induced metric from $`M`$, then Lemma 1 of Section 2.1 covers “most” fibers; For $`\omega `$ a Ricci-flat metric on $`M^{}`$ we study one family of examples in Section 2.2.
We will also give examples of SLag fibrations on $`M^{}`$, including some cases there $`H^1(M^{},)0`$.
In the second part of this paper we investigate Calabi-Yau hypersurfaces in $`M`$. We assume that the anti-canonical bundle of $`M`$ is ample and has a $`T`$-invariant holomorphic section $`\eta `$ (with a corresponding meromorphic section $`\sigma =\eta ^1`$ of the canonical bundle). We also assume that near the smooth part $`D^{}`$ of the zero set $`D`$ $`\eta `$ is transversal to zero. This gives a natural trivialization $`\phi ^{}`$ of the canonical bundle of $`D^{}`$. Consider the induced metric $`\omega `$ from $`M`$ on $`D^{}`$. We assume that the $`T`$-action on each connected component $`D_0`$ of $`D^{}`$ has a finite stabilizer. For $`D_0`$ one can choose an angle $`\theta _0`$ s.t. orbits of the $`T`$-action are SLag submanifolds of $`(D_0,\omega ,e^{i\theta _0}\phi ^{})`$. Consider transversal sections $`\eta _j`$, which converge to $`\eta `$. Their zero sets $`D_j`$ are smooth Calabi-Yau hypersurfaces. Near compact, $`T`$-invariant subsets $`N`$ of $`D_0`$ the manifolds $`D_j`$ converge to $`N`$ as $`\eta _j\eta `$. We will show (see Theorem 3) that one can perturb our SLag fibration on $`N`$ to SLag fibrations on the corresponding subsets of $`D_j`$ (all fibrations are SLag for metrics, induced from $`M`$).
We will see that this construction applies for $`M=P^4`$ and $`D_j`$ being quintics; and for $`M=G(2,4)`$ (the Grassmanian of 2-planes in $`^4`$) and $`D_j`$ being Calabi-Yau hypersurfaces in $`M`$.
We will also show that this can generalized from hypersurfaces to complete intersection Calabi-Yau submanifolds of a Kahler manifold $`M`$ with a higher complexity torus action (see Corollary 1 and Section 3.4).
Acknowledgments : This paper is a part of author’s work towards his Ph.D. at MIT. The author is grateful to his advisor, Tom Mrowka, for continuing support; and to Shing-Tung Yau for a helpful discussion.
Research is partially supported by an NSERC PGS B Award
## 2 Algebraic Complexity one Spaces
### 2.1 Fibration on M’ and it’s compactification
Let $`M^{2n}`$ be a compact algebraic manifold with an effective algebraic $`T^s`$-action. The fact the the action is algebraic is equivalent to saying that there is a very ample line bundle $`L`$ on $`M`$, and $`T`$ acts on it’s total space. We can use $`L`$ to construct a meromorphic section of the canonical bundle $`K`$ of $`M`$:
###### Theorem 1
For $`k`$ large enough there is a holomorphic section $`\sigma _1`$ of $`KL^k`$ and $`\sigma _2`$ of $`L^k`$ s.t. the section $`\sigma =\sigma _1\sigma _2^1`$ is a meromorphic, $`T`$-invariant section of $`K`$.
Proof: Consider the representation $`\rho `$ of $`T`$ on $`H^0(L,M)`$. We split it into a direct sum of irreducible representations, i.e. we find sections $`\alpha _1,\mathrm{},\alpha _d`$ of $`L`$ s.t. $`T`$ acts on $`span(\alpha _i)`$ with a character $`\xi _i`$. W.l.o.g we can assume that $`\xi _1`$ is a trivial character (otherwise we divide the action of $`T`$ on $`L`$ by $`\xi _1`$).
We can view $`\xi _j`$ as living in the character lattice of the dual Lie algebra $`𝒢^{}`$ of $`T`$. The condition that the $`T`$-action is effective on $`M`$ implies the $`\xi _i`$ $``$-span the character lattice of $`𝒢^{}`$. Indeed suppose $`\xi _i`$ do not span the character lattice of $`𝒢^{}`$. Then we can find an element $`e1`$ of $`T`$ s.t. $`\xi _i(e)=1`$, i.e. $`e`$ acts trivially on $`H^0(L,M)`$. But then $`e`$ acts trivially on $`M`$. Indeed suppose $`mM`$ and $`e(m)=m^{}m`$. Since $`L`$ is very ample, there is a section $`\beta `$ of $`L`$ s.t. $`\beta (m)=0,\beta (m^{})0`$. But then $`e`$ doesn’t act trivially on $`\beta `$\- a contradiction.
For $`l`$ large enough $`L^lK`$ has a holomorphic section. So $`H^0(L^lK,M)`$ is nontrivial. We can find a section $`\lambda `$ of $`L^lK`$ s.t. $`V=span(\lambda )`$ is $`T`$-invariant and $`T`$ acts on $`V`$ with a character $`\varphi `$. We can write $`\varphi =\mathrm{\Sigma }a_i\xi _i,a_i`$. We will consider only $`a_i0`$ and we divide them into subsets $`a_i^+`$ and $`a_i^{}`$ of positive and negative values. We will have a section $`\sigma _1^{}=\lambda \alpha _i^{a_i^{}}`$ of $`L^{k_1}K`$ and a section $`\sigma _2^{}=\alpha _i^{a_i^+}`$ of $`L^{k_2}`$. The $`T`$-actions leave $`V_i=span(\sigma _i^{})`$ invariant and $`T`$ acts on $`V_1`$ and $`V_2`$ with the same character. Also $`T`$ preserves $`\alpha _1`$. So tensoring $`\sigma _1^{}`$ and $`\sigma _2^{}`$ with powers of $`\alpha _1`$ we deduce that for $`k`$ large enough we can find sections $`\sigma _1`$ of $`L^kK`$ and $`\sigma _2`$ of $`L^k`$ s.t. $`T`$ acts on $`\sigma _i`$ with the same character. Q.E.D.
From now on we assume that $`s=n1`$, i.e. $`M`$ is an algebraic complexity one space. We refer the reader to and for discussion of complexity one spaces. Let $`\sigma `$ be any $`T`$-invariant meromorphic section of the canonical bundle of $`M`$. It defines a divisor $`D`$. Let $`M^{}=MD`$. Take $`\omega `$ to be any $`T`$-invariant Kahler form on $`M^{}`$ s.t. the $`T`$-action is Hamiltonian. We assume that $`H^1(M^{},)=0`$. From we get the following result (a similar result was also obtained in ):
###### Theorem 2
() $`M^{}`$ has a calibrated fibration $`\alpha `$ over an open subset of $`^n`$. Moreover for a generic point $`p`$ (outside of a countable union of $`(n2)`$-planes in $`^n`$) the fiber $`\alpha ^1(p)`$ is a smooth SLag submanifold. Connected components of smooth fibers are diffeomorphic to $`T^{n1}\times `$. If $`mM^{}`$ is a regular point of the torus action (i.e. the differential of the action is injective at $`m`$), then $`m`$ is a smooth point of the fiber $`L_m`$ through $`m`$ and the fiber $`L_m`$ is a SLag submanifold near $`m`$. Singular fibers have singularities of codimension at least 2 and near singular points they are diffeomorphic to a product of a cone with a Euclidean ball.
The fibration is defined as follows: Let $`e_1,\mathrm{},e_{n1}`$ be a basis for the Lie Algebra of $`T^{n1}`$ and $`X_1,\mathrm{},X_{n1}`$ be the corresponding flow vector fields on $`M`$. Let $`\sigma ^{}=i_{X_1}\mathrm{}i_{X_{n1}}\sigma `$ be the $`(1,0)`$-form on $`M^{}`$, obtained by contraction of $`\sigma `$ by vector fields $`X_1,\mathrm{},X_{n1}`$. Then from we know that $`\sigma ^{}`$ is a holomorphic $`(1,0)`$-form on $`M^{}`$ and $`\sigma =df`$ for a $`T`$-invariant, holomorphic function $`f`$ on $`M^{}`$. Let $`\mu `$ be a moment map for the $`T`$-action on $`M^{}`$. Then the fibers of the SLag fibration are given as level sets of the function $`\alpha =(\mu ,Imf)`$.
The closure of a connected component $`L^{}`$ of any smooth SLag fiber $`L`$ on $`M^{}`$ must intersect $`D`$ (for otherwise $`L^{}`$ would be compact, but it is diffeomorphic to $`T^{n1}\times `$). In order to further understand how the fibers of our SLag fibration on $`M^{}`$ compactify in $`M`$ we will have to make assumptions on the Kahler form $`\omega `$. First we treat the case then $`\omega `$ is induced from a Kahler form on $`M`$.
We decompose $`D`$ into $`D=D_+D_{}`$ corresponding to the meromorphic and holomorphic parts of $`\sigma `$. Consider the function $`f`$ near a point $`d`$ in $`DD_{}`$. Then the differential $`df=\sigma `$ is bounded, and so will be $`f`$. By Riemann extension theorem, $`f`$ extends to a holomorphic function near $`m`$. So $`f`$ extends to a holomorphic function on $`MD_{}`$. Also $`\mu `$ is smooth on $`MD_{}`$ and so our fibration extends to a fibration by level sets of $`(\mu ,Imf)`$ on $`MD_{}`$. Finally we wish to understand how the fibers compactify near $`D_{}`$. We have the following Lemma:
###### Lemma 1
Let $`\nu `$ be a value of the moment map s.t. $`T`$ acts freely on $`D_\nu =\mu ^1(\nu )D_{}`$. Consider a fiber $`L_t=(\mu =\nu ,Imf=t)`$ of the fibration over $`MD_{}`$ for some $`t`$. Then the boundary of $`L_t`$ in $`M`$ is $`D_\nu `$
Proof: Obviously the boundary of $`L_t`$ is contained in $`D_\nu `$. To prove the other inclusion we first note that the $`T^{n1}`$ action on $`M`$ induces the complex torus $`T^c=(^{})^{n1}`$ action on $`M`$. Indeed the flow vector fields $`X_v`$ give rise to the vector fields $`JX_v`$. The flow of $`X_u`$ preserves $`X_v`$ and commutes with $`J`$, hence it preserves $`JX_v`$, i.e. $`[X_u,JX_v]=0`$. Finally the vector field $`[X_uiJX_u,X_viJX_v]`$ is of type $`(1,0)`$, but it’s equal to $`[JX_u,JX_v]`$. So $`[JX_u,JX_v]=0`$ and all the vector fields $`X_u,JX_v`$ commute, and this induces the $`T^c`$-action on $`M`$ (with the real torus $`T^{n1}`$ being the product $`S^1\times \mathrm{}\times S^1`$ of unit circles in $`T^c`$). We have
$$_{JX_v}\sigma =d(i_{JX_v}\sigma )=id(i_{X_v}\sigma )=i_{X_v}\sigma =0$$
So the $`T^c`$-action preserves $`\sigma `$, and so it preserves $`D,D_+,D_{}`$. Also since $`f`$ is holomorphic then $`JX_v(f)=iX_v(f)=0`$, i.e. $`T^c`$ preserves $`f`$.
Let now $`dD_\nu `$. Since $`T`$ acts freely on $`D_\nu `$ then in particular the differential of the $`T`$-action on $`d`$ is injective. So the orbit $`L^c`$ of $`T^c`$-action on $`d`$ is a smooth complex submanifold of $`D_{}`$ of complex dimension $`n1`$. Since $`T^c`$-action preserves $`D_{}`$ then it is clear that near $`d`$ $`D_{}`$ coincides with $`L^c`$. Also the differential of $`\mu |_{L_c}`$ is surjective, so near $`d`$ $`D_\nu `$ coincides with the orbit $`L`$ of the $`T^{n1}`$-action on $`d`$.
Consider the level set $`\mathrm{\Sigma }_\nu `$ of $`\mu `$ on $`M`$. Then near $`D_\nu `$ $`\mathrm{\Sigma }_\nu `$ is smooth. Also $`T^{n1}`$ acts freely on $`D_\nu `$, hence it acts freely on $`\mathrm{\Sigma }_\nu `$ near $`D_\nu `$. We consider the symplectic reduction $`M_{red}`$ of $`\mathrm{\Sigma }_\nu `$ in a $`T`$-invariant neighbourhood $`U`$ of $`D_\nu `$. Then $`M_{red}`$ will be a smooth complex 1-dimensional manifold. Let $`\pi :\mathrm{\Sigma }_\nu M_{red}`$ be the quotient map. Then $`\pi (L)=d^{}`$ is one point in $`M_{red}`$ (here $`L`$ is the $`T^{n1}`$-orbit of $`d`$). Also in a sufficiently small neighbourhood $`U^{}`$ of $`d^{}`$ in $`M_{red}`$, $`\pi ^1(U^{}d^{})D_{}=\mathrm{}`$.
The function $`f`$ descends to a holomorphic function on $`U^{}`$. Suppose that $`f`$ has a singularity in $`d^{}`$. Then the image of $`f`$ covers a neighbourhood of $`\mathrm{}`$ in $``$. Hence the image of $`Imf`$ assumes all real values on $`U^{}d`$, in particular it assumes the value $`t`$. Hence the SLag fiber $`L_t`$ in $`MD_{}`$ intersects the neighbourhood $`\pi ^1(U^{})`$. From this we easily deduce that $`d`$ is in the closure of $`L_t`$ in $`M`$.
So we need to prove that $`f`$ has a singularity at $`d^{}`$. Suppose not. Then $`|f|`$ is bounded by some constant $`C`$. So $`|f|<C`$ on $`\mathrm{\Sigma }_\nu D_{}`$ near $`d`$. Now the $`T^c`$-action preserves $`f`$. The orbit of $`\mathrm{\Sigma }_\nu `$ under the $`T^c`$-action fills a neighbourhood $`W`$ of $`d`$ in $`M`$. On $`WD_{}`$ we have $`|f|<C`$. By Riemann extension theorem $`f`$ extends to a holomorphic function on $`W`$. So the norm of it’s differential $`|df|<C^{}`$ for some other constant $`C^{}`$, i.e. $`|\sigma ^{}|<C^{}`$. But $`\sigma ^{}=i_{X_1}\mathrm{}i_{X_{n1}}\sigma `$. We choose a local section $`\phi `$ of the canonical bundle of $`M`$ near $`d`$ s.t. $`|\phi |=1`$. Then $`\sigma =g\phi `$ and $`|g|`$ is not bounded since $`\sigma `$ has a singularity at $`d`$. But the vector fields $`X_i`$ are linearly independent and $`\omega `$-orthogonal to each other. Hence one easily deduces that $`|i_{X_1}\mathrm{}i_{X_{n1}}\phi |`$ is uniformly bounded from below - a contradiction. So $`f`$ is singular at $`d^{}`$ and we are done. Q.E.D.
Remark 1: The conclusions of Lemma 1 apply only for those values of $`\mu `$ on $`D_{}`$, on whose level sets $`T^{n1}`$ acts freely. In example 1 in Section 2.2 we’ll have that for all non-regular values $`\nu `$ on $`D_{}`$ the level set $`D_\nu `$ will not intersect the closure of any SLag fiber in $`MD_{}`$.
Ricci-Flat metrics Another interesting class of metrics on $`M^{}`$ are $`T`$-invariant, complete Ricci-flat metrics. We refer the reader to , for some existence results. If $`H^1(M^{},)=0`$ then from Theorem 2 we get a SLag fibration on $`M^{}`$. In this case the compactification of the fibers in $`M`$ is quite different from the case of induced metrics. We will consider one case of this setup in example 1 in Section 2.2.
### 2.2 Examples
1) K(N) () We will demonstrate a family of examples for the construction in Section 2.1. For those examples we will write down Ricci-flat metrics on $`M^{}`$ and we show how fibers of the SLag fibration on $`M^{}`$ compactify in $`M`$. We will also write down some metrics on $`M^{}`$, induced from metrics on $`M`$. We will show that for those metrics the conclusions of Lemma 1 do not apply for non-regular values of the moment map $`\mu `$ on $`D_{}`$ and thus Lemma 1 is sharp.
Let $`N^{2n}`$ be a toric Kahler-Einstein manifold with positive scalar curvature $`t`$. Consider $`K(N)`$ to be the total space of it’s canonical bundle. Then $`K(N)`$ is a Calabi-Yau manifold with a natural holomorphic $`(n+1,0)`$-form $`\phi `$. We refer to for definitions and properties of all relevant structures on $`K(N)`$. We will however give a Calabi construction of Kahler metrics on $`K(N)`$, since we will need them in further discussion. Let $`r^2:K(N)_+`$ be the square of the length of elements in $`K(N)`$. Take any positive function $`u`$ on $`_+`$ s.t. $`u^{}>0`$. Define the metric $`\omega _u`$ on $`K(N)`$ as follows: Let $`\omega `$ be the K-E metric on $`N`$. The connection on $`K(N)`$ induces a horizontal distribution for the projection $`\pi :K(N)N`$. We define the horizontal and the vertical distributions to be orthogonal. On the horizontal distribution we define the metric to be $`u(r^2)\pi ^{}(\omega )`$. On the vertical distribution the metric is $`t^1u^{}(r^2)\omega ^{}`$. Here $`\omega ^{}`$ is the induced metric on the linear fibers. The K-E condition ensures that $`\omega _u`$ is closed. If we choose $`u(r^2)=(tr^2+l)^{1/n+1}`$ for some positive constant $`l`$, then $`\omega _u`$ is Ricci-flat.
Let now $`W`$ be a 2-plane bundle over $`N`$, obtained by adding a trivial bundle to $`K(N)`$. Let $`M`$ be the projectivization of $`W`$. The $`T^n`$-action on $`N`$ induces an action on $`W`$ and on $`M`$. $`K(N)`$ naturally sits inside of $`M`$ and $`M`$ is obtained by adding to $`K(N)`$ a copy $`N_{\mathrm{}}`$ of $`N`$ at infinity. The holomorphic $`(n+1,0)`$-form $`\phi `$ on $`K(N)`$ becomes a meromorphic $`T^n`$-invariant section of the canonical bundle of $`M`$, which defines a divisor $`D=D_{}=N_{\mathrm{}}`$. $`\phi `$ has singularity of order 2 at $`D`$ and we get $`M^{}=MD=K(N)`$.
In we considered the SLag fibration on $`K(N)`$ with respect to the Ricci-flat metric on $`K(N)`$ (i.e. then $`u(r^2)=(tr^2+l)^{1/n+1}`$). We have shown that all fibers are asymptotic at infinity to a certain conical fiber $`L_0`$. Moreover the boundary in $`M`$ of each fiber is a certain minimal Lagrangian submanifold $`L_{\mathrm{}}`$ of $`N_{\mathrm{}}`$.
We wish to consider a metric on $`K(N)`$, induced from a metric on $`M`$. This can be done by choosing a different function $`u`$ : Consider a function $`w(x)=x^2u^{}(1/x)`$. This is a smooth positive function for $`x>0`$. Suppose it extends to a smooth positive function $`w`$ for $`x0`$. So in particular $`u^{}(y)`$ is asymptotic to $`1/y^2`$ at infinity and $`u^{}`$ integrates on $`(0,\mathrm{})`$. Thus we have a limit $`u_{\mathrm{}}>0`$ of $`u`$ at infinity. One can easily show that the metric $`\omega _u`$ compactifies to a smooth $`T^n`$-invariant metric $`\omega _u`$ on $`M`$.
In we have computed a specific moment map $`\mu `$ for the $`T^n`$-action on $`N`$. Moreover we have shown that the moment map of the action on $`K(N)`$ was $`\mu ^{}=u(r^2)\pi ^1(\mu )`$. This of course compactifies to a moment map on $`M`$. Let $`\nu `$ be a regular value of the moment map $`\mu `$ on $`N`$. Then $`\nu ^{}=u_{\mathrm{}}\nu `$ is a regular value of the moment map on $`D=N_{\mathrm{}}`$. Lemma 1 tells us that the level set $`D_\nu ^{}`$ of $`\mu ^{}`$ on $`D`$ is a boundary in $`M`$ of any SLag submanifold of the form $`(\mu ^{}=\nu ^{},Imf=c)`$ on $`K(N)`$. Let $`\nu `$ now be a value on the boundary of the moment polytope of $`N`$. Then we know from that $`0`$ is in the open part of the moment polytope. Hence by convexity of the polytope we deduce that a multiple $`\lambda \mu `$ is not in the moment polytope for any $`\lambda >1`$. Consider now the value $`\nu ^{}=u_{\mathrm{}}\nu `$ of $`\mu ^{}`$ on $`D`$. Then $`\nu ^{}`$ is not attained by $`\mu ^{}`$ on $`K(N)`$. Hence the level set $`D_\nu ^{}`$ does not intersect the closure of any element of our SLag fibration on $`K(N)`$. Hence the statement of Lemma 1 is sharp.
2) Toric Varieties: In section 2.1 we assumed that $`H^1(M^{},)=0`$. We will now demonstrate a class of examples there this condition doesn’t hold but we nevertheless have a SLag fibration on $`M^{}`$.
Let $`M`$ be toric i.e. $`T^{n1}T^n`$ and there is a $`T^n`$-action on $`M`$. From Theorem 1 we know that there is a section $`\sigma `$ of $`K(N)`$, which is invariant under $`T^n`$-action. First we prove the following Lemma, which will also be useful later:
###### Lemma 2
Let $`N^{2n}`$ be a connected complex manifold with a (non-zero) holomorphic $`(n,0)`$ form $`\sigma `$ and a Kahler form $`\omega `$. Suppose that we have a holomorphic $`T^c=(^{})^n`$-action on $`N`$ s.t. the action of the real torus $`T^n`$ is effective, preserves $`\sigma `$ and Hamiltonian with respect to $`\omega `$. Then the $`T^c`$-action is free and $`N`$ is biholomorphic to $`T^c`$ under the action. $`\sigma `$ is non-vanishing and equal to a constant multiple of the form $`dz_i/z_i`$ on $`T^c`$. Moreover for a choice of $`\theta `$ we have that orbits of the $`T^n`$-action give a SLag fibration on $`(N,\omega ,e^{i\theta }\sigma )`$.
Proof: Let $`v_1,\mathrm{},v_n`$ be a basis for the Lie algebra of $`T^n`$ and let $`X_i`$ be the flow vector field of $`v_i`$ on $`N`$. Let $`g=i_{X_1}\mathrm{}i_{X_n}\sigma `$. As we saw in , $`dg=0`$, i.e. $`g`$ is locally constant, and since $`N`$ is connected $`g`$ is constant. Since the action is effective, the differential of the action is injective at some point $`p`$, in which $`\sigma (p)0`$. The vectors $`X_1(p),\mathrm{},X_n(p)`$ span a Lagrangian plane, so $`g=g(p)=g_00`$. So $`\sigma `$ is non-vanishing and the differential of the $`T^n`$-action is everywhere injective. We can choose $`\theta `$ s.t. $`e^{i\theta }g_0`$ is real, and we get that orbits of the $`T^n`$-action are SLag submanifolds of $`(N,\omega ,e^{i\theta }\sigma )`$.
The differential of the $`T^c`$-action is an isomorphism everywhere. Orbits of the $`T^c`$-action are open, and since $`N`$ is connected, there is one orbit. So if $`H`$ is a stabilizer of some $`pN`$ under $`T^c`$-action, then $`H`$ is the stabilizer of $`N`$ under this action. We wish to prove $`H`$ is trivial. Let $`1hH`$. Then $`h=(z_1,\mathrm{},z_n)`$ with $`z_j=a_j+ib_j0`$. Since the $`T^n`$-action is effective, not all $`b_j`$ are equal to $`0`$. The action of $`h`$ is given by a time 1 flow of the vector field $`X=\mathrm{\Sigma }a_jX_j+b_jJX_j`$. Let $`\mu `$ be the moment map of the $`T^n`$-action. Consider a function $`\kappa =\mathrm{\Sigma }b_j\mu (v_j)`$. Then one easily shows that $`X(\kappa )>0`$, hence the time 1 flow of $`X`$ cannot return to the same point- a contradiction. Hence $`H`$ is trivial, the action is free and $`T^c`$ is biholomorphic to an $`N`$ under the action.
In the proof of Lemma 1 we showed that $`\sigma `$ is invariant under the $`T^c`$-action since it is invariant under the $`T^n`$-action. Consider the form $`\sigma _0=dz_j/z_j`$ on $`T^c`$. Then both the pullback of $`\sigma `$ under the action and $`\sigma _0`$ are $`T^c`$-invariant. Hence one is a constant multiple of the other and we are done. Q.E.D.
We return now to the toric manifold $`M`$. Consider $`N=MD_{}`$ with a holomorphic $`(n,0)`$-form $`\sigma `$ on it. We have a $`T^c`$-action on $`M`$, and it leaves $`MD_{}`$ invariant, hence it induces a $`T^c`$-action on $`N`$. From Lemma 2 we deduce that $`\sigma `$ is non-vanishing, i.e. $`D_+=\mathrm{}`$. Also for a choice of $`\theta `$ orbits of the $`T^n`$-action give a SLag fibration on $`(MD_{},\omega ,e^{i\theta }\sigma )`$.
We remark that $`M^{}`$ coincides with the regular points of the $`T`$-action. Indeed points in $`MD_{}`$ are regular points of $`T^n`$-action. Also every point $`dD_{}`$ is a singular point of the $`T`$-action (since otherwise the $`T^c`$-orbit of $`d`$ is open in $`M`$, and it cannot be contained in $`D_{}`$).
3) The Grassmanian G(2,4): This is yet another example, for which $`H^1(M^{},)0`$ but we can construct a SLag fibration on $`M^{}`$. Consider the Grassmanian $`G(2,4)`$ of complex 2-planes in $`^4`$, which we identify with a quadric hypersurface $`M=(z_1z_2+z_3z_4+z_5z_6=0)`$ in $`P^5`$ (see and ). The fourth power $`(\gamma ^{})^4`$ of the hyperplane bundle on $`P^5`$ restricted to $`M`$ is the anti-canonical bundle $`\overline{K}(M)`$. Thus polynomials of degree 4 give rise to holomorphic sections of $`\overline{K}(M)`$. There is a complex 3-torus $`T^c=(^{})^3`$-action on $`M`$ given by
$$(\lambda _1,\mathrm{},\lambda _3)(z_1,\mathrm{},z_6)=(\lambda _1z_1,\lambda _1^1z_2,\mathrm{},\lambda _3z_5,\lambda _3^1z_6)$$
This action of course contains the action of the real torus $`T^3(^{})^3`$. We would like to find a homogeneous polynomial $`p`$ of degree 4 on $`^6`$ s.t. $`p`$ defines a $`T^c`$-invariants section of $`\overline{K}(M)`$. To do that we will carefully set up the $`T`$-equivariant identifications between various bundles we use.
First on $`P^5`$ there is a constant bundle $`C=^6`$. It contains a universal bundle $`\gamma `$ as a sub-bundle. The tangent bundle to $`P^5`$ is isomorphic to $`\gamma ^{}(C/\gamma )`$. There is a short exact sequence
$$\gamma \gamma ^{}C\gamma ^{}(C/\gamma )\gamma ^{}$$
Taking the canonical bundles of the elements in the sequence we get $`K(P^5)K(C)\gamma ^6`$. Of course the bundle $`K(C)`$ is trivial, but the isomorphism $`K(P^5)\gamma ^5`$ is not $`GL(6,)`$-equivariant, but it is $`SL(6,)`$-equivariant. Since our torus $`T^c`$ is in $`SL(6,)`$, we are fine.
Next consider the quadric $`M`$, which can be viewed as a zero set of a section $`\eta `$ of $`(\gamma ^{})^2`$. The canonical bundle of $`M`$ is isomorphic to $`NK(P^5)`$. Here $`N`$ is the normal bundle to $`M`$ in $`P^5`$ and the isomorphism is given by $`v\phi i_v\phi |_M`$. Also $`N`$ is isomorphic to $`(\gamma ^{})^2`$ with an isomorphism given by $`v_v\eta `$. Since $`\eta `$ is $`T^c`$-invariant, this isomorphism is $`T^c`$-equivariant. So overall we get that $`K(M)NK(P^5)(\gamma ^{})^2\gamma ^6\gamma ^4`$. Dually $`\overline{K}(M)(\gamma ^{})^4`$ and this isomorphism is $`T^c`$-equivariant.
Consider a polynomial $`p=(z_1z_2)^2+(z_3z_4)^2(z_5z_6)^2`$. Then it defines a $`T^c`$-invariant holomorphic section of $`\overline{K}(M)`$, and dually it defines a $`T^c`$-invariant meromorphic section $`\sigma `$ of $`K(M)`$, which defines a divisor $`D=D_{}`$, which is the zero set of $`p`$. Let $`M^{}=MD`$ and pick a $`T^3`$-invariant Kahler metric $`\omega `$ on $`M^{}`$ s.t. the action is Hamiltonian (e.g one can pick the induced metric from $`P^5`$). We would like to get a SLag fibration for $`(M^{},\omega ,\sigma )`$. We will see that $`H^1(M^{},)0`$, but we can still do that if we replace $`\sigma `$ by a multiple $`e^{i\theta }\sigma `$.
Let $`e_1,e_2,e_3`$ be a basis for Lie algebra of $`T^3`$ and $`X_i`$ be corresponding flow vector fields on $`M`$. Let $`\sigma ^{}=i_{X_1}\mathrm{}i_{X_3}\sigma `$. Then from we know that $`\sigma ^{}`$ is a holomorphic $`(1,0)`$-form on $`M^{}`$. Moreover suppose that for some choice of $`\theta `$ we have $`Im(e^{i\theta }\sigma ^{})`$ is an exact 1-form on $`M^{}`$, i.e. it is equal to $`dh`$ for some function $`h`$. Then we saw in that $`(M^{},\omega ,e^{i\theta }\sigma )`$ has a SLag fibration on it with fibers given as level sets of $`(\mu ,h)`$ (here $`\mu `$ is a moment map for the $`T^3`$-action on $`M^{}`$).
To find this $`\theta `$ consider the following map $`\beta :M^{}P^1`$ given by $`\beta (z_1,\mathrm{},z_6)=(z_1z_2,z_3z_4)`$. We will show that $`\beta `$ is well-defined and compute it’s image. First we claim that on $`M^{}`$ we can’t have $`z_1z_2=0`$ or $`z_3z_4=0`$. Indeed suppose $`z_1z_2=0`$. Then $`z_3z_4=z_5z_6`$, hence $`p(z_1,\mathrm{},z_6)=(z_1z_2)^2+(z_3z_4)^2(z_5z_6)^2=0`$ and we are on $`D`$ and not $`M^{}`$. So we deduce that $`\beta `$ is well-defined and it’s image lies in $`C^{}=P^1((1,0)(0,1))`$.
Next we show that $`\beta `$ is surjective onto $`C^{}`$. Indeed $`\beta `$ has local inverses $`\alpha (a,b)=(a,1,b,1,ab,1)`$ into $`M^{}`$ (here $`(a,b)`$ is in $`C^{}`$). $`p(a,1,b,1,ab,1)=2ab0`$ and so indeed $`(a,1,b,1,ab,1)M^{}`$.
The map $`\beta `$ is $`T^c`$-invariant. Take now $`(a,b)C^{}`$ s.t. $`a+b0`$ i.e. $`(a,b)(1,1)`$. If $`(z_1,\mathrm{},z_6)\beta ^1((a,b))`$ then all $`z_i0`$. One easily deduces that $`\beta ^1((a,b))`$ is in fact the orbit of $`\alpha (a,b)`$ under $`T^c`$-action. Moreover because of the local inverse $`\alpha `$ the differential of $`\beta `$ is surjective at all points of $`\beta ^1((a,b))`$. So we get a principal fiber bundle $`M^{\prime \prime }=\beta ^1(C^{}(1,1))`$ over $`C^{}(1,1)`$ with the fiber being $`T^c`$.
Now the form $`\sigma ^{}`$ is $`T^c`$ -invariant. Moreover the tangent space to the orbits of $`T^c`$ is in the kernel of $`\sigma ^{}`$. From this we easily deduce that there is a holomorphic $`(1,0)`$-form $`\sigma ^{\prime \prime }`$ on $`C^{}(1,1)`$ s.t. $`\beta ^{}(\sigma ^{\prime \prime })=\sigma ^{}`$. We claim that $`\sigma ^{\prime \prime }`$ compactifies to a holomorphic $`(1,0)`$-form on $`C^{}`$. Indeed we have a local inverse $`\alpha `$ near $`(1,1)`$ and $`\alpha ^{}(\sigma ^{})=\sigma ^{\prime \prime }`$.
Now $`C^{}`$ is $`^{}`$, so $`H^1(C^{},)=`$. So one can find an angle $`\theta `$ s.t. $`Im(e^{i\theta }\sigma ^{\prime \prime })`$ is an exact 1-form on $`C^{}`$. We claim that also $`Im(e^{i\theta }\sigma ^{})`$ is an exact 1-form on $`M^{}`$. Indeed take a loop $`\gamma `$ on $`M^{}`$. Then one easily sees that
$$_\gamma Im(e^{i\theta }\sigma ^{})=_{\beta \gamma }Im(e^{i\theta }\sigma ^{\prime \prime })=0$$
and we are done.
## 3 Calabi-Yau hypersurfaces near a large complex limit Hypersurface
### 3.1 Deformation of SLag fibrations
In this section we wish to study Calabi-Yau hypersurfaces in an algebraic complexity one space $`M`$. We will assume that the anti-canonical bundle $`\overline{K}`$ of $`M`$ is ample and has a holomorphic section $`\eta `$, invariant under the $`T`$-action. Let $`D^{}`$ be a smooth part of the divisor $`D`$ of $`\eta `$. We will assume that $`\eta `$ is transversal to $`0`$ at $`D^{}`$. This defines a natural $`T^{n1}`$-invariant trivialization $`\phi ^{}`$ of the canonical bundle of $`D^{}`$. Indeed the canonical bundle of $`D^{}`$ is naturally isomorphic to $`KN`$. Here $`N`$ is a normal bundle to $`D^{}`$ and the isomorphism is explicitly given by $`\phi vi_v\phi |_D`$. Also the normal bundle $`N`$ is naturally isomorphic to $`\overline{K}`$, with isomorphism given by $`v_v\eta `$.
Take a connected component $`D_0`$ of $`D^{}`$. We assume that the $`T^{n1}`$-action on $`D_0`$ has a finite stabilizer $`H`$. We have an action of the quotients $`T_0=T^{n1}/H`$ and of $`T_0^c=T^c/H`$ on $`D_0`$. But $`T_0`$ is isomorphic to $`T^{n1}`$. Also the isomorphism between $`T_0`$ and $`T^{n1}`$ induces a biholomorphic isomorphism between $`T_0^c`$ and $`T^c`$. From these isomorphisms we get a holomorphic $`T^c`$-action on $`D_0`$, and this action preserves $`\phi ^{}`$. Also the corresponding $`T^{n1}`$-action is effective and it is Hamiltonian with respect to the metric $`\omega `$, induced from $`M`$. From Lemma 2 we deduce that the $`T^c`$-action is free and $`D_0`$ is biholomorphic to $`T^c`$. Also we can choose an angle $`\theta _0`$ s.t. orbits of the $`T^{n1}`$-action give a SLag fibration on $`D_0`$ for $`(\omega ,e^{i\theta _0}\phi ^{})`$. We have the following theorem :
###### Theorem 3
Let $`\eta ,D,D^{},D_0`$ be as above. Take $`N`$ to be any $`T`$-invariant compact subset of $`D_0`$ and $`U`$ a tubular neighbourhood of $`N`$ in $`M`$. Then there is a neighbourhood $`U_N`$ of $`\eta `$ in $`H^0(\overline{K}(M),M)`$ s.t. for $`\eta _pU_N`$ there is a SLag fibration of the neighbourhood $`U_p`$ in $`D_p`$ with respect to $`(\omega _p,e^{i\theta _p}\phi _p^{})`$.
Here $`D_p`$ is the zero set of $`\eta _p`$, $`U_p`$ contains $`D_pU`$, $`\omega _p`$ is the induced metric on $`D_p`$ from $`M`$, $`\phi _p^{}`$ is the natural trivialization of the canonical bundle of $`D_p`$, given by the section $`\eta _p`$ and $`\eta _p\theta _p`$ is a smooth function.
Proof: Choose a compact, $`T`$-invariant $`ND_0`$ and pick a compact, $`T`$-invariant neighbourhood $`N^{}`$ of $`N`$. Let $`U^{}`$ be a tubular neighbourhood of $`N^{}`$ in $`M`$.
For any $`\eta _pH^0(\overline{K}(M),M)`$ let $`D_p`$ be it’s zero set. The part of $`D_p`$ in $`U^{}`$ has a trivialization $`\phi _p^{}`$ of it’s canonical bundle, induced from the section $`\eta _p`$. Also one can choose a neighbourhood $`V`$ of $`\eta `$ in $`H^0(\overline{K}(M),M)`$ s.t. there is a map $`\alpha :N^{}\times VU^{}`$ so that $`D_pU^{}=\alpha (N,\eta _p)`$ for $`\eta _pV`$.
We wish to deform the SLag torus fibration on $`N^{}`$ to a SLag torus fibration on $`D_pU`$ for a tubular neighbourhood $`UU^{}`$ of $`N`$. To set up the deformation theory we consider $``$\- the connected component of the moduli-space of of $`C^{2,\alpha }`$ embeddings of $`T`$ into $`N^{}`$, which contains embeddings given by orbits of the $`T`$-action on $`N^{}`$. This is a Banach manifold, whose local chart at a particular (smooth) embedding are $`C^{2,\alpha }`$ sections of the normal bundle of the embedded $`T`$. We consider the space $`^{}=\times V`$, which can be thought as a space of embeddings of $`L`$ into various $`D_p`$ (with $`(f,\eta _p)\alpha (f,\eta _p)`$ for an embedding $`f`$).
So let $`(f,\eta _p)^{}`$ and $`\alpha (f,\eta _p)`$ be the corresponding embedding of $`T`$ into $`D_p`$. Then for a fixed $`\eta _p`$, various embeddings corresponding to different $`f`$’s are all isotopic in $`D_p`$, so they all carry the same isotopy class $`L_p`$ in $`D_p`$. The cohomology class $`[\omega ]`$ of $`\omega `$ restricted to $`D_p`$ is the cohomology class induced from $`\omega `$ on $`M`$. Also the image of $`\alpha (f,\eta _p)`$ is isotopic in $`M`$ to an orbit of the $`T`$-action on $`N^{}`$ and $`\omega `$ restricts to $`0`$ on the orbits. So $`[\omega ]|_{L_p}=0`$. Also on $`N`$ $`\phi ^{}`$ evaluated on orbits was not zero. By continuity $`\phi _p^{}(L_p)0`$. So we can choose $`e^{i\theta _p}`$ s.t.
$$Im(e^{i\theta _p}\phi _p^{}(L_p))=0$$
This $`e^{i\theta _p}`$ is defined up to a sign, and we can choose it to be a smooth function of $`\eta _p`$.
We want to consider the subset $`SLag^{}`$ of those embeddings $`(f,\eta _p)`$ s.t. the forms $`\omega `$ and $`Im(e^{i\theta _p}\phi _p^{})`$ restrict to $`0`$ on $`\alpha (f,\eta _p)`$. Let $`\pi :^{}V`$ be the projection onto the second factor. We also have a space $`S_0`$ of $`T`$-orbits in $`N^{}`$, and we will think of $`S_0\pi ^1(\eta )SLag`$. We have the following:
###### Lemma 3
The space $`SLag`$ near $`S_0`$ is a smooth manifold $`SLag_0`$ of dimension $`n1+2dim(H^0(\overline{K}(M),M))`$. Moreover the differential of the projection $`\pi `$, restricted to $`SLag_0`$ is surjective.
We claim that the statement of the Theorem follows from this Lemma. Indeed $`SLag_0`$ will be a fibration over $`V^{}`$ for sufficiently small neighbourhood $`V^{}`$ of $`\eta `$. If we choose any metric on $`SLag_0`$ (for instance the one induced from the product metric on $``$ and $`V^{}`$), then we have a horizontal distribution for the projection $`\pi `$ on $`SLag_0`$. For $`\eta _pV^{}`$ we can look at a line $`t\eta +t(\eta _p\eta )`$. This line induces a flow $`\rho _p(t)`$ on $`SLag_0`$, s.t $`\rho _p(1)`$ sends $`\pi ^1(\eta )`$ to $`\pi ^1(\eta _p)`$. It is an easy exercise in differential topology that one can choose a smaller neighbourhood $`U_N`$ of $`\eta `$ in $`H^0(\overline{K}(M),M)`$ s.t for $`\eta _pU_N`$ the image of $`S_0`$ under $`\rho _p(1)`$ gives a SLag fibration, which covers the neighbourhood $`UD_p`$.
Now we prove the Lemma: Since $`S_0`$ is compact, it is obviously enough to prove our claim near a point $`LS_0`$ (here we think of $`L`$ as an orbit of the $`T`$-action in $`N^{}`$). A neighbourhood $`Y`$ of $`L`$ in $``$ can be thought as a small ball in the space of $`C^{2,\alpha }`$ normal vector fields to $`L`$. There is a Banach vector bundle over $`^{}`$, given by the direct sum of exact $`C^{1,\alpha }`$ 2-forms with exact $`C^{1,\alpha }`$ $`(n1)`$-forms on $`T^{n1}`$. There is a section $`\sigma `$ of this bundle over $`Y\times V`$, given by $`\sigma (f,\eta _p)=(\alpha (f,\eta _p)^{}(\omega ),\alpha (f,\eta _p)^{}(Im(e^{i\theta _p}\phi _p^{})))`$. The space $`SLag`$ precisely corresponds to the zero set of this section.
We have shown in that one can slightly generalize the argument of to prove that the differential of this $`\sigma `$ is already surjective then restricted to the tangent space $`T_LT_L^{}`$ and the kernel of $`d\sigma `$ in $`T_L`$ is of dimension $`n1`$ (in our case the kernel in $`T`$ corresponds to vectors fields, which realize deformations of the orbits). So one deduces that $`SLag`$ is smooth of dimension $`n1+2dim(H^0(\overline{K}(M),M))`$ near $`L`$ and the differential of $`\pi `$ restricted to $`SLag`$ is surjective. Q.E.D.
Remark 2: Suppose $`D^{}`$ has several connected components. For each connected component we have a SLag fibration on a corresponding neighbourhood of $`D_p`$. But those fibrations are with respect to holomorphic volume forms $`e^{i\theta _p}\phi _p^{}`$ on $`D_p`$, there $`\theta _p`$ are potentially different for various connected components of $`D^{}`$. Suppose $`D_p`$ is smooth. Then $`\phi _p^{}`$ is a holomorphic volume form on $`D_p`$. It is clear from the proof that $`\theta _p`$ will coincide if the corresponding isotopy classes $`L_p`$ coming from different connected components of $`D^{}`$ define the same homology class in $`D_p`$. We will see two examples (3.2 and 3.3), in which this holds. Thus in particular we will produce SLag fibrations on subsets of $`D_p`$, whose complements have arbitrary small volume in $`D_p`$.
We would like to point out that Theorem 3 generalizes to complete intersections on a Kahler manifold $`M`$ with a $`T`$-action of complexity larger then 1. Indeed let $`(M^{2n},\omega )`$ be a Kahler manifold. Suppose that the anticanonical bundle $`\overline{K}(M)`$ is a tensor product
$$\overline{K}(M)L_1\mathrm{}L_d$$
(1)
Suppose we have a $`T^{nd}`$-action on $`M`$ s.t $`T^{nd}`$ also acts on the total space of each $`L_i`$. Moreover we assume that the action is equivariant with respect to the isomorphism (1). Suppose we have sections $`\eta _i`$ of $`L_i`$, which are invariant under the $`T`$-action. Let $`D`$ be their common zero set and $`D^{}`$ be the smooth part of $`D`$. We assume that $`(\eta _1,\mathrm{},\eta _d)`$ is transversal to $`0`$ along $`D^{}`$. Thus we get a trivialization $`\phi ^{}`$ of the canonical bundle of $`D^{}`$ and $`\phi ^{}`$ is invariant under the $`T`$ and the $`T^c`$-actions (here $`T^c`$ is a complex torus of dimension $`nd`$).
Let $`D_0`$ be a connected component of $`D^{}`$. Assume that the $`T`$-action on $`D_0`$ has a finite stabilizer $`H`$. We have, as before, the actions of the quotients $`T/H`$ and $`T^c/H`$ on $`D_0`$. The $`T^c/H`$ action will be free, $`D_0`$ is biholomorphic to $`T^c`$ and we can choose an angle $`\theta _0`$ s.t. orbits of the $`T`$-action will be SLag submanifolds of $`(D_0,\omega ,e^{i\theta _0}\phi ^{})`$. We will also have the following result, which generalizes Theorem 3:
###### Corollary 1
Let $`\eta =(\eta _1,\mathrm{},\eta _d),D,D^{},D_0`$ be as above. Let $`N`$ be a compact, $`T`$-invariant subset of $`D_0`$ and $`U`$ be it’s tubular neighbourhood in $`M`$. Then we can choose a neighbourhood $`U_N`$ of $`\eta `$ in $`H^0(L_1,M)\times \mathrm{}\times H^0(L_d,M)`$ s.t for any $`\eta ^p=(\eta _1^p,\mathrm{},\eta _d^p)U_N`$ we have a SLag torus fibration on a neighbourhood $`U_p`$ of $`D_p`$ with respect to $`(\omega ,e^{i\theta _p}\phi _p^{})`$.
Here $`D_p`$ is the common zero set of $`(\eta _1^p,\mathrm{},\eta _d^p)`$, $`U_p`$ contains $`UD_p`$, $`\omega `$ is the induced metric on $`D_p`$, $`\phi _p^{}`$ is the trivialization of the canonical bundle of $`D_p`$, given by the section $`\eta _p`$ and $`\eta _p\theta _p`$ is smooth.
The proof of the Corollary is analogous to the proof of Theorem 3.
In Sections 3.2-3.4 we consider some examples there Theorem 3 and Corollary 1 apply.
### 3.2 Fermat type quintics in $`P^4`$
Consider $`P^n`$ with a $`T^c=(^{})^{n1}`$-action given by
$$(c_1,\mathrm{},c_{n1})(z_1,\mathrm{},z_{n+1})=(c_1z_1,\mathrm{},c_{n1}z_{n1},c_j^1z_n,z_{n+1})$$
This action of course contains a $`T^{n1}T^c`$-action. Consider the following polynomial $`f_\lambda =z_j\lambda (z_{n+1})^{n+1}`$. Then for every $`\lambda `$ this polynomial defines a section of the anti-canonical bundle $`\overline{K}(P^n)`$, which is invariant under the $`T^c`$-action. There are essentially 2 distinct cases $`\lambda =0`$ or $`\lambda =1`$.
$`\lambda =\mathrm{𝟎}`$: $`D`$ has $`n+1`$ connected components $`D_j`$, given by $`D_j=(z_j=0)`$. Let $`D_j^{}`$ be the smooth part of $`D_j`$ in $`D`$, i.e. $`D_j^{}=(z_j=0,z_i0forij)`$. Consider the following 1-parameter family $`f^t=f_0+t\mathrm{\Sigma }z_j^{n+1}`$. Then $`n=4`$ the zero sets $`D_t`$ are Fermat type quintics. From each component $`D_j^{}`$ we get a SLag fibration on a corresponding part of $`D_t`$, but the fibrations are SLag for potentially different holomorphic volume forms on $`D_t`$ for different components $`D_j^{}`$ (see Remark 2). We will show that the isotopy classes $`L_t^j`$, coming from components $`D_j^{}`$, give the same homology class in $`D_t`$. Thus SLag fibrations will be with respect to one holomorphic volume form $`\phi _t`$ on $`D_t`$.
W.l.o.g it is enough to show that $`L_t^1`$ and $`L_t^2`$ define the same homology class in $`D_t`$. To show this consider a different $`T^{n1}`$-action on $`P^n`$, given by
$$(e^{i\theta _1},\mathrm{},e^{i\theta _{n1}})(z_1,\mathrm{},z_{n+1})=(z_1,z_2,e^{i\theta _1}z_3,\mathrm{},e^{i\theta _{n1}}z_{n+1})$$
This action has a moment map $`\mu =(\mu _1,\mathrm{},\mu _{n1})`$ with
$$\mu _j=\frac{|z_{j+2}|^2}{\mathrm{\Sigma }|z_i|^2}$$
Consider the level set $`\mathrm{\Sigma }=\mu ^1(1/2n,\mathrm{},1/2n)`$ restricted to $`D`$. It’s intersection with each of $`D_1^{}`$ and $`D_2^{}`$ it will be 1 smooth orbit $`L_i`$ in the interior ($`i=1,2`$). Also $`\mathrm{\Sigma }`$ will not intersect $`D_j`$ for $`j>2`$. Also $`L_i`$ are easily shown to coincide with the orbits of the original $`T^{n1}`$-action on $`D`$.
Consider now $`D_t`$ and $`\mu |_{D_t}`$. Then the level set $`\mathrm{\Sigma }_t=\mu ^1(1/2n,\mathrm{},1/2n)`$ on $`D_t`$ will also have 2 smooth components. Also clearly those components will be in the isotopy class of $`L_t^1`$ and $`L_t^2`$. Now the level set of any map to a non-compact manifold carries a trivial homology class. So $`L_t^1`$ and $`L_t^2`$ (with corresponding orientations) will have the same homology class in $`D_t`$ and we are done.
Remark 3: We have constructed Special Lagrangian fibrations on large parts of quintics near the large complex limit quintic. If one forgets about the Special condition and studies Lagrangian fibrations, then one can say more. In fact W.-D. Ruan has constructed in and Lagrangian tori fibrations on general quintics and also constructed symplectic mirrors to those fibrations.
$`\lambda =\mathrm{𝟏}`$: The singular points of $`D`$ are points there $`z_{n+1}=0`$ and $`z_1\mathrm{}z_n=0`$. The set of smooth points $`D^{}`$ will have 2 connected components: $`D_1^{}=(z_{n+1}=0,z_j0forj<n+1)`$ and $`D_2^{}=(z_j0)`$. $`D_1^{}`$ will be the orbit $`(1,\mathrm{},1,0)`$ under the $`T^c`$-action; $`D_2^{}`$ will be the orbit of $`(1,\mathrm{},1,1)`$ under $`T^c`$-action.
We have a moment map $`\mu =(\mu _1,\mathrm{},\mu _{n1})`$ for the $`T^{n1}`$-action on $`P^n`$ with $`\mu _j=\frac{|z_j|^2|z_n|^2}{\mathrm{\Sigma }|z_j|^2}`$. One easily sees that the preimage $`\mu ^1(0,\mathrm{},0)`$ doesn’t intersect the singular points of $`D`$. Also it intersects each component $`D_j^{}`$ in one smooth orbit of the $`T^{n1}`$-action. So we use the same trick as before to show that the isotopy classes $`L_t^1`$ and $`L_t^2`$ define the same homology class on $`D_t`$. Thus from Theorem 3 we get SLag fibrations on large parts of Calabi-Yau hypersurfaces near the complex limit hypersurface $`D`$.
### 3.3 Calabi-Yau Hypersurfaces in the Grassmanian $`G(2,4)`$
Consider the Grassmanian $`M=G(2,4)`$ of 2-planes in $`^4`$, which we continue to identify with a quadric hypersurface $`z_1z_2+z_3z_4+z_5z_6=0`$ in $`P^5`$. We have a $`T^3`$ and a $`T^c=(^{})^3`$-action on $`M`$ as in Section 2.2. The polynomial $`f=(z_1z_2)^2+(z_3z_4)^2+(z_5z_6)^2`$ defines a $`T^c`$-invariant section of the anti-canonical bundle $`\overline{K}(M)`$.
The singular part of the zero set $`D`$ of $`f`$ corresponds to points $`z_{2j1}z_{2j}=0`$ for $`j=1,2,3`$. Let $`D^{}`$ be the smooth part of $`D`$ and $`(z_1,\mathrm{},z_6)D^{}`$. We can assume w.l.o.g. that $`z_1z_2=1`$. Let $`z_3z_4=a`$. Then $`z_5z_6=a1`$ and $`a`$ satisfies $`a^2+a+1=0`$. We get 2 values $`a_i=\frac{1\pm \sqrt{3}i}{2}`$. For each $`a_i`$ we have a point $`b_i=(1,1,1,a_i,1,1a_i)`$ on $`D^{}`$. The orbits of $`T^c`$ through $`b_i`$ give 2 connected components $`D_1^{}`$ and $`D_2^{}`$ of $`D^{}`$ (explicitly $`D_i^{}`$ are given by $`z_3z_4/z_1z_2=a_i`$).
The $`T^3`$-action has a moment map $`\mu =(\mu _1,\mu _2,\mu _3)`$ with $`\mu _i=\frac{|z_{2i1}|^2|z_{2i}|^2}{\mathrm{\Sigma }|z_j|^2}`$. The preimage $`\mu ^1(0,0,0)`$ doesn’t intersect the singular part of $`D`$. Also it intersects each of $`D_j^{}`$ at one smooth orbit. We use the same argument as in Section 3.2 to prove that the isotopy classes $`L_p^1`$ and $`L_p^2`$ on $`D_p`$ give the same homology class. Thus we can use Theorem 3 to obtain SLag fibrations on large parts of Calabi-Yau hypersurfaces $`D_p`$, which are zero sets of polynomials $`f_p`$ of degree 4 for $`f_p`$ sufficiently close to $`f`$.
Remark 4: In we gave an example of a SLag submanifold on a Calabi-Yau hypersurface in $`G(2,4)`$.
### 3.4 Complete intersection of two degree 3 hypersurfaces in $`P^5`$
In this section we want to illustrate an application of Corollary 1. Let $`M`$ be $`P^5`$. We decompose it’s anticanonical bundle as
$$\overline{K}(M)L_1L_2$$
(2)
Here $`L_1=L_2=(\gamma ^{})^3`$. We have a $`T^3`$-action on $`M`$, given by
$$(e^{i\theta _1},\mathrm{},e^{i\theta _3})(z_1,\mathrm{},z_6)=(e^{i\theta _1}z_1,e^{i\theta _2}z_2,e^{i\theta _3}z_3,e^{i(\theta _1+\theta _2)}z_4,e^{i(\theta _1+\theta _2\theta _3)}z_5,e^{i(\theta _1+\theta _2)}z_6)$$
Then $`T`$ acts on $`\overline{K}(M),L_1,L_2`$. Since the linear action of $`T`$ on $`^6`$ is in $`SL(6,)`$, the $`T`$-action is equivariant with respect to the isomorphism (2) (see example (3) in Section 2.2). We have 4 monomials $`g_1=z_1z_2z_4,g_2=z_3z_4z_5,g_3=z_1z_2z_6,g_4=z_3z_5z_6`$, which can be viewed as $`T`$-invariant sections of $`L_1`$ and $`L_2`$. We pick $`\eta _1`$ and $`\eta _2`$ to be some of their linear combinations s.t. conditions of Corollary 1 hold. Thus we get SLag fibrations on large parts of complete intersections of 2 hypersurfaces of degree 3 in $`P^5`$ near $`(\eta _1,\eta _2)`$.
Massachusetts Institute of Technology
E-Mail : egold@math.mit.edu |
warning/0003/astro-ph0003072.html | ar5iv | text | # Achieving a wide field near infrared camera for the Calar Alto 3.5m telescope
## 1 SCIENTIFIC MOTIVATION
There are numerous scientific projects which would benefit from large area infrared surveys. Most fit into the category of discovering and characterising new objects. An example is a survey for very low mass stars, brown dwarfs and free floating giant planets in open clusters and star forming regions. All of these objects are cool ($`\mathrm{T}_{\mathrm{eff}}`$$`\genfrac{}{}{0pt}{}{_<}{^{}}\mathrm{\hspace{0.25em}3000}`$) and have a significantly larger J than R or I band flux. They can thus be detected with their optical–infrared colour, or even their Z–J colours obtainable on the same HgCdTe detector.
Other areas of science which would benefit from surveys include: the initial mass function of star forming regions; the dark matter content and the age of the Galaxy from cool white dwarfs in the Galactic disk and halo; Galactic structure traced via K and M giants (which, due to the lower extinction in the near infrared, can be traced to larger distances); galaxy surveys at low Galactic latitudes; quasar surveys; star formation history and damped Lyman alpha galaxies; high redshift galaxies; gravitational lenses and the cosmological constant.
To serve these scientific goals, we plan to build a wide field near infrared (0.8–2.5$`\mu `$m) imaging camera for doing large area surveys. Given the nature of the Calar Alto Observatory as a resource for German and Spanish astronomers, this camera (Omega 2000) is intended for use in common-user mode rather than undertaking pre-defined surveys.
Some projects are less concerned with area coverage than with volume, in which case deep ‘pencil-beam’ surveys are more suitable. This may be the appropriate strategy when searching for objects at a range of distances, and in some cases may be more efficient than shallower wide field surveys . A list of current and future near infrared survey facilities is given in Table 1.
The rest of this paper is as follows. After giving the important characteristics of the detector, we discuss the general issues influencing the design of the instrument. Much attention is paid to baffling schemes to minimise thermal radiation from the telescope structures. The chosen optical design for Omega 2000 is then presented, a long with a brief discussion of the electronic and software systems and readout modes required for this high data rate instrument.
## 2 INFRARED DETECTOR
The detector is a major factor influencing the design of an infrared camera, largely because of the limited choice and high cost of the available arrays. The largest near infrared science grade array currently available are 1K$`\times `$1K. These are in use in a number of existing cameras, such as the Omega Prime and Omega Cass cameras at Calar Alto. Omega 2000 will use the next generation HAWAII-2 2K$`\times `$2K HgCdTe array from Rockwell . This is very similar to the HAWAII array, except that it has slightly smaller pixels (18.0$`\mu `$m instead of 18.5$`\mu `$m) and 32 rather than 4 outputs, allowing a faster readout (Table 2). The arrays are expected to have a very high filling factor (i.e. essentially no ‘dead area’ between the pixels). Due to the background limited observing conditions, the detector will be operated with a 1 V reset voltage to increase the full well capacity to about 200,000 electrons.
## 3 DESIGN ISSUES
### 3.1 General Considerations
The primary science goal which Omega 2000 must address is wide field imaging, and for a fixed size array this demands that we should have the largest feasible pixel scale (acrseconds per pixel). “Feasible” in this case means (1) compatible with one of the focal stations (prime or Cassegrain) on the 2.2m or 3.5m telescopes at Calar Alto; (2) able to produce high quality images; (3) producible within an acceptable timescale ($`<`$2 years) and budget ($`\genfrac{}{}{0pt}{}{_<}{^{}}`$1 million DM). The timescale is set by the need to bring new technology into prompt scientific use and the delivery time of the science grade array (summer 2001).
Optical design problems aside, the upper limit on the pixel scale is set by requiring a sufficient sampling of the PSF to permit accurate photometry. Too large a pixel scale means that the PSF is undersampled, resulting in increased photometric errors. A low filling factor of the arrays would increase errors further. The size of the PSF is set by the seeing rather than diffraction at these wavelengths and this aperture size. Only limited seeing statistics are available for the 3.5m on Calar Alto: over the period 1993 to 1995, the median near infrared seeing was 1.0<sup>′′</sup>, and only better than 0.8<sup>′′</sup> 22% of the time. (These values may be slightly optimistic as they rely on integration times of less than 3s.) Sampling theory specifies that at least two pixels should span the FWHM of the PSF. However, in order to remove the variable background level common to infrared imaging, it is usually necessary to take multiple images of a given field with non-integer pixel offsets (dithers) between them. This enables a reconstruction of the PSF through sub-pixel sampling, using methods such as those employed with the undersampled WFPC2 camera on HST (e.g. the ‘drizzle’ algorithm ). Thus pixel scales of between 0.4$`{}_{}{}^{\prime \prime }/`$pix and 0.5$`{}_{}{}^{\prime \prime }/`$pix were considered for Omega 2000 on the grounds that they would only give a slight undersampling for a small fraction of the time.
The financial constraints of this project require that the instrument make use of existing telescope optics. Furthermore, as this will only be one of several instruments on the telescope, modifications to the telescope itself are not permitted. The design of our system is therefore limited to ‘conventional’ prime or Cassegrain focus solutions. If these constraints are relaxed, more sophisticated optical designs which allow good optical quality across a large (c. 1 deg.) field-of-view become possible. An example is a three-mirror telescope combined with a Schmidt-type corrector plate, as will be used by the wide field near infrared camera for UKIRT (WFCAM, see Atad-Ettedgui et al., these proceedings) and the dedicated infrared/optical survey telescope VISTA (Table 1). Although the former has made use of an existing telescope it requires a new f/9 secondary mirror, and has some other drawbacks (see section 3.2).
The dominant noise source for ground-based infrared imaging cameras is usually photon noise from the background. A significant part is from the bright sky, which (at night) is predominantly OH airglow (below 3$`\mu `$m) or thermal emission (above 3$`\mu `$m) . However, thermal radiation from the warm (0–15 C) telescope mirrors, structure and dome is significant longwards of about 2.2$`\mu `$m, i.e. for the K band . Therefore, as much of the optics as possible should be enclosed in a cold environment (a dewar). Often, the pupil is reimaged onto a cold Lyot stop to minimise the amount of radiation reaching the detector from structures around the mirror surfaces. However, the complex optics required for pupil reimaging can degrade the optical quality of the image to unacceptable levels as the field of view and/or pixel scale is increased. Nonetheless, a wide field of view can be achieved by dispensing with the reimaging optics and working at low f-ratio, for example at prime focus, but at the expense of a larger background in the K band.
There is, therefore, a trade-off between K band sensitivity, field-of-view (or pixel scale) and optical quality. This trade-off can be seen with the Omega Prime and Omega Cass infrared cameras at Calar Alto (Table 1). The former at prime focus has 1.8 times the field of view of the latter at Cassegrain focus, but is less sensitive in K. We stress that this trade-off is only relevant for the K band. In the J and H bands, thermal radiation from the warm surfaces is negligible compared to the OH emission from the sky. Thus if we only wanted a wide field imaging camera which operated up to about 2$`\mu `$m, prime focus would be the best option.
Weighing up these opposing factors is difficult, especially for a common-user instrument which will be used for a whole range of science projects. However, the scientific emphasis is on wide field imaging, and provided the K band problem is not too severe (see section 3.2) a large field was felt to be a more significant requirement. Furthermore, a reimaging Cassegrain focus solution would require more optics (minimum of 8 aspherical lenses) and probably a non-standard dewar, increasing project time and cost by about 50%.
We also investigated a non-reimaging Cassegrain focus design. As with the prime focus design, extra background from around the pupil can be seen by the detector. If the secondary mirror is undersized then the secondary is the pupil, and the extra background comes from the cold sky. This is a much lower background flux than the warm floor/dome, although the beam from the primary is vignetted by the undersized secondary. If an exact-sized secondary mirror is used the beam is not vignetted, but the camera would see some floor/dome via the secondary. However, in both cases the reduction of the beam from f/10 (from the existing secondary mirror) to f/2.35 (to achieve 0.45$`{}_{}{}^{\prime \prime }/`$pix) gave extremely poor optical quality, so this design had to be rejected.
### 3.2 Baffling a Non-reimaging Camera
The remaining design choices for Omega 2000 are therefore a Cassegrain focus camera with a cold Lyot stop and a no-cold stop prime focus camera. In this section we investigate the relative sensitivities and survey speeds of such cameras and as well as different baffling schemes for prime focus cameras.
Each detector pixel receives all the light from the $`2\pi `$ steradians (hemisphere) in the direction of the dewar window. Most of this is cold, dark dewar which contributes negligible radiation. With a Lyot stop, the rest of the light is from the pupil, but when it is omitted, radiation from the warm surfaces (telescope structure, floor/dome etc.) around the pupil reaches the detector, and the noise in this radiation lowers the instrument’s K band sensitivity.
This radiation can be reduced in a prime focus design by placing a cold annular baffle between the detector and primary mirror (Fig. 1). The further this baffle is from the detector (and the nearer it is to the pupil), the smaller is the solid angle subtended by warm surfaces at the detector. The ideal place for the baffle is at the primary mirror, but a cold baffle of 3.5m radius is not feasible! The only practical place is therefore inside the detector dewar, and its maximum distance from the detector is set by one or more of the following:
1. not vignetting the beam reaching the primary mirror from the sky,
2. the optical quality, mechanical stability and availability of the large dewar window,
3. the cost of and heat load on the increasingly larger dewar.
Typically one would use a baffle which is just large enough so as not to obstruct the view of the primary mirror from any point on the detector. This condition we call critical vignetting. Note that the larger the field-of-view, the larger this baffle must be in order to not vignet, and so the larger the solid angle of warm floor/dome which can be seen. This puts a sensitivity penalty on large pixel scales at prime focus. As the baffle is decreased in radius from its critically vignetting size, the floor/dome solid angle is reduced but at the expense of vignetting the light from the primary mirror (Fig. 2). Note that the amount of primary mirror which can be seen by a given pixel depends on the location of that pixel on the detector. This gives rise to differential vignetting across the field of view. As the baffle is made smaller still, there comes a point at which no floor/dome can be seen from any point on the detector and there is nothing gained in reducing the baffle radius further. This condition we refer to as super vignetting, and has the advantage of uniform vignetting across the whole field of view.
We have done some simulations to see how the sensitivity of an instrument on the 3.5m varies between these two extreme vignetting conditions. With a 0.45$`{}_{}{}^{\prime \prime }/`$pix prime focus camera and a cold baffle fixed 0.8m from the detector, its inner radius is varied between 160mm (critical vignetting) and 82mm (super vignetting). The calculation assumes that the background is due to thermal radiation from the primary mirror (r=1.75m, $`ϵ`$=0.06) and its central hole (r=0.325m, $`ϵ`$=1.0), and the signal is proportional to the area of the primary which can be seen (minus its hole). Sources are assumed to be black bodies at ambient temperatures. The cold baffle is in the dewar so contributes negligible radiation. Radiation from the primary support structure has also been neglected. We further assume that OH emission from the sky contributes an equal<sup>1</sup><sup>1</sup>1This is what is expected in summer. In winter the thermal radiation from the optical surfaces in a perfectly baffled camera only accounts for about 20% of the total background, the rest being sky OH emission. Thus in winter the performance of the prime focus camera relative to a Cassegrain focus one will be better than that shown here. amount of background radiation as the thermal emission from the primary mirror (McCaughran et al., unpublished). Thermal emission from the sky at these wavelengths is negligible compared to the OH emission. Fig. 3 shows how the sensitivity of the instrument relative to a perfectly baffled Cassegrain focus camera changes with baffle size. This perfect camera has only thermal emission from the primary mirror and hole, plus thermal emission from an exact-sized secondary mirror with an assumed emissivity of 0.03. As a reimaging Cassegrain focus camera would have several more lenses, it would have a slightly lower throughput, so the sensitivities of the prime focus systems relative to this are slightly conservative.
These simulations show that the prime focus camera is almost always less sensitive than a perfectly baffled Cassegrain focus camera in the K band. Interestingly, however, the performance of the prime focus camera is improved (sometimes considerably) by using an undersized baffle. The optimum baffling radius is different for different points on the detector, because the projection of the baffle onto the primary mirror is different for each detector point.
The situation can be improved by adding a warm baffle further away from the detector. The ideal case would be to put a 3.5m diameter high reflectivity ring around the primary mirror which looks through the dome slit at the cold sky. But as it would need an outer radius of almost 4.6m, it would be very expensive and highly impractical. Hence the warm baffle must hang as far below the dewar window as is possible without vignetting the beam from the sky, which is about 1.7m for the Calar Alto 3.5m (see section 4.1). As in Omega Prime, the warm baffle is a polished ellipsoidal annulus which looks into the dewar window to avoid collecting radiation from warm surfaces . Gold coating will further lower its emissivity, but we conservatively assume $`ϵ`$=0.10 here. Fig. 4 shows how the sensitivity varies as we change the size of this baffle. Note that a critically vignetting cold baffle is also included to reduce the amount of the warm baffle seen by the detector. Because the warm baffle is further from the detector, the variation in sensitivity both for different pixel positions and baffle sizes is less than in Fig. 3. Moreover, the sensitivity relative to the perfect Cassegrain focus camera is better. Uniform sensitivity across the whole field-of-view is usually desirable for a survey instrument, so in practice the super vignetting case may be preferred to intermediate vignetting.
As mentioned earlier, a larger field of view requires a larger baffle for any given vignetting mode, and therefore a reduced camera sensitivity. To get a quantitative idea of this, we repeated the above calculations for the super and critical vignetting modes but for two different pixel scales, 0.40$`{}_{}{}^{\prime \prime }/`$pix and 0.45$`{}_{}{}^{\prime \prime }/`$pix, and with baffles at various distances. The results are shown in Table 3 and indicate that for a given pixel scale and baffle distances, super vignetting provides better sensitivity than critical vignetting. Having the cold baffle 1.0m rather than 0.8m from the detector gives only a relatively small improvement, and the extra demands this places on the system (e.g. larger dewar window) are probably not worth it. When super vignetting with the warm baffle, the 0.40$`{}_{}{}^{\prime \prime }/`$pix camera is 1.18 times faster for a single shot than the 0.45$`{}_{}{}^{\prime \prime }/`$pix camera. However the 0.45$`{}_{}{}^{\prime \prime }/`$pix scale covers a large area 1.27 times faster, with the net result that the 0.45$`{}_{}{}^{\prime \prime }/`$pix scale is slightly faster for K band surveys. Of course, in the J and H bands super vignetting should never be used, so implementation of super vignetting requires a baffle which can be rapidly (and automatically) moved. Note that the 0.45$`{}_{}{}^{\prime \prime }/`$pix camera in super vignetting mode should have the same K band sensitivity as the current Omega Prime camera (which has a critically vignetting warm baffle). Interestingly, Omega Prime camera should be about 15% faster if it had a super vignetting warm baffle. Such a baffle has been made, but attempts to test it on the telescope have so far been foiled by bad weather.
The 0.45$`{}_{}{}^{\prime \prime }/`$pix camera is still 2.2 times slower than a perfectly baffled Cassegrain focus camera. However, as a pixel scale of only about 0.35$`{}_{}{}^{\prime \prime }/`$pix with good optical quality could be achieved at Cassegrain focus, the survey speed of the perfect Cassegrain focus camera is only about 1.3 times faster in K, and 1.7 times slower in J and H. It should be emphasised that these figures are rather sensitive to the assumptions laid out above, in particular with regard to the emissivity of the optical surfaces and the fraction of OH emission from the sky. Lowering the emissivity or temperature of the warm baffle in particular leads to much better sensitivity.
A three-mirror Cassegrain focus camera similar to that which will be used by WFCAM on UKIRT (3.8m f/2.5 primary) would be possible for the 3.5m at Calar Alto. However, although the current WFCAM design has a cold Lyot stop, it has a vignetting transmission of 73% in all wavebands (WFCAM web pages, ATC, Edinburgh). This is very similar to the K band transmission of Omega 2000 when using the super vignetting warm baffle. Additionally, WFCAM has a center-to-corner distortion of 0.5%, or 42 pixels. However, the super vignetting Omega 2000 design does not appear to be extendible to the much larger field-of-view obtainable with WFCAM.
## 4 OMEGA 2000 DESIGN
### 4.1 Optical Design
Following all of the considerations described in section 3.1, we decided to place Omega 2000 at the prime focus of the 3.5m telescope with a 0.45$`{}_{}{}^{\prime \prime }/`$pix scale. The optical system consists of four lenses (Fig. 5) and has excellent optical quality (Figs. 6, 7 and 8). Three filter wheels are envisaged, each with slots for six 3 inch (7.6cm) diameter filters. The parameters for the system are given in Table 4.
The very low distortion of $`<`$0.06% is less than 1 pixel center-to-corner, permitting images to be overlaid with a simple x,y shift. As demonstrated in the previous section, the sensitivity will be very similar to Omega Prime: with a two minute integration in 1<sup>′′</sup> seeing, the central pixel of a point source is 5$`\sigma `$ above the background noise at limits of J=19.2, H=18.1 and K’=17.5 magnitudes . This compares with a mean sky brightness of approximately J=15.2, H=13.6 and K’=13.0 (MAGIC instrument web pages, Calar Alto).
### 4.2 Mechanical Design
All of the optics will be enclosed in a dewar cooled with liquid nitrogen. The design will be similar to that of Omega Prime with two nested tanks to ensure a stable temperature of 77 K. The optimal baffling was discussed in section 3.2. A fixed critically vignetting cold baffle is on the inside of the dewar window. The window has a radius of 350mm. The cylindrical dewar will be about 0.6m in diameter and between 1.5m and 2.0m long: the exact figure will depend on the final cryogenic design. The heat load on the dewar is large, approximately 300W. A second warm baffle shaped as an oblate ellipsoid sits 1.1m from the dewar window (i.e. 1.7m from the detector). This distance corresponds with the bottom of the existing light shield attached to the bottom of the front ring. This shield actually slightly vignets the beam from the sky, but as it is required for other instruments it cannot be removed. Although the exact design has not yet been worked out, the warm baffle will be moveable between the critical and super vignetting modes (or may consist of two separate baffles). The former mode will be used for the J and H bands, and the latter for the K band. The projection of the super vignetting baffle onto the primary mirror (as seen from the detector) has a diameter of 3.0m, so does not threaten to diffraction limit the instrument at the longest wavelength (2.5$`\mu `$m).
### 4.3 Electronics and Readout Modes
A general readout electronic system is being developed in house. This will be used for a number of different instruments currently under development at MPIA, including Omega 2000 and the MIDI interferometer for the VLT. Only small modifications are required for different detector types.
The HAWAII-2 detector works either with 1 signal channels per quadrant or with 8 signal channels per quadrant simultaneously. All four quadrants can be operated in parallel. The pixels can be clocked at up to 1 MHz which will deliver 1 Megapixels per second per signal channel, with each pixel represented as 16 bits. The minimum integration time is given by the frame readout time, which is the normal state for arrays with integrate-while-read working mode. When working with all 32 channels, the data rate will be 64 Mbyte/sec and the full frame rate is about 8Hz. With only 4 channels (1 per quadrant) the data rate is 4 Mbyte/sec and the full frame rate is about 1Hz.
The background limit will be reached in broadband imaging with the HAWAII-2 array on Omega 2000 in a few seconds. Many images (“repeats”) are therefore required to achieve sufficient SNR, so it is very important that the array can be read out quickly with minimum dead time. It is also not necessary to use readout modes which greatly reduce the read noise. There are three different double correlated readout (DCR) modes suitable for detectors with non-destructive readout modes, such as the HAWAII detectors, and each give very different readout efficiencies:
1. Double correlated read. This is the most conservative mode, in that it passes throught the whole array three times to achieve a single image. The first time the whole array is reset, in the second the whole array is read to determine the pixel bias (offset) values, and finally the third read measures the integrated (light-exposed) values. The difference between the integrated and bias values for each pixel is the signal. At minimum integration time the efficiency (as measured by the ratio of integration time to total cycle time required to achieve an image) is 33%. To achieve an integration time longer than the minumum, there is a pause after the second reading of the whole array.
2. Double correlated read with fast reset. This is the current standard readout mode for MPIA instruments. In this mode the whole array must be clocked twice to achieve a single image. First, the whole array is clocked line by line, in which a fast reset is applied to the line followed immediately by a read of that line to determine the pixel biases. Then the whole array is clocked line by line a second time at the same speed as the first array clocking to read the integrated pixel values. The efficiency is 50% for the minimum integration time.
3. Full MPIA mode. The problem with the previous mode is that as soon as the array has been clocked the second time to read the integrated pixel values, those pixels continue to integrate: they are only reset once the first clocking (reset-read) of the next cycle is started. Full MPIA improves upon this by having only a single type of pass through the whole array. The array is clocked line by line starting with a read of the integrated pixels of a line, followed by a fast reset of that line, and then a read of the pixel bias values of that line. These bias values are those which will be subtracted from the integrated pixel values obtained the next time the whole array is read. Once the whole array has been read, there is a pause (if the integration time is above the minimum) and the cycle repeated for as many images (“repeats”) as required. Essentially no time is now wasted, and for any integration time the efficiency of this mode is almost 100%.
Whereas the first mode is a completely array-oriented approach, the full MPIA mode is a completely line-oriented approach. DCR with fast reset is essentially a mix of the two. Full MPIA mode is more efficient because each line is reset and the bias values read immediately after reading the integrated pixel values: DCR mode with fast reset waits until the whole array has been read before reseting. To obtain just a single image, DCR with fast reset takes the same time as full MPIA mode, but for a sequence of many repeats, full MPIA mode is much quicker (Table 5). The relative speeds for different integration times of the three modes has been empirically determined at MPIA using HAWAII arrays (Fig. 9). Note that the minimum integration time with full MPIA mode is twice that of DCR mode with fast reset (0.25 rather than 0.125s for the HAWAII-2), but this is not a problem for Omega 2000.
### 4.4 Data Acquisition and Control Software
The MPIAs GEneric InfraRed camera Software (GEIRS) will also be used for Omega 2000. The new challenge is the 26 times higher data rate compared to the 1K$`\times `$1K HAWAII detector on existing instruments. The data acquisition software must be able to sustain this data rate continuously, apply necessary preprocessing tasks, have data visualization control and save the data in time to non-volatile storage disks. The data are acquired by either one or two 30 MHz 16bit parallel interfaces.
The handling of the increased data rate can only be organised in a current state of the art symmetrical multiprocessor system (e.g. a SPARC Ultraserver 450), where nearly the full memory bandwidth is consumed for the basic tasks needed to control the camera system and human interfacing. The software package is already optimised for minimal usage of memory by sharing data memory between all processes. To keep the jobs synchronised without loss of efficiency, further parallelization making optimal use of the multiprocessing capabilities has to be implemented. The granularity of the shared memory buffers for the parallel processing also has to be reorganised to make more efficient use of the increased internal working memory. Although this can be done with standard computer architecture, there remains the option of using distributed processing on standard hardware with distributed memory (the so-called “multiple instructions on multiple data” organization). |
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