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warning/0507/hep-th0507006.html | ar5iv | text | # Braneworlds
## 1 Introduction
During recent years, cosmology has become one of the most successful fields in physics. The precise measurements of the anisotropies in the cosmic microwave background have confirmed a simple ’concordance model’: The Universe is spatially flat. Its energy density is dominated by vacuum energy (or a cosmological constant) which contributes about 70% to the expansion of the Universe, $`\mathrm{\Omega }_\mathrm{\Lambda }0.7`$. The next important contribution is cold dark matter (CDM) with $`\mathrm{\Omega }_{CDM}0.3`$. The expansion velocity of the Universe is given by the Hubble constant $`H_0=100h`$km/sMpc, with $`h0.7`$. Baryons only contribute a small portion of $`\mathrm{\Omega }_bh^20.02`$. Massive neutrinos contribute similarly or less.
The structures in the Universe (galaxies, clusters, voids and filaments) have formed out of small initial fluctuations which have been generated during inflation and have an almost scale-invariant spectrum, $`n=1\pm 0.1`$. There is probably also a small amount of tensor fluctuations (gravity waves) generated during inflation, which however has not yet been detected.
All the above numbers are accurate to a few percent and will be measured even more precisely with ongoing and planned experiments. This situation is unprecedented in cosmology. About twenty years ago, these numbers where known at best within a factor of two or even only by their order of magnitude. The concordance model is in agreement with most cosmological data, most notably the CMB anisotropy measurements, supernova type Ia distances (see contribution by Varun Shani), statistical analysis of the galaxy distribution, constraints from cosmic nucleosynthesis, cluster abundance and evolution etc.
However, on a theoretical level our understanding has remained poor. We have no satisfactory answers to the questions:
* What is dark matter ?
* What is dark energy?
* What is the ’inflaton’? Or what is, more precisely, the physics of inflation?
* How can we resolve the Big Bang and other singularities of classical general relativity?
There is justified hope that the last question could be resolved within a theory of quantum gravity, which is anyway needed if we want to put all fundamental interactions on a common footing. At present, the most successful attempt towards a theory of quantum gravity is string theory. This theory is based on the assumption that the ’fundamental objects’ are not particles but one dimensional strings. Particles then manifest as excitations, proper modes, of strings. It would lead us much too far to give an introduction to string theory at this point. The interested student will have to study the two volumes of Polchinski Pol .
In this course we next give an introduction to braneworlds, or, more generally, to physical effects of extra dimensions. In Section 3, we derive the gravitational equations for braneworlds with one co-dimension from Einstein’s equations in the bulk. We then discuss in detail the Randall–Sundrum II model, its background and its perturbations. In Section 5, we give an introduction to 4-dimensional cosmological perturbation theory and, especially to the CMB anisotropy spectrum. Only after this we are ready for braneworld cosmology in Section 6. We write down the most general brane cosmology in an empty bulk and we investigate some of its modifications w.r.t. 4-dimensional cosmology. In particular, we discuss the modification of the slow roll parameters during braneworld inflation. We also present one example of the modifications in the evolution of cosmological perturbations which are relevant in braneworlds. We end with some conclusions.
Notation: We use capital Latin indices $`A,B,\mathrm{}`$ to denote bulk coordinates, lower case Greek indices $`\mu ,\nu ,\mathrm{}`$ for coordinates on a four dimensional brane, and lower case Latin indices, $`i,j,\mathrm{}`$ for spatial 3-dimensional quantities. We sometimes also use bold symbols to denote 3d spatial vectors. We use the metric signature $`(,+,\mathrm{},+)`$. The 4-dimensional Minkowski metric is denoted by $`(\eta _{\mu \nu })`$.
Throughout we set $`c=\mathrm{}=k_{\mathrm{Boltzmann}}=1`$ so that time and length scales are measured in inverse energies (usually GeV’s), and mass and temperature correspond to an energy. The four dimensional Newton constant is then given by $`G_4=0.67\times 10^{38}`$GeV<sup>-2</sup>. Useful relations in this set of units are $`1=0.2`$GeV fm and 1 eV $`=1.16\times 10^4`$K. Here fm $`=`$ femtometer $`=10^{15}`$m.
## 2 Basics of Braneworlds
### 2.1 What are Braneworlds?
The interest of string theory lies in the fact that it may provide a unified description of gauge interactions and gravity. Its weak point is, that it is extremely hard to make predictions from string theory which are testable at energies available in experiments. The reason for that is that string theory probably fully manifests itself only at very high energies of the order the Planck scale. The observed 4-dimensional Planck scale is given by Newton’s constant, $`G_4`$. In our units with $`\mathrm{}=c=1`$ the Planck scale is $`E_4=M_4=1/\sqrt{4\pi G_4}3\times 10^{18}`$GeV. This energy scale cannot be achieved by far at terrestrial accelerators (the LHC presently under construction at CERN will achieve about 7000GeV).
Nevertheless, string theory makes some relatively firm predictions which might lead to observational consequences at low energy. First of all, it predicts that spacetime is ten-dimensional with one time and nine spatial dimensions. Since the observed world has only four dimensions, one usually assumes that the other six are compact and very small, so that they cannot be resolved by any physical experiment available to us so far.
Furthermore, string theory predicts the existence of so called $`p`$-branes, $`p+1`$-dimensional sub-manifolds of the ten dimensional spacetime on which open strings end. Gauge fields and gauge fermions which correspond to string end points can only move along these $`p`$-branes, while gravitons which are represented by closed strings (loops) can propagate in the full spacetime, the ’bulk’.
This basic fact of string theory has led to the idea of braneworlds: it may be that our $`3+1`$-dimensional spacetime is such a 3-brane. If this is so, only gravity can probe the bulk and the additional dimensions can be much larger than the smallest length scale which we have probed so far, which is of the order of $`(200`$GeV$`)^110^{18}`$m. Actually, Newton’s law has been tested only down to scales of about 0.1mm micGra . Hence, in the braneworld picture where only gravity can probe the extra-dimensions, these can be as large as 0.1mm$`=10^3`$m. In the next subsection we show how this fact can be employed to address the hierarchy problem.
### 2.2 Lowering the fundamental Planck scale
The fact that the 4-dimensional Planck scale, $`M_410^{19}`$GeV is so much larger than the fundamental scale in elementary particle physics, the electroweak scale $`E_{ew}10^3`$GeV is called the hierarchy problem. Apart from it seeming unnatural to have two so widely separated scales to describe fundamental physics, a more serious problem is the fact that as soon as we have a unified quantum theory which describes also gravity, the scale $`M_4`$ will enter in quantum corrections of all electroweak scale quantities which are not especially protected e.g. by symmetries and it will therefore completely spoil the so successful low energy standard model.
Here we show that within the braneworld picture, it is possible that the 4-dimensional Planck scale is not fundamental but only an effective scale which can become much larger than the fundamental Planck scale $`M_P`$ if the extra-dimensions are much large than $`M_P^1`$. Our argument goes back to Arkani-Hamed , Dimopoulos and Dvali (1998) ADD .
Let $`M_P`$ be the fundamental Planck scale and $`L`$ the size of $`n`$ extra dimensions. In addition there are 3 large spatial dimensions (and time). For simplicity we assume the $`n`$ extra-dimensions to be rolled up as a cylinder, $`(S^1)^n^3`$. The gravitational constants $`G_{(n+4)}`$ and $`G_4`$ are defined by the force laws of gravity. Due to the Gauss constraint these must have the forms
$$F_{(n+4)}=G_{(n+4)}\frac{m_1m_2}{r^{n+2}},\text{ and }F_4=G_4\frac{m_1m_2}{r^2}.$$
On small scales, $`rL`$, an observer on the brane sees $`n+4`$ dimensional gravity, while on large scales, $`rL`$, the cylinder can no longer be resolved and simple 4 dimensional gravity is observed. In order to relate the constants $`G_4`$ and $`G_{(n+4)}`$, we de-compactify the compact dimensions leading to an n-dimensional lattice of masses which looks from far like a hypersurface of mass density $`m/L^n`$ (see Fig. 1). Around the mass distribution, we now form a cylinder $`C`$ of dimension $`n+2`$, length $`d`$ and radius $`r`$ (see Fig. 1). To satisfy Gauss’ law, we require
$$_CF_{}𝑑\sigma =S_{(2+n)}G_{(4+n)}\times (\text{mass in }C).$$
Here, $`F_{}`$ is the component of the force normal to the surface and $`S_j`$ is the volume of a $`j`$ dimensional sphere, $`S_j=2\pi ^{(j+1)/2}/\mathrm{\Gamma }([j+1]/2)`$, where $`\mathrm{\Gamma }`$ denotes the Gamma-function, $`\mathrm{\Gamma }([j+1]/2)=([j1]/2)!`$ for integer values of $`(j+1)/2`$.
The first integral is $`4\pi r^2d^nF(r)`$ while the mass inside the cylinder is $`md^n/L^n`$. With the 4-dimensional force law this implies
$$G_4=\frac{S_{(2+n)}}{4\pi }\frac{G_{(4+n)}}{L^n}.$$
(1)
In order to relate the gravitational constant to the Planck mass, we express Newtonian gravity in terms of an action principle. For a static weak gravitational field $`\varphi `$ and a mass density $`\rho `$ in $`4+n`$ dimensions, we can obtain the Poisson equation by varying the action
$$I=d^{(3+n)}x\left[\frac{M_P^{(2+n)}}{2}\varphi _{(3+n)}^2\varphi +\rho ^{(4+n)}\varphi +\mathrm{}\right].$$
Integrating out $`\varphi `$, we obtain again the Newtonian force law and the relation
$$M_P^{(2+n)}=\frac{G_{(4+n)}^1}{S_{(2+n)}}.$$
(2)
Integrating the Lagrangian over the extra dimensions relates the 4- and $`4+n`$-dimensional Planck masses by
$$M_4^2=M_P^{(2+n)}L^n.$$
Together with Eq. (2) this reproduces the result (1). For a sufficiently large length scale $`L`$, $`M_43\times 10^{18}`$GeV can therefore be much larger than the fundamental higher dimensional Planck scale $`M_P`$.
Experimental ’micro-gravity’ bounds micGra require $`L<0.1`$mm. For $`n=2`$ and $`L1`$mm the fundamental Planck scale can be of the order of the electroweak scale, $`M_P(110)`$TeV . This Planck scale seems to be in agreement with most other bounds (cooling of supernovae, evolution of the Universe, etc) but leads to the very interesting prospective that effects from string theory might be observable at the Large Hadron Collider (LHC) presently under construction at CERN bounds .
Using ’large’ extra-dimensions, $`LM_P^1`$, the fundamental Planck scale can therefore be of the same order as the electroweak scale. On the other hand, there is no explanation for the length scale $`L0.1`$mm$`10^3`$eV. So the hierarchy problem has not actually been solved, but it has been moved from an energy hierarchy to an unexplained length scale. The hope is, that there would exist solutions of string theory which lead to such a scale dynamically.
### 2.3 New Physics from higher dimensions
##### Kaluza-Klein modes
We denote the four brane coordinates by $`x^\mu `$ and the additional $`n`$ bulk coordinates by $`y^a`$. For simplicity, we consider a bulk where the extra dimensions are rolled up in a cylinder of circumference $`L`$. In the general situation where the compact extra dimensions form a non-flat Calabi-Yau manifold, $`𝒞`$, the exponentials below have to be replaced by the corresponding eigen-functions of the Laplacian on $`𝒞`$. The case of a warped geometry, where the extra dimensions may even be non-compact, will be discussed separately in Section 4.
Be now $`\varphi `$ a massless scalar field in the ’bulk’, $`\varphi (x,y)`$ with $`\varphi (\mathrm{},y^a+L\mathrm{})=\varphi (\mathrm{},y^a,\mathrm{})`$. We can expand the y dependence of $`\varphi `$ in Fourier series
$$\varphi (x,y)=\underset{j}{}\varphi _j(x)\mathrm{exp}(i2\pi jy/L)$$
Since $`\varphi `$ satisfies the massless wave equation, $`_{4+n}^2\varphi =0`$, in the bulk, for a 4d observer which cannot resolve the scale $`L`$, the modes $`j0`$ will become massive fields,
$$(_4^2+m_j^2)\varphi _j=0\text{with }m_j^2=\frac{(2\pi j)^2}{L^2}.$$
These modes of the gravitational potential give raise to exponential corrections to Newton’s law,
$$V(r)=\frac{G_4}{4\pi r}[1+\frac{e^{r/L}}{r^2}+\mathrm{}].$$
If the extra dimensions are large, the first few masses can be very low, but if the graviton is the only bulk field, its weak coupling to other bulk modes leaves the theory nevertheless viable bounds .
As we shall see in Section 4, the mass spectrum can even be continuous, $`0<m<\mathrm{}`$, if there are non-compact extra dimensions.
##### Higher dimensional spin modes
The spin states of massless particles in $`d`$ space time dimensions are characterized by an irreducible representation of $`SO(d2)`$. For $`d=4`$, massless particles always carry a 1–dimensional representation of $`SO(2)`$ <sup>1</sup><sup>1</sup>1The full little group for massless particles is actually $`ISO(2)`$, the group of two dimensional Euclidean motions and rotations. But for finite dimensional representations of $`ISO(2)`$ the translations are acting trivially. The reason why not all representations of the universal covering group of $`SO(2)`$, , have to be considered, namely $`e^{ikx},k`$ is rather subtle. On the classical level, where $`SO(2)`$ is the relevant group, this problem disappears. A full discussion of the finite dimensional representations of the Poincaré group for massless particles in $`d=4`$ can be found e.g. in Wein .. Taking into account also parity, this leads to the two helicity modes of all massless particles in 4 dimensions, independent of their spin. This situation changes drastically if we allow for extra dimensions. Let us consider $`d=5`$: a massless particle of spin $`s`$ in $`d=5`$ spacetime dimensions, carries the representation $`D^s`$ of $`SO(3)`$ and thus has $`2s+1`$ spin states, like a massive particle of spin $`s`$ in $`4`$ dimensions.
Let us, for example, consider the graviton. In $`4+1`$ dimensions it has the $`5`$ helicity states of the tensor representation of SO(3). Projected onto a $`3+1`$ brane, two of them become the usual spin 2 graviton, two are a spin 1 particle, a gravi-vector and one has spin 0, the gravi-scalar.
The gravi-vector couples to the $`\mu 4`$ components of the energy momentum tensor. Interpreting these as the electromagnetic current $`J^\mu `$ and the gravi-vector as the electromagnetic potential $`A^\mu `$, the five-dimensional Einstein equations lead to Maxwell’s equations for $`A^\mu `$ and $`J^\mu `$. This is the so called Kaluza–Klein miracle, which is also true if any, non–Abelian gauge group replaces the one-dimensional torus which plays here the role of the electrodynamic gauge group $`U(1)`$. This finding of Kaluza and Klein KK has evoked an interest in extra-dimensions in the 20ties, long before string theory.
The positive aspects of the vector sector are, however, over shaded by the problems coming from the gravi-scalar. It couples to the four-dimensional energy momentum tensor and modifies gravity. It leads to a scalar-tensor theory of gravity with several observable consequences. For example, one can calculate the modification in the slowing down of the binary pulsar bipul PSR1913+16 due to the radiation of gravi-scalars. In Ref. DK it is shown that in the simple case of a 5-dimensional cylindrical bulk, this leads to a modification of the quadrupole formula by about 20%, while observations agree with the quadrupole formula to better than $`\frac{1}{2}`$%. If there are more than one extra-dimensions, there are several gravi-scalars and this problem is only enhanced.
Clearly, a modification of higher dimensional gravity is necessary to address the problem. Usually, one gives a mass to the gravi–scalar. There are several ways to do this and the resulting four dimensional theory in general depends on this choice. One proposal is the Goldberger–Wise mechanism GW .
Another solution is offered by non-compact extra-dimensions. As we shall see, in curved spacetimes, extra-dimensions can even be infinite. Due to a so called ’warp factor’ they become very small when seen from the brane. For infinite extra dimensions it can happen that the gravi-scalar represents a non-normalizable and therefore unphysical mode. This is precisely what happens in the Randall–Sundrum model and we shall discuss it in this context in Section 4.
## 3 Geometry of five dimensional Braneworld geometry
From now on we restrict ourselves to five dimensional braneworlds, i.e. braneworlds with only one extra-dimension. The idea here is still that spacetime has 10 (or for M-theory 11) dimensions, but 5 (or 6) of them are compactified to a static manifold of about Planck scale, while one of them is large. This picture is motivated mainly from 11-dimensional M-theory, e.g. the Horava-Witten model HOWI , but we shall not try to implement a realization of this model here. Nevertheless, it is important to note, that one co-dimension differs significantly from more than one. The main point is that the 3-brane splits space into two parts, the ’left’ and the ’right’ hand side of the brane. We shall see, that in this case it is possible to determine the gravitational equations on the brane by simply postulating Einstein’s equations in the bulk. This is no longer possible in the case of two or more extra-dimensions.
In this section we derive and discuss these brane gravity equations. In the next section we shall apply them to the Randall–Sundrum model, which we consider as beeing so far the most promising braneworld model.
To determine the gravitational equations on the brane, we start from the basic hypothesis that string theory predicts Einstein gravity in the bulk,
$$G_{AB}=\kappa _5T_{AB}.$$
(3)
We want to discuss in detail the situation of a $`3`$\- brane in a 5–dimensional bulk. We denote the bulk coordinates by $`(x^A)=(x^\mu ,y)`$, where $`(x^\mu )`$ are coordinates along the brane and $`y`$ is a transverse coordinate. We denote the brane position by $`y=y_b`$; in general $`y_b`$ depends on the point on the brane, $`y_b=y_b(x^\mu )`$. A more general embedding of the brane will be discussed below. Very often we consider an energy momentum tensor of the form
$$T_{AB}((x^C)=(x^\lambda ,y))=\frac{\mathrm{\Lambda }_5}{\kappa _5}g_{AB}+\delta _A^\mu \delta _B^\nu T_{\mu \nu }((x^\lambda ))\delta (yy_b).$$
(4)
Here $`\mathrm{\Lambda }_5`$ is a bulk cosmological constant, $`T_{\mu \nu }`$ is the energy momentum tensor on the brane and $`\kappa _5=6\pi ^2G_5`$ is the five-dimensional gravitational coupling constant. This is the most general ansatz if we do not allow for any matter fields in the bulk.
### 3.1 The second fundamental form
As above, $`g_{AB}`$ is the bulk metric. We denote the projection operator onto the brane by $`q_\mu ^A`$. The induced metric on the brane, also called the first fundamental form, is then given by
$$g_{\mu \nu }(x^\lambda )=q_\mu ^A(x^\lambda )q_\nu ^B(x^\lambda )g_{AB}(x^\lambda ,y_b(x^\lambda )).$$
(5)
We denote the covariant derivative in the bulk by $`{}_{}{}^{b}_{A}^{}`$ and covariant derivative on the brane (with respect to the induced metric) by $`_\mu `$. We also introduce the brane normal $`n`$, a vector field defined on the brane which is normal to all vectors parallel to the brane. Clearly, for an arbitrary vector field $`X=X^\mu _\mu `$ along the brane, $`_\mu X{}_{}{}^{b}_{\mu }^{}X`$. The difference of these two covariant derivatives is given by the extrinsic curvature also called the second fundamental form which we now introduce.
Be $`X=X^\mu _\mu `$ and $`Y=Y^\mu _\mu `$ two vector fields on the brane. Their covariant derivative on the brane is given by
$$_YX=(Y^\mu _\mu X^\nu +\mathrm{\Gamma }_{\mu \beta }^\nu X^\mu Y^\beta )_\nu ,$$
while the covariant derivative in the bulk is
$${}_{}{}^{b}_{Y}^{}X=(Y^\mu _\mu X^\nu +\mathrm{\Gamma }_{\mu \beta }^\nu X^\mu Y^\beta )_\nu +\mathrm{\Gamma }_{\mu \beta }^4_4.$$
Therefore, there exists a bi-linear form $`K_{\mu \nu }`$ on the brane such that
$${}_{}{}^{b}_{Y}^{}X=_YX+K(X,Y)n.$$
(6)
Since $`{}_{}{}^{b}_{Y}^{}X{}_{}{}^{b}_{X}^{}Y=[Y,X]`$ is tangent to the brane, $`K(X,Y)K(Y,X)=0`$, hence $`K`$ is symmetric. $`K`$ is called the ’second fundamental form’ or ’extrinsic curvature’ of the brane. Its sign is not uniquely defined in the literature; we shall use Eq. (6) as its definition.
Since $`n`$ is the unit normal of the brane, $`g(n,{}_{}{}^{b}_{Y}^{}X)=K(Y,X)`$. But as $`n`$ is normal to the brane vector field $`X`$ we have $`0={}_{}{}^{b}_{Y}^{}(g(n,X))=g({}_{}{}^{b}_{Y}^{}n,X)+g(n,{}_{}{}^{b}_{Y}^{}X)`$, so that
$$K(Y,X)=g({}_{}{}^{b}_{Y}^{}n,X)=\frac{1}{2}[g({}_{}{}^{b}_{Y}^{}n,X)+g({}_{}{}^{b}_{X}^{}n,Y)].$$
In components, $`K(X,Y)=K_{\mu \nu }X^\mu Y^\nu `$, this becomes
$$K_{\mu \nu }=\frac{1}{2}[{}_{}{}^{b}_{\mu }^{}n_\nu +{}_{}{}^{b}_{\nu }^{}n_\mu ].$$
(7)
Close to the brane we can choose coordinates such that
$$g_{AB}dx^Adx^Bds_b^2=g_{\mu \nu }dx^\mu dx^\nu +dy^2.$$
In these so called ’Gaussian normal coordinates’, $`n=_y`$ and $`{}_{}{}^{b}_{\mu }^{}n_\nu =\mathrm{\Gamma }_{\mu \nu }^4=(1/2)g_{\mu \nu },_4`$. The second fundamental form then becomes simply
$$K_{\mu \nu }=\frac{1}{2}_yg_{\mu \nu }.$$
For general coordinates $`(z^\mu )`$ on the brane we have to define a brane parameterization $`x^A=X_b^A(z^\mu )`$. The vector fields $`(e_\mu )=(_\mu X_b^A(z)_A)`$ then form a basis of tangent vectors on the brane. In terms of these one obtains by means of Eq. (7)
$$K_{\mu \nu }=\frac{1}{2}\left[g_{AB}(e_\mu ^A_\nu n^B+e_\nu ^A_\mu n^B)+e_\nu ^Ae_\mu ^Bn^Cg_{AB,C}\right],$$
(8)
where a comma denotes an ordinary derivative, $`f_{,C}=\frac{f}{x^C}`$.
### 3.2 The junction conditions
Einstein’s equations with a thin hyper-surface of matter become singular since there is a $`\delta `$\- function in the energy momentum tensor. Integrating them once across the brane leads to the so called junction or jump conditions (of Israel, Lancos, Darmois, Misner) Is ; La ; Da ; Mi .
Before we come to the algebraically more complicated situation of general relativity, let us first recall the well known junction conditions of electrostatics: we consider a conducting boundary surface (e.g. capacitor plate) with surface charge density $`\rho `$ and current density $`𝒋`$ along the surface.
The homogeneous Maxwell equations require that the tangential part of the electric field $`E_{}`$ and the normal component of the magnetic field, $`B_{}`$, are continuous across the boundary. This is usually derived by the following argument: denoting the two sides of the capacitor plate by the super-scripts $`+`$ and $``$ the homogeneous Maxwell equations imply for some surface $`S`$ spanning from one side to the other of our boundary or for a volume $`V`$ encompassing a little of both sides of the boundary (see Fig. 2)
$$\begin{array}{ccccc}0\hfill & =\hfill & _S(E)𝑑\sigma =\hfill & _SE_{}𝑑s=\hfill & L(E_{}^+E_{}^{})\text{ hence}\hfill \\ [E_{}]\hfill & \hfill & E_{}^+E_{}^{}=\hfill & 0,\hfill & \text{ and}\hfill \\ 0\hfill & =\hfill & _V(B)𝑑v=\hfill & _VB𝑑\sigma =\hfill & S(B_{}^+B_{}^{})\text{ hence}\hfill \\ [B_{}]\hfill & \hfill & B_{}^+B_{}^{}=\hfill & 0.\hfill & \end{array}$$
(9)
In the same manner, integrating the inhomogeneous Maxwell equations implies for the normal component of $`𝑬`$ and the tangential component of $`𝑩`$
$`[E_{}]E_{}^+E_{}^{}`$ $`=`$ $`4\pi \rho ,`$ (10)
$`[B_{}]B_{}^+B_{}^{}`$ $`=`$ $`4\pi jn.`$ (11)
Similar junction conditions exist also for Einstein’s equations. (Lanczos, 1922, Darmois 1927, Misner & Sharp 1964, Israel 1966, see Refs Is ; La ; Da ; Mi ), the so called junction conditions. To obtain them, we have to split the geometrical quantities into components parallel and transverse to a given hyper–surface. This is often also done in 4d gravity (3+1 formalism or ADM formalism), where one wants to study the time evolution of the metric on a 3d spatial hypersurface. Examples are numerical relativity, or canonical quantization of gravity where the canonical fields are the spatial metric components $`q_{ij}`$ and their canonical momenta are given by the extrinsic curvature, $`\pi _{ij}=K_{ij}`$.
There one considers spacelike hypersurfaces, i.e. , hypersurfaces with timelike normal $`n`$, $`g(n,n)<0`$. In braneworlds we have a timelike hypersurface with spacelike normals, $`g(n,n)>0`$.
Using a slicing of spacetime into 4d hyper surfaces, one can express the 5d Riemann curvature in terms of the 4d one and the extrinsic curvature. The equations are relatively simple if we write them in Gaussian coordinates (Gauss-Codazzi-Mainardi formulas, see e.g. Mi )
$`{}_{}{}^{5}R_{\nu \alpha \beta }^{\mu }`$ $`=`$ $`R_{\nu \alpha \beta }^\mu +K_{\nu \alpha }K_\beta ^\mu K_{\nu \beta }K_\alpha ^\mu \text{(Gauss formula)}`$ (12)
$`{}_{}{}^{5}R_{\mu \nu \alpha }^{4}`$ $`=`$ $`_\nu K_{\mu \alpha }_\alpha K_{\mu \nu }\text{(Codazzi formula)}`$ (13)
$`{}_{}{}^{5}R_{\mu 4\nu }^{4}`$ $`=`$ $`_yK_{\mu \nu }+K_{\mu \beta }K_\nu ^\beta \text{(Mainardi formula)}`$ (14)
$`{}_{}{}^{5}R`$ $`=`$ $`{}_{}{}^{5}R_{}^{\mu \nu }{}_{\mu \nu }{}^{}+2{}_{}{}^{5}R_{}^{4\mu }{}_{4\mu }{}^{}=R+2_yKK^2K_{\alpha \beta }K^{\alpha \beta }.`$ (15)
For the last equation we have used $`_yg^{\mu \nu }=2K^{\mu \nu }`$ so that
$$g^{\mu \nu }_yK_{\mu \nu }=_y(g^{\mu \nu }K_{\mu \nu })(_yg^{\mu \nu })K_{\mu \nu }=_yK2K^{\mu \nu }K_{\mu \nu },K=K_\mu ^\mu .$$
From these we can determine the 5-dimensional Einstein tensor,
$`{}_{}{}^{5}G_{4}^{4}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[R+K^2K^{\mu \alpha }K_{\alpha \mu }\right]`$ (16)
$`{}_{}{}^{5}G_{\mu }^{4}`$ $`=`$ $`_\mu K_\nu K_\mu ^\nu `$ (17)
$`{}_{}{}^{5}G_{\nu }^{\mu }`$ $`=`$ $`G_\nu ^\mu +2K_\alpha ^\mu K_\nu ^\alpha KK_\nu ^\mu +_yK_\nu ^\mu \delta _\nu ^\mu _yK+{\displaystyle \frac{1}{2}}g_{\mu \nu }(K^2+K_{\alpha \beta }K^{\alpha \beta }).`$ (18)
The derivatives wrt $`y`$ indicate that in order to determine the 5-dimensional Einstein tensor, it is not sufficient to know the second fundamental form on the brane itself, but we also have to know it on both sides of the brane.
We now derive equations on the brane from the bulk Einstein equation, $`G_{AB}=\kappa _5T_{AB}`$ which we assume to be valid as a low energy consequence from string theory. To identify the energy momentum tensor on the brane which contains a delta-function in $`y`$-direction, we define
$$S_B^A=\underset{ϵ0}{lim}_{y_bϵ}^{y_b+ϵ}T_B^A𝑑y,$$
so that
$$\underset{ϵ0}{lim}_{y_bϵ}^{y_b+ϵ}{}_{}{}^{5}G_{B}^{A}𝑑y=\kappa _5S_B^A.$$
The 4d metric is continuous and also $`K_{\mu \nu }`$ has no delta-function in $`y`$, but possibly a jump across the brane. This means that $`{}_{}{}^{5}G_{4}^{4}`$ and $`{}_{}{}^{5}G_{\mu }^{4}`$ have no delta-function. Only $`{}_{}{}^{5}G_{\mu \nu }^{}`$ may have one stemming from the term $`_yK_{\mu \nu }g_{\mu \nu }_yK`$ if $`K_{\mu \nu }`$ has a jump. Hence as consequence from the junction conditions, we obtain the following relations for $`S_{AB}`$
$`0`$ $`=`$ $`S_4^4`$ (19)
$`0`$ $`=`$ $`S_\mu ^4\text{ and}`$ (20)
$`\kappa _5S_\mu ^\nu `$ $`=`$ $`\left[K_\mu ^\nu \right]\delta _\mu ^\nu \left[K\right]\text{ or}`$ (21)
$`\left[K_\mu ^\nu \right]`$ $`=`$ $`\kappa _5(S_\mu ^\nu {\displaystyle \frac{1}{3}}\delta _\mu ^\nu S)\text{ where }S=S_\mu ^\mu =S_A^A.`$ (22)
### 3.3 $`Z_2`$ symmetry
In addition to Gaussian normal coordinates (which one can always choose, at least locally) we now assume $`Z_2`$ symmetry: the two sides of the brane are mirror images. For the metric components this implies
$`g_{\mu \nu }(x^\lambda ,y_b+y)`$ $`=`$ $`g_{\mu \nu }(x^\lambda ,y_by)`$ (23)
$`g_{\mu 4}(x^\lambda ,y_b+y)`$ $`=`$ $`g_{\mu 4}(x^\lambda ,y_by)`$ (24)
$`g_{44}(x^\lambda ,y_b+y)`$ $`=`$ $`g_{44}(x^\lambda ,y_by)`$ (25)
The same symmetry is required for $`T_{AB}`$ and any other bulk tensor field.
Under this condition we have $`K_+=K_{}`$ so that the junction conditions (22) reduce to
$$2K_\nu ^\mu =\kappa _5(S_\nu ^\mu \frac{1}{3}\delta _\nu ^\mu S).$$
(26)
If the braneworld satisfies $`Z_2`$ symmetry, the brane energy momentum tensor determines the second fundamental form.
However, we now show that even with $`Z_2`$ symmetry, knowing the brane energy momentum tensor is not sufficient to determine the brane Einstein tensor. To demonstrate this we first rewrite the 4d Einstein tensor, for a general coordinate system in terms of the 5d Riemann tensor and the extrinsic curvature:
$`G_{\mu \nu }`$ $`=`$ $`{}_{}{}^{5}G_{AB}^{}q_\mu ^Aq_\nu ^B+^5R_{AB}n^An^Bg_{\mu \nu }K_\mu ^\alpha K_{\alpha \nu }+KK_{\mu \nu }`$ (27)
$`+{\displaystyle \frac{1}{2}}g_{\mu \nu }(K^2K_{\alpha \beta }K^{\alpha \beta })\stackrel{~}{E}_{\mu \nu }.`$
Here $`{}_{}{}^{5}R_{AB}^{}n^An^B`$ corresponds to $`{}_{}{}^{5}R_{44}^{}={}_{}{}^{5}R_{4\mu 4}^{\mu }`$ in Gaussian coordinates and $`\stackrel{~}{E}_{\mu \nu }{}_{}{}^{5}R_{ABCD}^{}n^An^Cq_\nu ^Dq_\mu ^B`$ corresponds to $`{}_{}{}^{5}R_{4\mu 4\nu }^{}`$. Eq. (27) in Gaussian coordinates is a simple consequence of the expressions (12,13,14) for the Riemann tensor. In a general coordinate system it can be found in Ref. Shiru ; RoyRev (careful, the sign for $`K_{\mu \nu }`$ is different there). In 5 dimensions the Weyl tensor is given by
$$R_{ABCD}=\frac{2}{3}(g_{A[C}R_{D]B}g_{B[C}R_{D]A})\frac{1}{6}g_{A[C}g_{D]B}R+C_{ABCD}.$$
(28)
Here $`[AB]`$ indicates anti-symmetrization in the indices $`A`$ and $`B`$, and $`C_{ABCD}`$ is the 5-dimensional Weyl tensor defined by Eq. (28). It is easy to verify that $`C_{ABCD}`$ is traceless and obeys the same symmetries as the Riemann tensor, $`R_{ABCD}`$.
Inserting the 5-dimensional Einstein Eq. (3) for $`{}_{}{}^{5}G_{AB}^{}`$ and Eq. (28) in the expression containing the Ricci tensor $`R_{AB}`$, as well as in $`\stackrel{~}{E}_{\mu \nu }`$, we obtain the 4d brane gravity equation
$`G_{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{2}{3}}\kappa _5\left[T_{AB}q_\mu ^Aq_\nu ^B+g_{\mu \nu }(T_{AB}n^An^B{\displaystyle \frac{1}{4}}T)\right]K_\mu ^\alpha K_{\alpha \nu }+KK_{\mu \nu }`$ (29)
$`+{\displaystyle \frac{1}{2}}g_{\mu \nu }\left(K_{\alpha \beta }K^{\alpha \beta }K^2\right)E_{\mu \nu }`$
where
$$E_{\mu \nu }=C_{ABCD}n^An^Cq_\nu ^Dq_\mu ^B$$
(30)
is the ’projection’ of the Weyl tensor along the brane normal. The Codazzi equation (13) gives in addition
$$\kappa _5T_{AB}n^Aq_\mu ^B=_\mu K_\nu K_\mu ^\nu .$$
(31)
Because of the last term in Eq. (29), it is not sufficient to know the bulk energy momentum tensor and initial conditions for $`n`$, $`g_{AB}`$ and $`K_{AB}`$ to solve the gravitational equations on the brane. In addition we need to know $`E_{\mu \nu }`$, components of the bulk Weyl tensor on the brane. The Weyl tensor, which is the part of the curvature which can be non-vanishing even if the energy momentum tensor vanishes, contains information on bulk gravity waves. Bulk gravity waves can flow onto, respectively be emitted from the brane and thereby affect its evolution. This information is encoded only in the full bulk initial conditions. Therefore, to determine the evolution of the brane matter and geometry, in principle we have to solve the full bulk equations! Only in situations with very special symmetries this can be avoided. However, as soon as we want to perturb such symmetric solutions we have to take into account all the bulk modes, and we do expect the solutions to differ significantly from the results of 4-dimensional perturbation theory.
### 3.4 Brane gravity with an empty bulk
We now exemplify the effect of the bulk Weyl tensor in the case of an empty bulk. This will be the situation which we study for the rest of these lectures. We assume that the bulk is empty up to a simple cosmological constant $`\mathrm{\Lambda }_5`$.
$$T_{AB}=\frac{\mathrm{\Lambda }_5}{\kappa _5}g_{AB}+q_A{}_{}{}^{\mu }q_{B}^{}{}_{}{}^{\nu }S_{\mu \nu }^{}\delta (yy_b).$$
(32)
where $`S_{\mu \nu }`$ is the energy momentum tensor on the brane. It consists of a brane tension $`\lambda `$ and the matter energy momentum tensor $`\tau _{\mu \nu }`$.
$$S_{\mu \nu }=\lambda g_{\mu \nu }+\tau _{\mu \nu }.$$
(33)
The junction conditions read
$`\left[g_{\mu \nu }\right]`$ $`=`$ $`0\text{ first junction condition,}`$ (34)
$`\left[K_{\mu \nu }\right]`$ $`=`$ $`\kappa _5\left(S_{\mu \nu }{\displaystyle \frac{1}{3}}g_{\mu \nu }S\right)\text{ second junction condition.}`$ (35)
$`Z_2`$ symmetry requires that
$$K_{\mu \nu }^+=K_{\mu \nu }^{}=\frac{\kappa _5}{2}\left(S_{\mu \nu }\frac{1}{3}g_{\mu \nu }S\right).$$
(36)
Inserting our ansatz for $`T_{AB}`$ in the brane gravity equation (29) and using the second junction condition to eliminate the second fundamental form, we obtain
$$G_{\mu \nu }=\mathrm{\Lambda }_4g_{\mu \nu }+\kappa _4\tau _{\mu \nu }+\kappa _5^2\sigma _{\mu \nu }E_{\mu \nu },$$
(37)
with
$`\mathrm{\Lambda }_4`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\mathrm{\Lambda }_5+{\displaystyle \frac{\kappa _5^2}{6}}\lambda ^2),`$ (38)
$`\kappa _4`$ $`=`$ $`\kappa _5^2\lambda /6=2/M_4^2`$ (39)
$`\sigma _{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{1}{4}}\tau _{\mu \alpha }\tau _\nu ^\alpha +{\displaystyle \frac{1}{12}}\tau \tau _{\mu \nu }+{\displaystyle \frac{1}{8}}g_{\mu \nu }\tau _{\alpha \beta }\tau ^{\alpha \beta }{\displaystyle \frac{1}{24}}g_{\mu \nu }\tau ^2,`$ (40)
and, as before $`E_{\mu \nu }`$ denotes the projected 5d Weyl tensor, evaluated on either side of the brane (but not exactly on the brane where it may be ill-defined). The quantities $`\kappa _4`$ and $`M_4`$ denote the 4–dimensional gravitational coupling constant and Planck mass respectively and $`\tau =\tau _\mu ^\mu `$ is the trace of the matter energy momentum tensor. The relation between the 4- and 5-dimensional Planck mass in the braneworld approach is now obtained as follows: using $`\kappa _5M_5^3`$ and $`\kappa _4M_4^2`$, Eq. (39) shows that $`M_4^2M_5^6/\lambda `$. Using that the 4-dimensional cosmological constant is small, $`\mathrm{\Lambda }_4|\mathrm{\Lambda }_5|`$, we have $`\lambda \sqrt{6\mathrm{\Lambda }_5}/\kappa _5\sqrt{6\mathrm{\Lambda }_5}M_5^3`$, so that $`M_4^2M_5^3\sqrt{\mathrm{\Lambda }_5}=M_5^3L`$ with $`L\sqrt{\mathrm{\Lambda }_5}`$.
In the limit $`\lambda \tau _{\mu \nu }\tau _{\mu \alpha }\tau _\nu ^\alpha `$ we recover the 4–dimensional Einstein equation if $`E_{\mu \nu }`$ is negligible. The existence of this limit depends crucially on the existence of a 4–dimensional brane tension. In order for the 4d gravitational coupling constant to be positive, the brane tension must be positive, $`\lambda >0`$. At high energy densities (in the early universe) the quadratic term $`\sigma _{\mu \nu }`$ can become dominant and modify the dynamics (the expansion law of the universe). In general, there is an additional part, $`E_{\mu \nu }`$ , carrying information from the bulk geometry and evolution, which can affect the brane evolution in a crucial way.
### 3.5 Energy momentum conservation
The Codazzi equation (13) together with $`Z_2`$ symmetry implies
$$_\mu K_\nu K_\mu ^\nu =\kappa _5T_\mu ^4=0.$$
(41)
With the Gauss equation (12) and $`Z_2`$ symmetry this ensures energy and momentum conservation on the brane,
$$_\nu \tau _\mu ^\nu =0.$$
(42)
From the 4-dimensional contracted Bianchi identities we obtain in addition
$$_\nu E_\mu ^\nu =\kappa _5^2_\nu \sigma _\mu ^\nu .$$
(43)
Hence, the longitudinal part of $`E_{\mu \nu }`$ is fully determined by the matter content of the brane, while the transverse traceless part is not specified:
$$E_{\mu \nu }=E_{\mu \nu }^{(TT)}+E_{\mu \nu }^{(L)}$$
(44)
where $`_\nu E_\mu ^{\nu (TT)}=0`$ and
$$E_{\mu \nu }^{(L)}=\frac{1}{2}(_\mu \theta _\nu +_\nu \theta _\mu )\text{ with }_\mu \theta ^\mu =0.$$
Inserting this in Eq. (43), we obtain
$$^2\theta _\mu =\frac{\kappa _5^2}{2}\left[\tau _{\alpha \beta }(_\mu \tau ^{\alpha \beta }+^\alpha \tau _\mu ^\beta )+\frac{1}{3}(_\alpha \tau )(\tau _\mu ^\alpha q_\mu ^\alpha \tau )\right].$$
(45)
For given initial conditions, this equation has always a unique solution $`\theta _\mu `$ on the brane which determines $`E_{\mu \nu }^{(L)}`$. However, the transverse part, $`E_{\mu \nu }^{(TT)}`$ is not determined by the brane energy momentum tensor; it comes from bulk gravity waves. Only if $`E_{\mu \nu }^{(TT)}=0`$ does the brane energy momentum tensor determine the brane Einstein tensor. As we shall see, already in quite simple situations this is not the case.
## 4 The Randall Sundrum model
We now consider an Anti-de Sitter (AdS) bulk, $`\mathrm{\Lambda }_5<0`$ and would like to obtain Minkowski space on the brane. Since AdS is conformally flat, $`E_{\mu \nu }=0`$. A Minkowski brane can be achieved by setting $`\tau _{\mu \nu }=0`$. If in addition the brane tension is related to the 5–dimensional coupling constant and the cosmological constant by
$$\lambda ^2\kappa _5^2/6=\mathrm{\Lambda }_5,$$
(46)
Eq. (37) implies $`G_{\mu \nu }=0`$. Eq. (46) is it the Randall–Sundrum (RS) fine tuning condition RS1 ; RS2 . Small deviations from the RS condition lead to an exponentially expanding/contracting brane. The 4-dimensional gravitational constant becomes
$$\kappa _4=\frac{\lambda \kappa _5^2}{6}=\frac{\mathrm{\Lambda }_5}{\lambda }>0,\text{ or, equivalently }\lambda =\frac{\mathrm{\Lambda }_5}{\kappa _4},\text{ and }\kappa _4=\kappa _5\sqrt{\frac{\mathrm{\Lambda }_5}{6}}.$$
(47)
### 4.1 The metric
The following coordinates for (a part of) Anti-de Sitter will be useful for us:
$`ds^2`$ $`=`$ $`e^{2|z|/\mathrm{}}\eta _{\mu \nu }dx^\mu dx^\nu +dz^2\text{ Gaussian coordinates}`$ (48)
$`ds^2`$ $`=`$ $`\left({\displaystyle \frac{\mathrm{}}{y}}\right)^2\left(\eta _{\mu \nu }dx^\mu dx^\nu +dy^2\right)|y|>\mathrm{}\text{ conformal coordinates.}`$ (49)
Einstein’s equations, $`G_{AB}=\mathrm{\Lambda }_5g_{AB}`$, give $`\mathrm{\Lambda }_5=\frac{6}{\mathrm{}^2}`$. The RS fine tuning requires $`\lambda =\sqrt{6\mathrm{\Lambda }_5/\kappa _5^2}`$.
In their first model RS1 (RS1 model) Randall and Sundrum propose two branes, the first positioned at $`z=0`$ is called the hidden brane, and the second positioned at $`z=k\mathrm{}`$ is called the visible brane and represents our Universe. The gravitational force on the visible brane is suppressed by the factor $`\mathrm{exp}(2k)`$ w.r.t. the hidden brane, leading to an enhancement by a factor $`\mathrm{exp}(k)`$ of the apparent Planck mass. However, also this two brane model contains a gravi-scalar (also called radion) which has to obtain a mass by some non-gravitational mechanism.
This problem is resolved in the second model RS2 (RS2 model). There, our universe is located on the brane at $`z=0`$ corresponding to $`y=\mathrm{}`$ and no second brane is present. The apparent 4-dimensional Planck mass as measured on the brane on scales much larger than $`\mathrm{}`$ is then again given by
$`\lambda `$ $`=`$ $`{\displaystyle \frac{\kappa _4}{\mathrm{\Lambda }_5}}={\displaystyle \frac{3M_4^2}{\mathrm{}^2}}.\text{ Furthermore,}`$ (50)
$`\kappa _4^2`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Lambda }}{6}}\kappa _5^2=\mathrm{}^2\kappa _5^2\text{ so that}M_4^2=M_5^3\mathrm{}.`$ (51)
Hence in the RS2 model, the AdS curvature scale $`\mathrm{}`$ enters in the same way as $`L`$ for cylindric Kaluza-Klein models. Since the scale $`\mathrm{}`$ is limited by present day micro gravity experiments which have not detected any deviation from Newton’s law micGra , we have
$$\mathrm{}<0.1mm\text{ implying }\lambda >(1TeV)^4,\text{ hence }M_5=(M_4^2/\mathrm{})^{1/3}>10^5TeV.$$
### 4.2 Gravity waves in the RS2 model
In order to see that the radion mode is absent in the non-compact RS2 model, we consider perturbations to the AdS metric. It is easy to show that one can always choose a gauge (local coordinate system) so that the perturbed metric is of the form
$`ds^2`$ $`=`$ $`\left({\displaystyle \frac{\mathrm{}}{y}}\right)^2[(1+2\mathrm{\Psi })dt^24\mathrm{\Sigma }_idtdx^i+((12\mathrm{\Phi })\delta _{ij}+2H_{ij})dx^idx^j+4\mathrm{\Xi }_idydx^i`$ (53)
$`4dtdy+(1+2𝒞)dy^2]`$
Here $`H_{ij}`$ and $`\mathrm{\Sigma }_i`$, $`\mathrm{\Xi }_i`$ are transverse (i.e. divergence free) and $`H_{ij}`$ is traceless. In other words, $`_i\mathrm{\Sigma }_i=_i\mathrm{\Xi }_i=_iH_{ij}=H_i^i=0`$. As we shall see there are 5 homogeneous modes in these 10 physical perturbation variables corresponding to the 5 gravity wave modes in 5 dimensions. We consider the homogeneous 5d wave equation in the bulk.
Since the 3-brane is homogeneous, scalar-, vector- and tensor degrees of freedom decouple and we can consider them in turn (for more details see Section 5).
#### 4.2.1 Tensor perturbations
We first consider only the tensor $`H_{ij}`$, so that the perturbed metric is given by
$$ds^2=\left(\frac{\mathrm{}}{y}\right)^2\left(dt^2+(\delta _{ij}+2H_{ij})dx^idx^j+dy^2\right).$$
(54)
We Fourier transform $`H_{ij}`$ in the 3-dimensional $`𝐱`$ coordinates and consider one mode with fixed wave vector $`𝐤`$, so that $`H_{ij}(t,y,𝐱)=H_{ij}(t,y)\mathrm{exp}(i𝐤𝐱)`$. Since spacetime is isotropic and homogeneous in $`𝐱`$, different $`𝐤`$–modes do not couple. The bulk Einstein equations, $`\delta G_{AB}=\mathrm{\Lambda }\delta g_{AB}`$, for the Fourier mode $`k`$ then give
$$\left(_t^2+k^2_y^2+\frac{3}{y}_y\right)H_{ij}=0.$$
(55)
The general solution to this equation is of the form
$`H_{ij}=h_me_{ij}\text{ with }h_m`$ $`=`$ $`e^{i\omega t}(my)^2\left[AJ_2(my)+BY_2(my)\right].`$ (56)
Here $`e_{ij}`$ is the polarization tensor, $`k^ie_{ij}=e_j^i=0`$ and $`\omega ^2=m^2+k^2`$. The separation constant $`m^2`$ is arbitrary and can, in principle also be negative. $`J_2`$ and $`Y_2`$ are the Bessel functions of order 2. They are oscillating and decaying. Bessel functions represent “$`\delta `$–function normalizable” perturbations like harmonic waves in flat space, in the sense that Csaki:2000fc ; Bozza:2001xt
$$_0^{\mathrm{}}h_mh_m^{}\frac{dy}{m^2y^3}=m\delta (mm^{}).$$
(57)
These are just the ordinary gravity modes of $`4`$-dimensional mass $`m`$ without a mass gap which are discussed in the original RS paper RS2 . To find the correct weight $`1/y^3`$, we use that $`h_m`$ satisfies
$$(\mathrm{}{}_{4}{}^{}+_y^2\frac{3}{y}_y)h_m=0,$$
(58)
and thus $`\stackrel{ˇ}{h}=h_m/y^{3/2}`$ satisfies the equation of motion of a scalar field in a flat 5-dimensional spacetime (with $`y`$–dependent mass term),
$$(\mathrm{}{}_{4}{}^{}+_y^2\frac{15}{4y^2})\stackrel{ˇ}{h}=0.$$
(59)
This mode has to be normalizable w.r.t the Minkowski metric (no additional weight).
As we have mentioned above, $`m^2`$ is arbitrary and can also be chosen negative. However, if $`m^2<0`$ and therefore $`m`$ is imaginary, it is more useful to decompose the two independent solutions in the form
$$h_m=e^{i\omega t}(|m|y)^2\left[CK_2(|m|y)+DI_2(|m|y)\right],$$
(60)
where $`K_2`$ and $`I_2`$ are the modified Bessel functions of order 2. Considering the behavior of the Bessel functions, one sees that $`I_2`$ grows exponentially (see Fig. 5) and is clearly not normalizable (i.e. not square integrable with some weight which is a power law in $`y`$). Therefore, this mode is unphysical and we have to set $`D=0`$. It is important to note that $`\omega ^2=k^2+m^2`$ can become negative in this case leading to $`\omega =\pm i|\omega |`$. In other words, negative mass solutions become exponentially growing ’tachyonic’ instabilities! It is still unclear whether these tachyonic modes are relevant for cosmological braneworlds.
In the limit $`m0`$ only the $`Y_2`$–mode survives and we obtain $`h_{(m=0)}=C\mathrm{exp}(\omega t)`$ independent of the coordinate $`y`$. This zero-mode is normalizable with respect to the measure $`dy/y^3`$.
The general solution for a tensor perturbation is of the form
$$h=h_0+_{\mathrm{}}^{\mathrm{}}h_m𝑑m^2.$$
(61)
At the brane position, $`y=y_\mathrm{b}=\mathrm{}`$, the perturbations must satisfy the junction condition (36). These represent boundary conditions for the perturbations $`H_{ij}`$ in the bulk. On the right hand side of Eq. (36), we can in principle have an arbitrary perturbation of the matter energy momentum tensor. However, the only non-vanishing term of the tensor contribution to $`\tau _{\mu \nu }`$ is the traceless part of $`\tau _{ij}`$, i.e. the anisotropic stress on the brane, $`\mathrm{\Pi }_{ij}^{(T)}`$. A short computation shows
$`\delta K_{ij}|_{y_\mathrm{b}}`$ $`=`$ $`\left({\displaystyle \frac{2}{\mathrm{}}}H_{ij}_yH_{ij}\right)|_{y_\mathrm{b}},\text{hence}`$
$`2(_yH_{ij})|_{y_\mathrm{b}}`$ $`=`$ $`\kappa __5\mathrm{\Pi }_{ij}^{(T)},`$ (62)
where $`\mathrm{\Pi }^{(T)}`$ are tensor–type anisotropic stresses on the brane.
Let us first consider the homogeneous case $`\mathrm{\Pi }^{(T)}0`$. For $`m^2>0`$, the solutions are of the form
$$h=\mathrm{exp}(\pm i\omega t)(my)^2\left[AJ_2(my)+BY_2(my)\right].$$
(63)
The junction condition (62) then requires
$$B=A\frac{J_1(m\mathrm{})}{Y_1(m\mathrm{})}\frac{\pi }{4}(m\mathrm{})^2A,$$
(64)
where the last expression is a good approximation for $`m\mathrm{}1`$. This is precisely the result of Randall and Sundrum RS2 . It is not modified even if we allow for the negative mass modes, $`m^2>0`$, because a physical solution has to be of the form
$$h=C\mathrm{exp}(\pm t\sqrt{m^2k^2})(|m|y)^2K_2(|m|y),\text{ with }_yh=|m|C\mathrm{exp}(\pm t\sqrt{m^2k^2})(|m|y)^2K_1(|m|y),$$
(65)
and since $`K_1`$ has no zero, the junction condition (62) requires $`C=0`$.
But in a realistic brane universe, $`\mathrm{\Pi }^{(T)}`$ is not exactly zero. In cosmology, it is typically just a factor 10 smaller than other perturbations of the energy momentum tensor on the brane. We therefore cannot require $`C0`$. However, as long as $`\mathrm{\Pi }^{(T)}`$ remains small, we do not expect the unstable modes to be present, so that $`C(k,m)=0`$ for $`k^2<m^2`$. Within the framework of first order perturbation theory, the $`\mathrm{\Pi }^{(T)}`$ modes satisfy a Minkowski equation of motion and therefore they do not grow exponentially. Hence in this case, the $`K`$–mode can only be excited for $`\omega ^2=k^2+m^2>0`$, i.e. $`k^2>m^2`$. However, it is not clear whether this remains true to second order, where the evolution of $`H`$ feeds back in the equation of motion for $`\mathrm{\Pi }^{(T)}`$. Actually, within a toy model, it has been shown that the full, non-linear evolution can be exponentially unstable even if the linear equations do not excite the unstable mode CD4 .
#### 4.2.2 Vector perturbations
We now consider vector perturbations only, so that, in generalized longitudinal gauge, the metric takes the form
$`ds^2`$ $`=`$ $`\left({\displaystyle \frac{\mathrm{}}{y}}\right)^2\left(dt^24\mathrm{\Sigma }_idtdx^i+\delta _{ij}dx^idx^j+4\mathrm{\Xi }_idydx^i+dy^2\right)`$ (66)
The bulk Einstein equations for a mode $`𝐤`$ of the vector perturbations $`𝚺`$ and $`𝚵`$ are
$`\left(_y^2{\displaystyle \frac{3}{y}}_y\right)𝚺`$ $`=\left(_t^2+k^2\right)𝚺,`$ (67)
$`\left(_y^2{\displaystyle \frac{3}{y}}_y+{\displaystyle \frac{3}{y^2}}\right)𝚵`$ $`=\left(_t^2+k^2\right)𝚵,`$ (68)
$`\left(_y{\displaystyle \frac{3}{y}}\right)𝚵`$ $`=_t𝚺,`$ (69)
where $`𝚺`$ and $`𝚵`$ are transverse vectors, $`(𝐤𝚺)=(𝐤𝚵)=0`$. The constraint equation (69) fixes the relative amplitudes of $`\mathrm{\Sigma }`$ and $`\mathrm{\Xi }`$, showing that there is only one independent vector perturbation in the bulk (the “gravi-photon”). One can check that these equations are consistent, e.g. with the master function approach of Ref. Mukohyama:2000ui .
As in the tensor case, the solutions are Bessel functions of order two (and one). Considering just one component $`\mathrm{\Sigma }=\mathrm{\Sigma }_i`$ one obtains the expected oscillatory modes for positive mass-squared, $`m^2>0`$,
$`\mathrm{\Sigma }`$ $`=\mathrm{exp}(\pm i\omega t)(my)^2\left[AJ_2(my)+BY_2(my)\right],`$ (70)
$`\mathrm{\Xi }`$ $`={\displaystyle \frac{\pm i\omega }{m}}\mathrm{exp}(\pm i\omega t)(my)^2\left[AJ_1(my)+BY_1(my)\right],`$ (71)
where $`\omega =\sqrt{m^2+k^2}`$. These solutions have been found in Ref. Bridgman:2000ih . For a negative mass-square, $`m^2<0`$, we obtain again tachyonic solutions. Like in the tensor case, the solution containing the modified Bessel function $`I_\nu `$ cannot be accepted as it is exponentially growing and thus represents a non-normalizable mode. However, the $`K_\nu `$-solution is exponentially decaying and perfectly acceptable. For tachyonic vector perturbations with $`\omega ^2=m^2+k^2<0`$ we have
$`\mathrm{\Sigma }`$ $`=C\mathrm{exp}(\pm |\omega |t)(|m|y)^2K_2(|m|y),`$ (72)
$`\mathrm{\Xi }`$ $`={\displaystyle \frac{\pm |\omega |}{|m|}}C\mathrm{exp}(\pm |\omega |t)(|m|y)^2K_1(|m|y).`$ (73)
For large enough scales, $`m^2>k^2`$, these solutions again grow exponentially.
The boundary conditions at the brane relate these perturbations to the brane energy momentum tensor. For the energy momentum tensor on the brane, the vector degrees of freedom are defined according to
$$\left(S_{\mu \nu }\right)=\left(\begin{array}{cc}0& V_j\\ V_i& \mathrm{\Pi }_{ij}^{(V)}\end{array}\right)\lambda \left(q_{\mu \nu }\right),$$
(74)
where $`V_i`$ and $`\mathrm{\Pi }_i^{(V)}`$ are divergence-free vector fields and $`\mathrm{\Pi }_{ij}^{(V)}[_j\mathrm{\Pi }_i^{(V)}+_i\mathrm{\Pi }_j^{(V)}]`$. The first junction condition simply requires that $`\mathrm{\Sigma }`$ be continuous at the brane, which it is since the (modified) Bessel functions of even index are even functions. The second junction condition results in (for a detailed derivation, see Ringeval:2003na )
$`_t\mathrm{\Xi }+_y\mathrm{\Sigma }`$ $`=\kappa __5V,`$ (75)
$`\mathrm{\Xi }`$ $`=\kappa __5\mathrm{\Pi }^{(V)},`$ (76)
$`_tV`$ $`=k^2\mathrm{\Pi }^{(V)}.`$ (77)
The last equation follows from (75) and (76) and the bulk equations (67)–(69). It represents momentum conservation on the brane, which is guaranteed as long as we have vanishing energy flux off the brane and $`Z_2`$–symmetry.
Like for tensor perturbations, we consider homogeneous solutions, setting $`\mathrm{\Pi }^{(V)}V0`$. This requires $`\mathrm{\Xi }(|m|\mathrm{})=0`$, hence
$`B`$ $`=A{\displaystyle \frac{J_1(m\mathrm{})}{Y_1(m\mathrm{})}}\text{ for }m^2>0,`$ (78)
$`C`$ $`0\text{ for }m^2<0.`$ (79)
Equation (75) is then identically satisfied.
However, it seems more realistic to allow a small but non-vanishing anisotropic stress contribution $`\mathrm{\Pi }^{(V)}`$ and corresponding vorticity $`V`$. In this case, again, we can have solutions with $`C0`$ which can grow exponentially in time; hence small initial data can lead to an exponential instability like for tensor perturbations.
Using the normalization condition (57) for the $`m=0`$ mode of the variable $`\mathrm{\Xi }y`$ (this is the one which enters as dynamical variable in the perturbed action, see Csaki:1999jh ), one finds that $`|\mathrm{\Xi }|^2/y^3𝑑y`$ diverges logarithmically. Contrary to the tensor case, the vector zero-mode is not normalizable. Therefore, on the brane there is only the ordinary massless spin–2 graviton, but there are a continuous infinity of massive spin–2 and spin–1 particles (the modes discussed here, with $`m0`$).
#### 4.2.3 Scalar perturbations
We now discuss the most cumbersome, the scalar sector. Scalar–type metric perturbations in the bulk are of the form
$`ds^2`$ $`={\displaystyle \frac{\mathrm{}^2}{y^2}}[(1+2\mathrm{\Psi })dt^24dtdy`$
$`+(12\mathrm{\Phi })\delta _{ij}dx^idx^j+(1+2𝒞)dy^2].`$ (80)
The bulk Einstein perturbation equations for the mode $`𝐤`$ become, after some manipulations and introducing the combination $`\mathrm{\Gamma }\mathrm{\Phi }+\mathrm{\Psi }`$ (see Ref. CD4 ),
$`\mathrm{\Phi }\mathrm{\Psi }`$ $`=𝒞,`$ (81)
$`\left(_y^2{\displaystyle \frac{3}{y}}_y\right)\mathrm{\Gamma }`$ $`=\left(_t^2+k^2\right)\mathrm{\Gamma },`$ (82)
$`\left(_y^2{\displaystyle \frac{3}{y}}_y+{\displaystyle \frac{4}{y^2}}\right)𝒞`$ $`=\left(_t^2+k^2\right)𝒞,`$ (83)
$`_y\mathrm{\Phi }+\left(_y{\displaystyle \frac{3}{y}}\right)𝒞`$ $`=_t,`$ (84)
$`{\displaystyle \frac{3}{y}}\left(_y{\displaystyle \frac{2}{y}}\right)𝒞`$ $`=3_t^2\mathrm{\Phi }+k^2(\mathrm{\Phi }+𝒞),`$ (85)
$`3_t\left(_y\mathrm{\Phi }{\displaystyle \frac{1}{y}}𝒞\right)`$ $`=k^2,`$ (86)
$`_t\left(2\mathrm{\Phi }𝒞\right)`$ $`=\left(_y{\displaystyle \frac{3}{y}}\right).`$ (87)
Clearly these equations are not all independent, Eqs. (86)–(87) are identically satisfied if Eqs. (81)–(85) are. The solutions are obtained as for tensor and vector perturbations. For a positive mass-square, $`m^2>0`$, we find ($`\omega =\sqrt{m^2+k^2}`$)
$`\mathrm{\Gamma }`$ $`=\mathrm{exp}(\pm i\omega t)(my)^2\left[A^{}J_2(my)+B^{}Y_2(my)\right],`$ (88)
$`𝒞`$ $`=\mathrm{exp}(\pm i\omega t)(my)^2\left[AJ_0(my)+BY_0(my)\right],`$ (89)
$`\mathrm{\Phi }`$ $`={\displaystyle \frac{1}{2}}\mathrm{exp}(\pm i\omega t)(my)^2[A^{}J_2(my)+B^{}Y_2(my)`$
$`+AJ_0(my)+BY_0(my)],`$ (90)
$`\mathrm{\Psi }`$ $`={\displaystyle \frac{1}{2}}\mathrm{exp}(\pm i\omega t)(my)^2[A^{}J_2(my)+B^{}Y_2(my)`$
$`AJ_0(my)BY_0(my)],`$ (91)
$``$ $`={\displaystyle \frac{\pm im^3y^2}{2\omega }}\mathrm{exp}(\pm i\omega t)\times `$
$`\left[(A^{}3A)J_1(my)+(B^{}3B)Y_1(my)\right],`$ (92)
$`\text{with }A^{}`$ $`=3A{\displaystyle \frac{m^2}{m^2+2\omega ^2}},\text{ and }B^{}=3B{\displaystyle \frac{m^2}{m^2+2\omega ^2}}.`$ (93)
For a negative mass-square, $`m^2<0`$, we obtain ($`\omega =\sqrt{m^2k^2}`$)
$`\mathrm{\Gamma }`$ $`=\mathrm{exp}(\pm \omega t)(|m|y)^2C^{}K_2(|m|y),`$ (94)
$`𝒞`$ $`=\mathrm{exp}(\pm \omega t)(|m|y)^2CK_0(|m|y),`$ (95)
$`\mathrm{\Phi }`$ $`={\displaystyle \frac{1}{2}}\mathrm{exp}(\pm \omega t)(|m|y)^2\times `$
$`\left[C^{}K_2(|m|y)+CK_0(|m|y)\right],`$ (96)
$`\mathrm{\Psi }`$ $`={\displaystyle \frac{1}{2}}\mathrm{exp}(\pm \omega t)(|m|y)^2\times `$
$`\left[C^{}K_2(|m|y)CK_0(|m|y)\right],`$ (97)
$``$ $`={\displaystyle \frac{\pm |m|^3y^2}{2\omega }}\mathrm{exp}(\pm \omega t)[C^{}+3C]K_1(|m|y),`$ (98)
$`\text{with }C^{}`$ $`=3C{\displaystyle \frac{|m|^2}{|m|^2+2\omega ^2}},`$ (99)
where we have already used that the $`I`$–mode is not normalizable and therefore cannot contribute. Like for vector and tensor perturbations, we find again tachyonic solutions with $`m^2<0`$ which represent an exponential instability for sufficiently small wave numbers $`k`$ (large scales).
Determining the boundary conditions via the first and second junction conditions now requires a bit more care. Since we have already fully specified our coordinate system by the adopted choice of perturbation variables, we must allow for brane bending. We cannot fix the brane at $`y_\mathrm{b}=\mathrm{}`$, but we must allow for $`y_\mathrm{b}^+=\mathrm{}+`$ and $`y_\mathrm{b}^{}=\mathrm{}`$, respectively. Fortunately, $``$ is a scalar quantity and brane bending therefore does not affect vector and tensor perturbations. The anti-symmetry $`y_\mathrm{b}^+=y_\mathrm{b}^{}`$ is an expression of $`Z_2`$–symmetry. The introduction of the new perturbation variable $`(x^\mu )`$ describing brane bending enters the exressions for the first and second fundamental forms. From Eq. (5), we obtain $`q_{\mu \nu }=g_{\mu \nu }`$ to first order, which implies that $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$, hence $`𝒞`$, have to be continuous. At the brane position, the perturbed components of the extrinsic curvature (8) are
$`\delta K_{00}`$ $`={\displaystyle \frac{1}{\mathrm{}}}\left[\mathrm{\Phi }3\mathrm{\Psi }+2{\displaystyle \frac{}{\mathrm{}}}\right]+_y\mathrm{\Psi }2_t+_t^2,`$ (100)
$`\delta K_{0j}`$ $`=_j\left(_t\right),`$ (101)
$`\delta K_{ij}`$ $`=\left[{\displaystyle \frac{1}{\mathrm{}}}\left(\mathrm{\Psi }3\mathrm{\Phi }2{\displaystyle \frac{}{\mathrm{}}}\right)+_y\mathrm{\Phi }\right]\delta _{ij}+_i_j.`$ (102)
For the energy momentum tensor on the brane, we parameterize the four degrees of freedom according to
$$\left(S_{\mu \nu }\right)=\left(\begin{array}{cc}\rho & v_j\\ v_i& p\delta _{ij}+\mathrm{\Pi }_{ij}^{(S)}\end{array}\right)\lambda \left(q_{\mu \nu }\right),$$
(103)
where $`v_i_iv`$ and $`\mathrm{\Pi }_{ij}^{(S)}\left(_i_j\frac{1}{3}\mathrm{\Delta }\delta _{ij}\right)\mathrm{\Pi }^{(S)}`$. With Eqs. (100)–(102), the second junction condition reads
$`{\displaystyle \frac{1}{\lambda }}\left(2\rho +3p\right)`$ $`=\mathrm{\Phi }\mathrm{\Psi }+\mathrm{}_t\left(_t2\right)+L_y\mathrm{\Psi },`$ (104)
$`{\displaystyle \frac{3}{\lambda \mathrm{}}}v`$ $`=_t,`$ (105)
$`{\displaystyle \frac{3}{\lambda \mathrm{}}}\mathrm{\Pi }^{(S)}`$ $`=,`$ (106)
$`{\displaystyle \frac{1}{\lambda }}\left[\rho \mathrm{\Delta }\mathrm{\Pi }^{(S)}\right]`$ $`=\mathrm{\Psi }\mathrm{\Phi }+\mathrm{}_y\mathrm{\Phi }.`$ (107)
Combining the time derivative of Eq. (105) with Eqs. (104), (84) and (107), we obtain momentum conservation on the brane,
$$_tv=\frac{2}{3}\mathrm{\Delta }\mathrm{\Pi }^{(S)}+p.$$
(108)
Similar manipulations imply energy conservation on the brane,
$$_t\rho =\mathrm{\Delta }v.$$
(109)
Like for tensor and vector perturbations, we look for solutions with vanishing brane matter. Setting $`\mathrm{\Pi }^{(S)}\rho Pv0`$ forbids brane bending, $`=0`$. Then Eq. (105) implies $`(m\mathrm{})=0`$, thus
$`B^{}3B`$ $`=(A^{}3A){\displaystyle \frac{J_1(m\mathrm{})}{Y_1(m\mathrm{})}}\text{ for }m^2>0,`$ (110)
$`C^{}+3C`$ $`=0\text{ for }m^2<0.`$ (111)
The other equations are all satisfied if we require separately
$`{\displaystyle \frac{B}{A}}`$ $`={\displaystyle \frac{B^{}}{A^{}}}={\displaystyle \frac{J_1(m\mathrm{})}{Y_1(m\mathrm{})}}\text{ for }m^2>0,`$ (112)
$`C`$ $`=C^{}0\text{ for }m^2<0.`$ (113)
Since $`B/A=B^{}/A^{}`$, equations (110) are (112) are equivalent.
As for vector perturbations, the $`m=0`$ scalar mode is not normalizable. Like for tensor and vector perturbations, we have found “scalar gravitons” which appear on the brane as massive particles. If the brane matter is unperturbed, only oscillating $`m^2>0`$ solutions are possible. However, if we allow for non-vanishing matter perturbations on the brane, we can have $`C0`$ and the tachyonic modes $`m^2<0`$ can appear exactly like in the tensor and vector sectors.
It is not surprising that the same instability appears in the scalar, vector and tensor sectors, because all modes describe the same bulk particle, the five-dimensional graviton.
### 4.3 Green’s function, correction to the Newtonian potential
We want to determine the modification to Newton’s law in the RS2 model. Since the extra dimension is not compact and there are massive (homogeneous) modes of all masses $`m^2>0`$, we expect a modification which is not exponentially suppressed.
The 5d Green’s function is defined by
$$^2G(x,x)=\delta ^5(xx),$$
where $`^2`$ is the 5d d’Alembertian in AdS spacetime and we have to glue together a $`Z_2`$–symmetric solution on both sides which satisfies the homogeneous junction condition. We can obtain the retarded Green’s function in the standard way from the homogeneous solutions of the equation (see e.g. CoHi ):
$$G_R(x,x^{})=\frac{d^4k}{(2\pi )^4}e^{ik_\mu (x^\mu x^\mu )}\left[\frac{y^2y^2\mathrm{}^3}{𝒌^2(\omega +iϵ)^2}+_0^{\mathrm{}}𝑑m\frac{u_m(y)u_m(y^{})}{m^2+𝒌^2(\omega +iϵ)^2}\right].$$
The first term comes from the $`m=0`$ solution and the functions $`u_m`$ are the properly normalized massive modes,
$$u_m(y)=\sqrt{\frac{m\mathrm{}}{2}}\frac{J_1(m\mathrm{})Y_2(my)Y_1(m\mathrm{})J_2(my)}{\sqrt{J_1(m\mathrm{})^2+J_1(m\mathrm{})^2}}.$$
The general retarded solution for a given energy momentum tensor $`\tau _{\mu \nu }`$ on the brane is now of the form
$$h_{\mu \nu }(x)=2\kappa _5d^4xG_R(x,x^{})S_{\mu \nu }(x^{}).$$
For a stationary matter distribution it is simpler to use the Green’s function of the spatial Laplacian which is related to $`G_R`$ via integration over time
$`G(𝐱,y,𝐱^{},y^{})`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t^{}G_R(x,x^{})={\displaystyle \frac{d^3k}{(2\pi )^3}e^{i𝐤(𝐱𝐱^{})}\left[\frac{y^2y^2\mathrm{}^3}{𝐤^2}+_0^{\mathrm{}}𝑑m\frac{u_m(y)u_m(y^{})}{m^2+𝐤^2}\right]},`$ (114)
$`=`$ $`{\displaystyle \frac{y^2y^2\mathrm{}^3}{4\pi r}}+{\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{\mathrm{}}}𝑑mu_m(y)u_m(y^{})\mathrm{exp}(mr).`$ (115)
On the brane, $`y=y^{}=\mathrm{}`$, the first term gives the usual $`1/r`$ behavior, $`r=|𝐱𝐱^{}|`$. Expanding the second term to lowest order in $`\mathrm{}/r`$ we obtain
$$G(𝐱,\mathrm{},𝐱^{},\mathrm{})\frac{1}{4\pi \mathrm{}r}\left[1+\frac{\mathrm{}^2}{2r^2}+\mathrm{}\right].$$
This determines the Newtonian potential of a point mass on the brane with mass $`M`$,
$$\kappa _5MG\frac{\kappa _4M}{4\pi r}\left[1+\frac{\mathrm{}^2}{2r^2}+\mathrm{}\right].$$
(116)
Since the extra dimension is non-compact, the correction is not exponentially suppressed but only as a power law. Away from the wall, the potential at large separation is given by
$$G(𝐱,\mathrm{},𝐱^{},\mathrm{}+y)\frac{\mathrm{}}{8\pi (\mathrm{}+y)^2}\frac{2r^2+3y^2}{(r^2+y^2)^{3/2}}$$
(117)
The equipotential lines are shown in Fig 6. This formulas have been derived in Ref. TG from which also Fig. 6 is drawn.
## 5 Cosmological perturbation theory in 4 dimensions
Before studying brane cosmology and 5d effects on cosmological perturbations, I present a brief introduction to 4d cosmological perturbation theory and some aspects of 4d cosmology. Much more details can be found e.g. in Rfund ; Dod . Some knowledge of 4d cosmology is however assumed (Friedmann equations etc. as they can be found in the first chapter of standard textbooks on cosmology like Refs. peebles or Dod ). Students who are familiar with this subject may skip this section.
### 5.1 Perturbation variables
The observed Universe is not perfectly homogeneous and isotropic. Matter is arranged in galaxies and clusters of galaxies and there are large voids in the distribution of galaxies. Let us assume, however, that these inhomogeneities lead only to small variations of the geometry which we shall treat in first order perturbation theory. For this we define the perturbed geometry by
$$g_{\mu \nu }=\overline{g}_{\mu \nu }+a^2h_{\mu \nu },\overline{g}_{\mu \nu }dx^\mu dx^\nu =a^2\left(d\eta ^2+\gamma _{ij}dx^idx^j\right).$$
(118)
Here $`\overline{g}_{\mu \nu }`$ is the unperturbed Friedmann metric, $`a(\eta )`$ is the scale factor, $`\eta `$ denotes conformal time and $`\gamma _{ij}`$ is the 3d metric for a space of constant curvature $`K`$. The perturbations are assumed to be small, $`|h_{\mu \nu }|1`$. The energy momentum tensor is given by
$$\begin{array}{cccc}T_\nu ^\mu =\overline{T}_\nu ^\mu +\theta _\nu ^\mu ,\hfill & \overline{T}_0^0=\overline{\rho },\hfill & \overline{T}_j^i=\overline{p}\delta _j^i\hfill & |\theta _\nu ^\mu |/\overline{\rho }1.\hfill \end{array}$$
(119)
The background energy density $`\rho `$ and pressure $`p`$ satisfy the Friedmann equations,
$`^2\left({\displaystyle \frac{\dot{a}}{a}}\right)^2+K`$ $`=`$ $`{\displaystyle \frac{8\pi G_4}{3}}a^2\overline{\rho }+{\displaystyle \frac{1}{4}}a^2\mathrm{\Lambda }_4`$ (120)
$`\dot{\rho }=3(\overline{\rho }+\overline{p}),`$ (121)
where an over-dot denotes the derivative w.r.t. conformal time $`\eta `$.
Without loss of generality we can choose the so-called longitudinal gauge so that perturbations of the metric are of the form
$$\left(h_{\mu \nu }\right)=\left(\begin{array}{cc}(1+2\mathrm{\Psi })& B_i\\ B_i& (12\mathrm{\Phi })\gamma _{ij}+H_{ij}\end{array}\right).$$
(122)
Here $`B_i`$ is an divergence free vector and $`H_{ij}`$ is a trace-free, divergence free tensor field. The scalar quantities $`\mathrm{\Psi }`$ and $`\mathrm{\Phi }`$ are called the Bardeen potentials. In the Newtonian approximation they are both equal and reduce to the Neewtonian potential.
We also decompose the perturbations into different ’Fourier modes’,
$$\mathrm{\Psi }(\eta ,𝐱)=Y_𝐤(𝐱)\mathrm{\Psi }(\eta ,𝐤),\mathrm{\Phi }(\eta ,𝐱)=Y_𝐤(𝐱)\mathrm{\Phi }(\eta ,𝐤),B^i(\eta ,𝐱)=Y_𝐤^{(V)i}(𝐱)B(\eta ,𝐤),H_{ij}(\eta ,𝐱)=Y_{𝐤ij}^{(T)}(𝐱)H(\eta ,𝐤).$$
In a Friedmann Universe with vanishing curvature, these are just ordinary Fourier modes, while in the general case, the functions $`Y_𝐤`$ are eigenfunctions of the spatial Laplacian with eigenvalue $`k^2`$. The functions $`Y_𝐤^{(V)i}`$ and $`Y_{𝐤ij}^{(T)}`$ correspondingly are vector- and tensor-type eigenfunctions of the spatial Laplacian with vanishing divergence. For later use we also define a scalar type vector and tensor as well as a vector-type tensor,
$`Y_{𝐤i}^{(S)}`$ $`=`$ $`k^1_iY_𝐤,Y_{𝐤ij}^{(S)}=k^2(_i_j{\displaystyle \frac{1}{3}}\delta _{ij}\mathrm{\Delta })Y_𝐤,`$ (123)
$`Y_{𝐤ij}^{(V)}`$ $`=`$ $`{\displaystyle \frac{k^1}{2}}(_iY_{𝐤j}^{(V)}+_jY_{𝐤i}^{(V)}).`$ (124)
Note that contrary to the vector-type vector field $`Y_i^{(V)}`$, the vector field $`Y_i^{(S)}`$ is not divergence free. The same is true for the tensor fields $`Y_{ij}^{(S)}`$ and $`Y_{ij}^{(V)}`$.
Let $`T_\nu ^\mu =\overline{T}_\nu ^\mu +\theta _\nu ^\mu `$ be the full energy momentum tensor. We define its energy density $`\rho `$ and its energy flux 4-vector $`u`$ as the time-like eigenvalue and eigenvector of $`T_\nu ^\mu `$:
$$T_\nu ^\mu u^\nu =\rho u^\mu ,u^2=1.$$
(125)
We then parameterize their perturbations by
$$\rho =\overline{\rho }\left(1+\delta \right),u=u^0_t+u^i_i.$$
(126)
$`u^0`$ is fixed by the normalization condition,
$$u^0=\frac{1}{a}(1\mathrm{\Psi }).$$
(127)
We further set
$$u^i=\frac{1}{a}v^i=\frac{1}{a}\left(VY^{(S)i}+V^{(V)}Y^{(V)i}\right).$$
(128)
Here $`\delta `$ is called the density contrast and $`(v^i)`$ is the peculiar velocity.
We define $`P_\nu ^\mu u^\mu u_\nu +\delta _\nu ^\mu `$, the projection tensor onto the part of tangent space normal to $`u`$ and the stress tensor
$$\tau ^{\mu \nu }=P_\alpha ^\mu P_\beta ^\nu T^{\alpha \beta }.$$
(129)
In the unperturbed case we have $`\tau _0^0=0,\tau _j^i=\overline{p}\delta _j^i`$. Including perturbations, to first order we still obtain
$$\tau _0^0=\tau _i^0=\tau _0^i=0.$$
(130)
But $`\tau _j^i`$ contains in general perturbations. We set
$$\tau _j^i=\overline{p}\left[\left(1+\pi _L\right)\delta _j^i+\mathrm{\Pi }_j^i\right],\text{with}\mathrm{\Pi }_i^i=0.$$
(131)
We decompose $`\mathrm{\Pi }_j^i`$ into scalar- vector- and tensor-type contributions,
$$\mathrm{\Pi }_j^i=\mathrm{\Pi }^{(S)}Y_j^{(S)i}+\mathrm{\Pi }^{(V)}Y_j^{(V)i}+\mathrm{\Pi }^{(T)}Y_j^{(T)i}.$$
(132)
Another important variable is
$$\mathrm{\Gamma }=\pi _L\frac{c_s^2}{w}\delta $$
(133)
where $`c_s^2\dot{p}/\dot{\rho }`$ is the adiabatic sound speed and $`wp/\rho `$ is the enthalpy. One can show that $`\mathrm{\Gamma }`$ is proportional to the divergence of the entropy flux of the perturbations. Adiabatic perturbations are characterized by $`\mathrm{\Gamma }=0`$.
We shall use also other perturbation variables describing the density contrast and peculiar velocity, which actually correspond to these perturbations in different coordinate systems (gauges). One can show that on sub-horizon scales, $`k`$, on which perturbations are actually measurable, they all coincide.
$`D`$ $``$ $`\delta +3(1+w)\left({\displaystyle \frac{\dot{a}}{a}}\right){\displaystyle \frac{V}{k}},`$ (134)
$`D_g`$ $``$ $`\delta 3(1+w)\mathrm{\Phi },`$ (135)
$`\mathrm{\Omega }`$ $``$ $`V^{(V)}B^{(V)},`$ (136)
$`\mathrm{\Omega }V^{(V)}`$ $`=`$ $`B^{(V)}\sigma ^{(V)}.`$ (137)
Here we use the customary name, $`\sigma ^{(V)}=B^{(V)}`$, for the vector-type metric perturbation. These variables can be interpreted nicely in terms of gradients of the energy density and the shear and vorticity of the velocity field Ellis .
### 5.2 Einstein’s equations
We do not derive the first order perturbations of Einstein’s equations. This can be done by different methods, for example with Mathematica. We just write down the results.
#### 5.2.1 Constraint equations
$`\begin{array}{cccc}\hfill 4\pi Ga^2\rho D& =& (k^23K)\mathrm{\Phi }\hfill & (00)\\ \hfill 4\pi Ga^2(\rho +p)V& =& k\left(\mathrm{\Psi }+\dot{\mathrm{\Phi }}\right)\hfill & (0i)\end{array}\}`$ $`(\mathrm{scalar})`$ (140)
$`8\pi Ga^2(\rho +p)\mathrm{\Omega }={\displaystyle \frac{1}{2}}\left(2Kk^2\right)\sigma ^{(V)}(0i)`$ $`(\mathrm{vector})`$ (141)
#### 5.2.2 Dynamical equations
$`k^2\left(\mathrm{\Phi }\mathrm{\Psi }\right)`$ $`=`$ $`8\pi Ga^2p\mathrm{\Pi }^{(S)}(ij)(\mathrm{scalar})`$ (142)
$`\left[\dot{\mathrm{\Psi }}+\left({\displaystyle \frac{^2\dot{}}{^2}}\mathrm{\Phi }+^1\dot{\mathrm{\Phi }}\right)^{}\right]+`$ (143)
$`\left(^2+2\dot{}\right)\left[\mathrm{\Psi }+{\displaystyle \frac{^2\dot{}}{^2}}\mathrm{\Phi }+^1\dot{\mathrm{\Phi }}\right]`$ $`=`$ $`4\pi Ga^2\left(c_s^2D_g+w\mathrm{\Gamma }{\displaystyle \frac{2}{3}}w\mathrm{\Pi }\right)(ii)(\mathrm{scalar})`$ (144)
$`k\left(\dot{\sigma }^{(V)}+2\left({\displaystyle \frac{\dot{a}}{a}}\right)\sigma ^{(V)}\right)`$ $`=`$ $`8\pi Ga^2p\mathrm{\Pi }^{(V)}(ij)(\mathrm{vector})`$ (145)
$`\ddot{H}^{(T)}+2\left({\displaystyle \frac{\dot{a}}{a}}\right)\dot{H}^{(T)}+\left(2K+k^2\right)H^{(T)}`$ $`=`$ $`8\pi Ga^2p\mathrm{\Pi }^{(T)}(ij)(\mathrm{tensor})`$ (146)
For perfect fluids, where $`\mathrm{\Pi }_j^i0`$, we have $`\mathrm{\Phi }=\mathrm{\Psi }`$, $`\sigma ^{(V)}1/a^2`$, and $`H^{(T)}`$ obeys a damped wave equation. The damping term can be neglected on small scales (over short time periods) when $`\eta ^2\stackrel{<}{}2K+k^2`$, so that $`H^{(T)}`$ represents a propagating gravitational wave. For vanishing curvature, the scales $`k\eta 1`$ are simply the sub-horizon scales. For $`K<0`$, waves oscillate with a somewhat smaller frequency, $`\omega =\sqrt{2K+k^2}`$, while for $`K>0`$ the frequency is somewhat higher than $`k`$.
#### 5.2.3 Energy momentum conservation
The conservation equations, $`_\nu T^{\mu \nu }T_{;\nu }^{\mu \nu }=0`$ lead to the following perturbation equations.
$`\begin{array}{c}\dot{D}_g+3\left(c_s^2w\right)\left(\frac{\dot{a}}{a}\right)D_g+(1+w)kV+3w\left(\frac{\dot{a}}{a}\right)\mathrm{\Gamma }=0\\ \dot{V}+\left(\frac{\dot{a}}{a}\right)\left(13c_s^2\right)V=k\left(\mathrm{\Psi }+3c_s^2\mathrm{\Phi }\right)+\frac{c_s^2k}{1+w}D_g\\ +\frac{wk}{1+w}\left[\mathrm{\Gamma }\frac{2}{3}\left(1\frac{3K}{k^2}\right)\mathrm{\Pi }\right]\end{array}\}`$ $`(\mathrm{scalar}),`$ (150)
$`\dot{\mathrm{\Omega }}+\left(13c_s^2\right)\left({\displaystyle \frac{\dot{a}}{a}}\right)\mathrm{\Omega }={\displaystyle \frac{p}{2(\rho +p)}}\left(k{\displaystyle \frac{2K}{k}}\right)\mathrm{\Pi }^{(V)}`$ $`(\mathrm{vector}).`$ (151)
These can of course also be obtained from the Einstein equations since they are equivalent to the contracted Bianchi identities.
For scalar perturbations we have 4 independent equations and 6 variables. For vector perturbations we have 2 equations and 3 variables, while for tensor perturbations we have 1 equation and 2 variables. To close the system we must add matter equations. The simplest prescription is to set $`\mathrm{\Gamma }=\mathrm{\Pi }_{ij}=0`$. These matter equations, which describe adiabatic perturbations of a perfect fluid give us exactly two additional equations for scalar perturbations and one each for vector and tensor perturbations.
Another example is a universe with matter content given by a scalar field. We shall discuss this case in the next section. More complicated examples are those of several interacting particle species of which some have to be described by a Boltzmann equation. This is the actual universe at late times, say $`z\stackrel{<}{}10^{10}`$.
#### 5.2.4 A special case
Here we rewrite the scalar perturbation equations for a simple but important special case. We consider adiabatic perturbations of a perfect fluid. In this case $`\mathrm{\Pi }=0`$ and $`\mathrm{\Gamma }=0`$. Eq. (142) implies $`\mathrm{\Phi }=\mathrm{\Psi }`$. Using the first equation of (140) and Eqs. (135,134) to replace $`D_g`$ in the second of Eqs. (150) by $`\mathrm{\Psi }`$ and $`V`$, finally replacing $`V`$ by (140) one can derive a second order equation for $`\mathrm{\Psi }`$, which is, the only dynamical degree of freedom
$$\ddot{\mathrm{\Psi }}+3(1+c_s^2)\dot{\mathrm{\Psi }}+[(1+3c_s^2)(^2K)(1+3w)(^2+K)+c_s^2k^2]\mathrm{\Psi }=0.$$
(152)
Another interesting example (especially when discussing inflation) is the scalar field case. There, as we shall see in Section 5.4, $`\mathrm{\Pi }=0`$, but in general $`\mathrm{\Gamma }0`$ since $`\delta p/\delta \rho \dot{p}/\dot{\rho }`$. Nevertheless, since this case again has only one dynamical degree of freedom, we can express the perturbation equations in terms of one single second order equation for $`\mathrm{\Psi }`$. In Section 5.4 we shall find the following equation for a perturbed scalar field cosmology
$$\ddot{\mathrm{\Psi }}+3(1+c_s^2)\dot{\mathrm{\Psi }}+[(1+3c_s^2)(^2K)(1+3w)(^2+K)+k^2]\mathrm{\Psi }=0.$$
(153)
The only difference between the perfect fluid and scalar field perturbation equation is that the latter is missing the factor $`c_s^2`$ in front of the oscillatory $`k^2`$ term. Note also that for $`K=0`$ and $`w=c_s^2=`$ constant, the time dependent mass term $`m^2(\eta )=(1+3c_s^2)(^2K)+(1+3w)(^2+K)`$ vanishes.
It is useful to define the variable Mukhanov:1992tc
$$u=a\left[4\pi G(^2\dot{}+K)\right]^{1/2}\mathrm{\Psi },$$
(154)
which satisfies the equation
$$\ddot{u}+(\mathrm{{\rm Y}}k^2\ddot{\theta }/\theta )u=0,$$
(155)
where $`\mathrm{{\rm Y}}=c_s^2`$ or $`\mathrm{{\rm Y}}=1`$ for a perfect fluid or a scalar field background respectively, and
$$\theta =\frac{3}{2a\sqrt{^2\dot{}+K}}.$$
(156)
Another interesting variable is
$$\zeta \frac{2(^1\dot{\mathrm{\Psi }}+\mathrm{\Psi })}{3(1+w)}+\mathrm{\Psi }.$$
(157)
For the rest of this section we set $`K=0`$ for simplicity. Using Eqs. (152) and (153) respectively one then finds
$$\dot{\zeta }=k^2\frac{\mathrm{{\rm Y}}}{^2\dot{}}\mathrm{\Psi }.$$
(158)
On super-horizon scales, $`k/1`$, this time derivative is suppressed by a factor $`(k/)^2(k\eta )^2`$ and this variable is (nearly) conserved on large scales.
The evolution of $`\zeta `$ is closely related to the canonical variable $`v`$ defined by
$$v=\frac{a\sqrt{^2\dot{}}}{\sqrt{4\pi G}\mathrm{{\rm Y}}}\zeta ,$$
(159)
which satisfies the equation
$$\ddot{v}+(\mathrm{{\rm Y}}k^2\ddot{z}/z)v=0,\text{ for }z=\frac{a\sqrt{^2\dot{}+\kappa }}{\mathrm{{\rm Y}}}.$$
(160)
### 5.3 Dust and radiation
Next we discuss two simple applications which are important to understand the anisotropies in the cosmic microwave background (CMB).
#### 5.3.1 The pure dust fluid for $`K=0,\mathrm{\Lambda }=0`$
’Dust’ is the cosmological term for non-relativistic particles for which we can neglect the pressure so that $`w=c_s^2=p=0`$ and $`\mathrm{\Pi }=\mathrm{\Gamma }=0`$. The Friedmann equation implies for dust $`a\eta ^2`$ so tha $`=2/\eta `$. Equation (152) then reduces to
$$\ddot{\mathrm{\Psi }}+\frac{6}{\eta }\dot{\mathrm{\Psi }}=0,$$
(161)
with the general solution
$$\mathrm{\Psi }=\mathrm{\Psi }_0+\mathrm{\Psi }_1\frac{1}{\eta ^5}$$
(162)
with arbitrary constants $`\mathrm{\Psi }_0`$ and $`\mathrm{\Psi }_1`$. Since the perturbations are supposed to be small initially, they cannot diverge for $`\eta 0`$, and we have therefore to keep only the ’growing’ mode, $`\mathrm{\Psi }_1=0`$. But also the $`\mathrm{\Psi }_0`$ mode is only constant. This fact led Lifshitz who was the first to analyze cosmological perturbations to the conclusions that linear perturbations do not grow in a Friedman universe and cosmic structure cannot have evolved by gravitational instability Lif46 . However, the important point to note here is that, even if the gravitational potential remains constant, matter density fluctuations do grow on sub-horizon scales, scales where $`k\eta 1`$ and hence structure can evolve on scales which are smaller than the Hubble scale.
Defining $`x=k\eta `$, we obtain for the velocity potential and the density contrast
$`V`$ $`=`$ $`\mathrm{\Psi }_0{\displaystyle \frac{x}{3}}`$ (163)
$`D_g`$ $`=`$ $`5\mathrm{\Psi }_0{\displaystyle \frac{1}{6}}\mathrm{\Psi }_0x^2,D=D_g+3\mathrm{\Psi }+{\displaystyle \frac{6}{x}}V={\displaystyle \frac{1}{6}}\mathrm{\Psi }_0x^2.`$ (164)
In the variable $`D`$ the constant term has disappeared and we have $`D\mathrm{\Psi }`$ on super-horizon scales, $`x1`$.
On sub-horizon scales, the density fluctuations grow like the scale factor $`ax^2`$. Nevertheless, Lifshitz’ conclusion Lif46 that pure gravitational instability cannot be the cause for structure formation has some truth: if we start from tiny thermal fluctuations of the order of $`10^{35}`$, they can only grow to about $`10^{30}`$ due to this mild, power law instability during the matter dominated regime. Or, to put it differently, if we want to form structure by gravitational instability, we need initial fluctuations of the order of at least $`10^5`$, much larger than thermal fluctuations. According to what we have said here, we need these fluctuations at the beginning of the matter dominated phase, but as we shall see below, perturbations do not grow at all during the radiation dominated era, so that really initial fluctuations with amplitudes $`10^5`$ are needed. One possibility to create such fluctuations is quantum particle production in the classical gravitational field during inflation. The rapid expansion of the universe during inflation quickly expands microscopic scales, at which quantum fluctuations are important, to cosmological scales where these fluctuations are then “frozen in” as classical perturbations in the energy density and the geometry. We will discuss the induced spectrum on fluctuations in Section 5.5.
#### 5.3.2 The pure radiation fluid, $`K=0,\mathrm{\Lambda }=0`$
In this limit we set $`w=c_s^2=\text{“sevenrm1}/\text{“sevenrm3}`$, and $`\mathrm{\Pi }=\mathrm{\Gamma }=0`$ so that $`\mathrm{\Phi }=\mathrm{\Psi }`$. We conclude from $`\rho a^4`$ and the Friedmann equation that $`a\eta `$. For radiation, the $`u`$–equation (155) becomes
$$\ddot{u}+(\frac{1}{3}k^2\frac{2}{\eta ^2})u=0,$$
(165)
with general solution
$$u(x)=A\left(\frac{\mathrm{sin}(x)}{x}\mathrm{cos}(x)\right)+B\left(\frac{\mathrm{cos}(x)}{x}\mathrm{sin}(x)\right),$$
(166)
where we have set $`x=k\eta /\sqrt{3}=c_sk\eta `$. For the Bardeen potential we obtain with (154), up to constant factors,
$$\mathrm{\Psi }(x)=\frac{u(x)}{x^2}.$$
(167)
We must set $`B=0`$ for perturbations to remain regular at early times. On super-horizon scales, $`x1`$, we then have
$$\mathrm{\Psi }(x)\frac{A}{3}.$$
(168)
For the density and velocity perturbations one finds
$$D_g=2A\left[\mathrm{cos}(x)\frac{2}{x}\mathrm{sin}(x)\right],V=\frac{\sqrt{3}}{4}D_g^{}.$$
(169)
In the super-horizon regime, $`x1`$, this yields
$$\mathrm{\Psi }=\frac{A}{3},D_g=2A\frac{A}{3\sqrt{3}}x^2,V=\frac{A}{2\sqrt{3}}x.$$
(170)
On sub-horizon scales, $`x1`$, we obtain oscillating solutions with constant amplitude and with frequency $`k/\sqrt{3}`$:
$$V=\frac{\sqrt{3}A}{2}\mathrm{sin}(x),D_g=2A\mathrm{cos}(x),\mathrm{\Psi }=A\mathrm{cos}(x)/x^2.$$
(171)
Note that also radiation perturbations outside the Hubble horizon are frozen to first order. Once they enter the horizon they start to collapse, but pressure resists the gravitational force and the radiation fluid fluctuations oscillate at constant amplitude. The perturbations of the gravitational potential oscillate and decay like $`1/a^2`$ inside the horizon.
#### 5.3.3 Adiabatic initial conditions
Adiabaticity requires that the perturbations of all contributions to the energy density are initially in thermal equilibrium. This fixes the ratio of the density perturbations of different components. There is no entropy flux and thus $`\mathrm{\Gamma }=0`$. Here we consider a mixture of non relativistic matter and radiation. Since the matter and radiation perturbations behave in the same way on super-horizon scales,
$$D_g^{(r)}=A+Bx^2,D_g^{(m)}=A^{}+B^{}x^2,V^{(r)}V^{(m)}x,$$
(172)
we may require a constant ratio between matter and radiation perturbations. As we have seen in the previous section, inside the horizon ($`x>1`$) radiation perturbations start to oscillate while matter perturbations keep following a power law. On sub-horizon scales a constant ratio can thus no longer be maintained. There are two interesting possibilities: adiabatic and isocurvature perturbations. Here we concentrate on adiabatic perturbations which seem to dominate the observed CMB anisotropies.
From $`\mathrm{\Gamma }=0`$ one easily derives that two components with $`p_i/\rho _i=w_i=`$constant, $`i=1,2`$, are adiabatically coupled if $`(1+w_1)D_g^{(2)}=(1+w_2)D_g^{(1)}`$. Energy conservation then implies that their velocity fields agree, $`V^{(1)}=V^{(2)}`$. This result is also a consequence of the Boltzmann equation in the strong coupling regime. We therefore require
$$V^{(r)}=V^{(m)},$$
(173)
so that the energy flux in the two fluids is coupled initially.
We restrict ourselves to a matter dominated background, the situation relevant in the observed universe after equality. We first have to determine the radiation perturbations during a matter dominated era. Since $`\mathrm{\Psi }`$ is dominated by the matter contribution (it is proportional to the background density of a given component), we have $`\mathrm{\Psi }\mathrm{const}=\mathrm{\Psi }_0`$. We neglect the contribution from the sub-dominant radiation to $`\mathrm{\Psi }`$. Energy momentum conservation for radiation then gives, with $`x=k\eta `$, and $`d/dx=`$
$`D_g^{(r)}`$ $`=`$ $`{\displaystyle \frac{4}{3}}V^{(r)}`$ (174)
$`V^{(r)}`$ $`=`$ $`2\mathrm{\Psi }+{\displaystyle \frac{1}{4}}D_g^{(r)}.`$ (175)
Here $`\mathrm{\Psi }`$ is just a constant given by the matter perturbations, and it acts like a constant source term. The general solution of this system is then
$`D_g^{(r)}`$ $`=`$ $`A\mathrm{cos}(c_sx){\displaystyle \frac{4}{\sqrt{3}}}B\mathrm{sin}(c_sx)+8\mathrm{\Psi }\left[\mathrm{cos}(c_sx)1\right]`$ (176)
$`V^{(r)}`$ $`=`$ $`B\mathrm{cos}(c_sx)+{\displaystyle \frac{\sqrt{3}}{4}}A\mathrm{sin}(c_sx)+2\sqrt{3}\mathrm{\Psi }\mathrm{sin}(c_sx),`$ (177)
where $`c_s=1/\sqrt{3}`$ is the sound speed of radiation. Our adiabatic initial conditions require
$$\underset{x0}{lim}\frac{V^{(r)}}{x}=V_0=\underset{x0}{lim}\frac{V^{(m)}}{x}<\mathrm{}.$$
(178)
Therefore $`B=0`$ and $`V=V_0x`$ with $`V_0=A/42\mathrm{\Psi }`$ on super horizon scales, $`x1`$. Using in addition $`\mathrm{\Psi }=3V_0`$ (see (163)) we obtain
$`D_g^{(r)}`$ $`=`$ $`{\displaystyle \frac{4}{3}}\mathrm{\Psi }\mathrm{cos}\left({\displaystyle \frac{x}{\sqrt{3}}}\right)8\mathrm{\Psi }`$ (179)
$`V^{(r)}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{3}}}\mathrm{\Psi }\mathrm{sin}\left({\displaystyle \frac{x}{\sqrt{3}}}\right)`$ (180)
$`D_g^{(m)}`$ $`=`$ $`\mathrm{\Psi }(5+{\displaystyle \frac{1}{6}}x^2)`$ (181)
$`V^{(m)}`$ $`=`$ $`{\displaystyle \frac{1}{3}}\mathrm{\Psi }x.`$ (182)
On super-horizon scales, $`x1`$ we have
$$D_g^{(r)}\frac{20}{3}\mathrm{\Psi }\text{ and }V^{(r)}\frac{1}{3}x\mathrm{\Psi },$$
(183)
note that $`D_g^{(r)}=(4/3)D_g^{(m)}`$ and $`V^{(r)}=V^{(m)}`$ as it is required for adiabatic initial conditions.
### 5.4 Scalar field cosmology
We now consider the special case of a Friedmann universe filled with self interacting scalar field matter. We keep spatial curvature $`K=0`$ in this section. The action is given by
$$S=\frac{1}{16\pi G}d^4x\sqrt{|g|}R+d^4x\sqrt{|g|}\left(\frac{1}{2}_\mu \phi ^\mu \phi W(\phi )\right)=S_g+S_m$$
(184)
where $`\phi `$ denotes the scalar field and $`W`$ is the potential. The energy momentum tensor is obtained by varying the matter part of the action, $`S_m`$ wrt the metric $`g^{\mu \nu }`$,
$$T_{\mu \nu }=_\mu \phi _\nu \phi \left[\frac{1}{2}_\lambda \phi ^\lambda \phi +W\right]g_{\mu \nu }$$
(185)
The energy density $`\rho `$ and the energy flux $`u`$ are defined by
$$T_\nu ^\mu u^\nu =\rho u^\mu .$$
(186)
For a homogeneous and isotropic universe, $`\phi =\phi (t)`$ and $`g_{\mu \nu }=a^2\eta _{\mu \nu }`$ we obtain
$$\rho =\frac{1}{2a^2}\dot{\phi }^2+W(u^\mu )=\frac{1}{a}(1,\mathrm{𝟎}).$$
(187)
The pressure is given by
$$T_j^i=p\delta _j^ip=\frac{1}{2a^2}\dot{\phi }^2W.$$
(188)
We now consider scalar field perturbations,
$$\phi =\overline{\phi }+\delta \phi .$$
(189)
Clearly, the scalar field only generates scalar-type perturbations (to first order). The perturbed metric is therefore given by $`ds^2=a^2(1+2\mathrm{\Psi })d\eta ^2+a^2(12\mathrm{\Phi })\delta _{ij}dx^idx^j`$. Inserting Eq. (189) in the definition of the energy velocity perturbation $`V`$,
$$(u^\mu )=\frac{1}{a}(1\mathrm{\Psi },V,_i)$$
(190)
and the energy density perturbation $`\delta \rho `$,
$$\rho =\overline{\rho }+\delta \rho ,$$
(191)
we obtain
$$\delta \rho =\frac{1}{a^2}\dot{\overline{\phi }}\delta \dot{\phi }\frac{1}{a^2}\dot{\overline{\phi }}^2\mathrm{\Psi }+W,_\phi \delta \phi $$
(192)
and
$$V=\frac{k}{\dot{\overline{\phi }}}\delta \phi .$$
(193)
From the stress tensor, $`T_{ij}=\phi ,_i\phi ,_j[\frac{1}{2}_\lambda \phi ^\lambda \phi +W]g_{ij}`$ we find
$$p\pi _L=\frac{1}{a^2}\dot{\overline{\phi }}\delta \dot{\phi }\frac{1}{a^2}\dot{\overline{\phi }}^2\mathrm{\Psi }W,_\phi \delta \phi \text{and }\mathrm{\Pi }=0.$$
(194)
Short calculations give
$`D_g`$ $`=`$ $`(1+w)\left[4\mathrm{\Psi }+2{\displaystyle \frac{\dot{a}}{a}}k^1Vk^1\dot{V}\right],`$ (195)
$`D_s`$ $`=`$ $`D_g+3(1+w)\mathrm{\Psi },`$ (196)
$`\mathrm{\Gamma }`$ $`=`$ $`{\displaystyle \frac{2W,_\phi }{p\dot{\rho }}}\left[\dot{\overline{\phi }}\rho D_s\dot{\rho }\delta \phi \right],`$ (197)
$`\mathrm{\Pi }`$ $`=`$ $`0.`$ (198)
The Einstein equations then lead to the following second order equation for the Bardeen potential which we have discussed above:
$$\ddot{\mathrm{\Psi }}+2(\ddot{\phi }/\dot{\phi })\dot{\mathrm{\Psi }}+(2\dot{}2\ddot{\phi }/\dot{\phi }+k^2)\mathrm{\Psi }=0$$
(199)
or, using the definition $`c_s^2=\dot{p}/\dot{\rho }`$,
$$\ddot{\mathrm{\Psi }}+3(1+c_s^2)\dot{\mathrm{\Psi }}+(2\dot{}+(1+3c_s^2)^2+k^2)\mathrm{\Psi }=0.$$
(200)
As already mentioned above, this equation differs from the $`\mathrm{\Psi }`$ equation for a perfect fluid only in the last term proportional to $`k^2`$. This comes from the fact that the scalar field is not in a thermal state with fixed entropy, but it is in a fully coherent state ($`\mathrm{\Gamma }0`$) and field fluctuations propagate with the speed of light. On large scales, $`k\eta 1`$, this difference is not relevant, but on sub–horizon scales it does play a certain role.
### 5.5 Generation of perturbations during inflation
So far we have simply assumed some initial fluctuation amplitude $`A`$, without investigating where it came from or what the $`k`$–dependence of $`A`$ might be. In this section we discuss the most common idea about the generation of cosmological perturbations, namely their production from quantum vacuum fluctuations during an inflationary phase. The treatment here is focused mainly on getting the correct result with as little effort as possible; we ignore several subtleties related, e.g. to the transition from quantum fluctuations of the field to classical fluctuations in the energy momentum tensor. The idea is of course that the source for the metric fluctuations are the expectation values of the energy momentum tensor operator of the scalar field.
The basic idea is simple: A time dependent gravitational field very generically leads to particle production, analogously to the electron positron production in a classical, time dependent, strong electromagnetic field.
Let us first fix our notation. Inflation is an era during which the expansion of the scale factor is accelerated, $`\frac{d^2a}{dt^2}>0`$. In terms of conformal time, $`\frac{d}{d\eta }=\dot{}`$ , this becomes
$$\frac{d^2a}{dt^2}=\frac{1}{a}\dot{}>0.$$
We shall only consider simple power law inflation, where $`a=(c\eta )^q`$ for some constants $`c`$ and $`q`$. For the scale factor to be positive and real we require $`c\eta >0`$. Expansion then happens when $`cq>0`$ and accelerated expansion when in addition $`q<0`$. Hence for power law inflation, the scale factor behaves like
$$a|\eta |^q$$
and $`\eta <0`$ as well as $`q<0`$. It is easy to see that de Sitter inflation, $`a\mathrm{exp}(Ht)`$, corresponds to $`q=1`$. In general, for a fluid with $`p=w\rho `$
$$q=\frac{2}{13w}.$$
Inflation therefore requires $`w<1/3`$. During scalar field inflation, the energy density must therefore be dominated by the potential, $`W>a^2\dot{\phi }^2`$. We suppose that the field is ’slowly rolling’ down the potential until at some later moment the condition $`w<1/3`$ breaks down and inflation stops. How far away a given moment is from this end of inflation can be cast in terms of the slow roll parameters $`ϵ_1`$ and $`ϵ_2`$ defined by
$`ϵ_1`$ $`=`$ $`{\displaystyle \frac{\dot{H}}{aH^2}},H={\displaystyle \frac{}{a}}\text{ is the Hubble parameter}`$ (201)
$`ϵ_2`$ $`=`$ $`{\displaystyle \frac{\frac{a^2d^2\phi }{dt^2}}{\dot{\phi }}}=\left[1{\displaystyle \frac{\ddot{\phi }}{\dot{\phi }}}\right]=\left[1+{\displaystyle \frac{a^2W^{}}{\dot{\phi }}}\right].`$ (202)
#### 5.5.1 Scalar perturbations
The main result of this subsection is the following: During inflation, the produced particles induce a gravitational field with a (nearly) scale invariant spectrum,
$$k^3|\mathrm{\Psi }(k,\eta )|^2=k^{n1}\times \mathrm{const}.\text{ with }n1.$$
(203)
The quantity $`k^3|\mathrm{\Psi }(k,\eta )|^2`$ is the squared amplitude of the metric perturbation at comoving scale $`\lambda =\pi /k`$. To ensure that this quantity is small over a broad range of scales, so that neither black holes form on small scales nor large deviation from homogeneity and isotropy on large scales appear, we must require $`n1`$. These arguments have been put forward by Harrison and Zel’dovich HZ (ten years before the advent of inflation), leading to the name ’Harrison-Zel’dovich spectrum’ for a scale invariant perturbation spectrum.
To derive the above result we consider a scalar field background dominated by a potential, hence $`a|\eta |^q`$ with $`q1`$. Developing the action of this system,
$$S=𝑑x^4\sqrt{|g|}\left(\frac{R}{16\pi G}+\frac{1}{2}(\phi )^2W\right),$$
to second order in the perturbations (see Mukhanov:1992tc ) around the Friedmann solution one obtains
$$\delta S=𝑑x^4\sqrt{|\overline{g}|}\frac{1}{2}(_\mu v)^2$$
(204)
up to a total differential. Here $`v`$ is the perturbation variable
$$v=\frac{a\sqrt{^2\dot{}}}{\sqrt{4\pi G}}\zeta $$
(205)
introduced in Eq. (159). Via the Einstein equations, this variable can also be interpreted as representing the fluctuations in the scalar field. Therefore, we quantize $`v`$ and assume that initially, on small scales, $`k|\eta |1`$, $`v`$ is in the (Minkowski) quantum vacuum state of a massless scalar field with mode function
$$v_{\mathrm{in}}=\frac{v_0}{\sqrt{k}}\mathrm{exp}(ik\eta ).$$
(206)
The pre-factor $`v_0`$ is a $`k`$-independent constant which depends on convention, but is of order unity. From (158) we can derive
$$(v/z)^{}=\frac{k^2u}{z},$$
where $`za`$ is defined in Eq. (160) and $`ua\eta \mathrm{\Psi }`$ is given in Eq. (154). On small scales, $`k|\eta |1`$, this results in the initial condition for $`u`$
$$u_{\mathrm{in}}=\frac{iv_0}{k^{3/2}}\mathrm{exp}(ik\eta ).$$
(207)
In the case of power law expansion, $`a|\eta |^q`$, the evolution equation for $`u`$, Eq. (155), reduces to
$$\ddot{u}+(k^2\frac{q(q+1)}{\eta ^2})u=0.$$
(208)
The solutions to this equation are of the form $`(k|\eta |)^{1/2}H_\mu ^{(i)}(k\eta )`$, where $`\mu =q+1/2`$ and $`H_\mu ^{(i)}`$ is the Hankel function of the $`i`$th kind ($`i=1`$ or $`2`$) of order $`\mu `$. The initial condition (207) requires that only $`H_\mu ^{(2)}`$ appears, so that we obtain
$$u=\frac{\alpha }{k^{3/2}}(k|\eta |)^{1/2}H_\mu ^{(2)}(k\eta ),$$
where again $`\alpha `$ is a constant of order unity. We define the value of the Hubble parameter during inflation, which is nearly constant by $`H_i`$. With $`H=/a1/(|\eta |a)`$, we then obtain $`a1/(H_i|\eta |)`$. With the Planck mass defined by $`4\pi G=M_4^2`$, Eq. (154) then gives
$$\mathrm{\Psi }=\frac{H_i}{2M_4}u\frac{H_i}{M_4}k^{3/2}(k|\eta |)^{1/2}H_\mu ^{(2)}(k\eta ).$$
(209)
On small scales this is a simple oscillating function while on large scales $`k|\eta |1`$ it can be approximated by a power law,
$$\mathrm{\Psi }\frac{H_i}{M_4}k^{3/2}(k|\eta |)^{1+q},\text{ for }k|\eta |1.$$
(210)
Here we have used $`\mu =1/2+q<0`$. This yields
$$k^3|\mathrm{\Psi }|^2\left(\frac{H_i}{M_4}(k|\eta |)^{1+q}\right)^2k^{n1},$$
(211)
hence $`n1`$ if $`q1`$. Detailed studies have shown that the amplitude of $`\mathrm{\Psi }`$ can still be somewhat affected by the transition from inflation to the subsequent radiation era, the spectral index, however, is very stable. Simple deviations from de Sitter inflation, like e.g. power law inflation, $`q>1`$, lead to slightly blue spectra, $`n\stackrel{>}{}1`$.
With a somewhat more careful treatment, one finds that both, the amplitude and the spectral index depend on scale via the slow roll parameters $`ϵ_1`$ and $`ϵ_2`$,
$`k^3|\mathrm{\Psi }|^2`$ $`=`$ $`{\displaystyle \frac{2H_i^2}{M_4^2ϵ_1}}(k\eta _f)^{n1},`$ (212)
$`n|_{k=a(\eta )H(\eta )}`$ $`=`$ $`14ϵ_1(\eta )2ϵ_2(\eta ).`$ (213)
Vector perturbations are not generated during standard inflation; and even if they are generated they only decay during subsequent evolution and we therefore do not discuss them any further. This may change drastically in braneworlds (see Ringeval:2003na )!
#### 5.5.2 Tensor perturbations
The situation is different for tensor perturbations. Again we consider the perfect fluid case, $`\mathrm{\Pi }_{ij}^{(T)}=0`$. Eq. (146) implies
$$\ddot{H}_{ij}+\frac{2\dot{a}}{a}\dot{H}_{ij}+k^2H_{ij}=0.$$
(214)
If the background has a power law evolution, $`a\eta ^q`$ this equation can be solved in terms of Bessel or Hankel functions. The less decaying mode solution to Eq. (214) is $`H_{ij}=e_{ij}x^{1/2\beta }J_{1/2q}(x)`$, where $`J_\nu `$ denotes the Bessel function of order $`\nu `$, $`x=k\eta `$ and $`e_{ij}`$ is a transverse traceless polarization tensor. This leads to
$`H_{ij}`$ $`=`$ $`\mathrm{const}\mathrm{for}x1`$ (215)
$`H_{ij}`$ $`=`$ $`{\displaystyle \frac{1}{a}}\mathrm{for}x\stackrel{>}{}1.`$ (216)
One may also quantize the tensor fluctuations which represent gravitons. Doing this, one obtains (up to log corrections) a scale invariant spectrum of tensor fluctuations from inflation: for tensor perturbations the canonical variable is simple given by $`h_{ij}=M_PaH_{ij}`$. The evolution equation for $`h_{ij}=he_{ij}`$ is of the form
$$\ddot{h}+(k^2+m^2(\eta ))h=0,$$
(217)
where $`m^2(\eta )=\ddot{a}/a`$. During inflation $`m^2=q(q1)/\eta ^2`$ is negative, leading to particle creation. Like for scalar perturbations, the vacuum initial conditions are given on scales which are inside the horizon, $`k^2|m^2|`$,
$$h_{\mathrm{in}}=\frac{1}{\sqrt{k}}\mathrm{exp}(ik\eta )\text{ for }k|\eta |1.$$
Solving Eq. (217) with this initial condition, gives
$$h=\frac{1}{\sqrt{k}}(k|\eta |)^{1/2}H_{q1/2}^{(2)}(k\eta ),$$
where $`H_\nu ^{(2)}`$ is the Hankel function of degree $`\nu `$ of the second kind. On super-horizon scales where we have $`H_{q1/2}^{(2)}(k\eta )(k|\eta |)^{q1/2}`$, this leads to $`|h|^2|\eta |(k|\eta |)^{2q1}`$. Using the relation between $`h_{ij}=he_{ij}`$ and $`H_{ij}`$ one obtains the spectrum of tensor perturbations generated during inflation. For exponential inflation, $`q1`$ one finds again a scale invariant spectrum for $`H_{ij}`$ on super-horizon scales,
$$k^3|H_{ij}H^{ij}|(2H_{\mathrm{in}}/M_4)^2(k\eta _f)^{n_T}\text{ with }n_T=2(q+1)0.$$
(218)
Again, a more careful treatment within the slow roll approximation gives
$$n_T=2ϵ_1.$$
(219)
A more detailed analysis also of the amplitudes of the scalar and tensor spectra leads to the consistency relation $`n_T=2A_T^2/A_S^2`$ of slow roll inflation. Here $`A_T`$ and $`A_S`$ are the amplitudes of tensor and scalar perturbations, respectively.
More details on inflation can be found in many cosmology books, e.g. Refs. Dod ; LiLy ; Linde .
### 5.6 Power spectra
The quantities which we have calculated in the previous subsection are not the precise values of e.g. $`\mathrm{\Psi }(𝐤,\eta )`$, but only expectation values $`|\mathrm{\Psi }(𝐤,\eta )|^2`$. In different realizations of the same inflationary model, the ’phases’ $`\alpha (𝐤,\eta )`$ given by $`\mathrm{\Psi }(𝐤,\eta )=\mathrm{exp}(i\alpha (𝐤))|\mathrm{\Psi }(k)|`$ are different. They are random variables. Since the process which generates the fluctuations $`\mathrm{\Psi }`$ is stochastically homogeneous and isotropic, these phases are uncorrelated (for different values of $`𝐤`$). However, the quantity which we can calculate for a given model and which then has to be compared with observation is the power spectrum. Power spectra are the “harmonic transforms” of the two point correlation functions<sup>2</sup><sup>2</sup>2The “harmonic transform” in usual flat space is simply the Fourier transform. In curved space it is the expansion in terms of eigenfunctions of the Laplacian on that space, e.g. on the sphere it corresponds to the expansion in terms of spherical harmonics. If the perturbations of the model under consideration are Gaussian, this is a relatively generic prediction from inflationary models, then the two-point functions and therefore the power spectra contain the full statistical information of the model.
Let us first consider the power spectrum of matter,
$$P_D(k)=\left|D_g(𝐤,\eta _0)\right|^2.$$
(220)
Here $``$ indicates a statistical average, ensemble average, over “initial conditions” in a given model. $`P_D(k)`$ is usually compared with the observed power spectrum of the galaxy distribution.
The spectrum we can both, measure and calculate to the best accuracy is the CMB anisotropy power spectrum. It is defined as follows: The fluctuations of the radiation temperature as observed in the sky, $`\mathrm{\Delta }T/T`$, is a function of our position $`𝐱_0`$, time $`\eta _0`$ and the photon direction $`𝐧`$. We develop the $`𝐧`$–dependence in terms of spherical harmonics. We will suppress the argument $`\eta _0`$ and often also $`𝐱_0`$ in the following calculations. All results are for today ($`\eta _0`$) and here ($`𝐱_0`$). By statistical homogeneity statistical averages over an ensemble of realizations (expectation values) are supposed to be independent of position. Furthermore, we assume that the process generating the initial perturbations is statistically isotropic. Then, the off-diagonal correlators of the expansion coefficients $`a_\mathrm{}m`$ vanish and we have
$$\frac{\mathrm{\Delta }T}{T}(𝐱_0,\eta _0,𝐧)=\underset{\mathrm{},m}{}a_\mathrm{}m(𝐱_0,\eta _0)Y_\mathrm{}m(𝐧),a_\mathrm{}ma_\mathrm{}^{}m^{}^{}=\delta _{\mathrm{}\mathrm{}^{}}\delta _{mm^{}}C_{\mathrm{}}.$$
(221)
The $`C_{\mathrm{}}`$’s are the CMB power spectrum.
The two point correlation function is related to the $`C_{\mathrm{}}`$’s by
$`{\displaystyle \frac{\mathrm{\Delta }T}{T}}(𝐧){\displaystyle \frac{\mathrm{\Delta }T}{T}}(𝐧^{})_{𝐧𝐧^{}=\mu }={\displaystyle \underset{\mathrm{},\mathrm{}^{},m,m^{}}{}}a_\mathrm{}ma_\mathrm{}^{}m^{}^{}Y_\mathrm{}m(𝐧)Y_\mathrm{}^{}m^{}^{}(𝐧^{})=`$
$`{\displaystyle \underset{\mathrm{}}{}}C_{\mathrm{}}\underset{\frac{2\mathrm{}+1}{4\pi }P_{\mathrm{}}(𝐧𝐧^{})}{\underset{}{{\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}Y_\mathrm{}m(𝐧)Y_\mathrm{}m^{}(𝐧^{})}}={\displaystyle \frac{1}{4\pi }}{\displaystyle \underset{\mathrm{}}{}}(2\mathrm{}+1)C_{\mathrm{}}P_{\mathrm{}}(\mu ),`$ (222)
where we have used the addition theorem of spherical harmonics for the last equality; the $`P_{\mathrm{}}`$’s are the Legendre polynomials.
For given metric perturbations and perturbations of the energy momentum tensor of the cosmic fluid, the temperature perturbations can be determined by following the oscillations in the radiation fluid before decoupling (see subsection 5.3.2) and by following the propagation of photons along geodesics in the perturbed spacetime after decoupling. Decoupling of photons and matter happens during recombination ($`T3000`$K, $`z1000`$), where electrons and protons recombine to neutral hydrogen. During that process, the number density of free electrons with which the photons can scatter drops drastically and finally becomes so low, that the mean free path of the photons grows larger than the Hubble scale. The surface of constant temperature, $`T_{\mathrm{dec}}=T(\eta _{\mathrm{dec}})`$, at which this happens is also called the ’last scattering surface’. After last scattering, the photons effectively cease to interact and move freely along geodesics (more details can be found e.g. in Dod ; Rfund ; myCMB ).
Clearly the $`a_{lm}`$’s from scalar-, vector- and tensor-type perturbations are uncorrelated,
$$a_\mathrm{}m^{(S)}a_\mathrm{}^{}m^{}^{(V)}=a_\mathrm{}m^{(S)}a_\mathrm{}^{}m^{}^{(T)}=a_\mathrm{}m^{(V)}a_\mathrm{}^{}m^{}^{(T)}=0.$$
(223)
Since vector perturbations decay, their contributions, the $`C_{\mathrm{}}^{(V)}`$, are negligible in models where initial perturbations have been laid down very early, e.g. , after an inflationary period. Tensor perturbations are constant on super-horizon scales and perform damped oscillations once they enter the horizon.
Let us first discuss in somewhat more detail scalar perturbations. We restrict ourselves to the case $`K=0`$ for simplicity. We suppose the initial perturbations to be given by a spectrum,
$$\left|\mathrm{\Psi }\right|^2k^3=A^2k^{n1}\eta _0^{n1}.$$
(224)
We multiply by the constant $`\eta _0^{n1}`$, the actual comoving size of the horizon, in order to keep $`A`$ dimensionless for all values of $`n`$. $`A`$ then represents the amplitude of metric perturbations at horizon scale today, $`k=1/\eta _0`$.
For adiabatic perturbations we have obtained on super-horizon scales,
$$\frac{1}{4}D_g^{(r)}=\frac{5}{3}\mathrm{\Psi }+𝒪((k\eta )^2),V^{(b)}=V^{(r)}=𝒪(k\eta ).$$
(225)
The dominant contribution to the temperature fluctuations on super-horizon scales (neglecting the integrated Sachs–Wolfe effect $`\dot{\mathrm{\Phi }}\dot{\mathrm{\Psi }}`$ ) comes from two terms: the first, $`2\mathrm{\Psi }`$, is the change of photon energy due to the gravitational potential at the last scattering surface, $`\eta =\eta _{\mathrm{dec}}`$, and the second, $`\text{“sevenrm1}/\text{“sevenrm4}D_g^{(r)}`$ represents the intrinsic temperature fluctuations (for more details see SW ; peebles ; myCMB ). With Eq. (225) this yields the famous Sachs-Wolfe formula
$$\frac{\mathrm{\Delta }T}{T}(𝐱_0,𝐧,\eta _0)=2\mathrm{\Psi }(x_{\mathrm{dec}},\eta _{\mathrm{dec}})+\frac{1}{4}D_g^{(r)}(x_{\mathrm{dec}},\eta _{\mathrm{dec}})=\frac{1}{3}\mathrm{\Psi }(x_{\mathrm{dec}},\eta _{\mathrm{dec}}).$$
(226)
The Fourier transform of (226) gives
$$\frac{\mathrm{\Delta }T}{T}(𝐤,𝐧,\eta _0)=\frac{1}{3}\mathrm{\Psi }(k,\eta _{\mathrm{dec}})e^{i\mathrm{𝐤𝐧}\left(\eta _0\eta _{\mathrm{dec}}\right)}.$$
(227)
Using the decomposition
$$e^{i\mathrm{𝐤𝐧}\left(\eta _0\eta _{\mathrm{dec}}\right)}=\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}(2\mathrm{}+1)i^{\mathrm{}}j_{\mathrm{}}(k(\eta _0\eta _{\mathrm{dec}}))P_{\mathrm{}}(\widehat{𝐤}𝐧),$$
where $`j_{\mathrm{}}`$ are the spherical Bessel functions, we obtain
$`{\displaystyle \frac{\mathrm{\Delta }T}{T}}(𝐱_0,𝐧,\eta _0){\displaystyle \frac{\mathrm{\Delta }T}{T}}(𝐱_0,𝐧^{},\eta _0)`$ (229)
$`=`$ $`{\displaystyle \frac{1}{V}}{\displaystyle d^3x_0\frac{\mathrm{\Delta }T}{T}(𝐱_0,𝐧,\eta _0)\frac{\mathrm{\Delta }T}{T}(𝐱_0,𝐧^{},\eta _0)}`$
$`=`$ $`{\displaystyle \frac{1}{(2\pi )^3}}{\displaystyle d^3k\frac{\mathrm{\Delta }T}{T}(𝐤,𝐧,\eta _0)\left(\frac{\mathrm{\Delta }T}{T}\right)^{}(𝐤,𝐧^{},\eta _0)}`$
$`=`$ $`{\displaystyle \frac{1}{(2\pi )^39}}{\displaystyle d^3k\left|\mathrm{\Psi }\right|^2\underset{\mathrm{},\mathrm{}^{}=0}{\overset{\mathrm{}}{}}(2\mathrm{}+1)(2\mathrm{}^{}+1)i^{\mathrm{}\mathrm{}^{}}}`$
$`j_{\mathrm{}}(k(\eta _0\eta _{\mathrm{dec}}))j_{\mathrm{}^{}}(k(\eta _0\eta _{\mathrm{dec}}))P_{\mathrm{}}(\widehat{𝐤}𝐧)P_{\mathrm{}^{}}(\widehat{𝐤}𝐧^{}).`$
In the second equal sign we have used the unitarity of the Fourier transformation. Inserting $`P_{\mathrm{}}(\widehat{𝐤}𝐧)=\frac{4\pi }{2\mathrm{}+1}_mY_\mathrm{}m^{}(\widehat{𝐤})Y_\mathrm{}m(𝐧)`$ and $`P_{\mathrm{}^{}}(\widehat{𝐤}𝐧^{})=\frac{4\pi }{2\mathrm{}^{}+1}_m^{}Y_\mathrm{}^{}m^{}^{}(\widehat{𝐤})Y_\mathrm{}^{}m^{}(𝐧^{})`$, the integration over directions $`d\mathrm{\Omega }_{\widehat{k}}`$ gives $`\delta _{\mathrm{}\mathrm{}^{}}\delta _{mm^{}}_mY_\mathrm{}m^{}(𝐧)Y_\mathrm{}m(𝐧^{})`$.
Using as well $`_mY_\mathrm{}m^{}(𝐧)Y_\mathrm{}m(𝐧^{})=\frac{2\mathrm{}+1}{4\pi }P_{\mathrm{}}(\mu )`$, where $`\mu =𝐧𝐧^{}`$, we find
$`{\displaystyle \frac{\mathrm{\Delta }T}{T}}(𝐱_0,𝐧,\eta _0){\displaystyle \frac{\mathrm{\Delta }T}{T}}(𝐱_0,𝐧^{},\eta _0)_{\mathrm{𝐧𝐧}^{}=\mu }=`$ (230)
$`{\displaystyle \underset{\mathrm{}}{}}{\displaystyle \frac{2\mathrm{}+1}{4\pi }}P_{\mathrm{}}(\mu ){\displaystyle \frac{2}{\pi }}{\displaystyle \frac{dk}{k}\frac{1}{9}|\mathrm{\Psi }|^2k^3j_{\mathrm{}}^2(k(\eta _0\eta _{\mathrm{dec}}))}.`$
Comparing this equation with Eq. (222) we obtain for adiabatic perturbations on scales $`2\mathrm{}`$ $``$ $`(\eta _0\eta _{\mathrm{dec}})/\eta _{\mathrm{dec}}`$ $`100`$
$$C_{\mathrm{}}^{(SW)}\frac{2}{\pi }_0^{\mathrm{}}\frac{dk}{k}\left|\frac{1}{3}\mathrm{\Psi }\right|^2k^3j_{\mathrm{}}^2\left(k\left(\eta _0\eta _{\mathrm{dec}}\right)\right).$$
(231)
If $`\mathrm{\Psi }`$ is a pure power law as in Eq. (224) and we set $`k(\eta _0\eta _{\mathrm{dec}})k\eta _0`$, the integral (231) can be performed analytically. For the ansatz (224) one finds
$$C_{\mathrm{}}^{(SW)}=\frac{A^2}{9}\frac{\mathrm{\Gamma }(3n)\mathrm{\Gamma }(\mathrm{}\frac{1}{2}+\frac{n}{2})}{2^{3n}\mathrm{\Gamma }^2(2\frac{n}{2})\mathrm{\Gamma }(\mathrm{}+\frac{5}{2}\frac{n}{2})}\text{ for }3<n<3.$$
(232)
Of special interest is the scale invariant or Harrison–Zel’dovich spectrum, $`n=1`$ (see Section 5.5). It leads to
$$\mathrm{}(\mathrm{}+1)C_{\mathrm{}}^{(SW)}=\mathrm{const}.\left(\frac{\mathrm{\Delta }T}{T}(\vartheta _{\mathrm{}})\right)^2,\vartheta _{\mathrm{}}\pi /\mathrm{}.$$
(233)
This is precisely (within the accuracy of the experiment) the behavior observed by the DMR experiment aboard the satellite COBE DMR and more recently with the WMAP satellite WMAP .
Inflationary models predict very generically a HZ spectrum (up to small corrections). The DMR discovery has therefore been regarded as a great success, if not a proof, of inflation. There are other models like topological defects report or certain string cosmology models dgsv which also predict scale–invariant, i.e. Harrison Zel’dovich spectra of fluctuations. These models do however not belong to the class investigated here, since in these models perturbations are induced by seeds which evolve non–linearly in time.
For gravitational waves (tensor fluctuations), a formula analogous to (232) can be derived,
$$C_{\mathrm{}}^{(T)}=\frac{2}{\pi }𝑑kk^2\left|_{\eta _{\mathrm{dec}}}^{\eta _0}𝑑\eta \dot{H}(\eta ,k)\frac{j_{\mathrm{}}(k(\eta _0\eta ))}{(k(\eta _0\eta ))^2}\right|^2\frac{(\mathrm{}+2)!}{(\mathrm{}2)!}.$$
(234)
To a very crude approximation we may assume $`\dot{H}=0`$ on super-horizon scales and $`𝑑\eta \dot{H}j_{\mathrm{}}(k(\eta _0\eta ))H(\eta =1/k)j_{\mathrm{}}(k\eta _0)`$. For a pure power law,
$$k^3\left|H(k,\eta =1/k)\right|^2A_T^2k^{n_T}\eta _0^{n_T},$$
(235)
one obtains
$`C_{\mathrm{}}^{(T)}`$ $``$ $`{\displaystyle \frac{2}{\pi }}{\displaystyle \frac{(\mathrm{}+2)!}{(\mathrm{}2)!}}A_T^2{\displaystyle \frac{dx}{x}x^{n_T}\frac{j_{\mathrm{}}^2(x)}{x^4}}`$ (236)
$`=`$ $`{\displaystyle \frac{(\mathrm{}+2)!}{(\mathrm{}2)!}}A_T^2{\displaystyle \frac{\mathrm{\Gamma }(6n_T)\mathrm{\Gamma }(\mathrm{}2+\frac{n_T}{2})}{2^{6n_T}\mathrm{\Gamma }^2(\frac{7}{2}n_T)\mathrm{\Gamma }(\mathrm{}+4\frac{n_T}{2})}}.`$
For a scale invariant spectrum ($`n_T=0`$) this results in
$$\mathrm{}(\mathrm{}+1)C_{\mathrm{}}^{(T)}\frac{\mathrm{}(\mathrm{}+1)}{(\mathrm{}+3)(\mathrm{}2)}A_T^2\frac{8}{15\pi }.$$
(237)
The singularity at $`\mathrm{}=2`$ in this crude approximation is not real, but there is some enhancement of $`\mathrm{}(\mathrm{}+1)C_{\mathrm{}}^{(T)}`$ at $`\mathrm{}2`$ see Fig. 7). Again, inflationary models (and topological defects) predict a scale invariant spectrum of tensor fluctuations ($`n_T0`$).
On intermediate scales, $`100<\mathrm{}<1000`$, the acoustic oscillations of radiation density fluctuations before decoupling (see subsection 5.3.2) lead to a characteristic series of peaks in the CMB power spectrum which is being measured in great detail and contains very important information on cosmological parameters JL . On small angular scales, $`\mathrm{}\stackrel{>}{}800`$, fluctuations are damped by collisional damping (Silk damping). This effect has to be discussed with the Boltzmann equation for photons, which goes beyond the scope of this introduction (see Dod ; JL ).
## 6 Braneworld cosmology
We now want to study cosmology of an expanding maximally symmetric braneworld. We still require the bulk to be empty and $`Z_2`$–symmetric. One can show that the most general empty bulk allowing for a homogeneous and isotropic brane is Schwarzschild-AdS (Sch-AdS). In the cosmological setting we allow the brane to move in the bulk. As we shall see, this can mimic cosmological expansion. The situation is as depicted in Fig 8.
The 5d metric of Sch-AdS is of the form
$$ds^2=F(R)dT^2+\frac{dR^2}{F(R)}+R^2\left(\frac{dr^2}{1Kr^2}+r^2d\mathrm{\Omega }^2\right)$$
(238)
where the function $`F`$ is determined by the AdS curvature radius $`\mathrm{}`$, the 5d mass $`C`$ and the curvature $`K`$ of 3d space (on the brane) via
$$F(R)=K+\frac{R^2}{\mathrm{}^2}\frac{C}{R^2}.$$
From this we can calculate the 5d Weyl tensor, $`C_{ABCD}`$, and its ’electric’ components defined in (30) with the result
$`E_{\mu \nu }`$ $`=`$ $`\rho ^Eu_\mu ^Eu_\nu ^E+\pi _{\mu \nu }^E\text{with }u=R^1(1,\mathrm{𝟎}),`$ (239)
$`E_{00}`$ $`=`$ $`{\displaystyle \frac{C}{a^4}}=\rho ^E\text{and }\pi _{\mu \nu }^E=0.`$ (240)
Here $`R(T)=a(t)`$, where $`t`$ is cosmic time on the brane and $`u`$ is the unit normal vector in direction of cosmic time on the brane. $`R(T)`$ is the brane position at time $`t(T)`$.
### 6.1 The modified Friedmann equations
From the brane gravity equations,
$$G_{\mu \nu }=\mathrm{\Lambda }_5g_{\mu \nu }+\kappa _4\tau _{\mu \nu }+\kappa _5^2\sigma _{\mu \nu }E_{\mu \nu }$$
with $`\mathrm{\Lambda }_4=\frac{1}{2}(\mathrm{\Lambda }_5+\kappa _5/6\lambda ^2)`$, $`\kappa _4=\kappa _5^2\lambda /6`$ and
$$\sigma _{\mu \nu }=\frac{1}{4}\tau _{\mu \alpha }\tau _\nu ^\alpha +\frac{1}{12}\tau \tau _{\mu \nu }\frac{1}{8}g_{\mu \nu }\tau _{\alpha \beta }\tau ^{\alpha \beta }\frac{1}{24}g_{\mu \nu }\tau ^2$$
with $`\tau _{\mu \nu }=(\rho +p)u_\mu u_\nu pg_{\mu \nu }`$ we obtain
$$H^2=\frac{\kappa _4}{3}\rho \left(1+\frac{\rho }{2\lambda }\right)+\frac{C}{a^4}+\frac{\mathrm{\Lambda }_4}{3}+\frac{K}{a^2},$$
(241)
where $`H`$ is the Hubble parameter. In this section we denote the derivative with respect to cosmological time $`t`$ determined by $`dt=ad\eta `$ by an over-dot, so that $`H=\dot{a}/a`$ and $`=\dot{a}`$.
The term $`\rho ^2/(2\lambda )`$ in Eq. (241) is a correction to the Friedmann equation which is important only at high energies and $`\rho ^E=C/a^4`$, comes from the Weyl tensor. It is called ’Weyl radiation’ since it scales like cosmic radiation, $`a^4`$. Observations (nucleosynthesis) tell us that latest at the temperature $`T1`$MeV these additional terms should be unimportant. More precisely, $`\rho ^E/\rho |_{(nuc)}\stackrel{<}{}0.1`$ and $`\lambda \stackrel{>}{}(1`$MeV$`)^4`$. For the 5d Planck mass this implies $`M_5\stackrel{>}{}3\times 10^4`$GeV. The conservation equation of 4-dimensional cosmology remains unchanged,
$$\dot{\rho }=3(\rho +p)\frac{\dot{a}}{a}.$$
(242)
Solutions for $`C=K=\mathrm{\Lambda }_4=0`$ are readily found. If the equation of state is given by $`p=w\rho `$ we find
$$a=a_0\left[t(t+t_\lambda )\right]^{\frac{1}{3(1+w)}},t_\lambda =\frac{M_4}{\sqrt{3\pi \lambda }(3+w)}<1\mathrm{s}\mathrm{e}\mathrm{c},$$
(243)
where we have used $`\lambda >1(`$Mev$`)^4`$ for the inequality. This is to be compared with the usual 4-dimensional behavior from general relativity (GR). There we have $`a=a_0t^{2/3(1+w)}`$, which corresponds to the above result in the limit $`\lambda \mathrm{}`$. Especially interesting is also the case of de Sitter inflation where $`p=\rho `$ and hence $`\rho =`$constant, so that $`a=a_0e^{Ht}`$ , with
$$H=\sqrt{\kappa _4\frac{\rho }{3}\left(1+\frac{\rho }{2\lambda }\right)}>\sqrt{\frac{\kappa _4\rho }{3}}=H_{GR}.$$
(244)
At low energies we recover the usual Friedmann equation while at high energies, $`\rho \lambda `$, the expansion law differs. During a radiation epoch, $`p=\rho /3`$ with $`\rho \lambda `$ we have $`at^{1/4}`$ (instead of the usual GR behavior, $`at^{1/2}`$). This comes from the fact that in braneworlds at high energies $`H\rho `$, where as in 4d GR we have $`H\sqrt{\rho }`$. If $`\lambda `$ is given by the electroweak scale, $`\lambda 1`$TeV<sup>4</sup>, the observed low energy cosmology, like nucleosynthesis which starts after $`\rho 1`$MeV<sup>4</sup> is not affected (if $`C`$ is sufficiently small).
However, perturbations will carry 5d effects which should in principle be observable in the fluctuation spectrum of the cosmic microwave background radiation (CMB) and in the large scale distribution of matter. These effects are still under investigation. In Section 6.3 I shall present some preliminary partial results.
### 6.2 Brane inflation
We now want to study scalar field inflation in braneworlds. Since energy momentum conservation is still valid, the scalar field evolution equation is also not modified,
$$\ddot{\phi }+3H\dot{\phi }+W^{}(\phi )=0.$$
In 4d GR, the condition for inflation, $`\ddot{a}>0`$ , is equivalent to $`\dot{\phi }^2<V`$ which corresponds to the violation of the strong energy condition
$$p=\frac{1}{2}\dot{\phi }^2W<\frac{1}{3}\rho =\frac{1}{3}\left(\frac{1}{2}\dot{\phi }^2+W\right),w=\frac{p}{\rho }<\frac{1}{3},\text{or}\dot{\phi }^2<W.$$
The braneworld Friedmann equation (244) leads to a stronger condition on $`w`$ for accelerated expansion. For branewords $`0<\frac{\ddot{a}}{a}=\dot{H}+H^2`$ requires
$$w<\frac{1}{3}\left[\frac{1+2\rho /\lambda }{1+\rho /\lambda }\right],\text{or}\dot{\phi }^2<W+\left[\frac{\frac{1}{2}\dot{\phi }^2+W}{\lambda }\left(\frac{5}{4}\dot{\phi }^2\frac{1}{2}V\right)\right]$$
(245)
for inflation to happen. This becomes the usual $`w<1/3`$ at low energy, $`\rho <<\lambda `$, but turns into $`w<2/3`$ at high energy.
If the slow roll approximation ($`\dot{\phi }^2W`$) is satisfied we have
$$H^2\frac{\kappa _4}{3}W\left[1+\frac{W}{2\lambda }\right],\dot{\phi }\frac{W^{}}{3H}.$$
Since the Hubble rate is increased, slow roll is maintained longer than in usual 4d inflation. Correspondingly, the slow roll parameters are reduced,
$$ϵ_1\frac{\dot{H}}{H}=\frac{M_4^2}{16\pi }\left(\frac{W^{}}{W}\right)^2\left[\frac{1+W/\lambda }{(1+W/2\lambda )^2}\right],ϵ_2\frac{\ddot{\phi }}{\dot{\phi }H}=\frac{M_4^2}{8\pi }\left(\frac{W^{\prime \prime }}{W}\right)\left[\frac{1}{(1+W/2\lambda )}\right].$$
(246)
In the high energy regime, $`V\lambda `$ they are reduced by factors $`4\lambda /V`$ and $`2\lambda /V`$, respectively. Hence, the universe may be inflating only because it is in the high energy regime and turn into kinetic dominated expansion, $`w1`$, as soon as $`V<\lambda `$. Such models are constrained since they induce a blue spectrum of gravity waves braneGW .
In standard 4d GR the perturbation spectrum induced by inflation is well known and the scalar spectrum agrees extremely well with the observed anisotropies in the CMB. This will most probably lead to the most stringent constraints for braneworlds, which however have not yet been explored in full generality so far. This is still a very active field of research.
### 6.3 Observable consequences from braneworld cosmology
So far, we have seen that it is conceivable that our Universe is a 3-brane. At least at low energy and disregarding perturbations, we cannot distinguish cosmological evolution due to 4d Einstein equations or the (so different!) brane gravity equations. We finally want to study ways to discover whether the braneworld idea is realized in nature. Are there tell-tale observational signatures which would betray whether we live on a brane?
As we have seen, the gravitational equations for braneworlds differ significantly from Einstein gravity. However, at low energy the Friedmann equations for a homogeneous and isotropic Universe are recovered. Hence there are two regimes in which deviations from Einstein gravity will be found:
* At high energy. This is especially interesting for the generation of inflationary perturbations which may be affected by high energy braneworld behavior. As long as we restrict ourselves on background effects the calculations are relatively straight forward and well under control.
* In the perturbations. Cosmological brane perturbation theory is not well under control and still a subject of active research. One main point are bulk perturbations which can affect the brane and act there like ’sources’. On the other hand, gravity wave perturbations generated on the brane can be emitted into the bulk.
Here we give examples of both aspects, how effects from braneworlds can enter cosmological perturbations, but we are by no means exhaustive (more details and especially references can be found in RoyRev ).
Let us first consider the high energy universe. During inflation scalar and tensor perturbations are generated. The spectrum of scalar perturbations is $`|\mathrm{\Psi }|^2k^3=A_S^2k^{2q2}=A_S^2(k/H_0)^{n_s1}`$. The slow roll approximation for braneworlds gives RoyRev
$$A_S^2\frac{512\pi }{75M_4^6}\frac{W^3}{W^2}\left[\frac{2\lambda +W}{2\lambda }\right]^3,n_S=14ϵ_12ϵ_2.$$
(247)
Similarly, for tensor perturbations, $`|H|^2k^3=A_T^2k_T^n`$ one obtains in the slow roll approximation LiSm
$`A_T^2`$ $``$ $`{\displaystyle \frac{8W}{75M_4^2}}F^2(H/\mu ),n_T=2ϵ_1,\text{ where }F(x)=\left[\sqrt{1+x^2}x^2\mathrm{sinh}^1(1/x)\right]^{1/2}`$ (248)
$`\mu `$ $`=`$ $`\sqrt{{\displaystyle \frac{4\pi }{3}}}{\displaystyle \frac{\sqrt{\lambda }}{M_4}}\text{and }x=H/\mu \left({\displaystyle \frac{W}{\lambda }}\right)^{1/2}`$ (249)
Combining these slow roll equations, at low energy one obtains the same consistency relation as for ordinary inflation,
$$\frac{A_T^2}{A_S^2}\frac{M_4^2W^2}{16\pi W^2}=ϵ_1=n_T/2.$$
(250)
At higher energies, however the relation is different. Furthermore, the tensor to scalar ratio $`R=(A_T/A_S)^2`$ and the spectral indices, both depend on the energy scale $`W`$. In Fig. 9 the behavior of the different quantities is indicated as function of the energy. Of course also the amplitude of the perturbations strongly depends on the parameter $`W/M_4^4`$.
In Fig. 10 two models are shown, quartic inflation with $`W=\alpha \phi ^4`$ and quadratic inflation with $`W=m^2\phi ^2`$. The lines of these models in the $`(R,n_s)`$ plane for varying $`\lambda `$ are drawn. The parapeters $`m`$ respectively $`\alpha `$ are chosen such that the scalar amplitude is compatible with the measured value, $`A_S^210^{10}`$ for each brane tension $`\lambda `$. The observational constraints from WMAP data WMAP are also indicated. It is clear, that quartic braneworld inflation fares even worse than ordinary quartic inflation. It is virtually excluded. Also quadratic inflation with strong braneworld effects fits the data somewhat less well than usual quadratic inflation, since it predicts too strong tensor contributions. But clearly, in lack of a concrete model of inflation (e.g., a given potential) there is little which can be said.
Discussing the effects on perturbations from braneworlds is opening Pandora’s box. There is a plethora of new phenomena some of which we don’t even know the sign. For example: during ordinary inflation, gravitational waves are generated. For a given inflationary potential, their amplitude can be calculated accurately. However, in the braneworld context, a fraction of these waves will be radiated into the bulk and thereby reduce the gravity wave amplitude. On the other hand, there is also gravity wave generation in the bulk, and some of these accumulate on the brane, increasing the amplitude of gravity waves on the brane. Therefore, depending on the precise realization, even the sign of the braneworld effect on a gravity wave background is unknown.
For a more concrete example, let us concentrate on scalar perturbations. We just take into account, that on the perturbative level the Weyl tensor $`E_{\mu \nu }`$ can no longer be neglected. Its energy density perturbation $`\delta \rho _E`$ acts like a radiation perturbation. In addition, however it can have an arbitrary amount of anisotropic stress, $`\mathrm{\Pi }_E`$. The latter induces a difference between the two Bardeen potentials, $`\mathrm{\Psi }\mathrm{\Phi }\mathrm{\Pi }_E`$. This affects mainly the Sachs–Wolfe term in the CMB fluctuations, hence the low multi-poles up to roughly the first peak.
In Fig. 11 we show the effect of a Weyl perturbation as function of an amplitude parameter
$$C_{\mathrm{dark}}\frac{\delta \rho _E/\rho _r}{4\zeta _m}.$$
(251)
Here $`\zeta _m`$ is the $`\zeta `$ variable defined in Eq. (160), due to ordinary matter (without the Weyl component). The anisotropic stress $`\mathrm{\Pi }_E`$ , on large scales, $`\mathrm{}1/k`$, can be determined as function of $`C_{\mathrm{dark}}`$ and $`\zeta _m`$. Confidence plots for the amplitude $`C_{\mathrm{dark}}`$ and several other cosmological parameters from the WMAP data are shown in Fig. 12.
There are many more effects which may come from perturbations in the bulk and the different perturbation equations on the brane which are presently under study. A systematic investigation is still lacking.
## 7 Conclusions
In these lectures we have studied the possibility that our Universe may represent a 3-brane in a higher dimensional space. This idea is motivated by string theory. We have especially investigated the case of one large extra dimension where the brane gravitational equations can be obtained purely from the bulk equations. Even though the resulting gravity on the brane differs strongly from Einstein gravity, we have seen that for an Anti–de Sitter bulk, Newton’s law is recovered at large distances and the Friedmann equations for the evolution of the Universe are obtained at low energy.
It is, however by no means clear, to which extent the different gravitational equations will spoil the successes of cosmological perturbation theory. This is still an open question and its answer will be crucial for braneworlds.
Acknowledgment: I am most grateful to the organizers for inviting me to this successful, traditional school in such a beautiful environment. I warmly thank also the students for their interest and active participation which was a most stimulating experience for me. I thank Marcus Ruser for carefully reading the manuscript. |
warning/0507/math0507558.html | ar5iv | text | # 0. Introduction
## 0. Introduction
Let $`G`$ be a connected reductive group over an algebraically closed field $`k`$, and $`W`$ the Weyl group of $`G`$. For a unipotent element $`uG`$, let $`_u`$ be the variety of Borel subgroups containing $`u`$. According to Springer \[Sp2\], Lusztig \[L1\], $`W`$ acts naturally on the $`l`$-adic cohomology group $`H^n(_u)=H^n(_u,\overline{𝐐}_l)`$, the so-called Springer representations of $`W`$. Assume that $`k=𝐂`$, or the characteristic $`p`$ of $`k`$ is good. Then it is known that $`H^{\mathrm{odd}}(_u)=0`$. We consider the graded $`W`$-module $`H^{}(_u)=_{n0}H^{2n}(_u)`$. Let $`L`$ be a Levi subgroup of a parabolic subgroup of $`G`$. Let $`W_L`$ be the Weyl group of $`L`$, which is naturally a subgroup of $`W`$. If $`uL`$, the variety $`_u^L`$ is defined by replacing $`G`$ by $`L`$, and we have a graded $`W_L`$-module $`H^{}(_u^L)`$.
Lusztig proved in \[L3\] an induction theorem for Springer representations, which describes the $`W`$-module structure of $`H^{}(_u)`$ in terms of the $`W_L`$-module structure of $`H^{}(_u^L)`$, in the case where $`uL`$. However in this theorem, the information on the graded $`W`$-module structure is eliminated. In this paper, we try to recover partly the graded $`W`$-module structure, i.e., for a fixed positive integer $`e`$, we consider the $`W`$-modules $`V_{e,k}=_{nkmode}H^{2n}(_u)`$ for $`k=0,\mathrm{},e1`$. Let $`G`$ be a simple group modulo center defined over $`𝐂`$. We show, under a certain choice of $`L`$, $`u`$ and $`e`$, that the $`W`$-module $`V_{e,k}`$ can be described in terms of the graded $`W_L`$-module $`H^{}(_u^L)`$ with some additional data. In particular, we see that $`dimV_{e,k}`$ is independent of the choice of $`k`$.
In the case where $`u=1`$, $`H^{}(_u)`$ is isomorphic, as a graded $`W`$-module, to the coinvariant algebra of $`W`$. In this case $`V_{e,k}`$ has been studied by many authors, by Stembridge \[St\] for $`e`$ corresponding to the regular elements in $`W`$, by Morita and Nakajima \[MN1\] for $`W=𝔖_n`$ with $`e`$ such that $`1en`$, and by Bonnafé, Lehrer and Michel \[BLM\] for complex reflection groups $`W`$ in the most general framework. Our result partly covers the result of \[BLM\]. For general $`u1`$, Morita and Nakajima \[MN2\] considered certain types of unipotent elements for $`G=GL_n`$, which is a special case of ours.
The proof of the induction theorem in \[L3\] is done by passing to the finite field $`𝐅_q`$, and using a certain specialization argument $`q1`$ together with the properties of Deligne-Lusztig’s virtual character $`R_T(1)`$. Our argument is a variant of that in \[L3\]. We use a specialization $`q\zeta `$, where $`\zeta `$ is a primitive $`e`$-th root of unity. Thus our argument is closely related to the values of Green functions at root of unity. In the case where $`u`$ is a regular unipotent element in $`L`$, we obtain an explicit formula for such values, which is regarded as a generalization of the result by Lascoux, Leclerc and Thibon \[LLT\] for the case of Green polynomials of $`GL_n`$.
## 1. The statement of the main result
1.1. Let $`k`$ be an algebraic closure of a finite field with ch$`(k)=p>0`$ or the complex number field $`𝐂`$. Let $`G`$ be a connected reductive group $`G`$ over $`k`$. Let $``$ be the variety of Borel subgroups of $`G`$, and $`W`$ the Weyl group of $`G`$. For any $`gG`$, put $`_g=\{B^{}gB^{}\}`$. We consider the Springer representations of $`W`$ on $`H^n(_g,\overline{𝐐}_l)`$ (or on $`H^n(_g,𝐂)`$ in the case where $`k=𝐂`$).
Let $`L`$ be a Levi subgroup of a parabolic subgroup $`P`$ of $`G`$. The Weyl group $`W_L`$ of $`L`$ is naturally identified with a subgroup of $`W`$. Let $`^L`$ be the variety of Borel subgroups of $`L`$. For a unipotent element $`uL`$, we consider $`_u^L=\{B^{}^LuB^{}\}`$. Thus we have a $`W_L`$-module $`H^n(_u^L,\overline{𝐐}_l)`$, and a $`W`$-module $`H^n(_u,\overline{𝐐}_l)`$. The induction theorem for Springer representations asserts that
(1.1.1)
$$\underset{n0}{}(1)^nH^n(_u,\overline{𝐐}_l)=\mathrm{Ind}_{W_L}^W\left(\underset{n0}{}(1)^nH^n(_u^L,\overline{𝐐}_l)\right)$$
as virtual $`W`$-modules.
Remark 1.2. The induction theorem was stated in \[AL\], with a brief indication of the proof, in the case where $`k=𝐂`$, and was proved in \[L3\] for any $`k`$. Note that if $`p`$ is good, the unipotent classes in $`G`$ are parametrized in the same way as the case of $`k=𝐂`$, independent of $`p`$. Moreover in that case, it is known that $`H^n(_u,\overline{𝐐}_l)=0`$ for odd $`n`$. Then the algorithm of computing Green functions implies that the $`W`$-module structure of $`H^n(_u,\overline{𝐐}_l)`$ is independent of $`p`$. Thus by a general principle $`H^n(_u,\overline{𝐐}_l)`$ is isomorphic to the $`W`$-module $`H^n(_u^{},𝐂)`$, where $`u^{},_u^{}`$ are the corresponding objects in the algebraic group $`G_𝐂`$ over $`𝐂`$. In what follows, we express $`H^n(_u,\overline{𝐐}_l)`$ or $`H^n(_u^{},𝐂)`$ by $`H^n(_u)`$ by abbreviation.
1.3. Assume that $`k=𝐂`$. We consider the following variant of the induction theorem. Let $`\mathrm{\Gamma }`$ be a cyclic group of order $`e`$ generated by $`a`$. Let $`\zeta `$ be a primitive $`e`$-th root of unity in $`𝐂`$. Let $`V=_{n0}V_n`$ be a graded $`W`$-module. Then $`V`$ turns out to be a $`\mathrm{\Gamma }\times W`$-module by defining the action of $`\mathrm{\Gamma }`$ on $`V`$ by $`ax=\zeta ^nx`$ for $`xV_n`$. We denote by $`V^{(\zeta )}`$ the thus obtained $`\mathrm{\Gamma }\times W`$-module $`V`$.
For $`uL`$, we consider the graded $`W_L`$-module $`H^{}(_u^L)=_{n0}H^{2n}(_u^L)`$, where the degree $`n`$ part is given by $`H^{2n}(_u^L)`$, and similarly we consider the graded $`W`$-module $`H^{}(_u)=_{n0}H^{2n}(_u)`$. Let $`\mathrm{\Gamma }`$ be as before. We choose $`\mathrm{\Gamma }`$ such that $`\mathrm{\Gamma }N_W(W_L)`$, and consider the semidirect product $`\stackrel{~}{W}_L=\mathrm{\Gamma }W_L`$. We assume that the $`W_L`$-module $`H^n(_u^L)`$ can be extended to a $`\stackrel{~}{W}_L`$-module for each $`n`$. (In the case where $`aZ_W(W_L)`$, we have $`\stackrel{~}{W}_L=\mathrm{\Gamma }\times W_L`$. In this case, one can choose a trivial extension to $`\stackrel{~}{W}_L`$, i.e., we may asssume that $`\mathrm{\Gamma }(\stackrel{~}{W}_L)`$ acts trivially on $`H^{}(_u^L)`$.) Then one can define a $`\mathrm{\Gamma }\times \stackrel{~}{W}_L`$-module $`H^{}(_u^L)`$ as above, replacing $`W_L`$ by $`\stackrel{~}{W}_L`$, which we denote by $`H^{}(_u^L)^{(\zeta )}`$. (When we need to distinguish the group $`\mathrm{\Gamma }`$ as the first factor of $`\mathrm{\Gamma }\times \stackrel{~}{W}_L`$ from the subgroup of $`\stackrel{~}{W}_L`$, we write the latter as $`\mathrm{\Gamma }_0`$.) $`\mathrm{\Gamma }\times W`$-module $`H^{}(_u)^{(\zeta )}`$ is defined as before. Put $`V^{(\zeta )}=H^{}(_u^L)^{(\zeta )}`$, and let $`V_n^{(\zeta )}`$ be the degree $`n`$-part of $`V^{(\zeta )}`$. Let us consider the induced $`W`$-module
$$\mathrm{Ind}_{W_L}^WV^{(\zeta )}=\underset{wW/W_L}{}wV^{(\zeta )}.$$
Then $`\mathrm{Ind}_{W_L}^WV^{(\zeta )}`$ turns out to be a $`\mathrm{\Gamma }\times W`$-module by defining the action of $`\mathrm{\Gamma }`$ by $`b(wx)=\zeta ^n(wb^1bx)`$ for $`b\mathrm{\Gamma }_0,xV_n^{(\zeta )}`$, which we denote by $`\mathrm{\Gamma }\mathrm{Ind}_{W_L}^WV^{(\zeta )}`$.
1.4. In the remainder of this paper, we assume that $`G`$ is simple modulo center. Let $`TB`$ be a pair of maximal torus and a Borel subgroup of $`G`$. Put $`W=N_G(T)/T`$. Let $`L`$ be a Levi subgroup of a parabolic subgroup $`P`$ of $`G`$ containing $`B`$ such that $`LT`$. We have $`W_L=N_L(T)/T`$. Let $`\mathrm{\Phi }X(T)`$ be a root system for $`G`$ with respect to $`T`$, with a simple root system $`\mathrm{\Pi }`$ (with respect to $`B`$), where $`X(T)`$ is the character group of $`T`$. We denote by $`\mathrm{\Phi }_L`$ the sub system of $`\mathrm{\Phi }`$ corresponding to $`L`$ with the simple root system $`\mathrm{\Pi }_L\mathrm{\Pi }`$. Let $`\mathrm{\Pi }^{}`$ be the set of simple roots which are orthogonal to $`\mathrm{\Pi }_L`$ with respect to the standard inner product on $`V=𝐑_𝐙X(T)`$. We denote by $`L^{}`$ the Levi subgroup containing $`T`$ corresponding to $`\mathrm{\Pi }^{}`$. Let $`W_L^{}=N_L^{}(T)/T`$ be the Weyl group of $`L^{}`$. Then we have $`WW_L\times W_L^{}`$, and so $`W_L^{}N_W(W_L)`$.
We recall here the notion of regular elements of reflection groups due to Springer \[Sp1\]. Let $`W`$ be a reflection group in $`GL(V)`$. A vector $`vV`$ is called regular if $`v`$ is not contained in any reflecting hyperplane in $`V`$. An element $`aW`$ is called regular if $`a`$ has an eigenvector $`v`$ which is a regular element in $`V`$. If $`av=\zeta v`$, with $`\zeta `$ a primitive $`e`$-th root of unity, then the order of $`a`$ is equal to $`e`$ (\[Sp1, 4.2\]). In particular, if $`a`$ is regular of order $`e`$, there exists an eigenvalue $`\zeta `$ which is a primitive $`e`$-th root of unity.
The regular elements $`aW`$ in the case of classical groups are given as follows (cf. \[Sp1\]).
Type $`A_{n1}`$. In this case $`W=𝔖_n`$ and there are two types of regular elements.
(a) $`e`$ is a divisor of $`n`$, and $`a`$ is an $`n/e`$-product of (disjoint) $`e`$-cycles in $`𝔖_n`$.
(b) $`e`$ is a divisor of $`n1`$, and $`a`$ is an $`(n1)/e`$-product of $`e`$-cycles in $`𝔖_n`$
Type $`B_n`$. There are two types of regular elements.
(a) $`e`$ is an odd divisor of $`n`$, and $`a`$ is an $`n/e`$-product of positive cycles of length $`e`$.
(b) $`e`$ is an even divisor of $`2n`$, and $`a`$ is a $`2n/e`$-product of negative cycles of length $`e/2`$.
Type $`D_n`$. In this case there are 4 types of regular elements.
(a) $`e`$ is an odd divisor of $`n`$, and $`a`$ is a product of positive cycles of length $`e`$.
(b) $`e`$ is an odd divisor of $`n1`$, and $`a`$ is a product of positive cycle of length 1 and $`(n1)/e`$ positive cycles of length $`e`$.
(c) $`n`$ is even, and $`e`$ is an even divisor of $`n`$. $`a`$ is a product of negative cycles of length $`e/2`$.
(d) $`e`$ is an even divisor of $`2n2`$, and $`a`$ is a product of $`(n1)/e`$ negative cycles of length $`e/2`$ and one cycle of length 1, which is positive or negative according as $`(2n2)/e`$ is even or odd.
Regular elements in the exceptional Weyl groups are listed in \[Sp1\].
Returning to the original setting, we consider the subgroups $`W_L,W_L^{}`$ of $`W`$. Let $`V^{}`$ be the subspace of $`V`$ generated by $`\mathrm{\Pi }_L^{}`$. $`W_L^{}`$ is realized as a reflection group on $`V^{}`$. Assume that $`a`$ is a regular element of $`W_L^{}`$ of order $`e`$. Let $`\zeta `$ be a primitive $`e`$-th root of unity, and $`V(a,\zeta )`$ the eigensubspace of $`a`$ in $`V`$ with eigenvalue $`\zeta `$. Since $`a`$ is regular, $`V(a,\zeta )`$ is not contained in any reflecting hyperplane $`H_\alpha `$ for $`\alpha \mathrm{\Phi }_L^{}`$. We say that $`a`$ is $`L`$-regular if $`V(a,\zeta )`$ is not contained in any $`H_\alpha `$ for $`\alpha \mathrm{\Phi }\mathrm{\Phi }_L`$. If $`L`$ is the torus $`T`$, all the regular elements are $`L`$-regular. But if $`LT`$, regular elements are not necessarily $`L`$-regular. For example, if $`L`$ is not simple modulo center, regular elements in $`W_L^{}`$ are not $`L`$-regular in many cases. In the case where $`L`$ is simple modulo center, $`L`$-regular elements are classified as follows.
###### Lemma 1.5.
Assume that $`L`$ is simple modulo center.
1. If $`W`$ is of type $`A_n,B_n,D_n`$, take $`L`$ such that $`W_L`$ is of the same type as $`W`$ of rank $`m`$, and $`W_L^{}`$ is of type $`A_{nm1}`$. Then a regular element of $`W_L^{}`$ of type (a) in 1.4 is $`L`$-regular.
2. If $`W`$ is of type $`G_2,F_4`$ or $`E_8`$, there does not exist $`L`$-regular elements for any $`LT`$.
3. Assume that $`W`$ is of type $`E_6`$ or $`E_7`$. Let $`\mathrm{\Pi }=\{\alpha _1,\mathrm{},\alpha _7\}`$ (resp. $`\{\alpha _1,\mathrm{},\alpha _6\}`$) be the set of simple roots in $`E_7`$ (resp. in $`E_6`$) as in the figure. Take $`\mathrm{\Pi }_L=\{\alpha _k,\alpha _{k+1},\mathrm{},\alpha _7\}`$ (resp. $`\{\alpha _k,\alpha _{k+1},\mathrm{},\alpha _6\}`$) for $`k3`$. Then $`W_L^{}`$ is of type $`A_j`$ or of type $`A_j+A_1`$ for some $`j`$ except the case where $`W`$ is of type $`E_7`$ and $`\mathrm{\Pi }_L=\{\alpha _7\}`$, in which case $`\mathrm{\Pi }_L^{}`$ is of type $`D_5`$. In the former case, we choose $`a`$ a regular element of type (a) for type $`A`$, and in the latter case, we choose $`a`$ a regular element of type (a) for type $`D`$ in 1.4, respectively. Then $`a`$ is $`L`$-regular.
###### Proof.
If there exists $`\beta \mathrm{\Phi }\mathrm{\Phi }_L`$ such that $`\beta `$ is orthogonal to $`V_L^{}`$, then any regular element in $`W_L^{}`$ cannot be $`L`$-regular. By direct inspections, one can find such $`\beta `$ unless $`L`$ is the type given in (i), (iii) of the lemma. Assume that $`L`$ is as in the lemma, and let $`a`$ be a regular element in $`W_L^{}`$. If $`W_L^{}`$ is of type $`A_j`$ or type $`A_j+A_1`$, then a regular vector $`vV^{}`$ can be written explicitly, and one can check the $`L`$-regularity by direct inspections. If $`W_L^{}`$ is of type $`D_5`$ (in the case where $`W`$ is of type $`E_7`$), $`a`$ must be of type (a) (otherwise it is easy to see that $`a`$ is not $`L`$-regular). But this element is nothing but the regular element in $`A_4`$, and the checking is reduced to the previous case. The details are omitted. ∎
1.6. In what follows we consider a specific cyclic group $`\mathrm{\Gamma }N_W(W_L)`$, and $`uL`$ according to the following two cases.
Case (a): $`W_L^{}\{1\}`$.
In this case, we assume that $`L`$ is simple modulo center. We choose an $`L`$-regular element $`aW_L^{}`$, and put $`\mathrm{\Gamma }=\mathrm{}a\mathrm{}`$. Let $`e`$ be the order of $`\mathrm{\Gamma }`$. Thus $`\mathrm{\Gamma }W_L^{}`$ and we have $`\mathrm{\Gamma }\times W_LW`$. We take any unipotent element $`uL`$.
Case (b): $`W_L^{}=\{1\}`$.
In this case, we assume that $`L`$ is of type $`X_0+e(A_{n_11}+\mathrm{}+A_{n_r1})`$ with $`X_0`$ irreducible. We further assume that any $`\beta \mathrm{\Phi }\mathrm{\Phi }_L`$ is not orthogonal to the root system $`e(A_{n_11}+\mathrm{}+A_{n_r1})`$. (Note: since $`W_L^{}=\{1\}`$, any irreducible component of the Dynkin diagram corresponding to $`\mathrm{\Pi }\mathrm{\Pi }_L`$ consists of 1 or 2 nodes. The latter condition is satisfied for type $`B_n`$ if all the irreducible components consist of one node, and for type $`A_n,D_n`$ if the number of irreducible components having two nodes is at most 1.)
We choose $`aW`$ so that $`a`$ permutes each component $`A_{n_i1}`$ in a cyclic way, and acts trivially on $`X_0`$. Thus $`a𝔖_{en_1}\times \mathrm{}\times 𝔖_{en_r}`$, and $`a`$ is a product of disjoint cycles of length $`e`$. In particular, $`\mathrm{\Gamma }=\mathrm{}a\mathrm{}N_W(W_L)`$, and the subgroup of $`W`$ generated by $`\mathrm{\Gamma }`$ and $`W_L`$ coincides with the semidirect product $`\mathrm{\Gamma }W_L`$. Now $`L`$ is isogenic to $`G_0\times G_1\times \mathrm{}\times G_r`$ modulo center, where $`G_0`$ is of type $`X_0`$, and $`G_iGL_{n_i}\times \mathrm{}\times GL_{n_i}`$ ($`e`$-factors). We choose a unipotent element $`uL`$ so that $`u`$ corresponds to $`(u_0,u_1,\mathrm{},u_r)`$, where $`u_0G_0`$ is arbitrary, and $`u_i`$ is a diagonal element in $`G_i`$, i.e., $`u_i=(v_i,\mathrm{},v_i)`$ with $`v_iGL_{n_i}`$ for $`i=1,\mathrm{},r`$.
We can state our main theorem, whose proof will be given in the next section.
###### Theorem 1.7.
Assume that $`G`$ is defined over $`𝐂`$. Let $`L`$ be a Levi subgroup in $`G`$. Assume that a cyclic subgroup $`\mathrm{\Gamma }`$ of order $`e`$ in $`N_W(W_L)`$ and $`uL`$ are given as in 1.4. Put $`\stackrel{~}{W}_L=\mathrm{\Gamma }W_L`$. Then the followings hold.
1. $`W_L`$-module $`H^{}(_u^L)`$ can be extended to a $`\stackrel{~}{W}_L`$-module so that $`\mathrm{\Gamma }\times \stackrel{~}{W}_L`$-module $`H^{}(_u^L)^{(\zeta ^{})}`$ is defined for any $`e`$-th root of unity $`\zeta ^{}`$.
2. There exists a primitive $`e`$-th root of unity $`\zeta `$ such that
(1.7.1)
$$\mathrm{\Gamma }\mathrm{Ind}_{W_L}^W\left(H^{}(_u^L)^{(\zeta )}\right)H^{}(_u)^{(\zeta )}$$
as $`\mathrm{\Gamma }\times W`$-modules.
Remarks 1.8. (i) The extension of $`W_L`$-module $`H^{}(_u^L)`$ to $`\stackrel{~}{W}_L`$-module is not unique. The theorem aaserts that the statement (ii) holds for some choice of extension.
(ii) The theorem asserts that (1.7.1) holds for some choice of primitive $`e`$-th root of unity $`\zeta `$, but then it holds for any choice of primitive root of unity $`\zeta ^{}`$. In fact, we can write $`\zeta ^{}=\zeta ^j`$ for some $`j`$ prime to $`e`$, and we have an automorphism $`\tau `$ on $`\mathrm{\Gamma }`$ such that $`\tau (a)=a^j`$. It follows from (1.7.1) that we have an isomorphism of $`\mathrm{\Gamma }\times W`$ modules, where the action of $`\mathrm{\Gamma }`$ is twisted by $`\tau `$. It is easy to check that the twisted $`\mathrm{\Gamma }\times W`$-module $`\mathrm{\Gamma }\mathrm{Ind}_{W_L}^W\left(H^{}(_u^L)^{(\zeta )}\right)`$ is isomorphic to $`\mathrm{\Gamma }\mathrm{Ind}_{W_L}^W\left(H^{}(_u^L)^{(\zeta ^{})}\right)`$, and similarly the twisted $`H^{}(_u)^{(\zeta )}`$ is isomorphic to $`H^{}(_u)^{(\zeta ^{})}`$. Thus (1.7.1) holds also for $`\zeta ^{}`$.
## 2. Proof of Theorem 1.7
2.1. In the case where $`e=1`$, Theorem 1.7 is nothing but the original induction theorem. So we assume that $`e2`$ in what follows. Since the structure of the $`W`$-module $`H^n(_u)`$ is independent of $`p`$ provided that $`p`$ is a good prime, it is enough to show the corresponding formula for an appropriate $`p`$. So, we assume that $`G`$ is defined over $`𝐅_p`$, of split type, with Frobenius map $`F`$. We assume that $`TB`$ are both $`F`$-stable, and that $`LP`$ are $`F`$-stable. Thus $`F`$ acts trivially on $`W`$ and on $`W_L`$. We first note that
###### Lemma 2.2.
Let $`aN_W(W_L)`$ and choose $`\dot{a}N_G(T)N_G(L)`$. Assume that $`\dot{a}Z_G(u)`$. Then $`\mathrm{ad}\dot{a}`$ stabilizes $`_u^L`$, and acts on $`H^{}(_u^L)`$ in such a way that $`\mathrm{ad}\dot{a}(w)=awa^1`$ for $`wW_L`$.
###### Proof.
Since $`\dot{a}N_G(L)`$, $`\dot{a}`$ acts on $`^L`$ by the adjoint action $`\mathrm{ad}\dot{a}`$, which stabilizes $`_u^L`$ since $`\dot{a}Z_G(u)`$. Hence $`\dot{a}`$ acts naturally on $`H^{}(_u^L)`$. In order to compare this action with the action of $`W_L`$, we shall recall the construction of Springer representations of $`W_L`$. Let
$$\stackrel{~}{L}=\{(x,gB)L\times ^Lg^1xgB\},$$
and $`\pi :\stackrel{~}{L}L`$ be the first projection. Let $`L_r`$ be the set of regular semisimple elements in $`L`$. Then $`\pi ^1(L_r)`$ is isomorphic to
$$\stackrel{~}{L}_r=T_r\times L/T,$$
where $`T_r=TL_r`$. Let $`\pi _0:\stackrel{~}{L}_rL_r`$ be the map defined by $`\pi _0:(t,gT)g^1tg`$, which coincides with the restriction of $`\pi `$ on $`\stackrel{~}{L}_r`$ under the identification $`\pi ^1(L_r)\stackrel{~}{L}_r`$. Then $`\pi _0`$ is an unramified Galois covering with group $`W_L`$, and for a constant sheaf $`\overline{𝐐}_l`$ on $`\stackrel{~}{L}_r`$, $`=\pi _{}\overline{𝐐}_l`$ is a $`W_L`$-equivariant local system on $`L_r`$. Thus $`K=\mathrm{IC}(L,)`$ is a $`W_L`$-equivariant complex on $`L`$, and it is known by Lusztig that $`K\pi _{}\overline{𝐐}_l`$. Thus for each $`uL`$, the stalk $`_u^i(K)`$ at $`u`$ of the $`i`$-th cohomology sheaf of $`K`$ gives rise to a $`W_L`$-module $`H^i(_u^L)`$.
Now $`\dot{a}`$ acts on $`\stackrel{~}{L}_r`$ (resp. on $`L_r`$) by $`\mathrm{ad}\dot{a}:(t,gT)(\dot{a}t\dot{a}^1,\dot{a}g\dot{a}^1T)`$ (resp. $`\mathrm{ad}\dot{a}:x\dot{a}x\dot{a}^1`$), and $`\pi _0`$ commutes with $`\mathrm{ad}\dot{a}`$. Hence $``$ becomes an $`\dot{a}`$-equivariant local system. Since $`\pi _0^1(t)=\{(wtw^1,wT)wW_L\}`$ for $`tT_r`$, the stalk $`_t`$ has a natural structure of the regular $`W_L`$-module. Then the isomorphism $`_{\dot{a}t\dot{a}^1}_t`$ is given by $`\mathrm{ad}\dot{a}^1`$ under the identification $`_x\overline{𝐐}_l[W_L]`$ for $`xL_r`$. It follows that $``$ is $`\mathrm{}\dot{a}\mathrm{}W_L`$-equivariant, where $`\mathrm{}\dot{a}\mathrm{}`$ is a cyclic group generated by $`\dot{a}`$, and $`\dot{a}`$ acts on $`W_L`$ by $`\mathrm{ad}\dot{a}(w)=awa^1`$. By the functoriality of $`\mathrm{IC}`$ functor, $`K`$ turns out to be a $`\dot{a}`$-equivariant complex on $`L`$ under the adjoint action of $`\dot{a}`$, which is regarded as a $`\mathrm{}\dot{a}\mathrm{}W_L`$-equivariant complex on $`L`$. Hence for $`uL`$ such that $`\dot{a}u\dot{a}^1=u`$, $`_u^i(K)`$ has a structure of $`\mathrm{}\dot{a}\mathrm{}W_L`$-module.
On the other hand, $`\dot{a}`$ acts naturally on $`\stackrel{~}{L}`$ and on $`L`$ by the adjoint action, which commute with $`\pi `$. Thus $`\pi _{}\overline{𝐐}_l`$ is $`\dot{a}`$-equivariant, which is isomorphic to $`K`$ as the complex with $`\dot{a}`$-action. Hence the action of $`\dot{a}`$ on $`_u^i(K)`$ coincides with the action on $`H^i(_u^L)`$ induced from the adjoint action of $`\dot{a}`$ on $`_u^L`$. The lemma follows from this. ∎
Next we show the following lemma.
###### Lemma 2.3.
There exists a representative $`\dot{a}N_G(T)N_G(L)Z_G(u)`$ such that $`\dot{a}`$ acts trivially on $`H^{}(_u)`$ and that $`\dot{a}^e`$ acts trivially on $`H^{}(_u^L)`$. In particular, $`H^{}(_u^L)`$ has a structure of $`\stackrel{~}{W}_L`$-module.
###### Proof.
First consider the case (a) in 1.6. Let $`H`$ be the subgroup of $`G`$ generated by $`U_\alpha `$ with $`\alpha \mathrm{\Phi }_L^{}`$, where $`U_\alpha `$ is the root subgroup corresponding to $`\alpha `$. Then $`H`$ is a connected reductive subgroup of $`L^{}`$ whose Weyl group coincides with $`W_L^{}`$. Since $`HZ_G(u)`$, we have $`HZ_G^0(u)`$. One can choose a representative $`\dot{a}N_H(T_1)`$ of $`aW_L^{}`$, where $`T_1`$ is a maximal torus of $`H`$ contained in $`T`$. Then $`\dot{a}Z_G^0(u)N_G(L)`$ and $`\dot{a}^eT_1`$. Since $`T_1Z_G(u)`$, we see that $`T_1Z_L^0(u)`$. Thus, $`\dot{a}^eZ_L^0(u)`$. Hence $`\dot{a}`$ satisfies the condition.
Next consider the case (b) in 1.6. Let $`L_1`$ be the Levi subgroup containing $`L`$ of type $`X_{n_0}+A_{en_11}+\mathrm{}+A_{en_r1}`$. We have a natural projection $`\pi :L_1\overline{L}_1=L_1/Z^0(L_1)`$, and an isogeny map $`\theta :\stackrel{~}{L}_1=G_0\times SL_{en_1}\times \mathrm{}\times SL_{en_r}\overline{L}_1`$, where $`G_0`$ is the simply connected semisimple group of type $`X_0`$. Put $`\overline{u}=\pi (u)\overline{L}_1`$. Now $`Z_{L_1}(u)`$ acts on $`H^{}(_u)`$. Since $`Z^0(L_1)`$ acts trivially on $`H^{}(_u)`$, we have an action of $`Z_{L_1}(u)/Z^0(L_1)=Z_{\overline{L}_1}(\overline{u})`$ on $`H^{}(_u)`$. Let $`\stackrel{~}{u}`$ be an element in $`\stackrel{~}{L}_1`$ such that $`\theta (\stackrel{~}{u})=\overline{u}`$. $`\stackrel{~}{u}=(u_0,u_1,\mathrm{},u_r)`$ can be chosen as given in 1.4. We choose $`\ddot{a}\stackrel{~}{L}_1`$ as follows; put $`\ddot{a}=(a_0,a_1,\mathrm{},a_r)`$ with $`a_0G_0`$, and $`a_iSL_{en_i}`$ for $`1ir`$. We put $`a_0=1`$ and choose $`a_1,\mathrm{},a_r`$ so that $`a_iZ_{SL_{en_i}}^0(u_i)`$ and that $`a_i^eZ(SL_{en_i})`$. Such a choice is always possible for $`u_i`$ of type $`(n_i,\mathrm{},n_i)`$. Thus $`\ddot{a}Z_{\stackrel{~}{L}_1}^0(\stackrel{~}{u})`$. It follows that $`\theta (\ddot{a})`$ is contained in a connected subgroup of $`Z_{\overline{L}_1}(\overline{u}_1)`$, and by the previous remark, $`\theta (\ddot{a})`$ acts trivially on $`H^{}(_u)`$. Now take $`\dot{a}Z_{L_1}(u)`$ such that $`\pi (\dot{a})=\theta (\ddot{a})`$. Then $`\dot{a}N_G(T)N_G(L)`$, and acts trivially on $`H^{}(_u)`$. On the other hand, similar to $`\pi ,\theta `$, we have a map $`\pi ^{}:L\overline{L}=L/Z^0(L)`$ and $`\theta ^{}:\stackrel{~}{L}=G_0\times (SL_{n_1})^e\times \mathrm{}\times (SL_{n_r})^e\overline{L}`$. Let $`\overline{u}=\pi ^{}(u)\overline{L}`$, and $`\stackrel{~}{u}\stackrel{~}{L}`$ such that $`\overline{u}=\theta ^{}(\stackrel{~}{u})`$. Then we have an isomorphism $`H^{}(_u^L)H^{}(_{\overline{u}}^{\overline{L}})H^{}(_{\stackrel{~}{u}}^{\stackrel{~}{L}})`$ compatible with the actions of $`Z_L(u),Z_{\overline{L}}(\overline{u})`$ and $`Z_{\stackrel{~}{L}}(\stackrel{~}{u})`$ with respect to $`\pi ^{},\theta ^{}`$. We have $`\ddot{a}^eZ(SL_{n_1})^e\times Z(SL_{n_2})^e\times \mathrm{}`$. Since the action of $`Z(SL_{n_1})^e\times Z(SL_{n_2})^e\times \mathrm{}`$ can be extended to an action of $`Z(GL_{n_1})^e\times Z(GL_{n_2})^e\times \mathrm{}`$ on $`H^{}(_{\stackrel{~}{u}}^{\stackrel{~}{L}})`$, $`\ddot{a}^e`$ acts trivially on $`H^{}(_{\stackrel{~}{u}}^{\stackrel{~}{L}})`$, and so $`\dot{a}^e`$ acts trivially on $`H^{}(_u^L)`$. ∎
2.4. Let $`𝒵=Z_L^0`$ be the identity component of the center of $`L`$. Put $`_𝒵=\{B^{}𝒵B^{}\}`$. Then $`_𝒵`$ is decomposed into connected components
$$_𝒵=\underset{dW_L\backslash W}{}_{𝒵,d},$$
where $`_{𝒵,d}=\{{}_{}{}^{xd}BxL\}`$, which is isomorphic to $`^L`$ under the map $`B^{}B^{}L`$.
Put
$$𝒵_{\mathrm{reg}}=\{z𝒵Z_G^0(z)=L\}.$$
Then for any $`t𝒵_{\mathrm{reg}}`$, we have $`_t=_𝒵`$ by Lemma 2.2 (c) in \[L3\], and so $`_{tu}=_u_t=_u_𝒵`$. It follows that
$$_{tu}=\underset{dW_L\backslash W}{}(_{𝒵,d}_u),$$
where $`_{𝒵,d}_u`$ is isomorphic to $`_u^L`$ under the map $`B^{}B^{}L`$. This implies that
(2.4.1)
$$H^{2n}(_{tu})\underset{d^1W/W_L}{}H^{2n}(_{𝒵,d}_u).$$
The right hand side of (2.4.1) has a natural structure of the induced $`W`$-module $`\mathrm{Ind}_{W_L}^WH^{2n}(_u^L)`$. It is proved in \[L3, Proposition 1.4\] that (2.4.1) is actually an isomorphism of $`W`$-modules. Let $`aW`$ be as in the theorem. Since $`\dot{a}N_G(L)`$, it stabilizes $`𝒵`$, and so $`\dot{a}`$ acts on $`_𝒵`$ via $`\mathrm{ad}\dot{a}`$. It is easy to see that $`\dot{a}`$ induces a permutation action on the components of $`_𝒵`$; $`\dot{a}:_{𝒵,d}_{𝒵,ad}`$. It follows that $`\dot{a}`$ induces an automorphism on $`H^{2n}(_{tu})`$, which maps the factor corresponding to $`d^1W/W_L`$ to $`d^1a^1W/W_L`$. Under the isomorphism $`H^{2n}(_{𝒵,d}_u)H^{2n}(_u^L)`$, the factor corresponding to $`d^1W/W_L`$ is written as $`d^1H^{2n}(_u^L)`$, and $`\dot{a}`$ maps $`d^1H^{2n}(_u^L)d^1a^1H^{2n}(_u^L)`$. On the other hand, by Lemma 2.3, $`\dot{a}^e`$ acts trivially on $`H^{2n}(_u^L)`$, and induces an action of $`\stackrel{~}{W}_L`$ on it. Hence $`\dot{a}`$ induces an action of $`\mathrm{\Gamma }_0`$ on $`H^{2n}(_{tu})\mathrm{Ind}_{W_L}^WH^{2n}(_u^L)`$, which is given by $`\dot{a}:d^1xd^1a^1\dot{a}x`$ for each factor $`d^1H^{2n}(_u^L)`$.
Now we define an action of $`\mathrm{\Gamma }`$ on $`H^{}(_{tu})`$ by $`a:x\zeta ^n\dot{a}x`$ for $`xH^{2n}(_{tu})`$, where $`\dot{a}x`$ is the action of $`\mathrm{\Gamma }_0`$ on $`H^{2n}(_{tu})`$ given as above. Since the action of $`\dot{a}G`$ commutes with that of $`W`$, $`H^{}(_{tu})`$ turns out to be a $`\mathrm{\Gamma }\times W`$-module, which we denote by $`H^{}(_{tu})^{[\zeta ]}`$. The following lemma is immediate from the above discussion.
###### Lemma 2.5.
There exists an isomorphism of $`\mathrm{\Gamma }\times W`$-modules
$$H^{}(_{tu})^{[\zeta ]}\mathrm{\Gamma }\mathrm{Ind}_{W_L}^WH^{}(_u^L)^{(\zeta )}.$$
In view of Lemma 2.5, in order to prove the theorem it is enough to show the following proposition.
###### Proposition 2.6.
Under an appropriate choice of (a good prime) $`p`$, there exists an isomorphism of $`\mathrm{\Gamma }\times W`$-modules for any $`t𝒵_r`$,
$$H^{}(_u)^{(\zeta )}H^{}(_{tu})^{[\zeta ]}.$$
2.7. The remainder of this section is devoted to the proof of the proposition. We shall prove it by modifying the arguments in \[L3\]. By \[Sh1\], \[Sh2\], \[BS\], the following fact is known; assume that $`G`$ is simple modulo center. Then for each unipotent class $`C`$ of $`G`$, there exists $`u_1C^F`$, called a split unipotent element, such that $`F`$ acts on $`H^{2n}(_{u_1})`$ as a scalar multiplication by $`p^n`$. (In the case where $`G`$ is of type $`E_8`$, we assume that $`p1(mod4)`$). Since the component group $`A_G(u_1)=Z_G(u_1)/Z_G^0(u_1)`$ is isomorphic to $`S_3,S_4,S_5`$ or $`(𝐙/2𝐙)^k`$ for some $`k`$, there exists a positive integer $`s_0`$ (independent of $`p`$) such that $`F^{s_0}`$ acts on $`H^{2n}(_u)`$ by a scalar multiplication by $`p^{s_0n}`$ for any unipotent element $`u`$ of $`G^F`$ (e.g., one can take $`s_0=|S_5|`$.) Similarly, $`F^{s_0}`$ acts on $`H^{2n}(_u^L)`$ by a scalar multiplication by $`p^{s_0n}`$ for any unipotent element $`uL^F`$. Note that the isomorphism in (2.4.1) is $`F`$-equivariant. Hence $`F^{s_0}`$ acts also as a scalar multiplication by $`p^{s_0n}`$ for $`H^{2n}(_{tu})`$.
Note that $`\dot{a}`$ acts trivially on $`H^{2n}(_u)`$ by Lemma 2.3. It follows that one can write
(2.7.1) $`\mathrm{Tr}((F^s\dot{a})^iw,H^{}(_u))`$ $`={\displaystyle \underset{n0}{}}a_n(w)p^{isn},`$
(2.7.2) $`\mathrm{Tr}((w,a^i),H^{}(_u)^{(\zeta )})`$ $`={\displaystyle \underset{n0}{}}a_n(w)\zeta ^{in},`$
for any $`wW,0ie1`$ and for any positive integer $`s`$ divisible by $`s_0`$, where $`a_n(w)=\mathrm{Tr}(w,H^n(_u))`$ are integers for each $`n0`$.
On the other hand, by the description of the action of $`F`$ and of $`\dot{a}`$ on $`H^n(_{tu})`$ in 2.4, together with Lemma 2.5, one can write
(2.7.3) $`\mathrm{Tr}((F^s\dot{a})^iw,H^{}(_{tu}))`$ $`={\displaystyle \underset{n0}{}}b_{n,i}(w)p^{isn},`$
(2.7.4) $`\mathrm{Tr}((w,a^i),H^{}(_{tu})^{[\zeta ]})`$ $`={\displaystyle \underset{n0}{}}b_{n,i}(w)\zeta ^{in},`$
for $`w,i,s`$ as above, where $`b_{n,i}(w)`$ are certain integers.
For an integer $`x`$ and a prime number $`l`$, we denote by $`m_l(x)`$ the multiplicative order of $`x`$ in $`𝐙/l𝐙`$, i.e., the smallest positive integer $`m`$ such that $`x^m1(modl)`$. The following is a key for the proof of Proposition 2.6.
###### Lemma 2.8.
Assume that $`p1(mod4)`$. Let $`s_0,e`$ be fixed positive integers coprime to $`p`$. Then there exist infinitely many prime numbers $`l`$ satisfying the following properties.
1. $`m_l(p^s)=e`$ for a certain integer $`s`$ divisible by $`s_0`$.
2. $`l1`$ is divisible by $`e`$.
###### Proof.
By our assumption, the image of $`s_0e`$ on $`𝐅_p=𝐙/p𝐙`$ is non-zero. Hence the map $`xs_0ex+1`$ induces a bijective map on $`𝐅_p`$. Thus there exists $`c𝐙`$ such that the image of $`s_0ec+1`$ in $`𝐅_p`$ is contained in $`𝐅_p^{}(𝐅_p^{})^2`$. Put $`\alpha =s_0ec+1`$. Then $`\alpha `$ is prime to $`p`$, and so $`(\alpha 1)p`$ and $`\alpha `$ are coprime each other. Then by Dirichlet’s theorem on arithmetic progression, there exist infinitely many prime numbers $`l`$ of the form $`l=n(\alpha 1)p+\alpha `$ for some positive integer $`n`$. It is enough to show that these $`l3`$ satisfy the assertion of the lemma. For an integer $`a`$ and a prime number $`p`$, let $`\left({\displaystyle \frac{a}{p}}\right)`$ be the Legendre symbol, i.e.,
$$\left(\frac{a}{p}\right)=\{\begin{array}{cc}1\hfill & \text{ if }x^2a(modp)\text{ for some }x𝐙,\hfill \\ 1\hfill & \text{ otherwise.}\hfill \end{array}$$
We show that
(2.8.1)
$$\left(\frac{p}{l}\right)=1.$$
In fact, by the quadratic reciprocity law (e.g., \[Se\]), we have
$$\left(\frac{p}{l}\right)\left(\frac{l}{p}\right)=(1)^{\frac{p1}{2}\frac{l1}{2}}=1.$$
The second equality follows from the assumption that $`p1(mod4)`$. Hence we have $`\left({\displaystyle \frac{p}{l}}\right)=\left({\displaystyle \frac{l}{p}}\right)`$. But $`l\alpha (modp)`$, and so $`\left({\displaystyle \frac{l}{p}}\right)=\left({\displaystyle \frac{\alpha }{p}}\right)=1`$ since the image of $`\alpha `$ is not contained in $`𝐅_p^2`$ by our choice of $`\alpha `$. Hence (2.8.1) holds.
Now (2.8.1) is equivalent to $`p^{(l1)/2}1(modl)`$. It follows that $`m_l(p)=l1`$. Since $`l1=s_0ec(np+1)`$, we see that $`m_l(p^s)=e`$ for $`s=s_0c(np+1)`$ and that $`l1`$ is divisible by $`e`$. Thus this $`l`$ satisfies the assertion of the lemma. The lemma is proved. ∎
2.9 For given integers $`s_01,e2`$, we choose a prime number $`p`$ such that $`p`$ is not a factor of $`e,s_0`$ and that $`p1(mod4)`$, and fix it once and for all. For a multiple $`s`$ of $`s_0`$, put $`F^{}=F^s\dot{a}`$ and $`q=p^s`$. Under the setting in 1.6, we shall describe the set $`𝒵_{\mathrm{reg}}`$ more precisely. As in \[L3, Lemma 2.2\], $`𝒵_{\mathrm{reg}}`$ can be written as $`𝒵_{\mathrm{reg}}=𝒵_\beta \mathrm{ker}(\beta |_𝒵)`$, where $`\beta `$ runs over all the roots in $`\mathrm{\Phi }\mathrm{\Phi }_L`$. ($`\beta |_𝒵`$ gives a non-trivial character of $`𝒵`$ for $`\beta \mathrm{\Phi }\mathrm{\Phi }_L`$).
First consider the case (a). Let $`L^{}\mathrm{der}`$ be the derived subgroup of $`L^{}`$, and $`S^{}`$ be the split maximal torus of $`L^{}\mathrm{der}`$ contained in $`T`$. Then $`S^{}𝒵`$. Put $`S_{\mathrm{reg}}^{}=S^{}𝒵_{\mathrm{reg}}`$. Now $`W_L^{}`$ leaves the set $`\mathrm{\Phi }\mathrm{\Phi }_L`$ invariant. For each $`\beta \mathrm{\Phi }\mathrm{\Phi }_L`$, put $`H_\beta =_{x\mathrm{\Gamma }}\mathrm{ker}(x(\beta )|_S^{})`$. Then $`H_\beta `$ is an $`F^{}`$-stable subgroup of $`S^{}`$, and we see that
(2.9.1)
$$S_{}^{}{}_{\mathrm{reg}}{}^{F^{}}=S_{}^{}{}_{}{}^{F^{}}\underset{\beta \mathrm{\Phi }\mathrm{\Phi }_L}{}H_\beta ^F^{}.$$
$`H_\beta `$ is a closed subgroup of $`S^{}`$, and we put $`e_\beta =|H_\beta /H_\beta ^0|`$ for each $`\beta \mathrm{\Phi }\mathrm{\Phi }_L`$.
Let $`𝒫^{}`$ be the set of all prime numbers $`l`$ satisfying the condition in Lemma 2.8. Thus $`𝒫^{}`$ is an infinite set. We denote by $`𝒫`$ the subset of $`𝒫^{}`$ consisting of $`l`$ such that $`l>|\mathrm{\Phi }\mathrm{\Phi }_L|`$ and that $`l`$ does not divide $`e_\beta `$ ($`\beta \mathrm{\Phi }\mathrm{\Phi }_L)`$. Thus $`𝒫`$ is an infinite set also.
Next we consider the case (b). We may assume that $`G`$ has a connected center of dimension 1, and that the derived subgroup of $`G`$ is simply connected, almost simple. Let $`k`$ be an algebraic closure of $`𝐅_q`$. We see that there exists a subtorus $`S`$ of $`𝒵`$ such that $`S(k^{})^c`$, where $`c`$ is the number of irreducible components of $`\mathrm{\Phi }_L`$. Since $`a`$ permutes the factors $`k^{}`$ in $`S`$, we see that $`S^F^{}(𝐅_{q^e}^{})^r\times (𝐅_q^{})^r^{}`$, where $`r^{}`$ is equal to 1 or 0 according to the cases where $`X_0`$ is non-empty or empty. Since $`\mathrm{\Gamma }N_W(W_L)`$, $`\mathrm{\Gamma }`$ preserves the set $`\mathrm{\Phi }\mathrm{\Phi }_L`$. For each $`\beta \mathrm{\Phi }\mathrm{\Phi }_L`$, put $`K_\beta =_{x\mathrm{\Gamma }}\mathrm{ker}(x(\beta )|_S)`$. Then $`K_\beta `$ is an $`F^{}`$-stable subgroup of $`S`$, and we have
(2.9.2)
$$S_{\mathrm{reg}}^F^{}=S^F^{}\underset{\beta \mathrm{\Phi }\mathrm{\Phi }_L}{}K_\beta ^F^{},$$
where $`S_{\mathrm{reg}}=S𝒵_{\mathrm{reg}}`$. $`K_\beta `$ is a closed subgroup of $`S`$, and put $`e_\beta =|K_\beta /K_\beta ^0|`$ for each $`\beta \mathrm{\Phi }\mathrm{\Phi }_L`$. Under the identification $`S^F^{}(𝐅_{q^e}^{})^r\times (𝐅_q^{})^r^{}`$, we see that $`K_\beta ^{0F^{}}(𝐅_{q^e}^{})^{r1}\times (𝐅_q^{})^r^{}`$ or $`K_\beta ^{0F^{}}𝐅_{q^e^{}}^{}\times (𝐅_{q^e}^{})^{r1}\times (𝐅_q^{})^r^{}`$, where $`e^{}`$ is a proper divisor of $`e`$. (Let $`S_i`$ be the subtorus of $`S`$ corresponding to the factor $`eA_{n_i1}`$ for $`i=1,\mathrm{},r`$. Then the former case occurs if $`\beta |_{S_i},\beta |_{S_j}`$ are non-trivial for some $`ij`$, and the latter case occurs if $`\beta |_{S_i}`$ is non-trivial for only one $`i`$. Note that by our assumption in 1.4, $`\beta `$ is non-trivial on $`S_1\times \mathrm{}\times S_r`$.)
Let $`𝒫^{}`$ be as in the case (a). We define a subset $`𝒫`$ of $`𝒫^{}`$ as the set of prime numbers $`l𝒫^{}`$ such that $`l>|\mathrm{\Phi }\mathrm{\Phi }_L|`$ and that $`l`$ does not divide $`e_\beta `$.
The next lemma is a variant of Lemma 3.4 in \[L3\].
###### Lemma 2.10.
Assume that $`l𝒫`$, and let $`s`$ be a multiple of $`s_0`$ such that $`m_l(p^s)=e`$ (see Lemma 2.8). Put $`F^{}=F^s\dot{a}`$. Then there exists $`t𝒵_{\mathrm{reg}}`$ such that $`F^{}(t)=t`$ and that $`t^l=1`$.
###### Proof.
First consider the case (a) in 1.6. It is enough to show, for each $`l𝒫`$, that there exists $`tS_{}^{}{}_{\mathrm{reg}}{}^{F^{}}`$ such that $`t^l=1`$. Note that $`a`$ is a regular element of order $`e`$ in $`W_L^{}`$. Put $`V=𝐑_𝐙X(S^{})`$. Thus $`W_L^{}`$ acts on $`V`$ as a reflection group. Let $`\zeta `$ be a primitive $`e`$-th root of unity, and let $`a(e)`$ be the dimension of the eigenspace $`V(a,\zeta )V`$ of $`a`$ with eigenvalue $`\zeta `$. We show that
(2.10.1)
$$\mathrm{}\{tS_{}^{}{}_{}{}^{F^{}}t^l=1\}=l^{a(e)}.$$
By a general formula, we have $`|S_{}^{}{}_{}{}^{F^{}}|=|det_V(qIa)|=P_a(q)`$, where $`P_a(x)`$ is the characteristic polynomial of $`aW_L^{}`$. Since $`a`$ is regular $`P_a(x)`$ can be written, by \[Sp1, 4.2\], as
$$P_a(x)=\mathrm{\Phi }_e(x)^{a(e)}\mathrm{\Phi }^{}(x),$$
where $`\mathrm{\Phi }_e(x)`$ is the cyclotomic polynomial of degree $`e`$, and $`\mathrm{\Phi }^{}(x)`$ is a product of cyclotomic polynomials $`\mathrm{\Phi }_e^{}(x)`$ with $`e^{}<e`$. By our assumption $`m_l(q)=e`$, $`\mathrm{\Phi }_e(q)`$ is divisible by $`l`$, and $`\mathrm{\Phi }^{}(q)`$ is not divisible by $`l`$. This means that each minimal $`F^{}`$-stable torus $`M`$ of $`S^{}`$ corresponding to the factor $`\mathrm{\Phi }_e(x)`$ contains an element of order $`l`$. Since $`\{tM^F^{}t^l=1\}𝐅_{q^e}^{}`$, $`M^F^{}`$ contains exactly $`l`$ elements $`t`$ such that $`t^l=1`$. Thus (2.10.1) is proved.
For $`\beta \mathrm{\Phi }\mathrm{\Phi }_L`$, let $`V_\beta `$ be the subspace of $`V`$ which is orthogonal to $`x(\beta )`$ for all $`x\mathrm{\Gamma }`$. Then $`V_\beta `$ can be identified with $`𝐑_𝐙X(H_\beta ^0)`$. $`\mathrm{\Gamma }`$ stabilizes $`V_\beta `$, and let $`V_\beta (a,\zeta )`$ be the eigenspace of $`a`$ on $`V_\beta `$ with eigenvalue $`\zeta `$. Since $`a`$ is $`L`$-regular, we have $`dimV_\beta (a,\zeta )<dimV(a,\zeta )=a(e)`$. It follows that the characteristic polynomial $`P_a^{}(x)`$ of $`a`$ on $`V_\beta `$ contains the factor $`\mathrm{\Phi }_e(x)`$ with multiplicity less than $`a(e)`$. By a similar argument as above, minimal $`F^{}`$-stable subtori of $`H_\beta ^0`$ corresponding to $`\mathrm{\Phi }_e(x)`$ only contain elements of order $`l`$. This implies that
$$\mathrm{}\{tH_\alpha ^F^{}t^l=1\}=\mathrm{}\{tH_\alpha ^{0F^{}}t^l=1\}l^{a(e)1}.$$
It follows, by (2.9.1), that
$`\mathrm{}\{tS_{}^{}{}_{\mathrm{reg}}{}^{F^{}}t^l=1\}`$ $`=\mathrm{}\{tS_{}^{}{}_{}{}^{F^{}}t^l=1,t{\displaystyle \underset{\beta \mathrm{\Phi }\mathrm{\Phi }_L}{}}H_\beta ^F^{}\}`$
$`l^{a(e)}Nl^{a(e)1}=l^{a(e)1}(lN),`$
where $`N=|\mathrm{\Phi }\mathrm{\Phi }_L|`$. Since $`l>N`$ by our assumption, there exists $`tS_{}^{}{}_{\mathrm{reg}}{}^{F^{}}`$ such that $`t^l=1`$. This proves the lemma in the case (a).
Next consider the case (b) in 1.6. It is enough to show, for each $`l𝒫`$, that there exists $`tS_{\mathrm{reg}}^F^{}`$ such that $`t^l=1`$. We note that $`q^e^{}1`$ is not divisible by $`l`$ for any divisor $`e^{}<e`$ of $`e`$ by the assumption $`m_l(q)=e`$. Since $`S^F^{}(𝐅_{q^e}^{})^r\times (𝐅_q^{})^r^{}`$ (cf. 2.9), we have
$$\mathrm{}\{tS^F^{}t^l=1\}=l^r.$$
We consider $`K_\beta `$ given in 2.9. By the discussion in 2.9, we have
$$\mathrm{}\{tK_\beta ^F^{}t^l=1\}=\mathrm{}\{tK_\beta ^{0F^{}}t^l=1\}=l^{r1}.$$
It follows, by (2.9.2), that
$`\mathrm{}\{tS_{\mathrm{reg}}^F^{}t^l=1\}`$ $`=\mathrm{}\{tS^F^{}t^l=1,t{\displaystyle \underset{\beta \mathrm{\Phi }\mathrm{\Phi }_L}{}}K_\beta ^F^{}\}`$
$`l^rNl^{r1}=l^{r1}(lN),`$
where $`N`$ is as before. Since $`l>N`$ by our assumption, the lemma holds also for the case (b). ∎
We need the following lemma due to Lusztig.
###### Lemma 2.11 (\[L3, Lemma 3.2\]).
Let $`H`$ be a finite group, and $`\varphi `$ a virtual character of $`H`$ (over a field of characteristic 0). Assume that $`\varphi `$ is integral valued. Let $`x,yH`$ be such that $`xy=yx`$ and $`y^l=1`$ for a prime number $`l`$. Then $`\varphi (xy)\varphi (x)l𝐙`$.
2.12. Let $`s_0`$ be as in 2.7, and $`𝒫`$ be as in 2.9. Let $`F^{}=F^s\dot{a}`$ be as in Lemma 2.10 for a fixed $`l𝒫`$. Let $`R_{w,i}=R_{T_w}(1)`$ be the Deligne-Lusztig’s virtual character of $`G^{F^i}`$ for $`i=1,\mathrm{},e`$, where $`T_w`$ is an $`F^i`$-stable maximal torus of $`G`$ corresponding to $`wWW(T_1)`$ (here $`W(T_1)=N_G(T_1)/T_1`$ for an $`F^{}`$-stable pair $`T_1B_1`$). Let us choose $`t𝒵_{\mathrm{reg}}`$ as in Lemma 2.10. Then we have
(2.12.1) $`\mathrm{Tr}(F^iw,H^{}(_u))`$ $`=\mathrm{Tr}(u,R_{w,i}),`$
$`\mathrm{Tr}(F^iw,H^{}(_{tu}))`$ $`=\mathrm{Tr}(tu,R_{w,i}).`$
We remark that (2.12.1) was proved in \[L2\] under the assumption that $`p^s`$ is large enough (which is determined only by the data of the Dynkin diagram of $`G`$). Thus if we replace $`s_0`$ in 2.7 by a suitable large number, the result in \[L2\] is applicable. One can also apply \[Sh3, Theorem 2.2\] instead of \[L2\], where the restriction on $`p^s`$ is removed.
Since $`R_{w,i}`$ are integral valued, one can apply Lemma 2.11 for $`H=G^{F^i}`$ and $`x=u,y=t`$. Hence we have
$$\mathrm{Tr}(u,R_{w,i})=\mathrm{Tr}(tu,R_{w,i})modl𝐙.$$
It follows from (2.12.1) that
(2.12.2)
$$\mathrm{Tr}(F^iw,H^{}(_u))=\mathrm{Tr}(F^iw,H^{}(_{tu}))modl𝐙.$$
Let $`\zeta _0`$ be a fixed primitive $`e`$-th root of unity in $`𝐂`$, and $`R`$ the ring of integers of the cyclotomic field $`𝐐(\zeta _0)`$. Let $``$ be the set of non-zero prime ideals $`𝔭`$ in $`R`$ such that $`𝔭`$ contains one of the numbers $`1\zeta _0^i`$ for $`i=1,\mathrm{},e1`$ and $`\zeta _0`$. Let $`\overline{}`$ be the set of prime numbers $`l`$ such that $`𝔭𝐙=l𝐙`$ for $`𝔭`$. Since $``$ is a finite set, $`\overline{}`$ is a finite set. So, $`𝒫\overline{}`$ is an infinite set. Let $`𝒥`$ be the set of prime ideals $`𝔭`$ of $`R`$ such that $`𝔭𝐙=l𝐙`$ with $`l𝒫\overline{}`$. Then $`𝒥`$ is an infinite set. Now $`R/𝔭`$ is a finite extension of $`𝐅_l`$. Let $`\overline{\zeta }_0`$ be the image of $`\zeta _0`$ in $`R/𝔭`$. Since $`l𝒫`$, $`l1`$ is divisible by $`e`$. Hence $`\overline{\zeta }_0𝐅_l^{}`$, which has order $`e`$ by our choice of $`𝔭`$. Since $`m_l(p^s)=e`$, the image of $`p^s`$ in $`𝐙/l𝐙`$ has order $`e`$. Hence there exists $`j`$ such that
(2.12.3)
$$p^s\zeta _0^j𝔭.$$
Note that the number $`j`$ is determined by the choice of $`𝔭`$, which we denote by $`j(𝔭)`$. For $`j=1,\mathrm{},e1`$, let $`𝒥_j`$ be the set of prime ideals $`𝔭`$ in $`𝒥`$ such that $`j(𝔭)=j`$. Thus $`𝒥=_j𝒥_j`$, and so there exists $`j_0`$ such that $`𝒥_0=𝒥_{j_0}`$ is an infinite set. We put $`\zeta =\zeta _0^{j_0}`$. By (2.12.3), $`\zeta `$ is a primitive $`e`$-th root of unity.
We remark that $`H^{}(_{tu})=H^{}(_u_𝒵)`$ is independent of the choice of $`t𝒵_{\mathrm{reg}}`$. Then in view of (2.7.1) $``$ (2.7.4), together with (2.12.3), we see that
$`\mathrm{Tr}((F^s\dot{a})^iw,H^{}(_u))`$ $`=\mathrm{Tr}((w,a^i),H^{}(_u)^{(\zeta )})mod𝔭,`$
$`\mathrm{Tr}((F^s\dot{a})^iw,H^{}(_u_𝒵))`$ $`=\mathrm{Tr}((w,a^i),H^{}(_u_𝒵)^{[\zeta ]})mod𝔭`$
for any $`𝔭𝒥_0`$. Combined with (2.12.2), we have
$$\mathrm{Tr}((w,a^i),H^{}(_u)^{(\zeta )})=\mathrm{Tr}((w,a^i),H^{}(_u_𝒵)^{[\zeta ]})mod𝔭$$
for $`𝔭𝒥_0`$. Since $`𝒥_0`$ is an infinite set, we conclude that
$$\mathrm{Tr}((w,a^i),H^{}(_u)^{(\zeta )})=\mathrm{Tr}((w,a^i),H^{}(_u_𝒵)^{[\zeta ]}).$$
Hence Proposition 2.6 is proved, and the theorem follows.
## 3. Applications
3.1. Let $`W_L`$ be the subgroup of $`W`$, and $`\mathrm{\Gamma }`$ the subgroup of $`W`$ generated by $`aN_W(W_L)`$ such that $`\mathrm{\Gamma }`$ and $`W_L`$ generate the semidirect product group $`\stackrel{~}{W}_L=\mathrm{\Gamma }W_L`$. Let $`V=V^{(\zeta )}`$ be the $`\mathrm{\Gamma }\times \stackrel{~}{W}_L`$-module as in 1.3. (We write $`\mathrm{\Gamma }`$ as $`\mathrm{\Gamma }_0`$ if it is regarded as a subgroup of $`\stackrel{~}{W}_L`$, cf. 1.3.) Then $`V`$ can be decomposed as $`V=_{i𝐙/e𝐙}V^{(i)}`$, where $`V^{(i)}`$ is the eigenspace of $`a\mathrm{\Gamma }`$ with eigenvalue $`\zeta ^i`$, which is a $`\stackrel{~}{W}_L`$-submodule of $`V`$. Then we have
$`\mathrm{Ind}_{W_L}^WV`$ $`={\displaystyle \underset{i}{}}{\displaystyle \underset{wW/\stackrel{~}{W}_L}{}}{\displaystyle \underset{j}{}}wa^jV^{(i)}`$
$`={\displaystyle \underset{i}{}}{\displaystyle \underset{wW/\stackrel{~}{W}_L}{}}{\displaystyle \underset{k}{}}wb_kV^{(i)},`$
where $`b_k=_j\zeta ^{jk}a^j𝐂[\mathrm{\Gamma }]`$ (the group ring of $`\mathrm{\Gamma }`$). For each $`i𝐙`$, let $`\psi ^{(i)}`$ the linear character of $`\mathrm{\Gamma }`$ defined by $`\psi ^{(i)}(a)=\zeta ^i`$. Then $`\mathrm{\Gamma }`$-module $`𝐂b_k`$ is afforded by $`\psi ^{(k)}`$. Let $`V_n^{(i)}`$ be the eigenspace of $`a\mathrm{\Gamma }_0`$ on the $`\stackrel{~}{W}_L`$-module $`V^{(i)}`$ with eienvalue $`\zeta ^n`$. Let $`(\mathrm{\Gamma }\mathrm{Ind}_{W_L}^WV)^{(k)}`$ be the eigensapce of $`a\mathrm{\Gamma }`$ with eigenvalue $`\zeta ^k`$. Then we have the following lemma.
###### Lemma 3.2.
1. Let the notations be as above. We have
(3.2.1)
$$(\mathrm{\Gamma }\mathrm{Ind}_{W_L}^WV)^{(k)}\underset{wW/\stackrel{~}{W}_L}{}\underset{j𝐙/e𝐙}{}\underset{0n<e}{}wb_{knj}V_n^{(j)}$$
as vector spaces. In particular, $`dim(\mathrm{\Gamma }\mathrm{Ind}_{W_L}^WV^{(\zeta )})^{(k)}`$ is independent of the choice of $`k𝐙/e𝐙`$, which is given by
(3.2.2)
$$dim(\mathrm{\Gamma }\mathrm{Ind}_{W_L}^WV)^{(k)}=[W:\stackrel{~}{W}_L]dimV.$$
2. Assume that $`\mathrm{\Gamma }`$ commutes with $`W_L`$. Then we have
(3.2.3)
$$(\mathrm{\Gamma }\mathrm{Ind}_{W_L}^WV)^{(k)}\underset{j𝐙/e𝐙}{}\mathrm{Ind}_{\mathrm{\Gamma }\times W_L}^W(\psi ^{(k+j)}V^{(j)})$$
as $`W`$-modules.
###### Proof.
Under the action of $`\mathrm{\Gamma }`$ on $`\mathrm{\Gamma }\mathrm{Ind}_{W_L}^WV`$, $`wb_kV_n^{(i)}`$ is contained in an eigenspace of $`a`$ with eigenvalue $`\zeta ^{k+n+j}`$. Then (i) follows easily from the discussion in 3.1. Now assume that $`\mathrm{\Gamma }`$ commutes with $`W_L`$. Then $`b_kV^{(i)}`$ has a structure of $`\mathrm{\Gamma }\times W_L`$-module given by $`\psi ^{(k)}V^{(i)}`$. (ii) follows from the formula (3.2.1) by noticing that $`V^{(i)}=V_0^{(i)}`$. The lemma is proved. ∎
We consider a Levi subgroup $`LG`$ and a unipotent element $`uL`$, and take $`\mathrm{\Gamma }=\mathrm{}a\mathrm{}N_W(W_L)`$ satisfying the condition in 1.6. We apply the preceding argument to the situation $`V^{(\zeta )}=H^{}(_u^L)^{(\zeta )}`$. Then as a corollary to Theorem 1.7, we have
###### Proposition 3.3.
Under the setting in Theorem 1.7, we have, for $`0ke1`$,
(3.3.1)
$$\underset{nkmode}{}H^{2n}(_u)\underset{wW/\stackrel{~}{W}_L}{}\underset{j𝐙/e𝐙}{}\underset{0n<e}{}wb_{knj}H^{2j}(_u^L)_n$$
as vector spaces, where $`b_i𝐂[\mathrm{\Gamma }]`$ and $`H^{2n}(_u^L)_n`$ is the eigenspace of $`a\mathrm{\Gamma }_0`$ with eigenvalue $`\zeta ^n`$. In particular, $`dim\left(_{nkmode}H^{2n}(_u)\right)`$ is independent of the choice of $`k`$. In the case (a) in 1.6, (3.3.1) can be made more precise as follows;
(3.3.2)
$$\underset{nkmode}{}H^{2n}(_u)\mathrm{Ind}_{\mathrm{\Gamma }\times W_L}^W\left(\underset{j𝐙/e𝐙}{}\psi ^{(k+j)}H^{2j}(_u^L)\right)$$
as $`W`$-modules.
###### Proof.
By Theorem 1.7 $`\mathrm{\Gamma }\mathrm{Ind}_{W_L}^WH^{}(_u^L)^{(\zeta )}`$ is isomorphic to $`H^{}(_u)^{(\zeta )}`$ as $`\mathrm{\Gamma }\times W`$-modules. Since $`(H^{}(_u)^{(\zeta )})^{(k)}=_{nkmode}H^{2n}(_u)`$, the corollary follows from Lemma 3.2. ∎
Remarks 3.4. (i) In the case where $`u=1`$, the cohomology ring $`H^{}(_u)=H^{}()`$ coincides with the coinvariant algebra $`R`$ of $`W`$. In the special case where $`G`$ is of type $`A_{n1}`$, i.e., $`W𝔖_n`$, we consider $`W_L𝔖_{nre}`$ for $`1en`$. Then $`W_L^{}𝔖_{re}`$, and if we choose a regular element $`aW_L^{}`$ as a product of disjoint cycles of length $`e`$, Proposition 3.3 can be applied. This recovers the formula obtained by Morita and Nakajima \[MN1\].
More generally, consider the Weyl group $`W`$ acting on the real vector space $`V`$ as the reflection module. For $`vV`$, let $`W_v`$ be the stabilizer of $`v`$ in $`W`$, and $`N_v`$ the stabilizer of the line $`𝐑v`$ in $`W`$. Note that $`W_v`$ is normal in $`N_v`$, and $`W_v`$ coincides with $`W_L`$ for a certain Levi subgroup in $`G`$. Then for any $`\mathrm{\Gamma }=\mathrm{}a\mathrm{}`$ such that $`\mathrm{\Gamma }N_v`$, Bonnafé, Lehrer and Michel \[BLM\] have proved a similar formula as in Proposition 3.3. So our formula (3.3.2) can be regarded as a special case of theirs. (Note that they treat a more general case, where $`W`$ is a complex reflection group and $`\mathrm{\Gamma }`$ is not necessarily cyclic, in a framework of coinvariant algebras.)
(ii) We consider a unipotent element $`uL`$ in the case where $`G=GL_n`$. $`uG`$ can be written as $`u=u_\mu `$ by a partition $`\mu =(1^{m_1},2^{m_2},\mathrm{},n^{m_n})`$ of $`n`$. Take a positive integer $`e2`$, and let $`I`$ be a subset of $`\{1,\mathrm{},n\}`$ such that $`em_i`$ for $`iI`$. We consider a Levi subgroup $`L`$ of type $`X_0+e_{iI}A_{i1}`$, where $`X_0=A_k`$ with $`k=_{iI}im_i+_{iI}i(m_ie)1`$. Then we have $`W_L^{}=\{1\}`$, and one can choose $`uL`$ so that it satisfies the assumption of the case (b) in 1.6. Thus Proposition 3.3 can be applied. This covers the results on the stability of dimensions obtained in \[MN2\], \[MN3\], where they considered the case $`|I|=1`$ or the case all the $`m_i`$ are divisible by $`e`$.
Returning to the general setup, we consider the case where $`u`$ is a regular unipotent element in $`L`$. Then $`H^{}(_u^L)=H^0(_u^L)𝐂`$ is a trivial $`W_L`$-module. Thus Proposition 3.3 implies the following.
###### Corollary 3.5.
Let $`G`$ be a simple algebraic group modulo center, and $`L`$ a Levi subgroup in $`G`$. Let $`u`$ be a regular unipotent element in $`L`$. Let $`\mathrm{\Gamma }=\mathrm{}a\mathrm{}`$ be a subgroup of $`N_W(W_L)`$ of order $`e`$ satisfying the conditions in 1.6. Then for $`k=0,\mathrm{},e1`$, we have
$$\underset{nkmode}{}H^{2n}(_u)\mathrm{Ind}_{\mathrm{\Gamma }W_L}^W\stackrel{~}{\psi }^{(k)}$$
as $`W`$-modules, where $`\stackrel{~}{\psi }^{(k)}`$ is the character of $`\mathrm{\Gamma }W_L`$ obtained as the pull back of $`\psi ^{(k)}`$ under the projection $`\mathrm{\Gamma }W_L\mathrm{\Gamma }`$.
###### Proof.
In the case (a), the assertion follows from (3.3.2). So we consider the case (b). In the setup of 3.1, $`V^{(i)}`$ is a trivial $`W_L`$-module $`𝐂`$ for $`i=0`$ and zero otherwise. Then we see that $`V^{(0)}=V_0^{(0)}`$, and $`b_kV^{(0)}`$ has a structure of $`\stackrel{~}{W}_L`$-module $`\stackrel{~}{\psi }^{(k)}`$. The assertion follows from the formula in 3.1. ∎
3.6. Let $`G`$ be a simple algebraic group defined over $`𝐅_q`$ with Frobenius map $`F`$. We assume that $`G^F`$ is of split type. The Green function $`Q_{T_w}`$ is defined as the restriction of the Deligne-Lusztig’s virtual character $`R_{T_w}(1)`$ to the set of unipotent elements in $`G^F`$. We assume that $`p=\mathrm{ch}𝐅_q`$ is good, and in the case where $`G`$ is of type $`E_8`$, we further assume that $`q1(mod4)`$. Then as explained in 2.7, for each unipotent class $`C`$ of $`G`$, there exists a split element $`uC^F`$. As in 2.12, we have
(3.6.1)
$$Q_{T_w}(u)=\underset{n0}{}\mathrm{Tr}(w,H^{2n}(_u))q^n.$$
Hence there exists a polynomial $`𝐐_{w,C}(x)𝐙[x]`$ such that $`Q_{T_w}(u)=𝐐_{w,C}(q)`$. Concerning the values of Green functions at root of unity, we have the following.
###### Proposition 3.7.
Suppose that $`G,L`$ and $`uL`$ are as in Corollary 3.5. Then we have
(3.7.1)
$$𝐐_{w,C}(\zeta ^j)=|W_L|^1\mathrm{}\{xWx^1wxa^jW_L\}$$
for $`j=0,\mathrm{},e1`$. In particular, the value $`𝐐_{w,C}(\zeta ^{})`$ is independent of the choice of a primitive $`e`$-th root of unity $`\zeta ^{}`$.
###### Proof.
Put $`c_i(w)=\mathrm{}\{xWx^1wxa^iW_L\}`$ for $`i=0,\mathrm{},e1`$. Then
$`\left(\mathrm{Ind}_{\mathrm{\Gamma }W_L}^W\stackrel{~}{\psi }^{(k)}\right)(w)`$ $`=|\mathrm{\Gamma }W_L|^1{\displaystyle \underset{i=0}{\overset{e1}{}}}{\displaystyle \underset{\begin{array}{c}xW\\ x^1wxa^iW_L\end{array}}{}}\stackrel{~}{\psi }^{(k)}(x^1wx)`$
$`=|\mathrm{\Gamma }W_L|^1{\displaystyle \underset{i=0}{\overset{e1}{}}}c_i(w)\zeta ^{ki}.`$
It follows, by (3.6.1) together with Corollary 3.5, that
$`𝐐_{w,C}(\zeta ^j)`$ $`={\displaystyle \underset{k=0}{\overset{e1}{}}}\zeta ^{kj}{\displaystyle \underset{nkmode}{}}\mathrm{Tr}(w,H^{2n}(_u))`$
$`=|\mathrm{\Gamma }W_L|^1{\displaystyle \underset{i=0}{\overset{e1}{}}}c_i(w){\displaystyle \underset{k=0}{\overset{e1}{}}}\zeta ^{(ji)k}`$
$`=|W_L|^1c_j(w).`$
Hence we obtain the formula (3.7.1). Let $`\zeta ^j`$ be a primitive $`e`$-th root of unity. There exists an element $`\tau \mathrm{Gal}(𝐐(\zeta )/𝐐)`$ such that $`\tau (\zeta )=\zeta ^j`$. By (3.6.1), we see that $`𝐐_{w,C}(\zeta )𝐐(\zeta )`$ and that $`\tau (𝐐_{w,C}(\zeta ))=𝐐_{w,C}(\zeta ^j)`$. But since $`𝐐_{w,C}(\zeta )𝐙`$ by (3.7.1), we conclude that $`𝐐_{w,C}(\zeta )=𝐐_{w,C}(\zeta ^j)`$. This proves the proposition. ∎
Remark 3.8. In the case where $`G=GL_n`$ and $`L`$ is of type $`A_{m1}+\mathrm{}+A_{m1}`$ ($`e`$-times) with $`n=em`$, take a regular unipotent element $`u`$ in $`L`$. Then $`u=u_\mu G`$ with $`\mu =(m^e)`$. For $`wW=𝔖_n`$, let $`\lambda (w)=(1^{l_1},2^{l_2},\mathrm{})`$ be the partition of $`n`$ corresponding to the cycle decomposition of $`w`$. Then one can show by a direct computation (cf. \[M, (6.2)\]) that
$$|W_L|^1\mathrm{}\{xWx^1wxaW_L\}=\{\begin{array}{cc}e^{l(\lambda (w))}\hfill & \text{ if }el_i\text{ for all }i,\hfill \\ 0\hfill & \text{ otherwise,}\hfill \end{array}$$
where $`l(\lambda )`$ is the number of parts for a partition $`\lambda `$. Thus we recover the formula in \[LLT, Theorem 3.2, Theorem 3.4\] concerning the values of Green polynomials of $`GL_n`$ at roots of unity.
3.9. We give some more examples where Proposition 3.3 can be applied.
(i) Assume that $`G`$ is of type $`B_n`$ and $`L`$ is a Levi subgroup of type $`B_m`$ with $`m<n`$. Then $`L^{}`$ is of type $`A_{nm1}`$. For any $`uL`$ and a divisor $`e`$ of $`nm`$, the proposition can be applied. Similar results hold also for $`C_n`$ or $`D_n`$.
(ii) Assume that $`G=Sp_{2n}`$. Then a unipotent element $`uG`$ can be written as $`u=u_\mu `$ as an element of $`GL_{2n}`$, where $`\mu =(1^{m_1},2^{m_2},\mathrm{})`$ is a partition of $`2n`$ such that $`m_i`$ is even for odd $`i`$. Take an even integer $`e2`$, and let $`I`$ be a subset of odd integers $`\{1,3,\mathrm{},2n1\}`$ such that $`em_i`$ for $`iI`$. We consider a Levi subgroup $`L`$ of type $`X_0+e_{iI}A_{i1}`$, where $`X_0`$ is of type $`C_k`$ with $`2k=_{iI}im_i+_{iI}i(m_ie)`$. Then $`W_L^{}=\{1\}`$, and as in Remarks 3.4 (ii), one can find $`uL`$ so that the case (b) in 1.6 can be applied. Similar results hold also for type $`B_n`$ and $`D_n`$.
(iii) Assume that $`G`$ is of type $`E_7`$, and choose $`L`$ of type $`A_2`$ so that $`L^{}`$ is of type $`A_4`$. Take any unipotent element $`uL`$. Then the proposition can be applied with $`e=5`$. |
warning/0507/astro-ph0507003.html | ar5iv | text | # Core and Conal Component Analysis of Pulsar B1237+25
## 1 Introduction
Pulsar B1237+25 has been widely studied and represents the classic example of a pulsar with five emission components, apparently produced by a sightline traverse which passes almost through the center of two concentric emission cones and a central core beam \[e.g.,Rankin (1993)\]. It is also known to exhibit a “normal” and “abnormal” emission mode, frequent “null” pulses, and regular 2.8-period/cycle subpulse modulation \[Backer (1970a-c,1973)\] in its “normal” mode. Each of these properties has been well investigated, but little is known about possible linkages between them. We have carried out a new study of B1237+25’s polarized pulse sequences (hereafter PSs) with the purpose of investigating a) the characteristics and activity of its intermittently active core component and b) the recent claim that the star exhibits a third weak inner cone \[Gupta & Gangadhara (2003)\]. The pulsar’s individual-pulse behavior is so complex and varied, particularly with the benefit of excellent new Arecibo 327-MHz observations, that we have had to question again what interpretation should be made of its five profile features.<sup>1</sup><sup>1</sup>1Indeed, some early profiles suggest six components, ostensibly due to use of one circular or linear feed. This understanding has persisted, though, even with better observations (e.g.,Bartel et al. (1982); Qiao et al. (2000)) and we will develop some basis for understanding why below. PS polarization has provided the foundation for our analyses. Following the published studies of total sequences at 430 MHz \[Backer & Rankin (1980)\] and 1400 MHz \[Stinebring et al. (1984)\], Bartel et al. (1982) investigated the specific characteristics of the star’s two modes. Of particular significance in our current context was the discovery that the “normal” and “abnormal” emission (or profile) modes represented differing proportions of primary- and secondary-polarization-mode (hereafter PPM and SPM) power. We also draw on the studies of Oster & Sieber (1977), Prószyński & Wolszczan (1986), Psaltis &Seiradakis (1996) and Rankin & Ramachandran (2003).
While B1237+25’s relatively weak core component (III) has influenced beaming models \[e.g.,Backer (1976)\], its actual characteristics remain somewhat obscure. Its steeper spectrum renders it ever weaker at higher frequencies, and even around 400 MHz its close proximity to the stronger trailing component (IV) makes it difficult to isolate for study. We therefore began our study by focussing on the population of single pulses where the core was active. We quickly relearned that the core is active episodically, continuously so in “abnormal”-mode sequences, but also strongly at fairly regular intervals within “normal”-mode sequences. We thus now conclude that the classical “normal” mode is actually an alternation of two modes, a “quiet normal” and a “flare-normal” mode.
§2 then describes our observations, §3 & 4 introduce the characteristics of the three modes. In §5 we present the results of our search for a third cone following Gupta & Gangadhara (2003). §6 gives an analysis of the star’s nulling behavior, and §7 explores how the three modes interact. In §8 & §9 we discuss our results regarding the geometry and polarization of B1237+25’s core component. §10 gives an analysis of the star’s emission height, and §11 provides a summary and gives our overall conclusions.
## 2 Observations
The observations used in our analyses were made using the 305-m Arecibo Telescope in Puerto Rico. The primary 327-MHz polarized PSs were acquired using the upgraded instrument together with the Wideband Arecibo Pulsar Processor (WAPP<sup>2</sup><sup>2</sup>2http://www.naic.edu/$``$wapp) on 2003 July 13, 14 and 21, comprising 2340, 5094 and 4542 pulses, respectively. The ACFs and CCFs of the channel voltages produced by receivers connected to orthogonal linearly polarized feeds were 3-level sampled. Upon Fourier transforming, 64 channels were synthesized across a 25-MHz bandpass with a 512-$`\mu `$s sampling time, providing a resolution of 0.133 longitude. The Stokes parameters have been corrected for dispersion, interstellar Faraday rotation, and various instrumental polarization effects. Similar observations in four 21-cm bands on 2003 August 4 used 100-MHz widths, recorded 2062 pulses, and were sampled at the same resolution. Older Arecibo observations at 430 and 111.5 MHz are used for comparison as noted below.
## 3 Integrated Profile Analysis
Figure 1 reiterates for ease of comparison the familiar total profile of B1237+25, created by averaging a set of pulses without distinguishing its modes. We see the five classical components (I-V), which have usually been interpreted as a central sightline traverse through two concentric emission cones and a core beam. The validity of this inference, based on simple average profiles like this one, must be questioned for many reasons. The profile form is far from symmetric in either total power (solid curve) or linear polarization (dashed curve), and we see neither the full antisymmetric circular signature (dotted curve) often associated with a core component nor the expected 180°-position angle (PA) excursion of a central sightline traverse. Rather, the core region is depolarized as clearly shown by the linear minimum. Indeed, one might conclude almost immediately from such a profile alone that it cannot be comprised of a single orthogonal-polarization-mode (hereafter OPM) population of individual pulses. We know from the papers quoted above that both a “normal” and “abnormal” mode have been identified and well studied in this pulsar—and moreover that they differ considerably in relative OPM power. We argue below that the classical “normal” mode in fact consists of an alternation of two new modes, a “flare-normal” mode characterized by intervals of strong core activity, and a “quiet normal” mode wherein the core emission is nearly undetectable.
In our attempt to identify populations of individual pulses in which the core component is particularly active, we found that such pulses occur either in “abnormal”-mode intervals or within fairly discrete “normal”-mode intervals. Further investigation showed us that these bright-core intervals typically persist for 5-10 pulses and that they recur quasi-periodically about every 60-80 periods (hereafter $`P_1`$). We thus begin by delineating the characteristics of this “flare-normal”, bright-core mode in relation to the “quiet normal” and classical “abnormal” modes.
### 3.1 Third Mode
Figure 2 represents the average profile of the “flare-normal”-mode PSs, identified by visual inspection of the full observations using colour displays such as that in Figure 5. As we will show below, the boundaries of “flare-normal” mode subsequences can be distinguished from both “abnormal”- and “quiet normal”-mode intervals in these high quality observations generally within a pulse or two. As expected, the “flare-normal”-mode profile shows a much brighter core component, and we see the antisymmetic circular polarization often associated with core components. Note also that the core component is no longer linearly depolarized as in the total profile or classical “normal” mode. It here exhibits a symmetrical $`>`$25% linear “component” just under the total-power core and this is accompanied by a strong negative excursion of the PA (lower panel).
Components I, II and IV appear almost unaltered in “flare-normal”-mode sequences, but comp. V increases markedly in intensity and appears to occur slightly earlier, thus partially overlapping comp. IV. This modal component V is also substantially less linearly polarized, showing the effect (as we will see below) of an increased level of SPM power. The subtle decrease in the spacing of components I and II also appears significant and reflects a slightly narrower overall profile width as compared with the total profile or the “quiet normal”-mode profile below.
### 3.2 “Quiet normal” Mode
Figure 3 gives an average profile for the “quiet normal”-mode sequences. It represents the sum of those intervals exhibiting little or no core activity and in practice it differs indistinguishably from an average of the residuum after “flare-normal” and “abnormal”-mode intervals have been identified and removed. Apart from the decreased core presence, this profile differs little from the typical total profile, because some 85% of “normal”-mode pulses are “quiet normal” and the “abnormal” mode is yet rarer.<sup>3</sup><sup>3</sup>3Occasional “abnormal”-mode-dominated total profiles do occur, however, because it can sometimes persist for hundreds, or perhaps even thousands, of pulses. We again see a sharp minimum of the linear polarization near the longitude origin, and perhaps (owing to less core influence) a somewhat larger portion of the mostly unresolved negative PA excursion near the profile center. It is also noteworthy that in individual pulses comp. II is usually very narrow (a few bins) and is 100% polarized in typically 90% of its occurrences. The “quiet normal”-mode PA very often exhibits a full negative, nearly 180°, excursion under the core component. We see this behavior in many individual pulses, but not always in modal average profiles. Weak core power—often hardly detectable in individual “quiet normal” pulses—at times still accrues to form a distinct profile component; and when present, this residual core emission is often sufficient to disrupt the “conal” PA traverse. The effect is most pronounced at the linear power minimum in the profile center which indeed coincides with the positive-going PA “jump” under it. Thus, in many individual pulses and modal averages, when the core power is low enough, this “jump” connects negatively and symmetrically.
### 3.3 Classical “abnormal” Mode
Figure 4 gives the average profile of a long, “abnormal”-mode sequence and, in using the word “classical”, we emphasize that our criteria for identifying pulses belonging to this mode are the same as those in published studies. Again, the profile alterations exhibited by the pulsar’s “abnormal” mode are so dramatic at meter wavelengths that at first sight we can wonder if we are observing the same pulsar. Each of the five components assumes a different relative configuration in the “abnormal” mode, and the emission minimum, usually seen just before the core component, is now observed to follow it. The linear-polarization minimum, however, is even more prominent here, and falls at almost exactly the same longitude as in the “quiet normal” and total profiles. Also, the here positive “jump” associated with the minimum connects negatively in other such modal profiles, resulting in an overall PA traverse of more than -210°! Something of this behavior can also be seen in Bartel et al. ’s (1982) fig. 7 (left).
The core component, of course, predominates in the “abnormal” mode and we will see below that it is almost continuously active in adjacent individual pulses. Note, however, that while the circularly polarized signature of the core component is perhaps strongest here, the linear polarization properties of the “abnormal”-mode core are very different than those of the “flare-normal” mode. Here, we see even larger fractional linear polarization and again a long (but here longer) negative PA traverse associated with it. Both the total intensity and linear peaks fall slightly later than in the “flare-normal” mode, such that the linear polarization appears to trail within the visible core component. Moreover, this linear is delineated by leading and trailing minima which coincide closely with the duration of the negative (RH) part of the circular signature. All this then leaves the leading part of the visible core component conspicuously linearly depolarized.
Note also that the overall modal profile width is again narrower, by nearly 2°, than the total profile. Components IV and V have merged and highly polarized emission is now found in the “empty” region between components II and the core. Over much of the width of the profile we see evidence of increased SPM power: the PA of comp. V is now SPM dominated as is a leading-edge region of comp. I, but also we see evidence of PPM power on the extreme outer edges. Only the regions under and following comp. II and just after the core remain PPM dominated—this in stark contrast to the “normal” or total profiles where the PPM dominates “conal” regions of the profile everywhere apart from the extreme leading and trailing edges. Only comp. II appears nearly unaltered in the “abnormal” mode, but also we note that its leading edge is now more completely polarized.
## 4 Modal Intensity and Polarization Dynamics
We have noticed that the three largely distinct modal behaviors described above can be distinguished by their polarization signatures at the individual pulse level—and indeed we find that they can be most accurately and reliably so identified. Figure 5 displays the full polarization state of a 100-pulse sequence of the 327-MHz observation of 2003 July 21. The particular PS in Figure 5 was chosen because, unusually, good examples all three modal behaviors fall within this short interval.
### 4.1 “Quiet-normal” Mode
This fundamental and most characteristic mode can be identified in total power by noting the intervals where the conal components exhibit their usual 2.8-$`P_1`$ modulation. This is most dramatically seen in the PA column, where the modal power on the profile outer edges alternates cyclically between some $``$10° (purple) and +80° (chartreuse). This outer-edge OPM modulation is also discussed for B1237+25 by Rankin & Ramachandran (2003); see their fig. 4 for a 430-MHz example. During these intervals the central core emission is weak or absent, though there is usually sufficient linear power in the profile center to define the PA. Note that this PA behaves just as one would expect for a nearly central sightline traverse: we see the $``$10° (purple-angled) linear polarization first rotate negatively toward $``$30° (magenta) and then nearly in the next sample become +60° (green) and then +30° (cyan)—corresponding to a $``$120° excursion—whereafter the PA rotates more gradually around to 0° (blue). Components II and IV exhibit high levels of fractional linear power, with II frequently showing nearly complete polarization. Finally, the fractional circular polarization peaks positively before and within the center of the profile, reaching some 60% at about the position of the core component. Note that the negative circular is much less consistent in this mode; at times there is weak RHC following the LHC peak, but often not as well. We see little order to the circular polarization under the conal components in this mode.
### 4.2 “Flare-normal” Mode
Primarily, the “flare-normal” mode differs from the “quiet normal” mode by the presence of bright core emission. Note that this core activity is not continuous, but fluctuates from pulse to pulse in a manner that at times appears to mimic the conal modulation. Though these “flare-normal”-mode apparitions are characteristically of short (5–15 $`P_1`$) duration, they suggest (as in Fig. 5) a somewhat longer $`P_3`$, perhaps of about 4 $`P_1`$/c. Note especially the contrasting PA modulation as compared to the “quiet-normal” mode. We see little distinction between the “normal” modes in $`L/I`$, but there are clear differences in their PA and $`V/I`$ properties. “Flare-normal”-mode intervals exhibit a strange but consistent single-pulse PA behavior: at about comp. II’s longitude we again find a $``$10° PA (purple-angled) which now first rotates positively through 0° (blue) to +30° (cyan), then reverses to rotate negatively through 0° (blue) to $``$30° (magenta), and then finally rotates positively again to 0° (blue). A transition from “quiet”- to “flare-normal”-mode PA traverse properties is seen clearly at pulse 1978 and one in the opposite sense occurs at pulse 1960. These PA sense reversals cannot be produced geometrically, so must have some other cause as we will discuss below. Finally, note that the fractional circular polarization of the “flare-normal” mode tends to be somewhat smaller, but more fully antisymmetric than in the “quiet normal” mode. Most of these circumstances are well demonstrated by modal averages as in Fig. 2, but we have found it necessary to carefully inspect the individual-pulse behavior to assess just how they come about.
Overall, it is instructive to examine the “quiet”- to “flare-normal”-mode transition at about pulse 1978. Note the pronounced narrowing of the emission window, particularly on the trailing side of the profile. Comp. V becomes as bright as comp. I, but its linear polarization decreases—apparently as a result of mode-mixing because the individual-pulse fractional-linear distribution seems to change little across the boundary. The RHC power under components I and V becomes more orderly, apparently contributing to the small residual seen in the modal profile.
### 4.3 “Abnormal” Mode
“Abnormal”-mode intervals are easily distinguished from the other behaviors using only the total power. Not only does the core “flare”, emitting almost continuously in successive pulses, but at meter wavelengths strong “abnormal” emission only occurs in the lefthand portion of the pulse window—and what remains of comps IV and V merge into a single weak “hump”. Moreover, all evidence of the usual modulation ceases and SPM power dominates over large leading and trailing portions of the modal profile.<sup>4</sup><sup>4</sup>4Psaltis & Seiradakis (1996) find evidence for weak conal modulation in the “abnormal” mode at 21 cms.; however, this might be explained by the inclusion of “quiet-normal”–mode intervals, which could not be easily distinguished at this frequency. In $`L/I`$ we see little change in the “abnormal” mode, except at the longitude of comp. IV. The PA and $`V/I`$ behaviour, however, show that the “flare-normal” and “abnormal” modes share common features, but also remain distinct. Both the absence of modulation and SPM dominance are well known characteristics and immediately clear from the PA column of Fig. 5. Note, however, that we again see the strange multiple PA reversals, but in the “abnormal” mode the PA tends to first rotate negatively (to $``$30°, magenta) as the core is approached, then jump positively (to +30°, cyan), rotate negatively at first steeply and then more gradually (to $``$60°, full red), and finally jump back to 0° (blue). It is useful to contrast this “abnormal” PA traverse with that of the “flare-normal” mode: the jump to a +30° angle (cyan) occurs at just the same longitude in the two modes, but the linear power at this point in “abnormal” PSs is very small—note that this point coincides with the $`L`$ minimum in Fig. 4. The “abnormal” PA rotation under the core rotates further and much of it occurs at later longitudes.
We note that the short “abnormal”-mode PS of Fig. 5 is not fully representative of the longer examples in our observations. Comp. II is usually more consistently strong, so that we see a marked contrast between it and comp. IV in both the total power and fractional linear polarization. Additionally, long “abnormal” mode PSs are typically more prominently circularly polarized with significant fractional LHC seen first under comp. II and rising to at least 60% at the leading half-power point of the visible core. At its trailing half-power point the RHC is typically only some 30%, but it is important to emphasize that the LHC and RHC excursions represent about equal (antisymmetric) contributions of power.
### 4.4 “Normal”-mode Fluctuation Properties
Figure 6 gives longitude-resolved fluctuation spectra for our longest observation, the 5094-PS from 2003 July 14. This PS exhibits little “abnormal”-mode emission, thus it effectively represents the classical “normal” mode. We see strong features near 0.35 c/$`P_1`$ as well as a broad weaker one around 0.25 c/$`P_1`$, corresponding to the usual $`P_3`$=2.8 $`P_1`$ as well perhaps as the nearly 4 $`P_1`$ modulation seen in the “flare-normal” mode.
Furthermore, we also see the low frequency feature reported earlier\[Taylor & Huguenin (1971), Taylor, et al. (1975), Backer (1973), Bartel et al. (1982) \]. The brighter of the two features implies a $`P_3`$ of some 37 $`P_1`$ and the adjacent response corresponds to about twice this value. One of us \[Rankin (1986)\] had tended to associate this modulation with the core component, but we surely now see here that its effect is more general. The feature modulates not only the core, but a broad region around it surely including comps. II & IV. We therefore conclude that the low frequency feature represents the amplitude modulation of the quasi-periodic “flare-normal”-mode apparitions. Visual inspection of the corresponding PS appears consistent with this interpretation.
## 5 New Components?
We have investigated whether Gupta & Gangadhara (2003) (hereafter G&G) are correct in asserting that B1237+25 has two additional minor components, just leading and trailing the core component, which together constitute a smaller “further in” emission cone. Their analysis is based on a “window-threshold” method \[Gangadhara & Gupta (2001)\], which itself remains somewhat controversial, because power is sought only in certain discrete longitude “windows”. It is worth noting that G&G give some analytical details only about the leading component of the new pair—at –2.1° as opposed to +1.6°, relative to the core component—and our results pertain largely to it as well.
G&G applied their method to a 318-MHz, 1915-pulse, total-power observation having an apparently undetermined mix of modes. A cursory examination of the total and modal profiles in Figs. 1–4 reiterates that the region leading the core by a little over 2° generally represents a minimum in profile power. Only in the “abnormal” mode does this region emit at a level comparable to the other bright components. We cannot learn, however, from the “abnormal” average in Fig. 4 whether this power between components II and III represents a distinct component or has some other origin. Using a version of G&G’s method, which accumulated only “abnormal”-mode pulses having significant power levels in a window some 2°–2.5° before the core, we computed a partial modal profile. Though it contained only 13 of some 91 “abnormal”-mode pulses in our 2003 July 21 observation, it is almost identical to the full “abnormal”-mode profile in Fig. 4. The feature on the leading edge of the core in both profiles does not resolve into a separate component. Even if we interpret it as a “component”, it falls considerably later than the one in G&G’s fig. 2b, which peaks about halfway between components II and III. Similar results were obtained for the “flare-normal” mode, with only 3 pulses qualifying of some 370, and nothing like a distinct new component appearing.
A different result was obtained when we applied G&G’s analysis to groups of “quiet normal” pulses, and this partial profile appears in Figure 7. Here, we do see a discrete feature at about the same longitude shown in their fig. 2b. Fig. 7 averages 31 individual pulses which exceeded 108 times the off-pulse noise level $`\sigma `$—this out of a total of 2654 “quiet normal” pulses in the 2003 July 21 PS. The feature is not so clearly resolved as in G&G’s fig. 2, but its position and upwardly slanting baseline suggest that we have succeeded in roughly verifying their published analysis.
Please note that our Fig. 7 differs markedly from G&G’s fig. 2. Ours more resembles a hybrid of the star’s “flare-normal”-mode and total profiles, whereas theirs would seem to be dominated by “flare-normal” and “abnormal”-mode power. Perhaps the difference can be understood in terms of the greater sensitivity of the AO observations, where more weaker “quiet normal”-mode pulses were found to exceed an appropriate threshold. G&G do not tell us how many pulses went into their partial profile, but both its shape and the narrowness of the “new feature” might be understood if it were very few.
This said, the partial profile of Fig. 7 is peculiar: we see that the “new feature” is accompanied by a relatively strong core—including its antisymmetric circular polarization—and that the “new feature” is as closely marked by the LHC as the core is by the RHC polarization. Moreover, the relatively strong core and pronounced leading/trailing asymmetry are reminiscent of “abnormal”-mode profiles, despite the fact that the qualifying pulses here were selected from a “quiet-normal” population having no clear active-core intervals. We will return to this finding below.
Perhaps unfortunately, our similar effort to verify G&G’s purported trailing feature has met in failure. The core component (III) has a broad base in the two “normal” modes, as can readily be seen in the last several figures, and thus fills most of the “gap” with comp. IV. Even in the “abnormal” mode, where this trailing region is relatively empty, it proved impossible to define the window so as to exclude power associated with the “tail” of the core. Indeed, given the substantial width of the core component, we find it difficult to understand how these authors could have solved this problem. Their Table 3 gives the new component’s position as +1.56° after the core peak, where comp. IV lies at 2.98°. If the core is present, its trailing edge will surely overlie the position of the putative new component. Only when the core is absent—or nearly so as in “quiet normal”-mode PSs—would one seem to have a chance at finding this new component, and our searches for it in just this situation were negative as well.
## 6 Nulling Analysis
As reported previously, some 6% of B1237+25’s pulses are “nulls” \[Ritchings (1976);Rankin (1986)\]. We are able to discriminate these nulls with reasonable accuracy using a threshold of about 10% of the mean pulse intensity, depending on observation quality. On this basis, null-length histograms show clearly that 1-pulse nulls are about 5 times more frequent than 2-pulse nulls, with the frequency declining steadily to a maximum of about 8–10 pulses. The preponderance of 1-pulse nulls suggests that a population of unobserved and partial nulls also occur; the former could be quite frequent (some 3%), but the star’s small duty cycle makes the latter rare (perhaps 1 in 1500 for this pulsar) and thus difficult to positively identify. We have also looked for connections between nulls and modes, but nulls seem to occur within each of the modes and also at modal boundaries (e.g., , see Fig. 5).
One major hint, however, comes from constructing partial profiles comprised of the last pulses prior to a null or the first pulses afterward. Overall, these partial-profile pairs strongly resemble each other and the total profile, and the differences (e.g., strong comp. I after nulls) inconsistent among our various PSs. One subtle change, however, is seen in all our 327-MHz observations: the core component following nulls is elongated in the leading direction and often exhibits a bifurcated shape with the trailing part falling at the usual core position and the leading at just the longitude of the “new component” discussed in the preceding section. There is also a tendency for this feature pair to be brighter just after nulls.
We have thus investigated where the 31 qualifying pulses in Fig. 7 fall and they appear mainly just before or after ”flare-normal”-mode sequences, as identified from the PA as in Fig. 5. This suggests that our mode-identification criteria may need to be extended to include such transitional pulses, given the ”flare-normal” mode total power and polarization characteristics of Fig. 7.
One other interesting aspect of B1237+25’s emission is a population of very weak pulses—surely qualifying as nulls according to the above noise threshold—which nonetheless exhibit polarized emission at angles appropriate for their longitude and position within the PS. Two examples, seen in Fig. 5, are pulses 1997 and 2012 which happened to fall within the plotted interval, but other instances of the same effect are seen in any similar 100-pulse display. This surely suggests that nulls represent a more complex phenomenon than the simple bi-state (on/off) received understanding, and we note with interest the recent paper by Janssen & van Leeuwen (2004) which attributes a sort of “sputtering” to nulls in pulsar B0818–13.
## 7 How Do the Three Modes Interact?
We have not been able to see clearly whether transitions between the three modes exhibit any particular pattern or preferred sequence. It is tempting to see the “flare-normal” mode as being intermediate between the extremes of “quiet normal” and “abnormal” behavior. Both the “flare-normal” and “abnormal” modes have stronger core emission, somewhat narrower profiles, and a tendancy for components IV and V to merge. The “normal”-mode pair share the key property of regular subpulse modulation in the outer-cone regions (though perhaps the “flare-normal” mode has a somewhat longer $`P_3`$) as well as enough regularity in their succession to give the low frequency fluctuation-spectral feature. One might also see the “quiet normal” mode as exhibiting the least SPM power, the “flare-normal” somewhat more, and the “abnormal” dominated by it.
The low frequency modulation feature associated with“flare-normal”-mode apparitions argues that the two “normal” modes exhibit a cyclic relationship. This conclusion is strengthened by the fact that neither mode persists for very long. “Flare-normal”-mode intervals rarely last more than a dozen pulses and “quiet-normal” sequences hardly 50. The same cannot be said about the “abnormal” mode; it can switch on for a few pulses (as we see in Fig. 5), and our observations include a sequence 250 pulses long. It would appear that the star does not often or easily change from the “quiet normal” to the “abnormal” mode. As shown in Fig. 5, the more frequent transition is between “flare-normal” and “quiet normal”, perhaps indicating that more drastic geometric or energetic changes are required for “abnormal” emission. Once in the “abnormal” mode, however, it can stably continue much longer than the other modes.
## 8 Core-component Width and Conal Emission Geometry
The width of B1237+25’s core component is of interest because of its close connection to the full angular size of the stellar polar cap, and therefore is an indicator of the star’s emission geometry (Rankin 1990, 1993). In many profiles the relative weakness of the core makes its width difficult to determine and, even when not so, measurements often give values smaller than the polar-cap size. Added to this, the core is seen to substantially lag the center of the profile—that is, the midpoint between components I and V. We will discuss both issues in turn.
Because B1237+25 exhibits a double cone and because we have presumed it to have a highly central sightline traverse (an issue we will verify below), its magnetic geometry can be computed reliably on the basis of its conal dimensions alone. This computation was carried out in Rankin (1993) showing that the star’s magnetic latitude $`\alpha `$ is some 53°. This value agrees well with Lyne & Manchester’s (1988) somewhat different analysis giving 48°. On this basis, we can then compute the 1-GHz core width expected from the polar-cap angular width as $`W_{core}`$=2.45° $`P^{1/2}/\mathrm{sin}\alpha `$ \[Rankin (1990); eq.(5)\], and this in turn gives a width value of 2.61°. This is rather more than is usually measured as can easily be verified from the foregoing profiles: the total profile of Fig. 1 shows the difficulty of determining the core width, but a generous estimate here might be 2.5°. However, the isolated “flare-normal” and “abnormal”-mode core components in Figs. 2 & 4 have widths of only some 1.55°–1.75°.
The spectral behavior of the pulsar’s core width is not fully known. The relative weakness of its core and the narrower overall profile above 1 GHz make it increasingly difficult to measure this width. At low frequencies some observations suggest that the core width is roughly independent of frequency, but most also suggest a width somewhat less than 2.6° \[e.g.,Hankins & Rankin (2005)\]. Indeed, the several published observations in the 110–150-MHz range (Backer 1970; Lyne et al. 1971; Hankins & Rankin 2005) all suggest a “narrow” core with a slow leading and steep trailing edge. More illuminating are 112-MHz Pushchino profiles made in 2000, each 153 pulses long and with a resolution of about a milliperiod. In some of these the core is nearly as bright as the other components, often shows a double form with a stronger following feature, and has an overall half-power width of hardly 2.8° (Suleymanova 2005). This surely tends to reiterate that the B1237+25 core is ”double” and that its width changes little, at least at meter wavelengths.
At 327 MHz, the core shows a similar shape; either we see evidence for some leading-edge extension as in Figs. 2 & 4, or for a leading-edge feature as in Fig. 7. We are thus left with a consistent impression that the star’s core is incomplete (or “absorbed”; e.g.,Rankin (1983)) on its leading edge. This interpretation is further supported by the analyses of Qiao et al. (2000, 2003), who through Gaussian-component fitting to average profiles argue for a sixth weak component just prior to the visible core.
Additional strong support for this interpretation comes from the core’s circularly polarized signature. In the modal profiles above, we see that the visible core component aligns closely with the trailing negative (RHC) peak; whereas, the positive (LHC) peak falls in the “extended” region at about the half-power point on the core’s leading edge. This asymmetry in the total-power shape of the core component together with the balanced antisymmetric circular polarization is very unlike the typical core behavior in other pulsars. Indeed, were we to take the “core width” not as the 3-db total power value, but rather as the interval between the respective leading and trailing half-power points on the circular signature, a value near 2.6° can be scaled from the profiles of Figs. 2 & 4. Or, to carry this one step further, we can measure the core’s half width in “flare-normal” and “abnormal”mode profiles between the circular zero-crossing point and the visible total-power component’s trailing half-power point—and this gives values around 1.3° as expected.
Finally, in Figure 8 we give a modal profile of a “quiet normal”-mode PS in which the core emission was particularly weak. Note that here we see no core feature at all in either the total intensity or the circular polarization. What we do see, however, is the full conal PA traverse! Apart from the OPM dominance “jumps” at the outer edges of the profile, here finally is the expected PA traverse which is completely in keeping with the single-vector model \[Radhakrishnan &Cooke (1969); Komesaroff (1970)\]. We have taken the longitude origin in this and the earlier figures to fall within half a sample of the PA-traverse symmetry point. Note that the PA traverse accumulates a full 180° in the longitude interval between components I and V and also that the PA rate near the center of the traverse is exceptionally steep as expected. We measure the maximum PA-rate value to be some $``$185$`\pm `$5 °/°, which may well be the largest ever determined—though this may be an underestimate owing to averaging over pulses which “jitter” slightly in their longitude centers (see Fig. 4). Using this value (rather than infinity) to compute the impact angle $`\beta `$ \[see Rankin (1993), Table 2\], we find that $`\beta `$ is 0.25° and $`\beta /\rho `$ 0.051, where $`\rho `$ is the outer cone radius to the outside 3-db point.
## 9 Core-component Linear Polarization
An old B1237+25 mystery is why its total profile fails to exhibit the expected 180° PA “swing” expected for a nearly central sightline traverse. We have been able to show in Fig. 8 that, indeed, the pulsar does exhibit such a traverse—but only in partial profiles restricted to the individual-pulse population having virtually no core emission. This clearly established, we are now in a position to begin to understand what core characteristics are responsible for the distorted PA traverse we generally observe. Or, said the other way around, this modal profile shows us how fragile is this full conal R&C PA traverse, because it is carried by pulses which not only have little emission at the longitude of the core but are also linearly depolarized—probably because the confounding core emission is never entirely absent.
It then follows that when the core is present, both most of the power and most of the polarization belong to the core component. This surely explains why the PA “disruption” is confined to a longitude region just under the core. Note then that in those profiles above where core power is present, the PA exhibits a “hook” entailing usually four sense reversals. (Only the “quiet-normal”-mode profile in Fig. 3 behaves differently—exhibiting two reversals—because core power is so weak here that the PA almost connects negatively between the two reversals.) Then, we can observe that the PA centerpoint between the 2nd and 3rd reversals aligns accurately with the PA value far from the core (the average, for instance, under components 2 & 4). Such an unusual PA signature—which we have seen is encountered in both core-active modes—has a very natural explanation: the core component must be dominated by linearly polarized power which is orthogonal to the conal power so clearly seen in Fig. 8. And if this “quiet normal”-mode conal power represents the primary polarization mode, it follows that the core is largely comprised of “secondary” polarization-mode emission.
This behavior can also clearly be seen in individual pulses. In the context of discussing Fig. 5 above, we pointed out the different PA behaviors exhibited by the three modes. “Quiet-normal”-mode pulses exhibit little evidence of PA reversals; whereas, both the “flare-normal” and “abnormal” modes show the same individual-pulse PA behavior as seen in their modal profiles. Returning to these respective modal profiles in Figs. 2 & 4, it is interesting to ponder the significance of their details. In both cases, the centers of their PA “hooks” fall about 0.4° later than the zero-crossing points of their circular polarization signatures—points which nearly coincide with the center of the conal PA traverse in Fig. 7. It is interesting then to see that the “flare-normal”-mode profile shows no depolarized OPM-dominance boundary, whereas the “abnormal” profile exhibits one on the leading edge of the visible core component and perhaps one on its trailing edge as well.
Finally, Figure 9 provides a simulation of the “hook” which we have found characteristic of the disrupted PA behavior under the core component. We have simply added a Gaussian-shaped, linearly cross-polarized component to the natural core-absent profile of Fig. 8—and this results in the four-fold-reversed PA signature in the bottom panel. This surely demonstrates that orthogonally polarized power must be a factor in both the depolarization and the nearly unique PA behavior seen in B1237+25. It does not, however, fully delineate what the source of this cross-polarized power might be. An obvious possibility is the SPM emission seen so prominently at other longitudes in this pulsar, but strictly SPM power is so far not known to be a property of core emission. Were this the source, we might expect to see strongly depolarized mode-dominance boundaries, which are apparently seen in the “abnormal”-mode profile of Fig. 4, but interestingly not so in the “flare-normal”-mode profile of Fig. 2. Emission-height variations such as suggested by Mitra & Seiradakis (2003) may also provide an explanation for the observed polarization effects. Detailed study of the depolarization will be required to distinguish these two alternatives.
## 10 Core-component Symmetry and Conal Emission Height
Blaskiewicz, Cordes & Wasserman (1991) and Malov & Suleymanova (1998; hereafter M&S) have attempted to account for core/cone asymmetries in pulsar emission profiles in terms of aberration and retardation effects and their analyses have led to useful methods for estimating the height of the emitting regions. Both approaches require that the longitude positions of conal component pairs be measured relative to a fiducial longitude associated with that of the magnetic axis; and the two methods use, respectively, the linear PA-traverse and core-component centers, to estimate this position. Extending their earlier work based on M&S (Gangadhara & Gupta 2001), G&G applied their analysis to B1237+25. As mentioned above, they reported a pair of new conal components comprising a third “further in” cone. We, however, have only been able to verify their leading feature, which we find is clearly associated with the core component. G&G thus took the visible core component as indicative of the profile center, but we have seen above that only a trailing portion of the B1237+25 core is emitted. For these reasons, their analysis using M&S’s technique is poorly founded for B1237+25 and thus their height estimates unreliable.
In fact, our analysis above interestingly provides two possible fiducial points within the profile, the center (inflection point) of the linear PA traverse and the zero-crossing point of the antisymmetric circular polarization signature. These two respective origins permit us to apply the methods of Blaskiewicz et al. (1991) and M&S in the same star. However, it is important to note that the two points do not occur simultaneously as one can be identified in “quiet normal” intervals and the other within “flare-normal” and “abnormal” PSs. We give the measured (total or “quiet normal”-mode) conal component positions relative to both the PA- and CP-defined origins in Table 1. It is configured precisely as was G&G’s Table 3 and we also use their definitions for the conal leading and trailing component-pair positions $`\varphi _l^i`$ and $`\varphi _t^i`$, the total aberration-retardation longitude shift $`\nu ^i`$, the magnetic azimuth angle $`\mathrm{\Gamma }^i`$, emission height $`r_{em}^i`$, and relative polar cap annulus $`s_L^i`$.
We thus also base our computations on Gangadhara & Gupta’s (2001) eqs.(1-8, 10-15)\] with one significant exception: Dyks et al. (2004) have provided a subtle correction of their eq.(9) so we use the latter paper’s eq.(7) instead—though in practice this correction results in only 10% smaller emission heights for B1237+25. Table 1 then shows that the outer and inner conal emission heights, relative to the PA traverse, are roughly the same, about 310$`\pm `$75 km, but that these cones are emitted on different field lines having their “feet” at some 0.78$`\pm `$0.07 and 0.53$`\pm `$0.05 of the polar cap radius, respectively. \[These values are substantially smaller and more external than the ones computed by G&G (2003) that is, 460 & 600 km at 41 & 59%, respectively.\] The zero-crossing point of the circular signature was determined by making modal profiles of short “flare-normal” or “abnormal” intervals, which were then compared with corresponding “quiet normal” intervals such as that shown in Fig. 8. The circular zero-crossing point generally precedes the center of the PA traverse by some 0.10$`{}_{}{}^{}\pm `$0.05 (though we occasionally see the opposite behavior). Computations based on this second definition of the profile center are also given in Table 1 where one can see there that the emission heights thereby obtained are some 60 km smaller and the annuli somewhat more exterior.
We give this second calculation as an example, rather than an attempt to be definitive, as a thorough analysis and interpretation of the relation between the PA traverse and circular signature is beyond the scope of this paper. On the scale of a few pulses with the large S/N and high resolution, the component shapes and positions fluctuate continuously, though some of the polarization variation appears to be systematic. Note, for instance, the apparent “drift” to later longitudes of the positive circular during the last 50 or so pulses in Fig. 5. Owing to these complexities, for example, we have not attempted here to interpret the narrowing of the profile in the core-active modes. We plan to continue this work in a subsequent paper.
## 11 Discussion and Conclusions
This paper provides new—and we hope somewhat clarifying—analyses of the famous five-component (M) pulsar B1237+25, based largely on the straightforward method of partial polarization profiles. Pulsars with such clear double-conal/core structure remain unexplained by pulsar theories. Not only has core emission thus far received little detailed analytical attention, but the properties of double cones also remain understudied. Thus this pulsar still provides a primary context for such investigation; not only is it perhaps the brightest of the M stars, but it also exhibits a nearly precise central sightline geometry, regular subpulse modulation, modes, OPM effects and nulls.
The study began with a reexamination of the star’s modes, subpulse modulation and nulls. In addition to studying the star’s “abnormal” mode, we find that its long known “normal” mode in fact consists of two distinct behaviors, which we have designated the “quiet normal”, where little or no core activity can be distinguished in individual pulses, and the “flare-normal” mode, where the core is nearly as bright as in the “abnormal” mode, but displays distinct properties. In many ways this “flare-normal” mode can be regarded as an intermediate state between the “quiet normal” and “abnormal” behaviors. In this “flare-normal” mode, the core is bright, but not so fully dominated by the SPM. The profile narrows and components IV and V partially merge, but not as much as in the “abnormal” mode. Subpulse modulation in the “flare-normal” mode is not entirely quenched (as in the “abnormal”), but exhibits less regularity and perhaps a $`P_3`$ value around 4 $`P_1`$/c. One can also view the “flare-normal” mode as exhibiting much more SPM than the “quiet normal” mode, but far less than the “abnormal” mode wherein it is dominant over most of the profile.
“Flare-normal” and “quiet normal” intervals alternate with each other quasi-periodically. It is this regular appearance of bright “flare-normal”-mode core power in the center of the profile that produces the long known “core”-associated fluctuation feature (e.g., Rankin 1986). Typically, “quiet normal” sequences persist for some 40–70 pulses and “flare-normal” intervals for 5–15 pulses, making a complete cycle of some 60–80 pulses in duration. We note that what appear to be short “abnormal”-mode intervals are often interspersed within this overall “quiet normal”/“flare-normal”-mode cycle; generally, however, these persist for only a few pulses. At unpredictable times, though, the “abnormal” mode appears to “catch”, and then it can persist without interruption for a few or many hundreds of pulses. “Abnormal”-mode sequences thus appear to reflect a distinct relatively stable“state” of the star’s magnetospheric emission; whereas, the “normal”-mode alternation of “quiet” and “flare” intervals represents a further cyclic “state”.
While it has almost uniformly been assumed that B1237+25 had a highly central sightline geometry, its average PA traverse does not at all simply bear this out. Its total profile exhibits neither the constant PA nor the steep, 180, central traverse expected for such a sightline geometry. Some “abnormal”-mode profiles (e.g., Bartel et al. 1982) exhibited more of the expected PA traverse, but it remained to be explained why the total profile did not. We have largely been able to resolve this issue. The most core-quiet intervals of “quiet normal”-mode PSs are found to exhibit an almost textbook-quality PA behavior, as shown above in Fig. 8. The center of this PA traverse has a sweep rate of at least 180$`{}_{}{}^{}/_{}^{}`$ and the PA in the “wings” of the profile becomes almost constant at about the same PA value (apart from OPM dominance effects on the extreme profile edges.) This, in turn, permits us to determine the star’s $`\beta `$ value from its PA sweep rate for the first time, confirming that $`\alpha `$ is 53$`{}_{}{}^{}\pm `$2 (Rankin 1993) and $`\beta `$ 0.25$`{}_{}{}^{}\pm `$0.05. In terms of the geometrical models of the above paper, this implies that $`\beta /\rho `$ is hardly 5%.
Another mystery has revolved around the behavior of B1237+25’s core component. First, its expected antisymmetric circularly polarized signature is weak at best in the total profile (e.g., Fig. 1), and second, the observed width of its core component at meter wavelengths is often or usually substantially less than the expected 2.45$`{}_{}{}^{}P_{}^{1/2}/\mathrm{sin}\alpha `$ (Rankin 1990). Our 327-MHz “abnormal”- (Fig. 4) and “flare-normal”-mode (Fig. 2) profiles show beautiful anti-symmetric circular signatures under the core component. However, it is noteworthy that this circular symmetry is retained despite the partial character of the total-power core component itself. Thus we see that the leading portion of the core is at least 50% circularly polarized, whereas the trailing portion (after the core peak) exhibits only 25–30% fractional circular polarization. We note also that the circular signatures in the “abnormal”- and “flare-normal”-mode profiles are nearly identical—that is, the two modal core components have about the same fractional circular polarization and both their total intensities and circular zero-crossings nearly overlie each other. The strong difference in their respective core linear polarization is then noteworthy: The linear peak falls late in the “abnormal” mode and is bounded on both sides by near minima; whereas, in the “flare-normal” mode the linear peak lies near the center of the component and the linear polarization appears continuous almost across the entire profile. This difference can also be observed in the two PA traverses: The “abnormal”-mode “hook” under the core component falls later within the component and rotates further negatively, as we have also seen in individual pulses in Fig. 5 above.
We reemphasize that the B1237+25 core component is dominated by cross-polarized power. The extent and manner in which this occurs varies from mode to mode. Even in “quiet normal” PSs (e.g., Figs. 3 & 5) the pronounced depolarization under the core appears to corroborate the contribution of orthogonally polarized power. Also, the relative contribution of SPM power becomes stronger in the “flare-normal” and “abnormal” modes, as can be seen very clearly in the colour PS polarisation display above. Note in particular the manner in which the SPM power can be seen in the PA column on the outer edges of the emission window. We have been able to model the PA “hook” by adding cross-polarized power in the form of the core to the “quiet normal” profile. We cannot yet be sure whether the cross-polarized power in the core is SPM power or, for instance, it results from differences in emission height as suggested by Mitra & Seiradakis (2003).
The width of the B1237+25 core component at meter wavelengths remains an issue. We find that the visible core in modal profiles is hardly 2, substantially less than the expected at least 2.6. However, this visible core is clearly late, as it aligns with the trailing, RHC peak of the antisymmetric circular signature. This trailing portion of the core seems to be complete as its width, measured between the circular zero-crossing and its trailing 3-db points, is just half of 2.6° as expected. We see little which aligns with the leading LHC peak in the modal profiles. However, this is just the point where G&G found the earlier of their two “new components” using their window-threshold method—and this is the feature which we have been able to verify as well using their method. G&G aside, we do occasionally see individual pulses which have peaks at the longitude of the LHC peak, about –1.
We therefore conclude on the basis of the various evidence that the core component in B1237+25 is asymmetric and incomplete. However, we find evidence for a complete core—with about the expected 2.5–3 width—in both the scale of the two respective peaks of the circular signature and also by the observed emission in some single pulses which align with the LHC peak. Some short PSs at lower frequencies seem to exhibit just this sort of core more clearly (Suleymanova 2005). It is then hardly surprising that when Qiao et al. (2003) fitted Gaussian functions to this star’s profile, they found it necessary to add a sixth component at just this “leading core” longitude. They also concluded that the core is hollow and we find no clear evidence from our work to contravene this conclusion—though neither can we fully verify it. B1237+25 is nearly unique in providing a traverse so close to the magnetic axis.
Also, we used the unusually clear definition of both the PA-traverse and circular symmetry points that our analysis provides to estimate the star’s emission heights using the recent reformulation of Dyks et al. (2004). We find respective outer and inner emission-cone heights of 340$`\pm `$79 km and 278$`\pm `$76 km relative to the PA-traverse inflection point, such that they lie on polar cap annuli of 78 and 53%. We also determine these heights with respect to the circular polarization zero-crossing point giving us values some 60 km smaller. An obvious interpretation of this difference might be that the core emission is emitted much closer to the stellar surface, but at a height of some 60 km.
Finally, although our study is focussed primarily on core characteristics. it is interesting to compare our results which those of Psaltis & Seiradakis (1996). We are not surprised that they found evidence for three emission rings in total power, and we believe that these should be associated with the inner cone as well as the PPM and SPM subrings which comprise the outer cone (Rankin & Ramachandran 2003). Our results further tend to support their conclusion that the conal modulation feature is comprised of several preferred values (see Fig. LABEL:fig6), and it is surely true that the low frequency feature is associated more with the inner conal components than the outer ones. Our present analysis has not been aimed at discerning the subbeam “carousel” structure, but it is difficult for us to understand how the low frequency “modal” modulation—in which the core participates—could be commensurate with the well known 2.8-$`P_1`$ modulation seen prominently in the outer components.
Surely, our new analyses of B1237+25 above have demonstrated that this marvelous star still has very much to teach us. For almost no other pulsar do we have the opportunity to study the characteristics of the emission so close to the magnetic axis—and the pulsar remains the paragon of the five-component (M) species. We plan to continue some of the lines of investigation reported here and, in particular, expect to report further on the star’s core cross-polarized emission in a subsequent paper.
## Acknowledgments
We thank Avinash Deshpande for help with the observations, Stephen Redman and Svetlana Suleymanova for assistance with aspects of our analysis, and Jaroslaw Dyks, Svetlana Suleymanova, and Geoff Wright for discussions. Portions of this work were carried out with support from US National Science Foundation Grants AST 99-87654 and 00-98685. Arecibo Observatory is operated by Cornell University under contract to the NSF. This work made use of the NASA ADS system. |
warning/0507/q-bio0507039.html | ar5iv | text | # A complex adaptive systems approach to the kinetic folding of RNA
## 1 Introduction
RNA takes part in a variety of important cellular activities, including protein synthesis, intron splicing, gene silencing, and genome rearrangement (Lee et al., 2002; Mochizuki et al., 2002; Gratias & Betermier, 2003; Yang et al., 2003). Considering the extensive functional repertoire of RNA molecules, it is of interest to determine their (functional) native structures and to understand the (kinetic) process by which they fold into such structures. The native structures of RNA molecules can be computed efficiently, at the functionally relevant secondary structure level, using free-energy minimization methods (Hofacker et al., 1994; Mathews et al., 2004) or, in cases where sufficient homologous sequences are available, by phylogenetic comparisons (Woese & Pace, 1993; Cannone et al., 2002).
Several algorithms have been developed for studying the kinetic process by which RNA molecules fold into their native secondary (Mironov & Kister, 1985; Morgan & Higgs, 1996; Flamm et al., 2000; Zhang & Chen, 2000; Wolfinger et al., 2003; Tang et al., 2004; Ndifon & Nkwanta, 2005) and tertiary (Abrahams et al., 1990; Gultyaev et al., 1990; Isambert, 2000; Xayaphoumine et al., 2003) structures. The majority of these algorithms (e.g., see Isambert, 2000; Xayaphoumine et al., 2003; Ndifon & Nkwanta, 2005) operate on a helix-based move-set, involving the formation and dissociation of entire RNA helices. On the other hand, a few of the algorithms (e.g., see Zhang & Chen, 2000; Wolfinger et al., 2003) operate on a pair-based move-set; they model the RNA folding process as a time-series of structural transitions, involving the formation and dissociation of individual RNA base pairs.
Flamm and colleagues (Flamm, 1998; Flamm et al., 2000) have extended the afore-mentioned pair-based move-set by introducing the concept of base pair shifting, which makes possible description of the biological process of defect-diffusion, believed to be an important feature of the in vivo folding kinetics of RNA (Poerschke, 1974a). The folding model developed in this paper implements this extended pair-based move-set and is inspired by the theory of complex adaptive systems (see Section 3). The applicability of the model is illustrated through several examples based on selected natural and synthetic RNAs (see Section 4). In particular, the folding kinetics of the yeast tRNA<sup>Phe</sup> is shown to be strongly influenced by modifications to specific hairpin loops. In addition, a characteristic optimal folding temperature $`T_{opt}`$ $`\left(313K\right)`$ of tRNA<sup>Phe</sup>, at which the native state exhibits maximal accessibility, is identified. Furthermore, estimates are obtained for the population dynamics of two alternative stable states of SV11, an RNA species that is replicated by $`Q\beta `$ replicase (Zamora et al., 1995).
The remainder of this paper is organized as follows. In Section 2, we present some concepts related to RNA secondary structures and folding kinetics. We introduce the theory of complex adaptive systems and discuss details of the new folding model in Section 3. In Section 4, we apply the model to some example problems and discuss other possible applications in Section 5.
## 2 Background information
### 2.1 RNA secondary structure
Let $`X`$ be an arbitrary RNA sequence of length $`n`$. We think of $`X`$ as a string $`X=x_1x_2\mathrm{}x_n`$ defined over the nucleotide alphabet $`\{A,C,G,U\}`$. The nucleotides or bases of $`X`$ have a propensity to pair (or form canonical and non-canonical bonds) with each other. A pair formed by the bases $`x_i`$ and $`x_j`$, $`i<j`$, is denoted by $`(i,j)`$. Two base pairs $`(i,j)`$ and $`(i,j)`$ are said to be compatible if either $`i<i<j<j`$ or $`i<j<i<j`$. If we let $`H`$ be the set of possible pairs that can be formed by the bases of $`X`$, then a secondary structure $`S`$ of $`X`$ can be thought of as a set of mutually compatible base pairs drawn from $`H.`$ The multiset consisting of all subsets of $`H`$, including the empty set (i.e., the open chain), forms the conformation space of $`X`$, denoted here by $`\zeta (X)`$. Note that incompatible base pairs form pseudoknots, which are prohibited from occurring in the folding model developed in this paper. The model can, however, be readily extended to allow the formation of pseudoknots once reliable thermodynamic parameters for such tertiary structural elements become available.
### 2.2 Kinetic folding
The kinetic folding of an RNA sequence $`X`$ at the coarse-grained secondary structure level can be thought of as a time-series of structural transitions, mediated by a set of operations called the move-set (Flamm, 1998). Each operation or move converts one secondary structure $`S_i\zeta \left(X\right)`$ into another $`S_j\zeta \left(X\right)`$. For each structure $`S_i\zeta \left(X\right)`$, the move-set defines a neighborhood $`N\left(S_i\right)`$ such that $`S_j`$ $`\zeta \left(X\right)`$ belongs to $`N\left(S_i\right)`$ if and only if $`d(S_i,S_j)d`$ $`^+`$, where $`d(S_i,S_j)`$ is the (Hamming) distance between structures $`S_i`$ and $`S_j`$ and $`d`$ is the ”move distance”. Only moves that convert $`S_i`$ into some $`S_jN\left(S_i\right)`$ are legal. The probability of a legal move is given by a rate equation, an example of which is the Metropolis rule (Metropolis et al., 1953):
$$k_{ij}=\{\begin{array}{c}e^{\frac{\left(G_jG_i\right)}{RT}}\text{, if }G_j<G_i\\ 1\text{, if }G_jG_i\end{array}$$
(1)
where $`P\left(S_iS_j\right)`$ is the probability of converting $`S_i`$ into $`S_j`$ $`N\left(S_i\right)`$ by a single move; $`G_i`$ and $`G_j`$ are, respectively, the free energies of $`S_i`$ and $`S_j`$, computed using a suitable choice of free-energy parameters (e.g., Mathews et al., 1999).
Conventional Monte Carlo RNA folding algorithms execute in each time step $`t+\delta t`$ a move that converts the nascent RNA structure $`S_t`$ into some structure $`S_{t+\delta t}`$, where $`S_t`$ denotes the secondary structure of $`X`$ at time $`t`$. The time-ordered series of structures $`\left\{S_t\right\}_{t0},`$ with $`S_t\zeta \left(X\right)`$ and $`S_{t=0}`$ the open chain, is called a folding trajectory of $`X`$. The folding time $`\tau _f`$ associated with a given folding trajectory is the minimum value of $`t`$ for which $`S_t`$ is the native structure of $`X`$. A folding trajectory satisfying the following condition is called a folding path (Flamm et al., 2000): $`S_{t_1}=S_{t_2}`$ if and only if $`t_1=t_2`$.
## 3 The RNA folding model
### 3.1 Complex adaptive systems
The RNA folding model presented below is inspired by basic ideas from the theory of complex adaptive systems. Specifically, a complex adaptive system (CAS) is characterized by the presence of a diverse ensemble of components that engage in local interactions and an autonomous process that selects a subset of those components for enhancement based on the results of the local interactions (Levin, 1998). From these component-level dynamics emerge important global (i.e., system-level) properties such as self-organization and nonlinearity. Self-organization refers to the emergence of order from local interactions between the components of a CAS. Self-organization tends to drive a CAS towards stable configurations or states. In addition, a CAS exhibits nonlinearity; the rules that govern local interactions between the components of a CAS change as the CAS evolves (Levin, 1998). Consequently, a CAS may evolve along any one of a multitude of trajectories and may attain alternative stable states (Levin, 1998), depending on its particular evolutionary trajectory. See Levin (1998) and the references therein for further information on CASs.
### 3.2 Details of the model
The kinetic folding of RNA sequences into secondary structures is viewed here as a time-series of structural rearrangements (SRs), involving the formation, dissociation, and shifting of individual RNA base pairs. It is modeled as a hierarchically-structured CAS; at the lowest level of the hierarchy are RNA bases and base pairs that engage in local stacking interactions. The results of these stacking interactions determine the probabilities (or fitnesses) of possible SRs. These probabilities are given by
$$P_f^{(i,j)}=e^{\frac{\mathrm{\Delta }G^{ij}}{2RT}},\text{ }$$
(2)
$$P_d^{(i,j)}=\frac{1}{P_f^{(i,j)}},\text{ and}$$
(3)
$$\text{ }P_s^{(i,j)(i,k)}=P_d^{(i,j)}P_f^{(i,k)},$$
(4)
where $`P_f^{(i,j)}`$, $`P_d^{(i,j)}`$, and $`P_s^{(i,j)(i,k)}`$ are, respectively, the probabilities of formation and dissociation of $`(i,j)`$, and of shifting $`(i,j)`$ into $`(i,k),`$ $`R`$ is the gas constant, $`T`$ is the absolute temperature, and $`\mathrm{\Delta }G^{ij}`$ is the stacking (including single-base stacking) energy associated with $`(i,j)`$. For an isolated base pair $`(i,j)`$,
$$\mathrm{\Delta }G^{ij}=\mathrm{\Delta }G^{ij}+c\mathrm{ln}\left(ji\right),c0,$$
(5)
Equation $`\left(\text{5}\right)`$ takes into account the entropy of a loop of size $`\left(ji\right)`$. A suitable value for the parameter $`c`$ is $`1.75`$ (Fisher, 1966). Stacking energy calculations are based on the Turner $`3.1`$ energy rules (Mathews et al., 1999), at temperature $`T=310.5K`$, and on the Turner $`2.3`$ energy rules (Freier et al., 1986), at other temperatures. Equations $`\left(\text{2}\right)`$, $`\left(\text{3}\right)`$, $`\left(\text{4}\right)`$ are based on the Kawasaki dynamics (1966); here, the dynamics involve transitions between states associated with specific local contexts of an RNA secondary structure.
The next level of the hierarchical structure is occupied by SRs. It is at this level that selection operates. An autonomous stochastic sampling process (Baker, 1987) periodically (i.e., from one time step to another) selects a subset of possible SRs for realization based on the fitnesses of the SRs. The ensemble of possible SRs changes from one time step to another thereby assuring its diversity (see below). Note that due to the inter-dependence of local stacking interactions, on which the fitnesses of SRs depend, it is necessary to ensure that the SRs that are selected for realization in the same time step be mutually independent. Therefore if there is an SR involving the base $`x_i`$, then no other SR that involves the nearest-neighbor bases and base pairs of $`x_i`$ can occur in the same time step. Furthermore, in order to prevent the formation of pseudoknots SRs may only involve accessible bases; two bases $`x_i`$ and $`x_j`$, $`i<j`$, are accessible if for $`k=i+1,\mathrm{},j1`$ there is no base pair $`(k,l)`$ $`\left(\text{resp}.(l,k)\right)`$ such that $`l>j`$ (resp. $`l<i`$).
To understand the model just described, consider the kinetic folding of an RNA sequence $`X`$. Denote by $`\mathrm{}`$ the set of possible SRs that are available for realization in a given time step. $`\mathrm{}`$ will contain as many elements as there are structures in $`N\left(S_t\right)`$, where $`S_t`$ is the nascent structure of $`X`$. Each element of $`\mathrm{}`$ is associated with specific bases and base pairs that belong to that element’s ”local context”. For instance, an SR that involves the formation of the base pair $`(i,j)`$ is associated with $`(i,j)`$ and all nearest-neighbor bases and base pairs of $`(i,j)`$. Stacking interactions between the bases and base pairs associated with a given SR determine that SR’s probability or fitness. This fitness is used periodically by an autonomous stochastic process to select a subset of SRs from $`\mathrm{}`$ for realization.
As SRs are realized (and removed from $`\mathrm{}`$), existing SRs may be become impractical while new SRs may become possible. For instance, the formation of $`(i,j)`$ in a given time step may make possible the shifting of $`(i,j)`$ into some $`(i,k)`$ in the next time step. Conversely, the dissociation of $`(i,j)`$ in a given time step will render impractical the shifting of $`(i,j)`$ into some other base pair in the next time step. New SRs that become possible are added to $`\mathrm{}`$ while those that become impractical are removed from $`\mathrm{}`$. This assures the diversity of possible SRs. Note that the idea that the selective enhancement of components (i.e., SRs in this case) of a given system leads to their removal (or elimination) as well as the elimination of other components from that system appears to be at odds with what takes place in most known CASs. In the present case, the goal of selection is to, indirectly, enhance the thermal stabilities of the local contexts associated with SRs.
We note here that the inter-dependence of local stacking interactions, on which the fitnesses of SRs depend, implies that a given SR may ”interact” with many other SRs. For instance, an SR that involves the base pair $`(i,j)`$ will ”interact” with all SRs that involve either of the bases $`x_i`$ and $`x_j`$. This implies a relatively high average degree of ”epistatic” interactions between SRs. Results from studies based on random Boolean networks predict that such a high degree of epistasis leads to rugged fitness landscapes with numerous attractors (Kauffman and Levin, 1987; Kauffman, 1989). This prediction is consistent with the well-known rugged nature of RNA folding energy landscapes.
We further note that SRs, as defined in the model, operate at much more local scales of space than is the case with most existing folding methods (e.g., see Flamm et al., 2000; Isambert et. al., 2000; Tang et al., 2004). The fitnesses of SRs depend exclusively on the stacking energies associated with specific local contexts of the nascent RNA structure, $`S_t`$, and not on the free-energies of structures found in $`N\left(S_t\right)`$. Therefore, there is no need for explicit computation of the free-energies of RNA structures in the present model, in contrast to, e.g., the folding method of Flamm et al. (2000). We expect the model to reproduce global characteristics of CASs such as self-organization and nonlinearity. In particular, RNA molecules are ”self-organizing” since, through their own internal dynamics, they tend to fold into thermodynamically favorable or stable states. Folding RNA molecules also exhibit nonlinear dynamics, as evinced by their attainment of alternative stable states. In Section 4, we will illustrate such nonlinear dynamics using an example based on SV11.
### 3.3 Computer implementation of the model
The above model has been implemented in the computer program kfold, which is available from the author upon request. For an input RNA sequence of length $`n`$, the program selects $`m=1`$ SR, if $`n30`$, and $`m=7`$ SRs, if $`n>30`$, for realization in each time step. The folding time is incremented in each time step by the reciprocal of the product of $`m`$ and the sum of the fitnesses of all possible SRs. The number of selected SRs $`m`$ can be adjusted by the user. Note that the choice of $`m`$ influences the computer time required to fold an input RNA sequence but has minimal effect on the qualitative folding kinetics of the sequence (see example in Table 1). Also note that in order to speed up folding simulations, the program currently only allows the formation of base pairs that can be stacked. Specifically, a base pair $`(i,j)`$ can be stacked if there exists complementary bases $`x_l`$ and $`x_k,`$ $`l<k,`$ such that either $`ik=lj=1`$ or $`ki=j`$ $`l=1`$.
| $`m`$ | Time steps | Fraction of A | Fraction of B |
| --- | --- | --- | --- |
| $`1`$ | $`249`$ | $`0.40`$ | $`0.60`$ |
| $`3`$ | $`89`$ | $`0.41`$ | $`0.59`$ |
| $`5`$ | $`134`$ | $`0.38`$ | $`0.62`$ |
| $`7`$ | $`169`$ | $`0.42`$ | $`0.58`$ |
| $`9`$ | $`241`$ | $`0.45`$ | $`0.55`$ |
| $`11`$ | $`>30000`$ | $`0.40`$ | $`0.60`$ |
> Table 1. Influence of $`m`$ on folding kinetics for the sequence
>
> $$GUCCUUGCGUGAGGACAGCCCUUAUGUGAGGGC,$$
>
> with $`n=33`$. It was folded with ((((((((((((((…..)))))))))))))) (A) and ((((((….)))))).((((((….)))))) (B) serving as target structures. The fraction of simulations that found either structure within the allowed time scale (i.e., $`4\times 10^4`$ time steps) is similar for different values of $`m`$. On the other hand, the number of time steps, which reflects the amount of computer time required for folding, decreases as $`m`$ increases from $`1`$ to $`3`$, and subsequently increases with $`m`$. $`50\%`$ of simulations failed to find either target structure for $`m=11`$. The number of time steps given for $`m=11`$ thus represents a lower bound of its actual value. Note that each data point was averaged from just $`500`$ folding simulations run at $`T=310.5K`$. Therefore, there may be errors in the data resulting from limited sampling of possible folding trajectories.
## 4 Example applications
We now use the new folding model, as implemented in kfold, to study the effects of base modifications and temperature on the folding kinetics of the yeast tRNA<sup>Phe</sup>. We also estimate the population dynamics of two alternative stable states of the synthetic SV11. These examples will demonstrate that the folding model qualitatively reproduces characteristic RNA folding dynamics. Note that in the following examples, the folding times (i.e., mean first passage times) were scaled using experimentally measured folding times (in $`\mu s`$) of the hairpin $`AAAAAACCCCCCUUUUUU`$ (Poerschke, 1974b). This was done in order to allow direct comparisons with folding times reported in Flamm (1998). Unless otherwise noted, all folding simulations were run at $`T=310.5K`$.
### 4.1 Influence of base modifications on folding kinetics
A number of tRNA sequences are known to contain base modifications. Such modifications, believed to be the consequence of evolutionary optimization, have been shown to improve the foldabilities of some tRNAs (Flamm, 1998; Flamm et al., 2000). We have introduced several base modifications to the individual hairpins of the yeast tRNA<sup>Phe</sup> sequence and studied the folding kinetics of the modified sequences. All modified bases were prohibited from engaging in bond formation and stacking interactions. The modified sequences are shown in Table 2.
> | Sequence | Modified Hairpins | Modified Sequence Positions |
> | --- | --- | --- |
> | $`seq1`$ | $`1,`$ $`2`$ $`\&`$ $`3`$ | $`15,17,19,37,38,55,56`$ $`\&`$ $`59`$ |
> | $`seq2`$ | $`1`$ $`\&`$ $`2`$ | $`15,17,19,37`$ $`\&`$ $`38`$ |
> | $`seq3`$ | $`1`$ $`\&`$ $`3`$ | $`15,17,19,55,56`$ $`\&`$ $`59`$ |
> | $`seq4`$ | $`2`$ $`\&`$ $`3`$ | $`37,38,55,56`$ $`\&`$ $`59`$ |
> | $`seq5`$ | $`None`$ | $`None`$ |
>
> Table 2. Modified tRNA sequences used in this example. Note that hairpins are labeled from left to right, with ”1” representing the left-most hairpin.
We found base modifications to elicit substantial improvements in tRNA foldabilities, in the form of drastic decreases in folding times (see Figure 1). For the unmodified sequence, $`seq5`$, the fraction of folded sequences (i.e., sequences that have found the native cloverleaf structure) increased relatively slowly, reaching about $`45\%`$ after the first $`600us`$, and $`100\%`$ within $`4500\mu s`$. For all modified sequences, on the other hand, the fraction of folded sequences increased rapidly, reaching $`100\%`$ within $`1500\mu s`$. Among these sequences, $`seq1`$ was the fastest folder with an estimated folding time of $`300\mu s`$, while $`seq4`$ was the slowest folder with a folding time of $`1500\mu s`$. The folding times of $`seq2`$ and $`seq3`$ were approximately equal (i.e., about $`800\mu s`$). These observed effects of base modifications on tRNA folding kinetics are consistent with experimental data, as well as with predictions made by kinfold (Flamm, 1998). Note that we were able to simulate much longer folding times for tRNA<sup>Phe</sup>, up to $`4500\mu s`$, than was done in Flamm (1998).
### 4.2 Temperature dependence of folding kinetics
We have used the folding model to study the temperature dependence of the folding time $`\tau _f`$ of the yeast tRNA<sup>Phe</sup>. We found a $`V`$-shaped temperature dependence of $`\tau _f`$ (see Figure 2), suggesting the existence of an optimal folding temperature $`T_{opt}313K`$. A possible explanation for the existence of $`T_{opt}`$ is as follows: At temperatures $`T>T_{opt}`$, there are numerous structures with similar free-energies as the native cloverleaf. The native cloverleaf is therefore relatively unstable at temperatures above $`T_{opt}`$ and may be associated with a much smaller basin of attraction in the energy landscape. It therefore takes the folding tRNA longer to ”find” the cloverleaf among the ensemble of nonnative states. On the other hand, at temperatures $`T<T_{opt}`$, the stability of nonnative structures increases leading to a growth in the number and, perhaps, sizes of nonnative basins of attraction in the energy landscape. These nonnative basins of attraction may decrease the folding tRNA’s chances of finding the native state. Both scenarios (i.e., $`T>T_{opt}`$ and $`T<T_{opt}`$) lead to suboptimal native state accessibilities. Meanwhile, the maximal accessibility of the native state that is evident at $`T=T_{opt}`$ suggests the existence at $`T_{opt}`$ of optimal balance between the thermal stabilities of native and nonnative states. Note that $`T_{opt}`$ is close to the optimal growth temperature range for many yeast species. Detailed analysis of the temperature dependence of folding kinetics for tRNAs and other functional RNAs from various organisms will allows us to determine if this observation is a consequence of evolutionary optimization or simply a chance occurrence.
### 4.3 (Meta)stable states
During kinetic folding, some RNA molecules may get trapped in long-lived, nonnative states called metastable states/conformations. Examples of such metastable RNA molecules include riboswitches that regulate gene expression in bacteria by switching between alternative stable conformations (Vitreshack et al., 2004). By adopting a repressing conformation, a riboswitch can elicit the premature termination of DNA transcription or the inhibition of protein translation. Detailed in silico analysis of the folding kinetics of a metastable RNA molecule can provide insight into its functionality. For instance, Nagel et al. (1999) used computer simulations to identify a metastable structure of the Hok mRNA that mediates apoptosis in plasmid R1-free cells. The predictions made by their simulations were in good agreement with experimental data (Nagel, 1999). See Higgs (2000) for other examples of how computer simulations have yielded important insight into the folding kinetics of metastable RNA molecules.
SV11 is a 115nt synthetic RNA species that exists in two alternative stable states, a rod-like stable conformation and a multi-component metastable conformation (see Figure 3). While the metastable conformation is a template for Q$`\beta `$ replicase, the stable conformation is not (Zamora et al., 1995). Several authors have previously performed in silico analysis of the folding kinetics of SV11. Some of the authors (Flamm, 1998; Flamm et al., 2000) successfully predicted the existence of the metastable conformation while others (e.g., see Gultyaev et al., 1990; Morgan & Higgs, 1996) could only do so if folding was constrained to occur in conjunction with transcription. However, none of the authors obtained detailed theoretical estimates of the population dynamics of the molecule’s two alternative stable states.
Here we report detailed estimates of the population dynamics of the stable and metastable conformations of SV11. As shown in Figure 4 the population of the metastable state increases steadily, reaching a maximum after about $`2000\mu s`$. In experiments performed using kinfold, Flamm (1998) reported that the fraction of the metastable state reached about $`16\%`$ after $`500\mu s`$. This estimate is consistent with the results shown in Figure 4. However, we were able to fold the molecule for much longer, up to $`1.5\times 10^4\mu s`$, than was done in Flamm (1998). This allowed us to estimate the population dynamics of not just the metastable conformation but also of the stable native conformation. The ratio of the fraction of simulations that found the native conformation to the fraction that found the metastable conformation in the time scale of the simulation was approximately $`3`$ to $`1`$. Note that the accuracy of these results can be tested directly, in the laboratory.
## 5 Discussion
In this paper, we introduced a new model for the kinetic folding of RNA sequences into secondary structures that was inspired by the theory of complex adaptive systems. In the folding model, RNA bases and base pairs engage in local stacking interactions that determine the probabilities (or fitnesses) of possible RNA structural rearrangements (SRs). Meanwhile, selection operates at the level of SRs; an autonomous stochastic sampling process periodically selects a subset of possible SRs for realization based on the fitnesses of the SRs. Several examples were used to illustrate the applicability of the model. In particular, certain base modifications were shown to substantially improve the foldability of tRNA<sup>Phe</sup>. In addition, a characteristic optimal folding temperature $`T_{opt}\left(313K\right)`$ of tRNA<sup>Phe</sup> was identified. Furthermore, the model was used to confirm previous experimental results (Zamora et al., 1995) regarding the existence of two alternative stable states of the Q$`\beta `$ variant SV11, and to obtain (experimentally verifiable) estimates of the population dynamics of those states. The above examples demonstrated, among other things, the emergence from (local) SRs of nonlinear RNA folding dynamics (i.e., the realization of alternative stable states). Other possible applications of the model are discussed below.
The analysis of properties of fitness landscapes is currently of interest to researchers in a wide range of fields including the life, computer, and social sciences (e.g., see Hadany & Beker, 2003; McCarthy, 2004; Skellett et al., 2005). A number of interesting general features of these landscapes such as positive correlations between the degree of epistasis, the number of local optima, and the expected value of the global optimum have thus far been elucidated in the context of adaptive walks on landscapes generated by random Boolean networks (RBNs) (Kauffman & Levin, 1987; Kauffman, 1989; Skellett et al., 2005) . The simplicity of RBNs and their accessibility to some degree of mathematical analysis make them convenient for use as generators of fitness landscapes. However, some of the general properties of such RBN-generated landscapes may differ from those of landscapes found in Nature. RNA folding kinetics, as modeled in this paper, could serve as a generator of, and therefore assist the analysis of properties of ”natural” fitness landscapes (i.e., RNA folding energy landscapes). Note that this proposed use of RNA in the investigation of fitness landscapes differs from the previous related use of the RNA sequence to structure mapping (Schuster & Stadler, 1994), which was not based on folding kinetics.
The model could also be used to study the RNA sequence to structure (or genotype to phenotype) mapping, from a kinetics perspective. Previous investigations, based on the thermodynamics of RNA folding, have made several important findings about the RNA sequence to structure mapping such as the existence of (1) extended neutral networks of sequences that fold into the same secondary structures (2) few common or ”typical” structures, realized with relatively high frequencies, and (3) many rare structures that have little or no evolutionary significance (Schuster et al., 1994; Schuster et al., 1998). These findings, together with results from thermodynamics-based RNA optimization experiments (Fontana & Schuster, 1998), have confirmed previous hypotheses on, and shed light into several important features of the process of molecular evolution, including the role of neutrality in adaptation and the existence of continuous/discontinuous transitions or punctuated equilibria in evolutionary trajectories. It would be interesting to determine how the nature of the RNA sequence to structure mapping, as obtained from thermodynamic folding experiments, changes when RNA folding kinetics is taken into account. It is also possible that an investigation of the genotype to phenotype mapping based on the kinetics of RNA folding will yield further insight into the process of molecular evolution.
## 6 Acknowledgements
This work was funded in part by DOE-ER63580 and by RCMI Grant no. RR017581. The author thanks Dr. Asamoah Nkwanta and anonymous reviewers for their useful comments on an earlier version of the manuscript.
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warning/0507/cond-mat0507549.html | ar5iv | text | # Electronic Transport in Fullerene C20 Bridge Assisted by Molecular Vibrations
## Abstract
The effect of molecular vibrations on electronic transport is investigated with the smallest fullerene C<sub>20</sub> bridge, utilizing the Keldysh nonequilibrium Green’s function techniques combined with the tight-binding molecular-dynamics method. Large discontinuous steps appear in the differential conductance when the applied bias-voltage matches particular vibrational energies. The magnitude of the step is found to vary considerably with the vibrational mode and to depend on the local electronic states besides the strength of electron-vibration coupling. On the basis of this finding, a novel way to control the molecular motion by adjusting the gate voltage is proposed.
preprint: APS/123-QED
Inelastic transport associated with local heating in nanoscale devices has been a growing interest in the fields of nanoscience and nanotechnology over the past few years rf:Smit ; rf:Agrait ; rf:Ness ; rf:Mont ; rf:Chen . Fullerenes are considered promising candidates for basic elements in nanoscale devices, therefore, electronic transport in fullerene bridges has received significant attention both experimentally rf:Park ; rf:Pasupathy ; rf:Joachim and theoretically rf:Nakanishi ; rf:Palacios . McEuen and co-workers measured current-voltage ($`I`$-$`V`$) characteristics of the C<sub>60</sub> and C<sub>140</sub> bridges rf:Park ; rf:Pasupathy and reported that coupling between electronic and vibrational degrees of freedom plays an important role in electronic transport.
The strength of the electron-vibration coupling is known to be enhanced with decreasing the diameter of the fullerene rf:Devos-L . Thus, among fullerenes, the smallest fullerene C<sub>20</sub> is expected to have the largest electron-vibration coupling. Recently, the C<sub>20</sub> fullerene was synthesized by Prinzbach et al rf:Prinzbach . To the best of our knowledge, the issue of inelastic transport through the C<sub>20</sub> bridge in the presence of electron-vibration scattering has not been addressed. Although elastic transport through the C<sub>20</sub> bridge has been theoretically investigated by several authors rf:Miyamoto ; rf:Roland ; rf:Otani , characteristic features remain to be elucidated.
The aim of this Letter is to clarify the role of molecular vibrations in electronic transport through the C<sub>20</sub> connected to Au electrodes. A semi-infinite one-dimensional (1D) chain is used as a simple and ideal model for the Au electrode, as shown in Fig. 4. In an effort to achieve the aim, molecular vibrations in the C<sub>20</sub> are first examined using the tight-binding molecular-dynamics (TBMD) method rf:Wang . The influence of molecular vibrations on electronic transport characteristics of the C<sub>20</sub> bridge is discussed in terms of the Keldysh nonequilibrium Green’s function (NEGF) method rf:Keldysh , a powerful method used to analyze inelastic transport in the presence of electron-vibration interactions rf:Caroli ; rf:Datta . Recently, the NEGF formalism has been successfully applied to inelastic transports in several nanostructures, such as atomic wires rf:Frederiksen1 ; rf:Frederiksen2 and molecular junctions rf:Zhu ; rf:Galperin ; rf:Pecchia ; rf:Asai .
Before studying the vibrational properties of the C<sub>20</sub> bridge, the stable structure of C<sub>20</sub> is determined using the TBMD method. The structure of C<sub>20</sub> is distorted from the highest possible $`I_h`$ symmetry to the $`D_{3d}`$ symmetry after the optimization, due to the Jahn-Teller effect. The optimized structure of the C<sub>20</sub> contains three types of nonequivalent atoms, labeled $`a`$, $`b`$, and $`c`$ in Fig. 4. There are four distinct bond lengths: $`a_{ab}=`$1.464 Å, $`a_{bc}=`$1.469 Å, $`a_{cc^{}}=`$1.519 Å, and $`a_{cc^{\prime \prime }}=`$1.435 Å. The obtained bond lengths differ by less than 1 $`\%`$ from those obtained using the first-principles calculations rf:Galli .
The vibrational energies of the isolated C<sub>20</sub> are first calculated by diagonalizing a 60$`\times `$60 force-constant matrix derived from the TBMD calculations. The results of the vibrational energies are classified using the irreducible representations of the $`D_{3d}`$ symmetry group, as listed in Table 1. Similarly, the vibrational energies of the C<sub>20</sub> connected by two springs to the Au electrodes are calculated. The springs are attached to the $`a`$-atoms on either side of the principal axis of $`D_{3d}`$ symmetry of the C<sub>20</sub>, as shown in Fig. 4. The spring constant was assumed to be $`4.37`$ eV/Å<sup>2</sup>, an experimental value for the C<sub>60</sub> attached to Au electrodes rf:Park . In the present calculations, atomic vibrations of Au atoms were neglected because an Au atom is considerably heavier than a C atom. The connection of the electrodes only influences the vibrational modes with $`A_{1g}`$ and $`A_{2u}`$ symmetries, because modes with other symmetries show no stretching of the springs between the C<sub>20</sub> and Au electrodes. The vibrational energies shifted by the connection are listed in the $`A_{1g}^{}`$ and $`A_{2u}^{}`$ columns in Table 1. It is important to note that the vibrational mode of 11.1 meV in the $`A_{2u}^{}`$ column corresponds to the shuttle motion of C<sub>20</sub> that goes back and forth between the two Au electrodes. Using these 55 vibrational modes, the Hamiltonian for the molecular vibrations of the C<sub>20</sub> bridge can be written as $`_{\mathrm{vib}}=_\lambda \mathrm{}\omega _\lambda (b_\lambda ^{}b_\lambda +\frac{1}{2})`$, where $`b_\lambda ^{}`$ ($`b_\lambda `$) is the creation (annihilation) operator of the vibrational quanta (i.e., phonon) with energy $`\mathrm{}\omega _\lambda `$.
The Hamiltonian for conduction electrons in the C<sub>20</sub> bridge by the tight-binding model within the Hückel approximation is subsequently described. The tight-binding Hamiltonian is expressed by the sum of five parts: $`_{\mathrm{el}}=_\mathrm{L}+_{\mathrm{LM}}+_\mathrm{M}+_{\mathrm{MR}}+_\mathrm{R}`$. $`_\mathrm{M}`$ represents the Hamiltonian for an extended molecule, Au-C<sub>20</sub>-Au, including the edge Au atoms in the semi-infinite 1D electrodes, $`_{\mathrm{L}/\mathrm{R}}`$ for the left/right electrodes without the edge atom, and $`_{\mathrm{LM}/\mathrm{MR}}`$ for the contact between them. The Hamiltonian for the extended molecule is expressed as $`_\mathrm{M}=_iϵ_ic_i^{}c_i+_{i,j}t_{ij}^0(c_i^{}c_j+\mathrm{h}.\mathrm{c}.)`$, where $`ϵ_i`$ is the on-site energy of $`\pi `$ orbitals for C atoms and of $`6s`$ orbitals for Au atoms, $`t_{ij}^0`$ is a hopping parameter between $`i`$th and $`j`$th orbitals in equilibrium, and $`c_i^{}`$ ($`c_i`$) is the creation (annihilation) operator of an electron on the $`i`$th orbital. In the present investigation, the on-site energy $`ϵ_{\mathrm{Au}}`$ for the Au electrode is assumed to lie in the center of the HOMO-LUMO gap of the isolated C<sub>20</sub>. The hopping parameters between Au and C atoms are chosen to be $`1.0`$ eV, whereas those between $`\pi `$ orbitals are determined using the TBMD calculations rf:Wang .
The electron-vibration interaction is given by
$`_{\mathrm{e}\text{-}\mathrm{v}}`$ $`=`$ $`{\displaystyle \underset{\lambda }{}}{\displaystyle \underset{ij}{}}g_{ij}^\lambda (c_i^{}c_j+\mathrm{h}.\mathrm{c}.)(b_\lambda ^{}+b_\lambda ),`$ (1)
where the coupling constant $`g_{ij}^\lambda `$ is given as
$`g_{ij}^\lambda ={\displaystyle \underset{\alpha =x,y,z}{}}{\displaystyle \frac{t_{ij}}{s_{i\alpha }}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _\lambda }}}\left[{\displaystyle \frac{u_{i\alpha }^\lambda }{\sqrt{M_i}}}{\displaystyle \frac{u_{j\alpha }^\lambda }{\sqrt{M_j}}}\right],`$ (2)
where $`u_{i\alpha }^\lambda `$ is an eigenvector of the dynamic matrix and $`dt_{ij}/ds_{i\alpha }`$ is the hopping modulation with an atomic displacement $`s_{i\alpha }`$ from equilibrium along the $`\alpha `$ direction of $`i`$th carbon atom. The modulations $`dt_{ij}/ds_{i\alpha }`$ within the C<sub>20</sub> are determined using the TBMD calculations, and those between the C<sub>20</sub> and Au electrodes are assumed to be $`1.0`$ eV/Å.
According to the NEGF formalism, the electrical current from the left electrode to a system is given by
$`I={\displaystyle \frac{2e}{h}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑ϵ\mathrm{Tr}\left[𝚺_L^<(ϵ)𝑮^>(ϵ)𝚺_L^>(ϵ)𝑮^<(ϵ)\right].`$ (3)
For the steady-state transport, the lesser and greater Green’s functions obey the steady-state Keldysh equation: $`𝑮^{<>}=𝑮^r(𝚺_\mathrm{L}^{<>}+𝚺_\mathrm{R}^{<>}+𝚺_{\mathrm{e}\text{-}\mathrm{v}}^{<>})𝑮^a`$, where the self-energy $`𝚺_{\mathrm{L}(\mathrm{R})}^{<>}`$ comes from the left (right) contacts and $`𝚺_{\mathrm{e}\text{-}\mathrm{v}}^{<>}`$ is related to electron-vibration interaction. The retarded (advanced) Green’s function is given by the Dyson equation: $`𝑮^{r(a)}=𝑮^{0,r(a)}+𝑮^{0,r(a)}𝚺_{\mathrm{e}\text{-}\mathrm{v}}^{r(a)}𝑮^{r(a)}`$, where $`𝑮^{0,r(a)}`$ is an unperturbed retarded (advanced) Green’s function and $`𝚺_{\mathrm{e}\text{-}\mathrm{v}}^{r(a)}`$ is the retarded (advanced) self-energy due to electron-vibration interaction. In the present calculation, the self-energies due to electron-vibration interaction are treated perturbatively using the Feynman diagram technique, and are expanded up to the lowest order of self-energy diagrams. On the other hand, the self-energies of semi-infinite electrodes can be calculated analytically in the case of the semi-infinite 1D electrode rf:Datta . The equation (3) is reduced to the well-known Landauer formula for the elastic current $`I_{\mathrm{el}}`$ by setting $`𝚺_{\mathrm{e}\text{-}\mathrm{v}}^<=0`$ in the Keldysh equation.
The $`I`$-$`V`$ characteristics of the C<sub>20</sub> bridges are shown in the inset in Fig. 4, where the solid and dashed curves represent the total current $`I`$ and the elastic current $`I_{\mathrm{el}}`$, respectively. From the calculated $`I`$-$`V`$ characteristics, the total current curve is found to deviate slightly upward from the elastic curve as the applied bias-voltage increases. The dissipation power into molecular vibrations is also estimated, up to $`V_{\mathrm{bias}}=100`$ mV, as of the order of $`1`$ nW, which is less than $`3\%`$ of the total power generated in the entire bridge. This means that local heating due to the molecular vibrations is not a serious bottleneck for transport characteristic of the C<sub>20</sub> bridge compared with that due to contact resistance.
The contribution from each vibrational mode to the transport characteristics can be seen clearly in the differential conductance, $`dI/dV_{\mathrm{bias}}`$, shown in Fig. 4. Solid and dashed curves in Fig. 4 represent the total differential conductance $`dI/dV_{\mathrm{bias}}`$ and its elastic part $`dI_{\mathrm{el}}/dV_{\mathrm{bias}}`$, respectively. In contrast to the smooth curve for the elastic differential conductance, large discontinuous jumps appear in the total differential conductance curve at particular bias-voltages, as indicated by arrows in Fig. 4. The contribution from vibrational modes other than those indicated by the arrows is increasingly small, despite the nonzero coupling constants $`g_{ij}^\lambda `$. A typical example is illustrated by the conductance step at $`11.1`$ mV, originating from the shuttle motion of C<sub>20</sub>, which is two orders of magnitude less than the step at $`29.9`$ mV. The shuttle motion will be discussed in further detail later.
The physical origin for the considerable variability of the magnitude of the discontinuous step in Fig. 4 with the vibrational mode, is now discussed from the viewpoint of the electronic structures of the C<sub>20</sub> bridge. Fermi’s golden rule gives us not only the scattering rate of electronic states, but also the magnitude of these steps in the conductance curve. According to Fermi’s golden rule, the magnitude of the steps is proportional to $`|_{i,j}g_{ij}^\lambda \mathrm{\Psi }_i^{}(ϵ\mathrm{}\omega _\lambda )\mathrm{\Psi }_j(ϵ)|^2`$ . Here $`\mathrm{\Psi }_i(ϵ)`$ represents a scattering electronic state at the $`i`$th atom, in the absence of an electron-vibration interaction. The scattering electronic states within the narrow bias-window $`[ϵ_FeV/2,ϵ_F+eV/2]`$ can be approximately replaced by $`\mathrm{\Psi }_i(ϵ_F)`$ at the Fermi level, therefore, the magnitude of the conductance steps can be estimated by $`S_\lambda \mathrm{}^2|_{i,j}g_{ij}^\lambda G_{ij}^{0,<}(ϵ_F)|^2`$, referred to as the scattering intensity hereafter, where $`G_{ij}^{0,<}(ϵ_F)=i\mathrm{}^1\mathrm{\Psi }_i^{}(ϵ_F)\mathrm{\Psi }_j(ϵ_F)`$ is the unperturbed lesser Green’s function. Figure 4 shows the scattering intensity $`S_\lambda `$ for the C<sub>20</sub> bridge. The intensity exhibits large peaks for particular modes indicated by the arrows. These mode energies coincide completely with the positions of the discontinuous steps indicated by the arrows in Fig 4. As seen in the expression of $`S_\lambda `$, the magnitude of the conductance steps is determined by the local electronic states near the Fermi level, as well as vibrational states of the C<sub>20</sub> bridge. $`G_{ij}^{0,<}(ϵ_F)`$ describes the correlation between electrons with the Fermi energy on the $`i`$th and $`j`$th atoms and the $`g_{ij}^\lambda `$ has a large value for strong stretching of interatomic bond, therefore, large conductance-steps are thought to appear when the bond between atoms with high electron-densities stretch strongly.
Thus, a novel way to control the motion of C<sub>20</sub> between two electrodes under current flow, utilizing the obtained knowledge from the correlation between scattering electronic states and molecular vibrations is proposed. The shuttle motion of C<sub>20</sub> is a key motion, however, its excitation rate is very low owing to the small $`S_\lambda `$ at $`\mathrm{}\omega _\lambda =11.1`$ meV in Fig. 4. As explained previously, the rate is expected to increase significantly when the electronic state localized at C atoms adjacent to Au electrodes lie close to the Fermi level, because only two springs between C<sub>20</sub> and Au electrodes stretch in the shuttle motion. However, such a localized state lies $`1.5`$ eV below the Fermi level ($`ϵ_F=0`$ eV), as shown in Fig. 4(a), and the scattering intensity of the shuttle motion exhibits a maximum peak at $`1.5`$ eV in Fig. 4(c). Therefore, its excitation rate can be enhanced by tuning the gate voltage to shift the localized state to the Fermi level. Of course, more precise analyses beyond the rigid-band picture assumed here will be necessary to quantify the gate-voltage effect on the Fermi level. The control mechanism of the shuttle motion of C<sub>20</sub> proposed herein is complementary with the Coulomb blockade mechanism for the recent experiment on electronic transports in the C<sub>60</sub> weakly suspended between Au electrodes rf:Park . The mechanism proposed in the current investigation is efficient, even for inelastic electronic transport through fullerenes strongly connected to the electrodes.
We finally emphasize the extention of the proposed idea to other molecular bridges. For the 1,4-benzenedithiol (BDT=C<sub>6</sub>H<sub>4</sub>S<sub>2</sub>) bridge, it has been theoretically predicted that a rotational motion of 1,4-BDT molecule is strongly coupled to scattering electronic states at low-bias voltages rf:Sergueev . By analogy with our results, the rotational motion will be switched on when a scattering state with high electron-densities on sulfur atoms in 1,4-BDT adjacent to the electrodes is shifted to the Fermi level by tuning the gate voltage.
In conclusion, the inelastic electronic transport assisted by molecular vibrations in the smallest fullerene C<sub>20</sub> bridge was investigated using the NEGF formalism combined with the TBMD method. The effect of the molecular vibrations on the transport characteristics of C<sub>20</sub> appear as discontinuous steps with various heights in the differential conductance curve. The stepwise behavior of the conductance, and the variation of magnitude of the step was clearly understood by analyzing the electronic structures based on Fermi’s golden rule, as well as the vibrational structures. Local heating caused by the molecular vibrations of the C<sub>20</sub> bridge were also found not to be a serious bottleneck for functions of fullerene-based nanodevices. Moreover, the idea that the shuttle motion of the fullerenes can be controlled by adequately tuning the gate voltage was proposed. The present study provides a clue toward a more complete understanding of inelastic electronic transport and the electromechanics in various molecular devices.
This work was supported in part by “Academic Frontier” Project of MEXT (2005-2010). Part of the numerical calculations were performed on the Hitachi SR8000s at ISSP, The University of Tokyo. |
warning/0507/math0507496.html | ar5iv | text | # The 𝑝-adic local monodromy theorem for fake annuli
## 1 Introduction
This paper proves a generalization of the $`p`$-adic local monodromy theorem of André , Mebkhout , and the present author . That theorem, originally conjectured by Crew as an analogue in rigid ($`p`$-adic) cohomology of Grothendieck’s local monodromy theorem in étale ($`\mathrm{}`$-adic) cohomology, asserts the quasi-unipotence of differential modules with Frobenius structure on certain one-dimensional rigid analytic annuli.
The $`p`$-adic local monodromy theorem has far-reaching consequences in the theory of rigid cohomology, particularly for curves . However, although one can extend it to a relative form \[15, Theorem 5.1.3\] to obtain some higher-dimensional results, for some applications one needs a version of the monodromy theorem which is truly higher-dimensional.
This theorem takes an initial step towards producing such a higher-dimensional monodromy theorem, by proving a generalization in which the role of the annulus is replaced by a somewhat mysterious space, called a *fake annulus*, described by certain rings of multivariate power series. When there is only one variable, the space is a true annulus, so the result truly generalizes the original monodromy theorem. Indeed, this paper has the side effect of giving an exposition of the original theorem, albeit one somewhat encumbered by extra notation needed for the fake case.
In the remainder of this introduction, we explain further the context in which the $`p`$-adic local monodromy theorem arises, introduce and justify the fake analogue, and outline the structure of the paper.
### 1.1 Monodromy of $`p`$-adic differential equations
Let $`K`$ be a field of characteristic zero complete with respect to a nonarchimedean absolute value, whose residue field $`k`$ has characteristic $`p>0`$. Suppose we are given a rigid analytic annulus over $`K`$ and a differential equation on the annulus, i.e., a module equipped with an integrable connection. We now wish to define the “monodromy around the puncture” of this connection, despite not having recourse to the analytic continuation we would use in the analogous classical setting. In particular, we would like to construct a representation of an appropriate étale fundamental group, whose triviality or unipotence amounts to the existence of a full set of horizontal sections or log-sections; the latter relates closely to the existence of an extension or logarithmic extension of the connection across the puncture (e.g., \[16, Theorem 6.4.5\]).
We can define a monodromy representation associated to a connection if we can find enough horizontal sections on some suitable covering space. In particular, we are mainly interested in connections which become unipotent, i.e., can be filtered by submodules whose successive quotients are trivial for the connection, on some cover of the annulus which is “finite étale near the boundary” (in a sense that can be made precise). One can then construct a monodromy representation, using the Galois action on horizontal sections, giving an equivalence of categories between such *quasi-unipotent* modules with connection and a certain representation category. (Beware that if the field $`k`$ is not algebraically closed, these representations are only semilinear over the relevant field, namely the maximal unramified extension of $`K`$. See \[13, Theorem 4.45\] for a precise statement.)
In order for such an equivalence to be useful, we need to be able to establish conditions under which a module with connection is forced to be quasi-unipotent. As suggested in the introduction to , a natural, geometrically meaningful candidate restriction (analogous to the existence of a variation of Hodge structure for a complex analytic connection) is the existence of a *Frobenius structure* on the connection. For $`K`$ discretely valued, the fact that connections with a Frobenius structure are quasi-unipotent is the content of the $`p`$-adic local monodromy theorem ($`p`$LMT) of André , Mebkhout , and this author .
Note that in this paper, we will not go all the way to the construction of monodromy representations. These appear directly in André’s proof of the $`p`$LMT (at least for $`k`$ algebraically closed); for a direct construction (of Fontaine type) assuming the $`p`$LMT, see . See Remark 6.2.7 for more details.
### 1.2 Fake annuli
The *semistable reduction problem* (or *global quasi-unipotence problem*) for overconvergent $`F`$-isocrystals, as formulated by Shiho \[23, Conjecture 3.1.8\] and reformulated in \[16, Conjecture 7.1.2\], is essentially to give a higher dimensional version of the $`p`$LMT. From the point of view of , this can be interpreted as proving a uniform version of the $`p`$LMT across all divisorial valuations on the function field of the original variety. This interpretation immediately suggests that one needs to exploit the quasi-compactness of the Riemann-Zariski space associated to the function field of an irreducible variety; this observation is developed in more detail in .
The upshot is that one must prove the $`p`$LMT uniformly for the divisorial valuations in a neighborhood (in Riemann-Zariski space) of an arbitrary valuation, not just a divisorial one. Ideally, one could proceed by first verifying whether the $`p`$LMT itself makes sense and continues to hold true when one passes from a divisorial valuation to a more general one. This entails replacing the annuli in the $`p`$LMT with some sort of “fake annuli” which cannot be described as rigid analytic spaces in the usual sense. Nonetheless, one can still sensibly define rings of analytic functions in a neighborhood of an irrational point, and thus set up a ring-theoretic framework in which an analogue of the $`p`$-adic local monodromy theorem can be formulated. (This allows us to get away with the linguistic swindle of speaking meaningfully about “$`p`$-adic differential equations on fake annuli” without giving the noun phrase “fake annulus” an independent meaning!)
Indeed, this framework fits naturally into the context of the slope filtration theorem of . That theorem, which gives a structural decomposition of a semilinear endomorphism on a finite free module over the Robba ring (of germs of rigid analytic functions on an open annulus with outer radius 1), does not make any essential use of the fact that the Robba ring is described in terms of power series. Indeed, as presented in , the theorem applies directly to our fake annuli; thus to prove the analogue of the $`p`$-adic monodromy theorem, one needs only analogize Tsuzuki’s unit-root monodromy theorem from . With a bit of effort, this can indeed be done, thus illustrating some of the power of the Frobenius-based approach to the monodromy theorem. Note that our definition of fake annuli will actually include true rigid analytic annuli, so the monodromy theorem given here will strictly generalize the $`p`$-adic local monodromy theorem.
Unfortunately, it is not so clear how to prove a form of the $`p`$LMT for arbitrary valuations; in this paper, we restrict to a somewhat simpler class. These are the monomial valuations, which correspond to monomial orderings in a polynomial ring. For instance, these include valuations on $`k(x,y)`$ in which the valuations of $`x`$ and $`y`$ are linearly independent over the rational numbers. There are many valuations that do not take this form, namely the *infinitely singular* valuations; the semistable reduction problem for these must be treated in a more roundabout fashion, which we will not discuss further here.
### 1.3 Structure of the paper
We conclude this introduction with a summary of the various sections of the paper.
In Section 2, we define the rings corresponding to fake annuli, and verify that they fit into the formalism within which slope filtrations are constructed in .
In Section 3, we define $`F`$-modules, $``$-modules, and $`(F,)`$-modules on fake annuli, and verify that the category of $`(F,)`$-modules is invariant of the choice of a Frobenius lift (Proposition 3.4.7).
In Section 4, we give a fake annulus generalization of Tsuzuki’s theorem on unit-root $`(F,)`$-modules (Theorem 4.5.2).
In Section 5, we invoke the technology of slope filtrations from (via ), and apply it to deduce from Theorem 4.5.2 a form of the $`p`$-adic local monodromy theorem for $`(F,)`$-modules on fake annuli (Theorem 5.2.4).
In Section 6, we deduce some consequences of the $`p`$-adic local monodromy theorem. Namely, we calculate some extension groups in the category of $`(F,)`$-modules, establish a local duality theorem, and generalize some results from and on the full faithfulness of overconvergent-to-convergent restriction.
### Acknowledgments
Thanks to Nobuo Tsuzuki for some helpful remarks on the unit-root local monodromy theorem. Thanks also to Francesco Baldassarri and Pierre Berthelot for organizing useful workshops on $`F`$-isocrystals and rigid cohomology in December 2004 and June 2005. The author was supported by NSF grant number DMS-0400727.
## 2 Fake annuli
In this section, we describe ring-theoretically the fake annuli to which we will be generalizing the $`p`$-adic local monodromy theorem, deferring to for most of the heavy lifting.
First, we put in some notational conventions that will hold in force throughout the paper; for the most part, these hew to the notational régime of (which in turn mostly follows ), with a few modifications made for greater consistency with .
###### Convention 2.0.1.
Throughout this paper, let $`K`$ be a complete *discretely valued* field of characteristic $`0`$, whose residue field $`k`$ has characteristic $`p>0`$. Let $`𝔬=𝔬_K`$ be the ring of integers of $`K`$, and let $`\pi `$ denote a uniformizer of $`K`$. Let $`w`$ be the valuation on $`𝔬`$ normalized so that $`w(\pi )=1`$. Let $`q`$ be a power of $`p`$, and assume extant and fixed a ring endomorphism $`\sigma _K:KK`$, continuous with respect to the $`\pi `$-adic valuation, and lifting the $`q`$-power endomorphism on $`k`$. Let $`K_q`$ be the fixed field of $`K`$ under $`\sigma _K`$, and let $`𝔬_q`$ be the fixed ring of $`𝔬`$ under $`\sigma _K`$. Finally, let $`I_n`$ denote the $`n\times n`$ identity matrix over any ring.
###### Remark 2.0.2.
As noted in the introduction, the restriction to $`K`$ discretely valued is endemic to the methods of this paper; see Remark 2.4.4 for further discussion.
### 2.1 Monomial fields
###### Definition 2.1.1.
Let $`k`$ be a field. A *nearly monomial field (of height $`1`$) over $`k`$* is a field $`E`$ equipped with a valuation $`v:E^{}`$ (also written $`v:E\{+\mathrm{}\}`$) satisfying the following restrictions.
1. The field $`E`$ is a separable extension of $`k`$, i.e., $`kE`$ and $`kE^p=k^p`$.
2. The image $`v(E^{})`$ of $`v`$ is a finitely generated $``$-submodule of $``$, and $`v(k^{})=\{0\}`$.
3. The field $`E`$ is complete with respect to $`v`$.
4. With the notations
$`𝔬_E`$ $`=\{xE:v(x)0\}`$
$`𝔪_E`$ $`=\{xE:v(x)>0\}`$
$`\kappa _E`$ $`=𝔬_E/𝔪_E,`$
the natural map $`k\kappa _E`$ is finite.
If $`\kappa _E=k`$, we say $`E`$ is a *monomial field* (or *fake power series field*) over $`k`$; in that case, $`k`$ is integrally closed in $`E`$. We define the *rational rank* of $`E`$ to be the rank of $`v(E^{})`$ as a $``$-module.
###### Remark 2.1.2.
One can also speak of monomial fields of height greater than 1, by allowing the valuation $`v`$ to take values in a more general totally ordered abelian group. The techniques used in this paper do not apply to that case, so we will ignore it; in the semistable reduction context, one can eliminate the case of height greater than 1 by an inductive argument \[17, Proposition 4.2.4\].
###### Example 2.1.3.
A monomial field of rational rank 1 is just a power series field, by the Cohen structure theorem. This characterization generalizes to arbitrary monomial fields; see Definition 2.1.9. Note also that nearly monomial fields are examples of *Abhyankar valuations*, i.e., valuations in which equality holds in Abhyankar’s inequality \[26, Théorème 9.2\].
###### Remark 2.1.4.
If $`E`$ is a nearly monomial field and $`E^{}/E`$ is a finite separable extension, then $`E^{}`$ is also nearly monomial: $`E^{}`$ is separable over $`k`$, the valuation $`v`$ extends uniquely to a valuation $`v^{}`$ on $`E^{}`$, $`E^{}`$ is complete with respect to $`v^{}`$, and the index $`[v^{}((E^{})^{}):v(E^{})]`$ and degree $`[\kappa _E^{}:\kappa _E]`$ are both finite since their product is at most $`\mathrm{deg}(E^{}/E)`$ \[26, Proposition 5.1\]. Conversely, every nearly monomial field can be written as a finite separable extension of a monomial field; see Definition 2.1.9.
###### Remark 2.1.5.
If $`E`$ is a nearly monomial field over $`k`$ and $`\kappa _E/k`$ is separable, then by Hensel’s lemma, the integral closure $`k^{}`$ of $`k`$ in $`E`$ is isomorphic to $`\kappa _E`$; in other words, $`E`$ is a monomial field over $`k^{}`$. In particular, if $`k`$ is perfect, then any finite extension of a monomial field over $`k`$ is a monomial field over some finite separable extension of $`k`$. This fails if $`k`$ is not perfect, even for finite separable extensions of the monomial field: the field $`k`$ is integrally closed in
$$k((t))[z]/(z^pzct^p)(ckk^p),$$
but the latter has residue field $`k(c^{1/p})k`$.
It will frequently be convenient to work with monomial fields in terms of coordinate systems.
###### Definition 2.1.6.
Let $`m`$ be a nonnegative integer. Let $`L`$ be a lattice in $`^m`$, i.e., a $``$-submodule of $`^m`$ which is free of rank $`m`$, and which spans $`^m`$ over $``$. Let $`L^{}(^m)^{}`$ denote the lattice dual to $`L`$:
$$L^{}=\{\mu (^m)^{}:\mu (z)zL\}.$$
Given a formal sum $`_{zL}c_z\{z\}`$, with the $`c_z`$ in some ring, define the *support* of the sum to be the set of $`zL`$ such that $`c_z0`$; define the support of a matrix of formal sums to be the union of the supports of the entries. If $`SL`$ and a formal sum or matrix has support contained in $`S`$, we also say that the element or matrix is “supported on $`S`$”. For $`R`$ a ring, let $`R[L]`$ denote the group algebra of $`L`$ over $`R`$, i.e., the set of formal sums $`_{zL}c_z\{z\}`$ with coefficients in $`R`$ and finite support.
###### Remark 2.1.7.
It is more typical to denote the class in $`R[L]`$ of a lattice element $`zL`$ by $`[z]`$, rather than $`\{z\}`$. However, we need to use brackets to denote Teichmüller lifts, so we will stick to braces for internal consistency.
###### Definition 2.1.8.
For any ring $`R`$ and any $`\lambda (^n)^{}`$, let $`v_\lambda `$ denote the valuation on $`R[L]`$ given by
$$v_\lambda \left(\underset{zL}{}c_z\{z\}\right)=\mathrm{min}\{\lambda (z):zL,c_z0\}.$$
Let $`P_\lambda L`$ denote the submonoid of $`zL`$ for which $`\lambda (z)0`$, and let $`R[L]_\lambda `$ denote the monoid algebra $`R[P_\lambda ]`$. Let $`RL_\lambda `$ and $`R((L))_\lambda `$ denote the $`v_\lambda `$-adic completions of $`R[L]_\lambda `$ and $`R[L]`$, respectively.
###### Definition 2.1.9.
Given a lattice $`L`$ and some $`\lambda (^m)^{}`$, we say $`\lambda `$ is *irrational* if $`L\mathrm{ker}(\lambda )=\{0\}`$. In this case, $`k((L))_\lambda `$ is a monomial field over $`k`$. Conversely, given a nearly monomial field $`E`$ over a field $`k`$, with valuation $`v`$, a *coordinate system* for $`E`$ is a sequence $`x_1,\mathrm{},x_m`$ of elements of $`E`$ such that $`v(x_1),\mathrm{},v(x_m)`$ freely generate $`v(E^{})`$ as a $``$-module. Given a coordinate system, put $`L=^m`$ with generators $`z_1,\mathrm{},z_m`$, and define $`\lambda (^m)^{}`$ by $`\lambda (z_i)=v(x_i)`$; then $`\lambda `$ is irrational, and the continuous map $`k((L))_\lambda E`$ given by $`\{z_i\}x_i`$ is injective. If we identify $`k((L))_\lambda `$ with its image in $`E`$, then $`E`$ is finite separable over $`k((L))_\lambda `$; if $`E`$ is monomial over $`k`$, then in fact $`E=k((L))_\lambda `$ by Proposition 2.1.10 below. This fact may be viewed as a monomial version of the Cohen structure theorem in equal characteristics.
###### Proposition 2.1.10.
Let $`L`$ be a lattice in $`^m`$, choose $`\lambda (^m)^{}`$ irrational, and let $`E`$ be a finite separable extension of $`k((L))_\lambda `$ with value group $`\lambda (L)`$ and residue field $`k`$. Then $`E=k((L))_\lambda `$.
###### Proof.
The claim is equivalent to showing that $`k((L))_\lambda `$ is inseparably defectless in the sense of Ostrowski’s lemma \[22, Théorème 2, p. 236\]. In particular, since a tamely ramified extension of $`E`$ is necessarily without defect, it suffices to check that there is no Artin-Schreier defect extension. If $`E=k((L))_\lambda [z]/(z^pzx)`$ were such an extension, we could rewrite it as $`k((L))_\lambda [z]/(z^pzy)`$ with the leading term of $`y`$ being either an element of $`pL`$ times a non-$`p`$-th power $`c`$ in $`k`$, or an element of $`LpL`$ times a nonzero element of $`k`$. But in the first case the residue field of $`E`$ would be $`k(c^{1/p})`$, and in the second case we would have $`[v(E):\lambda (L)]=p`$; in either case, we would contradict the assumption that $`E`$ is a defect extension. This contradiction yields the claim. ∎
###### Remark 2.1.11.
The term “monomial field” is modeled on the use of the term “monomial valuation”, e.g., in , to refer to a valuation $`v`$ of the sort considered in Definition 2.1.1. (Such valuations, each of which endows the lattice $`L`$ with a total ordering, are more common in mathematics than one might initially realize: for example, they are used to define highest weights in the theory of Lie algebras, and they are sometimes used to construct term orders in the theory of Gröbner bases.) In a previous version of this paper, the term “fake power series field” was used instead; we have decided that it would be better to save this term for describing the completion of a finitely generated field extension of $`k`$ with respect to *any* valuation of height 1. (See Remark 2.3.7 for some reasons why we are not considering such valuations here.)
### 2.2 Witt rings and Cohen rings
We now enter the formalism of \[14, § 2\].
###### Definition 2.2.1.
Let $`K^{\mathrm{perf}}`$ be the completion of the direct limit $`K\stackrel{\sigma _K}{}K\stackrel{\sigma _K}{}\mathrm{}`$ for the $`\pi `$-adic topology; this is a complete discretely valued field of characteristic 0 with residue field $`k^{\mathrm{perf}}`$, so it contains $`W(k^{\mathrm{perf}})`$ by Witt vector functoriality. For $`E`$ a perfect field of characteristic $`p`$ containing $`k`$, put $`\mathrm{\Gamma }^E=W(E)_{W(k^{\mathrm{perf}})}𝒪_{K^{\mathrm{perf}}}`$; note that the valuation $`w`$ extends naturally to $`\mathrm{\Gamma }^E`$.
###### Definition 2.2.2.
For $`E`$ a perfect field of characteristic $`p`$ containing $`k`$, complete for a valuation $`v`$ trivial on $`k`$, define the *partial valuations* $`v_n`$ on $`\mathrm{\Gamma }^E[\pi ^1]`$ as follows. Given $`x\mathrm{\Gamma }^E[\pi ^1]`$, write $`x=_i[\overline{x_i}]\pi ^i`$, where each $`\overline{x_i}E`$ and the brackets denote Teichmüller lifts. Set
$$v_n(x)=\underset{in}{\mathrm{min}}\{v(\overline{x_i})\}.$$
As in \[14, Definition 2.1.5\], the partial valuations satisfy some useful identities (here using $`\sigma `$ to denote the $`q`$-power Frobenius):
$`v_n(x+y)`$ $`\mathrm{min}\{v_n(x),v_n(y)\}`$ $`(x,y\mathrm{\Gamma }^E[\pi ^1],n)`$
$`v_n(xy)`$ $`\underset{m}{\mathrm{min}}\{v_m(x)+v_{nm}(y)\}`$ $`(x,y\mathrm{\Gamma }^E[\pi ^1],n)`$
$`v_n(x^\sigma )`$ $`=qv_n(x)`$ $`(x\mathrm{\Gamma }^E[\pi ^1],n)`$
$`v_n([\overline{x}])`$ $`=v(\overline{x})`$ $`(\overline{x}E,n0).`$
In the first two cases, equality holds whenever the minimum is achieved exactly once. Define the *levelwise topology* (or *weak topology*) on $`\mathrm{\Gamma }^E`$ by declaring that a sequence $`\{x_i\}`$ converges to zero if and only if for each $`n`$, $`v_n(x_i)\mathrm{}`$ as $`i\mathrm{}`$.
###### Definition 2.2.3.
For $`r>0`$, write $`v_{n,r}(x)=rv_n(x)+n`$; for $`r=0`$, write conventionally
$$v_{n,0}(x)=\{\begin{array}{cc}n\hfill & v_n(x)<\mathrm{}\hfill \\ \mathrm{}\hfill & v_n(x)=\mathrm{}.\hfill \end{array}$$
Let $`\mathrm{\Gamma }_r^E`$ be the subring of $`\mathrm{\Gamma }^E`$ for which $`v_{n,r}(x)\mathrm{}`$ as $`n\mathrm{}`$; then $`\sigma `$ sends $`\mathrm{\Gamma }_r^E`$ to $`\mathrm{\Gamma }_{r/q}^E`$. Define the map $`w_r`$ on $`\mathrm{\Gamma }_r^E`$ by
$$w_r(x)=\underset{n}{\mathrm{min}}\{v_{n,r}(x)\};$$
then $`w_r`$ is a valuation on $`\mathrm{\Gamma }_r`$ by \[14, Lemma 2.1.7\], and $`w_r(x)=w_{r/q}(x^\sigma )`$. Put
$$\mathrm{\Gamma }_{\mathrm{con}}^E=_{r>0}\mathrm{\Gamma }_r^E;$$
this is a henselian discrete valuation ring with maximal ideal $`\pi \mathrm{\Gamma }_{\mathrm{con}}^E`$ and residue field $`E`$ (see discussion in \[14, Definition 2.2.13\]).
###### Convention 2.2.4.
For $`E`$ a not necessarily perfect field complete for a valuation $`v`$ trivial on $`k`$, we write $`E^{\mathrm{perf}}`$ and $`E^{\mathrm{alg}}`$ for the *completed* (with respect to $`v`$) perfect and algebraic closures of $`E`$. When $`E`$ is to be understood, we abbreviate $`\mathrm{\Gamma }^{E^{\mathrm{perf}}}`$ and $`\mathrm{\Gamma }^{E^{\mathrm{alg}}}`$ to $`\mathrm{\Gamma }^{\mathrm{perf}}`$ and $`\mathrm{\Gamma }^{\mathrm{alg}}`$, respectively. Note that this is consistent with the conventions of but *not* with those of , where the use of these superscripts is taken not to imply completion.
Since we are interested in constructing $`\mathrm{\Gamma }^E`$ for $`E`$ a monomial field, which is not perfect, we must do a bit more work, as in \[14, § 2.3\].
###### Definition 2.2.5.
Let $`E`$ be a nearly monomial field over $`k`$ with valuation $`v`$. Let $`\mathrm{\Gamma }^E`$ be a complete discrete valuation ring of characteristic $`0`$ containing $`𝔬`$ and having residue field $`E`$, such that $`\pi `$ generates the maximal ideal of $`\mathrm{\Gamma }^E`$. Suppose that $`\mathrm{\Gamma }^E`$ is equipped with a *Frobenius lift*, i.e., a ring endomorphism $`\sigma `$ extending $`\sigma _K`$ on $`𝔬_K`$ and lifting the $`q`$-power Frobenius map on $`E`$. We may then embed $`\mathrm{\Gamma }^E`$ into $`\mathrm{\Gamma }^{\mathrm{perf}}`$ by mapping $`\mathrm{\Gamma }^E`$ into the first term of the direct system $`\mathrm{\Gamma }^E\stackrel{\sigma }{}\mathrm{\Gamma }^E\stackrel{\sigma }{}\mathrm{}`$, completing the direct system, and mapping the result into $`\mathrm{\Gamma }^{\mathrm{perf}}`$ via Witt vector functoriality. In particular, we may use this embedding to induce partial valuations and a levelwise topology on $`\mathrm{\Gamma }^E`$, taking care to remember that these depend on the choice of $`\sigma `$. If $`E^{}`$ is a finite separable extension of $`E`$, and we start with a suitable $`\mathrm{\Gamma }^E`$ equipped with a Frobenius lift $`\sigma `$, we may form the unramified extension of $`\mathrm{\Gamma }^E`$ with residue field $`E^{}`$; this will be a suitable $`\mathrm{\Gamma }^E^{}`$, and carries a unique Frobenius lift extending $`\sigma `$.
###### Definition 2.2.6.
Let $`E`$ be a nearly monomial field over $`k`$, and fix a pair $`(\mathrm{\Gamma }^E,\sigma )`$ as in Definition 2.2.5. Write
$$\mathrm{\Gamma }_{\mathrm{con}}^E=\mathrm{\Gamma }^E\mathrm{\Gamma }_{\mathrm{con}}^{\mathrm{perf}},$$
with the intersection taking place within $`\mathrm{\Gamma }^{\mathrm{perf}}`$. For $`r>0`$, we say that $`\mathrm{\Gamma }^E`$ *has enough $`r`$-units* if $`\mathrm{\Gamma }^E\mathrm{\Gamma }_r^{\mathrm{perf}}`$ contains units lifting all nonzero elements of $`E`$. We say that $`\mathrm{\Gamma }^E`$ *has enough units* (or more properly, the pair $`(\mathrm{\Gamma }^E,\sigma )`$ has enough units) if $`\mathrm{\Gamma }^E`$ has enough $`r`$-units for some $`r>0`$; this implies that $`\mathrm{\Gamma }_{\mathrm{con}}^E`$ is a henselian discrete valuation ring with maximal ideal $`\pi \mathrm{\Gamma }_{\mathrm{con}}^E`$ and residue field $`E`$. If $`\mathrm{\Gamma }^E`$ has enough units, then so does $`\mathrm{\Gamma }^E^{}`$ for any finite separable extension $`E^{}`$ of $`E`$ \[14, Lemma 2.2.12\].
### 2.3 Toroidal interpretation
The condition of having enough units is useful in the theory of slope filtrations, but is not convenient to check in practice. Fortunately, it has a more explicit interpretation in terms of certain “naïve” analogues of the functions $`v_n`$ and $`w_r`$, as in \[11, § 2\] or \[14, § 2.3\].
###### Definition 2.3.1.
Let $`L`$ be a lattice in $`^m`$ and let $`\lambda (^m)^{}`$ be an irrational linear functional. Let $`\mathrm{\Gamma }^\lambda `$ denote the $`\pi `$-adic completion of $`𝔬((L))_\lambda `$; its elements may be viewed as formal sums $`_{zL}c_z\{z\}`$ with $`w(c_z)\mathrm{}`$ as $`\lambda (z)\mathrm{}`$. Define the *naïve partial valuations* on $`\mathrm{\Gamma }^\lambda [\pi ^1]`$ by the formula
$$v_n^{\mathrm{naive}}\left(c_z\{z\}\right)=\mathrm{min}\{\lambda (z):zL,w(c_z)n\},$$
where the minimum is infinite if the set of candidate $`z`$’s is empty. These functions satisfy the identities
$`v_n^{\mathrm{naive}}(x+y)`$ $`\mathrm{min}\{v_n^{\mathrm{naive}}(x),v_n^{\mathrm{naive}}(y)\}`$ $`(x,y\mathrm{\Gamma }^\lambda [\pi ^1])`$
$`v_n^{\mathrm{naive}}(xy)`$ $`\underset{m}{\mathrm{min}}\{v_m^{\mathrm{naive}}(x)+v_{mn}^{\mathrm{naive}}(y)\}`$ $`(x,y\mathrm{\Gamma }^\lambda [\pi ^1]),`$
with equality in each case if the minimum is achieved only once. Define the *naïve levelwise topology* (or *naïve weak topology*) on $`\mathrm{\Gamma }^\lambda `$ by declaring that a sequence $`\{x_i\}`$ converges to zero if and only if for each $`n`$, $`v_n^{\mathrm{naive}}(x_i)\mathrm{}`$ as $`i\mathrm{}`$.
###### Definition 2.3.2.
For $`r>0`$ and $`n`$, write
$$v_{n,r}^{\mathrm{naive}}(x)=rv_n^{\mathrm{naive}}(x)+n;$$
extend the definition to $`r=0`$ by setting
$$v_{n,0}^{\mathrm{naive}}(x)=\{\begin{array}{cc}n\hfill & v_n^{\mathrm{naive}}(x)<\mathrm{}\hfill \\ \mathrm{}\hfill & v_n^{\mathrm{naive}}(x)=\mathrm{}.\hfill \end{array}$$
Let $`\mathrm{\Gamma }_r^{\mathrm{naive}}`$ be the set of $`x\mathrm{\Gamma }^\lambda `$ such that $`v_{n,r}^{\mathrm{naive}}(x)\mathrm{}`$ as $`n\mathrm{}`$. Define the map $`w_r^{\mathrm{naive}}`$ on $`\mathrm{\Gamma }_r^{\mathrm{naive}}`$ by
$$w_r^{\mathrm{naive}}(x)=\underset{n}{\mathrm{min}}\{v_{n,r}^{\mathrm{naive}}(x)\};$$
then $`w_r^{\mathrm{naive}}`$ is a valuation on $`\mathrm{\Gamma }_r^{\mathrm{naive}}[\pi ^1]`$, as in \[14, Lemma 2.1.7\]. Put
$$\mathrm{\Gamma }_{\mathrm{con}}^{\mathrm{naive}}=_{r>0}\mathrm{\Gamma }_r^{\mathrm{naive}}.$$
###### Remark 2.3.3.
The ring $`\mathrm{\Gamma }_r^{\mathrm{naive}}`$ is a principal ideal domain; this will follow from \[14, Proposition 2.6.5\] in conjunction with Definition 2.3.6 below.
We may view $`\mathrm{\Gamma }^\lambda `$ as an instance of the definition of $`\mathrm{\Gamma }^E`$ in the case $`E=k((L))_\lambda `$; this gives sense to the following result.
###### Proposition 2.3.4.
Let $`\sigma `$ be a Frobenius lift on $`\mathrm{\Gamma }^E=\mathrm{\Gamma }^\lambda `$ for $`E=k((L))_\lambda `$. Then for $`r>0`$, the following are equivalent.
1. $`\sigma `$ is continuous for the naïve levelwise topology (i.e., that topology coincides with the levelwise topology induced by $`\sigma `$), and for each $`zL`$ nonzero, $`\{z\}^\sigma /\{z\}^q`$ is a unit in $`\mathrm{\Gamma }_r^{\mathrm{naive}}`$.
2. For $`s(0,qr]`$, $`n0`$, and $`x\mathrm{\Gamma }^E`$,
$$\underset{jn}{\mathrm{min}}\{v_{j,s}(x)\}=\underset{jn}{\mathrm{min}}\{v_{j,s}^{\mathrm{naive}}(x)\}.$$
(2.3.4.1)
3. $`\mathrm{\Gamma }^E`$ has enough $`qr`$-units, and for each $`zL`$ nonzero, $`\{z\}`$ is a unit in $`\mathrm{\Gamma }_{qr}^E`$.
In particular, in each of these cases, for $`s(0,qr]`$, $`\mathrm{\Gamma }_s^{\mathrm{naive}}=\mathrm{\Gamma }_s^E`$ and $`w_s(x)=w_s^{\mathrm{naive}}(x)`$ for all $`x\mathrm{\Gamma }_s`$.
###### Proof.
Given (a), for $`s(0,qr]`$, we have
$$\underset{jn}{\mathrm{min}}\{v_{j,s}^{\mathrm{naive}}(x)\}=\underset{jn}{\mathrm{min}}\{v_{j,s/q}^{\mathrm{naive}}(x^\sigma )\}$$
(2.3.4.2)
for each $`n0`$ and each $`x\mathrm{\Gamma }`$, as in the proof of \[14, Lemma 2.3.3\]. We then obtain (b) as in the proof of \[14, Lemma 2.3.5\], from which (c) follows immediately.
Given (c), the equation (2.3.4.1) holds for $`x=\{z\}`$ for any $`zL`$, since the minima both occur for $`j=0`$. For $`x=c_z\{z\}`$ a finite sum, we have
$$\underset{jn}{\mathrm{min}}\{v_{j,s}^{\mathrm{naive}}(x)\}=\underset{jn}{\mathrm{min}}\{\underset{zL}{\mathrm{min}}\{v_{j,s}^{\mathrm{naive}}(c_z\{z\})\}\}$$
(2.3.4.3)
and so the left side of (2.3.4.1) is greater than or equal to the right side. On the other hand, if $`j`$ is taken to be the smallest value for which the outer minimum is achieved on the right side of (2.3.4.3), then the inner minimum is achieved by a unique value of $`z`$. Thus we actually may deduce equality in (2.3.4.1), again for $`x=c_z\{z\}`$ a finite sum. For general $`x`$, we may obtain the desired equality by replacing $`x`$ by a finite sum $`x^{}`$ such that $`xx^{}=y+z`$ for some $`y\mathrm{\Gamma }^E`$ with $`w(y)`$ greater than $`n`$, and some $`z\mathrm{\Gamma }_{qr}^E`$ with $`w_{qr}(z)`$ greater than either side of (2.3.4.1). Hence (c) implies (b).
Finally, note that (b) implies (a) straightforwardly. ∎
###### Remark 2.3.5.
Note that in Proposition 2.3.4, conditions (a) and (c) may be checked for $`z`$ running over a basis of $`L`$. Note also that Proposition 2.3.4 implies that for $`E=k((L))_\lambda `$, if $`\mathrm{\Gamma }^E`$ has enough units, then $`\mathrm{\Gamma }^E`$ is isomorphic to $`\mathrm{\Gamma }^\lambda `$.
###### Definition 2.3.6.
By the *standard extension* of $`\sigma _K`$ to $`\mathrm{\Gamma }^\lambda `$, we will mean the Frobenius lift $`\sigma `$ defined by
$$\underset{zL}{}c_z\{z\}\underset{zL}{}c_z^{\sigma _K}\{z\}^q.$$
(We will also refer to such a $`\sigma `$ as a *standard Frobenius lift*.) When equipped with a standard Frobenius lift, $`\mathrm{\Gamma }^\lambda `$ has enough $`r`$-units for every $`r>0`$; by Proposition 2.3.4, it follows that $`v_n(x)=v_n^{\mathrm{naive}}(x)`$ for all $`n`$ and all $`x\mathrm{\Gamma }^\lambda [\pi ^1]`$. Thus many of the results of \[14, § 2\], proved in terms of the Frobenius-based valuations, also apply verbatim to the naïve valuations.
###### Remark 2.3.7.
For applications to semistable reduction, one would also like to consider a similar situation in which the residue field $`k((L))_\lambda `$ is replaced by the completion of a finitely generated field extension of $`k`$ with respect to an arbitrary valuation of height (real rank) 1, at least in the case where the transcendence degree over $`k`$ is equal to 2. This would require a slightly more flexible set of foundations: one must work only with finitely generated $`k`$-subalgebras of the complete field, so that one has hope of having enough units. A more serious problem is how to perform Tsuzuki’s method (a/k/a Theorem 4.5.2) in this context.
### 2.4 Analytic rings
We now introduce “analytic rings”, citing into for their structural properties.
###### Definition 2.4.1.
Let $`E`$ be a nearly monomial field over $`k`$, or the completed perfect or algebraic closure thereof. In the first case, suppose that $`\mathrm{\Gamma }^E`$ has enough $`r_0`$-units for some $`r_0>0`$ (otherwise take $`r_0=\mathrm{}`$). Let $`I`$ be a subinterval of $`[0,r_0)`$ bounded away from $`r_0`$ (i.e., $`I`$ is a subinterval of $`[0,r]`$ for some $`r(0,r_0)`$). Let $`\mathrm{\Gamma }_I^E`$ denote the Fréchet completion of $`\mathrm{\Gamma }_{r_0}^E[\pi ^1]`$ under the valuations $`w_s`$ for $`sI`$; this ring is an integral domain \[14, Lemma 2.4.6\]. If $`I`$ is closed, then $`\mathrm{\Gamma }_I^E`$ is a principal ideal domain \[14, Proposition 2.6.9\]. Put
$$\mathrm{\Gamma }_{\mathrm{an},r}^E=\mathrm{\Gamma }_{(0,r]}^E;$$
this ring is a Bézout ring, i.e., a ring in which every finitely generated ideal is principal \[14, Theorem 2.9.6\]. Put $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}^E=_{r>0}\mathrm{\Gamma }_{\mathrm{an},r}^E`$; then $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}^E`$ is also a Bézout ring. The group of units in $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}^E`$ consists of the nonzero elements of $`\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$ \[14, Corollary 2.5.12\]. For $`E^{}`$ finite separable over $`E`$, $`E^{}=E^{\mathrm{perf}}`$, or $`E^{}=E^{\mathrm{alg}}`$, by \[14, Proposition 2.4.10\] one has
$$\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}^E^{}=\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}^E_{\mathrm{\Gamma }_{\mathrm{con}}^E}\mathrm{\Gamma }_{\mathrm{con}}^E^{}.$$
###### Remark 2.4.2.
It is likely that $`\mathrm{\Gamma }_I^E`$ is a Bézout ring for any $`I`$ as above. However, this statement is not verified in , and we will not need it anyway, so we withhold further comment on it.
###### Remark 2.4.3.
If $`E=k((t))`$ is a power series field, then the ring $`\mathrm{\Gamma }_I^E`$ is the ring of rigid analytic functions on the annulus $`w(t)I`$ in the $`t`$-plane. Thus our construction of fake annuli includes “true” one-dimensional rigid analytic annuli over $`K`$, and most of our results on fake annuli (like the $`p`$-adic local monodromy theorem) generalize extant theorems on true annuli. On the other hand, if $`E=k((L))_\lambda `$ and $`\mathrm{rank}(L)>1`$, then the ring $`\mathrm{\Gamma }_I^E`$ is trying to be the ring of rigid analytic functions on a subspace of the rigid affine plane in the variables $`\{z_1\},\mathrm{},\{z_m\}`$ for some basis $`z_1,\mathrm{},z_m`$ of $`L`$, consisting of points for which there exists $`rI`$ with $`w(\{z_i\})=r\lambda (z_i)`$ for $`i=1,\mathrm{},m`$. If $`I=[r,r]`$, then this space is an affinoid space in the sense of Berkovich, but otherwise it is not (because one can only cut out an analytic subspace of the form $`w(x)=\alpha w(y)`$ for $`\alpha `$ rational). Indeed, as far as we can tell, this space is not a $`p`$-adic analytic space in either of the Tate or Berkovich senses, despite the fact that it has a sensible ring of analytic functions; hence the use of the adjective “fake” in the phrase “fake annulus”, and the absence of an honest definition of that phrase.
###### Remark 2.4.4.
Since one can sensibly define rigid analytic annuli over arbitrary complete nonarchimedean fields, Remark 2.4.3 suggests the possibility of working with fake annuli over more general complete $`K`$. However, the algebraic issues here get more complicated, and we have not straightened them out to our satisfaction. For example, the analogue of the ring $`\mathrm{\Gamma }_r^{\mathrm{naive}}`$ fails to be a principal ideal ring if the valuation on $`K`$ is not discrete; it probably still has the Bézout property (that finitely generated ideals are principal), but we have not checked this. In any case, the formalism of completely breaks down when $`K`$ is not discretely valued, so an attempt here to avoid a discreteness hypothesis now would fail to improve upon our ultimate results; we have thus refrained from making such an attempt.
## 3 Frobenius and connection structures
We now introduce a notion which should be thought of as a $`p`$-adic differential equation with Frobenius structure on a fake annulus. We start with some notational conventions.
###### Convention 3.0.1.
Throughout this section, assume that $`E`$ is a monomial field and that $`\mathrm{\Gamma }^E`$ is equipped with a Frobenius lift such that $`\mathrm{\Gamma }^E`$ has enough $`r_0`$-units for some $`r_0>0`$; we view $`\mathrm{\Gamma }^E`$ as being equipped with a levelwise topology via the choice of a coordinate system. (This choice does not matter, as the topology can be characterized as the coarsest one under which the $`v_{n,r}`$ are continuous for all $`n`$ and all $`r(0,r_0)`$.) We suppress $`E`$ from the notation, writing $`\mathrm{\Gamma }`$ for $`\mathrm{\Gamma }^E`$, $`\mathrm{\Gamma }_{\mathrm{con}}`$ for $`\mathrm{\Gamma }_{\mathrm{con}}^E`$, and so on.
###### Convention 3.0.2.
When a valuation is applied to a matrix, it is defined to be the minimum value over entries of the matrix.
We also make a definition of convenience.
###### Definition 3.0.3.
Under Convention 3.0.1, we will mean by an *admissible ring* any one of the following topological rings.
* The ring $`\mathrm{\Gamma }`$ or $`\mathrm{\Gamma }[\pi ^1]`$ with its levelwise topology.
* The ring $`\mathrm{\Gamma }_r`$ or $`\mathrm{\Gamma }_r[\pi ^1]`$ with the Fréchet topology induced by $`w_s`$ for all $`s(0,r]`$, for $`r(0,r_0)`$. Note that for $`\mathrm{\Gamma }_r`$, this coincides with the topology induced by $`w_r`$ alone.
* The ring $`\mathrm{\Gamma }_{\mathrm{con}}`$ or $`\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$ topologized as the direct limit of the $`\mathrm{\Gamma }_r`$ or $`\mathrm{\Gamma }_r[\pi ^1]`$.
* The ring $`\mathrm{\Gamma }_I`$ with the Fréchet topology induced by the $`w_s`$ for $`sI`$, for some $`I[0,r_0)`$ bounded away from $`r_0`$.
* The ring $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$ topologized as the direct limit of the $`\mathrm{\Gamma }_{\mathrm{an},r}`$.
By a *nearly admissible ring*, we mean one of the above rings with $`E`$ replaced by a finite separable extension.
### 3.1 Differentials
###### Definition 3.1.1.
Let $`S/R`$ be an extension of topological rings. A *module of continuous differentials* is a topological $`S`$-module $`\mathrm{\Omega }_{S/R}^1`$ equipped with a continuous $`R`$-linear derivation $`d:S\mathrm{\Omega }_{S/R}^1`$, having the following universal property: for any topological $`S`$-module $`M`$ equipped with a continuous $`R`$-linear derivation $`D:SM`$, there exists a unique morphism $`\varphi :\mathrm{\Omega }_{S/R}^1M`$ of topological $`S`$-modules such that $`D=\varphi d`$. Since the definition is via a universal property, the module of continuous differentials is unique up to unique isomorphism if it exists at all.
Constructing modules of continuous differentials is tricky in general (imitating the usual construction of the module of Kähler differentials requires a topological tensor product, which is a rather delicate matter); however, for fake annuli, the construction is straightforward.
###### Definition 3.1.2.
By a *coordinate system* for $`\mathrm{\Gamma }`$, we will mean a lattice $`L`$ in some $`^m`$, an irrational linear functional $`\lambda (^m)^{}`$, an isomorphism $`\mathrm{\Gamma }^\lambda \mathrm{\Gamma }`$ carrying $`zL`$ to a unit in $`\mathrm{\Gamma }_{r_0}`$ for each nonzero $`zL`$, and a basis $`z_1,\mathrm{},z_m`$ of $`L`$. Such data always exist thanks to Proposition 2.3.4.
###### Definition 3.1.3.
For the remainder of this subsection, choose a coordinate system for $`\mathrm{\Gamma }`$, and let $`\mu _1,\mathrm{},\mu _mL^{}`$ denote the basis dual to $`z_1,\mathrm{},z_m`$. For $`\mu L^{}`$ and $`S`$ an admissible ring, let $`_\mu `$ be the continuous derivation on $`S`$ defined by the formula
$$_\mu \left(\underset{z}{}c_z\{z\}\right)=\underset{z}{}\mu (z)c_z\{z\};$$
note that $`\mu (z)`$, so it may sensibly be viewed as an element of $`𝔬`$. (The continuity of $`_\mu `$ is clear in terms of naïve partial valuations, so Proposition 2.3.4 implies continuity in terms of the Frobenius-based valuations.) For $`\mu =\mu _i`$, write $`_i`$ for $`_{\mu _i}`$. Define $`\mathrm{\Omega }_{S/𝔬}^1`$ to be the free $`S`$-module $`Sd\{z_1\}\mathrm{}Sd\{z_m\}`$, equipped with the natural induced topology and with the continuous $`𝔬`$-linear derivation $`d:S\mathrm{\Omega }_{S/𝔬}^1`$ given by
$$dx=\underset{i=1}{\overset{m}{}}_i(x)d\mathrm{log}\{z_i\}$$
(where $`d\mathrm{log}(f)=df/f`$).
###### Proposition 3.1.4.
The module $`\mathrm{\Omega }_{S/𝔬}^1`$ is a module of continuous derivations for $`S`$ over $`𝔬`$. In particular, the construction does not depend on the choice of the coordinate system.
###### Proof.
This is a straightforward consequence of the fact that one of $`𝔬[\{z_i\}^{\pm 1}]`$ or $`𝔬[\pi ^1,\{z_i\}^{\pm 1}]`$ is dense in $`S`$. ∎
###### Remark 3.1.5.
Note that Proposition 3.1.4 also allows us to construct the module of continuous differentials $`\mathrm{\Omega }_{S/𝔬}^1`$ when $`S`$ is only nearly admissible.
###### Remark 3.1.6.
For $`\mathrm{rank}(L)=1`$ and $`\mu L`$ nonzero, the image of $`_\mu `$ is closed; however, this fails for $`\mathrm{rank}(L)>1`$, because bounding $`\lambda (z)`$ does not in any way limit the $`p`$-adic divisibility of $`z`$ within $`L`$. This creates a striking difference between the milieux of true and fake annuli, from the point of view of the study of differential equations. On true annuli, one has the rich theory of $`p`$-adic differential equations due to Dwork-Robba, Christol-Mebkhout, et al. On fake annuli, much of that theory falls apart; the parts that survive are those that rest upon Frobenius structures, whose behavior differs little in the two settings.
### 3.2 $``$-modules
###### Definition 3.2.1.
Let $`S`$ be a nearly admissible ring. Define a *$``$-module* over $`S`$ to be a finite free $`S`$-module $`M`$ equipped with an integrable $`𝔬`$-linear connection $`:MM\mathrm{\Omega }_{S/𝔬}^1`$; here integrability means that, letting $`_1`$ denote the induced map
$$M_S\mathrm{\Omega }_{S/𝔬}^1\stackrel{1}{}M_S\mathrm{\Omega }_{S/𝔬}^1_S\mathrm{\Omega }_{S/𝔬}^1\stackrel{1}{}M_S_S^2\mathrm{\Omega }_{S/𝔬}^1,$$
the composite map $`_1`$ is zero. We say $`𝐯M`$ is *horizontal* if $`(𝐯)=0`$.
###### Definition 3.2.2.
Suppose $`S`$ is admissible, and fix a coordinate system for $`\mathrm{\Gamma }`$. Given a $``$-module $`M`$ over $`S`$, for $`\mu L^{}`$, define the map $`\mathrm{\Delta }_\mu :MM`$ by writing $`(𝐯)=_{i=1}^m𝐰_id\mathrm{log}\{z_i\}`$ with $`𝐰_iM`$, and setting
$$\mathrm{\Delta }_\mu (𝐯)=\underset{i=1}{\overset{m}{}}\mu (z_i)𝐰_i.$$
Also, write $`\mathrm{\Delta }_i`$ for $`\mathrm{\Delta }_{\mu _i}`$.
###### Remark 3.2.3.
The maps $`\mathrm{\Delta }_\mu `$ satisfy the following properties.
* The map $`L^{}\times MM`$ given by $`(\mu ,𝐯)\mathrm{\Delta }_\mu (𝐯)`$ is additive in each factor.
* For all $`\mu L^{}`$, $`sS`$, and $`𝐯M`$, we have the Leibniz rule
$$\mathrm{\Delta }_\mu (s𝐯)=s\mathrm{\Delta }_\mu (𝐯)+_\mu (s)𝐯.$$
* For $`\mu _1,\mu _2L^{}`$, the maps $`\mathrm{\Delta }_{\mu _1},\mathrm{\Delta }_{\mu _2}`$ commute.
Conversely, given a finite free $`S`$-module $`M`$ equipped with maps $`\mathrm{\Delta }_\mu :MM`$ for each $`\mu L^{}`$ satisfying these conditions, one can uniquely reconstruct a $``$-module structure on $`M`$ that gives rise to the $`\mathrm{\Delta }_\mu `$.
###### Remark 3.2.4.
Note that for true annuli (i.e., $`\mathrm{rank}(L)=1`$), the integrability restriction is empty because $`\mathrm{\Omega }_{S/𝔬}^1`$ has rank 1 over $`S`$. However, for fake annuli, integrability is a real restriction: even though the ring theory looks one-dimensional, the underlying “fake space” is really $`m`$-dimensional, inasmuch as $`\mathrm{\Omega }_{S/𝔬}^1`$ has rank $`m`$ over $`S`$.
###### Definition 3.2.5.
Let $`M`$ be a $``$-module over $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}^E^{}`$, for $`E^{}`$ a finite separable extension of $`E`$. We say $`M`$ is:
* *constant* if $`M`$ admits a horizontal basis (a basis of elements of the kernel of $``$);
* *quasi-constant* if there exists a finite separable extension $`E^{\prime \prime }`$ of $`E^{}`$ such that $`M\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}^{E^{\prime \prime }}`$ is constant;
* *unipotent* if $`M`$ admits an exhaustive filtration by saturated $``$-submodules, whose successive quotients are constant;
* *quasi-unipotent* if $`M`$ admits an exhaustive filtration by saturated $``$-submodules, whose successive quotients are quasi-constant.
We extend these definitions to $`(F,)`$-modules by applying them to the underlying $``$-module.
###### Remark 3.2.6.
If $`M`$ is quasi-unipotent, then there exists a finite separable extension $`E^{\prime \prime }`$ of $`E^{}`$ such that $`M\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}^{E^{\prime \prime }}`$ is unipotent. The converse is also true: if $`M\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}^{E^{\prime \prime }}`$ is unipotent, then the shortest unipotent filtration of $`M\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}^{E^{\prime \prime }}`$ is unique, so descends to $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}^E^{}`$.
### 3.3 Frobenius structures
###### Definition 3.3.1.
Let $`S`$ be a nearly admissible ring stable under $`\sigma `$; for instance, $`\mathrm{\Gamma },\mathrm{\Gamma }_{\mathrm{con}},\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$ are permitted, but $`\mathrm{\Gamma }_r`$ is not. Define an *$`F`$-module* over $`S`$ (with respect to $`\sigma `$) to be a finite free $`S`$-module $`M`$ equipped with a $`S`$-module homomorphism $`F:\sigma ^{}MM`$ which is an isogeny, i.e., which becomes an isomorphism upon tensoring with $`S[\pi ^1]`$. We typically view $`F`$ as a $`\sigma `$-linear map from $`M`$ to itself; we occasionally view $`M`$ as a left module for the twisted polynomial ring $`S\{\sigma \}`$. Given an $`F`$-module $`M`$ over $`S`$ and an integer $`c`$, which must be nonnegative if $`\pi ^1S`$, define the *twist* $`M(c)`$ of $`M`$ to be a copy of $`M`$ with the action of $`F`$ multiplied by $`\pi ^c`$.
###### Definition 3.3.2.
Let $`S`$ be a nearly admissible ring stable under $`\sigma `$. Define an *$`(F,)`$-module* over $`S`$ to be a finite free $`S`$-module $`M`$ equipped with the structures of both an $`F`$-module and a $``$-module, which are compatible in the sense of making the following diagram commute:
###### Remark 3.3.3.
We may regard $`\mathrm{\Omega }_{S/𝔬}^1`$ itself as an $`F`$-module via $`d\sigma `$, in which case the compatibility condition asserts that $`:MM\mathrm{\Omega }_{S/𝔬}^1`$ is an $`F`$-equivariant map. The fact that $`_\mu (f^\sigma )0(mod\pi )`$ for any $`f\mathrm{\Gamma }_{\mathrm{con}}`$ means that $`\mathrm{\Omega }_{\mathrm{\Gamma }_{\mathrm{con}}/𝔬}^1`$ is isomorphic as an $`F`$-module to $`N(1)`$, for some $`F`$-module $`N`$ over $`\mathrm{\Gamma }_{\mathrm{con}}`$. In the language of , this means that the generic HN slopes of $`\mathrm{\Omega }_{\mathrm{\Gamma }_{\mathrm{con}}/𝔬}^1`$ are positive \[14, Proposition 5.1.3\].
###### Definition 3.3.4.
For $`a`$ a positive integer, define an *$`F^a`$-module* or *$`(F^a,)`$-module* as an $`F`$-module or $`(F,)`$-module relative to $`\sigma ^a`$. Given an $`F`$-module $`M`$, viewed as a left $`S\{\sigma \}`$-module, define the $`F^a`$-module $`[a]_{}M`$ to be the left $`S\{\sigma ^a\}`$-module given by restriction along the inclusion $`S\{\sigma ^a\}S\{\sigma \}`$; in other words, replace the Frobenius action by its $`a`$-th power. Given an $`F^a`$-module $`N`$, viewed as a left $`S\{\sigma ^a\}`$-module, define the $`F`$-module $`[a]^{}M`$ to be the left $`S\{\sigma \}`$-module
$$[a]^{}M=S\{\sigma \}_{S\{\sigma ^a\}}M;$$
then the functors $`[a]^{}`$ and $`[a]_{}`$ are left and right adjoints of each other. See \[14, § 3.2\] for more on these operations.
### 3.4 Change of Frobenius
The category of $`(F,)`$-modules over $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$ relative to $`\sigma `$ turns out to be canonically independent of the choice of $`\sigma `$, by a Taylor series argument (as in \[25, § 3.4\]).
###### Convention 3.4.1.
Throughout this subsection, fix a coordinate system on $`\mathrm{\Gamma }`$. Given an $`m`$-tuple $`J=(j_1,\mathrm{},j_m)`$ of nonnegative integers, write $`J!=j_1!\mathrm{}j_m!`$; if $`U=(u_1,\mathrm{},u_m)`$, write $`U^J=u_1^{j_1}\mathrm{}u_m^{j_m}`$. Also, define the “falling factorials”
$$^{\underset{¯}{J}}=\underset{i=1}{\overset{m}{}}\underset{l=0}{\overset{j_i1}{}}(_il)$$
$$\mathrm{\Delta }^{\underset{¯}{J}}=\underset{i=1}{\overset{m}{}}\underset{l=0}{\overset{j_i1}{}}(\mathrm{\Delta }_il),$$
with the convention that $`^{\underset{¯}{0}}`$ and $`\mathrm{\Delta }^{\underset{¯}{0}}`$ are the respective identity maps. (The use of falling factorial notation is modeled on .)
###### Lemma 3.4.2.
Let $`M`$ be a $``$-module over $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$. Then for any $`r\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$, any $`𝐯M`$, and any $`m`$-tuple $`J`$ of nonnegative integers,
$$\frac{1}{J!}\mathrm{\Delta }^{\underset{¯}{J}}(r𝐯)=\underset{J_1+J_2=J}{}\left(\frac{1}{J_1!}^{\underset{¯}{J_1}}(r)\right)\left(\frac{1}{J_2!}\mathrm{\Delta }^{\underset{¯}{J_2}}(𝐯)\right).$$
###### Proof.
Since
$$_i(_i1)\mathrm{}(_ij+1)=\{z_i\}^j(\{z_i\}^1_i)^j$$
and similarly for $`\mathrm{\Delta }_i`$, this amounts to a straightforward application of the Leibniz rule. ∎
###### Lemma 3.4.3.
For any $`u_1,\mathrm{},u_m\mathrm{\Gamma }_{\mathrm{con}}`$ with $`w(u_i)>0`$ for $`i=1,\mathrm{},m`$, and any $`x\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$, the series
$$f(x)=\underset{j_1,\mathrm{},j_m=0}{\overset{\mathrm{}}{}}\frac{1}{J!}U^J^{\underset{¯}{J}}(x)$$
converges in $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$, and the map $`xf(x)`$ is a continuous ring homomorphism sending $`\{z_i\}`$ to $`u_i`$.
###### Proof.
Pick $`r>0`$ such that $`u_1,\mathrm{},u_m,x\mathrm{\Gamma }_{\mathrm{an},r}`$ and $`w_r(u_i)>0`$ for $`i=1,\mathrm{},m`$. Write $`x=_{zL}c_z\{z\}`$; note that for each $`J`$,
$$\frac{1}{J!}^{\underset{¯}{J}}(x)=\underset{zL}{}\left(\underset{i=1}{\overset{m}{}}\left(\genfrac{}{}{0pt}{}{\mu _i(z)}{j_i}\right)\right)c_z\{z\},$$
so that $`w_s(^{\underset{¯}{J}}(x)/J!)w_s(x)`$ for $`s(0,r]`$. This yields the desired convergence, as well as continuity of the map $`xf(x)`$. Moreover, $`f`$ is a ring homomorphism on $`𝔬[\{z_1\},\mathrm{},\{z_m\}]`$ by Lemma 3.4.2, so must be a ring homomorphism on $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$ by continuity; the fact that it sends $`\{z_i\}`$ to $`u_i`$ is apparent from the formula. ∎
###### Lemma 3.4.4.
Let $`M`$ be a $``$-module over $`\mathrm{\Gamma }_{\mathrm{an},r}`$ for some $`r>0`$. Suppose that for some positive integer $`h`$, $`M`$ admits a basis $`𝐞_1,\mathrm{},𝐞_n`$ such that the $`n\times n`$ matrices $`N_1,\mathrm{},N_m`$ defined by $`\mathrm{\Delta }_i(𝐞_l)=_j(N_i)_{jl}𝐞_j`$ satisfy $`w_r(N_i)>w((p^h)!)`$ for $`i=1,\mathrm{},m`$. For $`J`$ an $`m`$-tuple of nonnegative integers, define the $`n\times n`$ matrix $`N_J`$ by
$$\mathrm{\Delta }^{\underset{¯}{J}}(𝐞_l)=\underset{j}{}(N_J)_{jl}𝐞_j.$$
Then
$$w_r(N_J)w(J!)w(p)(j_1+\mathrm{}+j_m)/(p^h(p1)).$$
###### Proof.
The condition that $`w_r(N_i)>w((p^h)!)`$ means that for any $`a`$ and any $`b\{0,\mathrm{},p^h1\}`$, if we write
$`𝐯`$ $`={\displaystyle \underset{j=1}{\overset{m}{}}}x_j𝐞_j`$ $`(x_j\mathrm{\Gamma }_{\mathrm{an},r})`$
$`(\mathrm{\Delta }_iap^h)(\mathrm{\Delta }_iap^h1)\mathrm{}(\mathrm{\Delta }_iap^hb)𝐯`$ $`={\displaystyle \underset{j=1}{\overset{m}{}}}y_{ijab}𝐞_j`$ $`(y_{ijab}\mathrm{\Gamma }_{\mathrm{an},r}),`$
then $`\mathrm{min}_j\{w_r(y_{ijab})\}\mathrm{min}_j\{w_r(x_j)\}+w(b!)`$ (i.e., the same bound as for the trivial connection with $`𝐞_1,\mathrm{},𝐞_n`$ horizontal). This gives the bound
$`w_r(N_J)`$ $`w(J!)+{\displaystyle \underset{i=1}{\overset{m}{}}}\left(w(j_i!)+j_i/p^hw((p^h)!)+w((j_ip^hj_i/p^h)!)\right)`$
$`w(J!){\displaystyle \underset{i=1}{\overset{m}{}}}w(p)j_i/(p^h(p1))`$
using the fact that $`w(j_i!)=_{g=1}^{\mathrm{}}w(p)j_i/p^g`$. This yields the claim. ∎
###### Lemma 3.4.5.
Let $`M`$ be an $`(F,)`$-module over $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$ or over $`\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$, and let $`𝐞_1,\mathrm{},𝐞_n`$ be a basis of $`M`$. For each nonnegative integer $`g`$, define the $`n\times n`$ matrices $`N_{g,1},\mathrm{},N_{g,m}`$ by $`\mathrm{\Delta }_i(F^g𝐞_l)=_j(N_{g,i})_{jl}(F^g𝐞_j)`$. Then there exist $`r_1(0,r_0)`$ and $`c>0`$ such that for each nonnegative integer $`g`$ and for each of $`i=1,\mathrm{},m`$, $`N_{g,i}`$ has entries in $`\mathrm{\Gamma }_{\mathrm{an},r_1q^g}`$ and
$$w_{rq^g}(N_{g,i})gc(r[r_1/q,r_1]).$$
Moreover, if $`M`$ is defined over $`\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$, we can also ensure that $`w(N_{g,i})gc`$.
###### Proof.
Define $`a_{hi}\mathrm{\Gamma }_{\mathrm{con}}`$ by the formula
$$_i(x^\sigma )=\underset{h=1}{\overset{m}{}}a_{hi}(_hx)^\sigma (x\mathrm{\Gamma }_{\mathrm{con}});$$
then $`w(a_{hi})1`$ as in Remark 3.3.3. In particular, we can choose $`r_1(0,r_0)`$ as in Proposition 2.3.4 such that for $`i=1,\mathrm{},m`$, $`a_i\mathrm{\Gamma }_{r_1}`$, $`w_{r_1}(a_{hi})1`$, and $`N_{0,i}`$ has entries in $`\mathrm{\Gamma }_{\mathrm{an},r_1}`$. Then the formula
$$N_{g+1,i}=\underset{h=1}{\overset{m}{}}a_{hi}N_{g,h}^\sigma $$
yields the claim for any $`c`$ with $`\mathrm{min}_i\{\mathrm{min}_{r[r_1/q,r_1]}\{w_r(N_{0,i})\}\}c`$ and (in case $`M`$ is defined over $`\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$) $`\mathrm{min}_i\{w(N_{0,i})\}c`$. ∎
###### Lemma 3.4.6.
Let $`M`$ be an $`(F,)`$-module over $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$ (resp. over $`\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$). Then for any $`u_1,\mathrm{},u_m\mathrm{\Gamma }_{\mathrm{con}}`$ with $`w(u_i)>0`$ for $`i=1,\mathrm{},m`$, and any $`𝐯M`$, the series
$$f(𝐯)=\underset{j_1,\mathrm{},j_m=0}{\overset{\mathrm{}}{}}\frac{1}{J!}U^J\mathrm{\Delta }^{\underset{¯}{J}}(𝐯)$$
converges for the natural topology of $`M`$, and the map $`𝐯f(𝐯)`$ is semilinear for the map defined by Lemma 3.4.3.
###### Proof.
Pick a basis $`𝐞_1,\mathrm{},𝐞_n`$ of $`M`$; for each nonnegative integer $`g`$, define the $`n\times n`$ matrices $`N_{g,1},\mathrm{},N_{g,m}`$ by $`\mathrm{\Delta }_i(F^g𝐞_l)=_j(N_{g,i})_{jl}(F^g𝐞_j)`$. By Lemma 3.4.5, we can choose $`r_1(0,r_0)`$ such that for some $`c>0`$, $`w_{rq^g}(N_{g,i})gc`$ for all nonnegative integers $`g`$ and all $`r[r_1/q,r_1]`$; moreover, if we are working over $`\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$, we can ensure that $`w(N_{g,i})gc`$.
Now choose a positive integer $`h`$ with $`w(p)/(p^h(p1))<1/2`$. Then by the previous paragraph, for each sufficiently small $`r>0`$, there exists a basis $`𝐯_1,\mathrm{},𝐯_n`$ of $`M`$ (depending on $`r`$) on which each $`\mathrm{\Delta }_i`$ acts via a matrix $`N_i`$ with $`w_r(N_i)>w((p^h)!)`$. By Lemma 3.4.4, the matrix $`N_J`$ defined by
$$\mathrm{\Delta }^{\underset{¯}{J}}(𝐯_l)=\underset{j}{}(N_J)_{jl}𝐯_j$$
satisfies $`w_r(N_J)w(J!)(j_1+\mathrm{}+j_m)/2`$.
On the other hand, since $`w(u_i)1`$ for $`i=1,\mathrm{},m`$, we have that $`w_r(u_i)>1/2`$ for $`r`$ sufficiently small. We conclude that for each sufficiently small $`r>0`$, there exists a basis $`𝐯_1,\mathrm{},𝐯_n`$ such that the series defining each of $`f(𝐯_1),\mathrm{},f(𝐯_n)`$ converges under $`w_r`$. By Lemma 3.4.2 and Lemma 3.4.3, for each $`\mathrm{\Gamma }_{\mathrm{an},r}`$-linear combination $`𝐯`$ of $`𝐯_1,\mathrm{},𝐯_n`$, the series defining $`f(𝐯)`$ converges under $`w_r`$. By the same token, in case $`M`$ is defined over $`\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$, for each $`\mathrm{\Gamma }_r[\pi ^1]`$-linear combination $`𝐯`$ of $`𝐯_1,\mathrm{},𝐯_n`$, the series defining $`f(𝐯)`$ converges under $`w`$. This yields the desired convergence of $`f`$; again, the semilinearity follows from Lemma 3.4.2 and Lemma 3.4.3. ∎
###### Proposition 3.4.7.
Let $`\sigma _1`$ and $`\sigma _2`$ be Frobenius lifts on $`\mathrm{\Gamma }`$ such that $`\mathrm{\Gamma }`$ has enough units with respect to each of $`\sigma _1`$ and $`\sigma _2`$, and for each $`zL`$ nonzero, $`\{z\}`$ is a unit in $`\mathrm{\Gamma }_{\mathrm{con}}`$ under both definitions. (By Proposition 2.3.4, it is equivalent to require that the definitions of $`\mathrm{\Gamma }_{\mathrm{con}}`$ with respect to $`\sigma _1`$ and to $`\sigma _2`$ coincide.) Then there is a canonical equivalence of categories between $`(F,)`$-modules over $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$ (resp. over $`\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$) relative to $`\sigma _1`$ and relative to $`\sigma _2`$, acting as the identity on the underlying $``$-modules.
###### Proof.
Put $`u_i=\{z_i\}^{\sigma _2}/\{z_i\}^{\sigma _1}1`$. Let $`M`$ be a $``$-module admitting a compatible Frobenius structure $`F_1`$ relative to $`\sigma _1`$. For $`𝐯M`$, define
$$F_2(𝐯)=\underset{j_1,\mathrm{},j_m=0}{\overset{\mathrm{}}{}}\frac{1}{J!}U^JF_1(\mathrm{\Delta }^{\underset{¯}{J}}(𝐯));$$
this series converges thanks to Lemma 3.4.6. Moreover, the result is $`\sigma _2`$-linear thanks to Lemma 3.4.2. ∎
###### Remark 3.4.8.
By tweaking the proof of Proposition 3.4.7, one can also obtain the analogous independence from the choice of $`\sigma `$ for the category of $`(F,)`$-modules over $`\mathrm{\Gamma }_{\mathrm{con}}`$. We will not use this result explicitly, though a related construction will occur in Subsection 4.2.
## 4 Unit-root $`(F,)`$-modules (after Tsuzuki)
In this section, we give the generalization to fake annuli of Tsuzuki’s unit-root local monodromy theorem , variant proofs of which are given by Christol and in the author’s unpublished dissertation . Our argument here draws on elements of all of these; its specialization to the case of true annuli constitutes a novel (if only slightly so) exposition of Tsuzuki’s original result.
###### Convention 4.0.1.
Throughout this section, let $`E`$ denote a nearly monomial field over $`k`$, viewed in a fixed fashion as a finite separable extension of a monomial field over $`k`$. We assume that any Frobenius lift $`\sigma `$ considered on $`\mathrm{\Gamma }=\mathrm{\Gamma }^E`$ is chosen so that $`\mathrm{\Gamma }`$ has enough units. In particular, $`\mathrm{\Gamma }=\mathrm{\Gamma }^E`$ and $`\mathrm{\Gamma }_{\mathrm{con}}=\mathrm{\Gamma }_{\mathrm{con}}^E`$ are nearly admissible in the sense of Definition 3.0.3.
### 4.1 Unit-root $`F`$-modules
###### Definition 4.1.1.
We say an $`F`$-module $`M`$ over $`\mathrm{\Gamma }^E`$ or $`\mathrm{\Gamma }_{\mathrm{con}}^E`$, with respect to some Frobenius lift $`\sigma `$, is *unit-root* (or *étale*) if the map $`F:\sigma ^{}MM`$ is an isomorphism (not just an isogeny). We say an $`(F,)`$-module over $`\mathrm{\Gamma }^E`$ or $`\mathrm{\Gamma }_{\mathrm{con}}`$ is unit-root if its underlying $`F`$-module is unit-root.
We will frequently calculate on such modules in terms of bases, so it is worth making the relevant equations explicit.
###### Remark 4.1.2.
Assume that $`E`$ is a monomial field, and fix a coordinate system for $`\mathrm{\Gamma }`$. Let $`M`$ be a $``$-module over $`\mathrm{\Gamma }`$ or $`\mathrm{\Gamma }_{\mathrm{con}}`$ with basis $`𝐞_1,\mathrm{},𝐞_n`$. Given $`\mu L^{}`$, define the $`n\times n`$ matrix $`N_\mu `$ by $`\mathrm{\Delta }_\mu (𝐞_l)=_j(N_\mu )_{jl}𝐞_j`$; if we identify $`𝐯=c_1𝐞_1+\mathrm{}+c_n𝐞_nM`$ with the column vector with entries $`c_1,\mathrm{},c_n`$, then we have
$$\mathrm{\Delta }_\mu (𝐯)=N_\mu 𝐯+_\mu (𝐯).$$
Given an $`F`$-module with the same basis $`𝐞_1,\mathrm{},𝐞_n`$, define the $`n\times n`$ matrix $`A`$ by $`F(𝐞_l)=_jA_{jl}𝐞_j`$; then with the same identification of $`𝐯`$ with a column vector, we have
$$F(𝐯)=A𝐯^\sigma .$$
In case the Frobenius lift $`\sigma `$ is standard, the compatibility of Frobenius and connection structures is equivalent to the equations
$$N_\mu A+_\mu (A)=qAN_\mu ^\sigma (\mu L^{});$$
of course it is only necessary to check this on a basis of $`L^{}`$.
###### Remark 4.1.3.
It is also worth writing out how the equations in Remark 4.1.2 transform under change of basis. First, if $`U`$ is an invertible $`n\times n`$ matrix, then
$$N_\mu A+_\mu (A)=0(U^1N_\mu U+U^1_\mu (U))(U^1A)+_\mu (U^1A)=0.$$
Second, in case $`\sigma `$ is standard, the equations
$$N_\mu A+_\mu (A)=qAN_\mu ^\sigma \text{and}N_\mu ^{}A^{}+_\mu (A^{})=qA^{}(N_\mu ^{})^\sigma $$
are equivalent for
$`N_\mu ^{}`$ $`=U^1N_\mu U+U^1_\mu (U)`$
$`A^{}`$ $`=U^1AU^\sigma .`$
### 4.2 Unit-root $`F`$-modules and Galois representations
We now consider unit-root $`F`$-modules over $`\mathrm{\Gamma }`$, obtaining the usual Fontaine-style setup.
###### Lemma 4.2.1.
Let $`\mathrm{}`$ be a separably closed field of characteristic $`p>0`$, and let $`\tau `$ denote the $`q`$-power Frobenius on $`\mathrm{}`$. Let $`A`$ be an invertible $`n\times n`$ matrix over $`\mathrm{}`$.
1. There exists an invertible $`n\times n`$ matrix $`U`$ over $`\mathrm{}`$ such that $`U^1AU^\tau `$ is the identity matrix.
2. For any $`1\times n`$ column vector $`𝐯`$ over $`\mathrm{}`$, there are exactly $`q^n`$ distinct $`1\times n`$ column vectors $`𝐰`$ over $`\mathrm{}`$ for which $`A𝐰^\tau 𝐰=𝐯`$.
###### Proof.
Part (a) is \[9, Proposition 1.1\]; part (b) is an easy corollary of (a). ∎
We next introduce a “big ring” over which unit-root $`F`$-modules over $`\mathrm{\Gamma }`$ can be trivialized.
###### Definition 4.2.2.
Let $`\stackrel{~}{\mathrm{\Gamma }}`$ be the $`\pi `$-adic completion of the maximal unramified extension of $`\mathrm{\Gamma }`$; then any Frobenius lift $`\sigma `$ on $`\mathrm{\Gamma }`$ extends uniquely to $`\stackrel{~}{\mathrm{\Gamma }}`$, and the derivation $`d`$ extends uniquely to a derivation $`d:\stackrel{~}{\mathrm{\Gamma }}(\mathrm{\Omega }_{\mathrm{\Gamma }/𝔬}^1_\mathrm{\Gamma }\stackrel{~}{\mathrm{\Gamma }})`$. Likewise, any $``$-module $`M`$ over $`\mathrm{\Gamma }`$ induces a connection $`:(M_\mathrm{\Gamma }\stackrel{~}{\mathrm{\Gamma }})(M_\mathrm{\Gamma }\mathrm{\Omega }_{\mathrm{\Gamma }/𝔬}^1_\mathrm{\Gamma }\stackrel{~}{\mathrm{\Gamma }})`$. Let $`\stackrel{~}{𝔬}_q`$ be the fixed subring of $`\stackrel{~}{\mathrm{\Gamma }}`$ under $`\sigma `$; this is a complete discrete valuation ring with residue field $`𝔽_q`$ and maximal ideal generated by $`\pi `$.
###### Proposition 4.2.3.
Let $`M`$ be a unit-root $`F`$-module over $`\stackrel{~}{\mathrm{\Gamma }}`$. Then $`M`$ admits an $`F`$-invariant basis.
###### Proof.
Applying Lemma 4.2.1(a) produces a basis which is fixed modulo $`\pi `$. Given a basis fixed modulo $`\pi ^n`$, correcting it to a basis fixed modulo $`\pi ^{n+1}`$ amounts to solving a set of vector equations of the form of Lemma 4.2.1(b). The resulting sequence of bases converges to the desired $`F`$-invariant basis. ∎
###### Definition 4.2.4.
Assume that $`𝔽_qk`$. Given a unit-root $`F`$-module $`M`$ over $`\mathrm{\Gamma }`$, let $`D_\mathrm{\Gamma }(M)`$ denote the set of $`F`$-invariant elements of $`M_\mathrm{\Gamma }\stackrel{~}{\mathrm{\Gamma }}`$; then $`D_\mathrm{\Gamma }(M)`$ is a finite free $`\stackrel{~}{𝔬}_q`$-module equipped with a continuous action of $`G=\mathrm{Gal}(E^{\mathrm{sep}}/E)`$. By Proposition 4.2.3, the natural map $`D_\mathrm{\Gamma }(M)_{\stackrel{~}{𝔬}_q}\stackrel{~}{\mathrm{\Gamma }}M_\mathrm{\Gamma }\stackrel{~}{\mathrm{\Gamma }}`$ is an isomorphism. Conversely, given a finite free $`\stackrel{~}{𝔬}_q`$-module $`N`$ equipped with a continuous action of $`G`$, let $`V(N)`$ denote the set of $`G`$-invariant elements of $`N_{\stackrel{~}{𝔬}_q}\stackrel{~}{\mathrm{\Gamma }}`$; by Galois descent, the natural map $`V(N)_\mathrm{\Gamma }\stackrel{~}{\mathrm{\Gamma }}N_{\stackrel{~}{𝔬}_q}\stackrel{~}{\mathrm{\Gamma }}`$ is an isomorphism. The functors $`D_\mathrm{\Gamma }`$ and $`V`$ thus exhibit equivalences of categories between unit-root $`F`$-modules over $`\mathrm{\Gamma }`$ and finite free $`\stackrel{~}{𝔬}_q`$-modules equipped with continuous $`G`$-action.
So far in this subsection, we have considered only unit-root $`F`$-modules over $`\stackrel{~}{\mathrm{\Gamma }}`$, without connection structure. The reason is that the connection does not really add any extra structure in this case.
###### Proposition 4.2.5.
Let $`M`$ be a unit-root $`F`$-module over $`\mathrm{\Gamma }`$ (resp. over $`\stackrel{~}{\mathrm{\Gamma }}`$). Then there is a unique integrable connection $`:MM\mathrm{\Omega }^1`$ compatible with $`F`$.
###### Proof.
We first check existence and uniqueness for the connection, without worrying about integrability. Let $`_0:MM\mathrm{\Omega }^1`$ be any connection (not necessarily integrable). Then a map $`:MM\mathrm{\Omega }^1`$ is a connection if and only if $`_0`$ is a $`\mathrm{\Gamma }`$-linear map from $`M`$ to $`M\mathrm{\Omega }^1`$, i.e., if it corresponds to an element of $`M^{}M\mathrm{\Omega }^1`$. Moreover, $``$ is compatible with $`F`$ if and only if
$$(_0)F(Fd\sigma )(_0)=(Fd\sigma )_0_0F;$$
in other words, we can write down a particular $`𝐰M^{}M\mathrm{\Omega }^1`$ such that $``$ is $`F`$-equivariant if and only if $`_0`$ corresponds to an element $`𝐯M^{}M\mathrm{\Omega }^1`$ with $`𝐯F𝐯=𝐰`$. By Remark 3.3.3, $`M^{}M\mathrm{\Omega }^1`$ can be written as a twist $`N(1)`$ for some $`F`$-module $`N`$ over $`\mathrm{\Gamma }`$ (resp. over $`\stackrel{~}{\mathrm{\Gamma }}`$); hence the series $`𝐰+F𝐰+F^2𝐰+\mathrm{}`$ converges in $`M^{}M\mathrm{\Omega }^1`$ to the unique solution of $`𝐯F𝐯=𝐰`$.
It remains to prove that the unique connection $``$ compatible with $`F`$ is in fact integrable. It is enough to check this over $`\stackrel{~}{\mathrm{\Gamma }}`$; moreover, it is enough to exhibit a single integrable connection compatible with $`F`$, as this must then coincide with the connection constructed above. To do this, we apply Proposition 4.2.3 to produce an $`F`$-invariant basis $`𝐞_1,\mathrm{},𝐞_n`$ of $`M`$, then set
$$(c_1𝐞_1+\mathrm{}+c_n𝐞_n)=𝐞_1dc_1+\mathrm{}+𝐞_ndc_n;$$
this map is easily seen to be an integrable connection compatible with $`F`$. ∎
###### Remark 4.2.6.
For true annuli, the construction of Definition 4.2.4 is due to Fontaine \[7, 1.2\]. In general, one consequence of the construction is that the categories of unit-root $`F`$-modules over $`\mathrm{\Gamma }`$ relative to two different Frobenius lifts are canonically equivalent, since the category of finite free $`\stackrel{~}{𝔬}_q`$-modules equipped with continuous $`G`$-action does not depend on the choice of the Frobenius lift. In fact, one may even change the choice of the underlying Frobenius lift $`\sigma _K`$, as long as it does not change what $`\stackrel{~}{𝔬}_q`$ is. Note that the same is true of unit-root $`(F,)`$-modules; more precisely, if a $``$-module $`M`$ over $`\mathrm{\Gamma }`$ admits a unit-root Frobenius structure for one Frobenius lift, it admits a unit-root Frobenius structure for any Frobenius lift. That is because the connection on $`M`$ can be recovered from $`D_\mathrm{\Gamma }(M)`$, by specifying that elements of $`D_\mathrm{\Gamma }(M)`$ are horizontal.
### 4.3 Positioning Frobenius
It will be useful to prove a positioning lemma for elements of $`k((L))_\lambda `$.
###### Lemma 4.3.1.
Assume that the field $`k`$ is algebraically closed. Suppose $`xk((L))_\lambda `$ cannot be written as $`a^qa`$ for any $`ak((L))_\lambda `$. Then there exists $`c>0`$ such that for any nonnegative integer $`i`$ and any $`yk((L))_\lambda `$ with $`xy+y^qk((q^iL))_\lambda `$, we have $`v_\lambda (xy+y^q)cq^i`$.
###### Proof.
Clearly there is no harm in replacing $`x`$ by $`xy_0+y_0^q`$ for any $`y_0k((L))_\lambda `$. In particular, write $`x=_{zL}c_z\{z\}`$, and let $`x_{},x_0,x_+`$ be the sum of $`c_z\{z\}`$ over those $`z`$ with $`\lambda (z)`$ negative, zero, positive, respectively. Since $`k`$ is algebraically closed, we have $`x_0=yy^q`$ for some $`yk`$. Since $`v_\lambda (x_+)>0`$, we have $`x_+=yy^q`$ for $`y=x_++x_+^q+x_+^{q^2}+\mathrm{}`$. We may thus reduce to the case $`x=x_{}`$; in particular, $`x`$ has finite support.
For $`zL`$ nonzero, let $`i(z)`$ denote the largest integer $`i`$ such that $`z/q^iL`$. Set
$$y_1=\underset{zL,\lambda (z)<0}{}\underset{i=1}{\overset{i(z)}{}}(c_z\{z\})^{q^i},$$
so that $`x_1=x+y_1y_1^q`$ is supported on $`LqL`$. We cannot have $`x_1=0`$, or else we could have written $`x=a^qa`$ for some $`ak((L))_\lambda `$. There must thus be a smallest (under $`\lambda `$) element $`z`$ of the support of $`x_1`$. For any nonnegative integer $`i`$ and any $`yk((L))_\lambda `$ with $`xy+y^qk((q^iL))_\lambda `$, the support of $`xy+y^q`$ must contain $`q^{i+j}z`$ for some nonnegative integer $`j`$, and so $`v_\lambda (xy+y^q)\lambda (z)q^i`$, as desired. ∎
### 4.4 Successive decimation
We now give a version of Tsuzuki’s construction for solving $`p`$-adic differential equations.
###### Convention 4.4.1.
Throughout this subsection, assume that $`k`$ is algebraically closed and that $`E`$ is a monomial field over $`k`$, and fix a coordinate system on $`\mathrm{\Gamma }`$. Also assume $`\sigma `$ is a standard Frobenius lift; we may then safely confound the naïve and Frobenius-based partial valuations.
###### Definition 4.4.2.
For $`\mu L^{}`$ nonzero, write $`L_\mu `$ for the sublattice of $`zL`$ for which $`\mu (z)p`$.
###### Lemma 4.4.3.
Suppose that $`A`$ is an invertible $`n\times n`$ matrix over $`\mathrm{\Gamma }`$ with $`w(AI_n)>0`$, that $`N_\mu `$ is an $`n\times n`$ matrix over $`\mathrm{\Gamma }`$ supported on $`L_\mu `$, and that $`N_\mu A+_\mu (A)=qAN_\mu ^\sigma `$. Then $`A`$ is supported on $`L_\mu `$.
###### Proof.
Suppose the contrary; write $`A=B+C`$ with $`B`$ supported on $`L_\mu `$ and $`h=w(C)`$ as large as possible, so in particular $`h>0`$. Write $`C=_{zL}C_z\{z\}`$. Since $`h`$ is as large as possible, there exists $`zLL_\mu `$ such that $`w(C_z)=h`$; the coefficient of $`\{z\}`$ in $`_\mu (C)`$ then also has valuation $`h`$. However, in the equality
$$_\mu (C)=(qBN_\mu ^\sigma N_\mu B_\mu (B))+(qCN_\mu ^\sigma N_\mu C),$$
the first term in parentheses is supported on $`L_\mu `$ while the second term has valuation strictly greater than $`h`$ (since $`w(AI_n)>0`$ forces $`w(N_\mu )>0`$). This contradiction yields the desired result. ∎
###### Lemma 4.4.4.
Pick $`\mu L^{}`$ nonzero. Given $`r>0`$, let $`N_\mu `$ be an $`n\times n`$ matrix over $`\mathrm{\Gamma }_r`$ such that $`w(N_\mu )>0`$ and $`w_r(N_\mu )>0`$. Then for any $`s(0,r)`$, there exists an invertible $`n\times n`$ matrix $`U`$ over $`\mathrm{\Gamma }_s`$ such that $`w_s(UI_n)>0`$, $`w(UI_n)w(N_\mu )`$, and $`U^1N_\mu U+U^1_\mu (U)`$ is supported on $`L_\mu `$.
###### Proof.
Define a sequence $`U_0,U_1,\mathrm{}`$ of invertible matrices over $`\mathrm{\Gamma }_r`$ with $`w(U_jI_n)w(N_\mu )`$ and $`w_r(U_jI_n)w_r(N_\mu )`$, as follows. Start with $`U_0=I_n`$. Given $`U_j`$, put $`N_j=U_j^1N_\mu U_j+U_j^1_\mu (U_j)`$. Write $`N_j=_{zL}N_{j,z}\{z\}`$, let $`X_j`$ be the sum of $`\mu (z)^1N_{j,z}\{z\}`$ over all $`zLL_\mu `$, and put $`U_{j+1}=U_j(I_nX_j)`$.
For $`j>0`$, if $`w(U_jI_n)<\mathrm{}`$, one sees that $`w(U_{j+1}I_n)>w(U_jI_n)`$. Hence the $`U_j`$ converge $`\pi `$-adically; since they all satisfy $`w_r(U_jI_n)w_r(N_\mu )>0`$, the $`U_j`$ converge under $`w_s`$ to a limit $`U`$ satisfying $`w_s(UI_n)w_s(N_\mu )`$. In particular, $`U`$ is invertible and $`U^1N_\mu U+U^1_\mu (U)`$ is supported on $`L_\mu `$. (Compare \[24, Lemma 6.1.4\] and \[10, Lemma 5.1.3\].) ∎
###### Lemma 4.4.5.
Let $`N_1,\mathrm{},N_m`$ be $`n\times n`$ matrices over $`\mathrm{\Gamma }_r`$, such that $`w(N_i)>0`$ and $`w_r(N_i)>0`$ for $`i=1,\mathrm{},m`$. Suppose that there exists an invertible matrix $`A`$ over $`\mathrm{\Gamma }`$ with $`w(AI_n)>0`$, such that
$$N_iA+_i(A)=qAN_i^\sigma (i=1,\mathrm{},m).$$
Then for any $`s(0,r)`$, there exists an invertible matrix $`U`$ over $`\mathrm{\Gamma }_s`$, such that $`w(UI_n)\mathrm{min}_i\{w(N_i)\}`$, $`w_s(UI_n)\mathrm{min}_i\{w_s(N_i)\}`$, and $`U^1AU^\sigma `$ is supported on $`pL`$.
###### Proof.
For $`i=0,\mathrm{},m`$, let $`S_i`$ be the sublattice of $`zL`$ such that $`\mu _j(z)p`$ for $`j=1,\mathrm{},i`$. Pick $`s_1,\mathrm{},s_{m1}`$ with $`0<s<s_{m1}<\mathrm{}<s_1<r`$, and put $`s_0=r`$ and $`s_m=s`$. Put $`U_0=I_n`$. Given $`U_i`$ invertible over $`\mathrm{\Gamma }_{s_i}`$ such that $`A_i=U_i^1AU_i^\sigma `$ is supported on $`S_i`$, note that $`M_i=U_i^1N_iU_i+U_i^1_i(U_i)`$ satisfies the equation $`M_iA_i+_i(A_i)=qA_iM_i^\sigma `$. We may then argue (as in the proof of Proposition 4.2.5) that $`M_i`$ is congruent to a matrix supported on $`S_i`$ modulo successively larger powers of $`\pi `$.
Since $`M_i`$ is supported on $`S_i`$, we may apply Lemma 4.4.4 to produce $`U_{i+1}=U_iV_i`$, with $`V_i`$ supported on $`S_i`$, such that $`U_{i+1}^1N_iU_{i+1}+U_{i+1}^1_i(U_{i+1})`$ is supported on $`S_{i+1}`$. Then $`U_{i+1}^1AU_{i+1}^\sigma =V_i^1A_iV_i^\sigma `$ is supported on $`S_i`$; by Lemma 4.4.3, $`U_{i+1}^1AU_{i+1}^\sigma `$ is also supported on $`S_{i+1}`$. Thus the iteration continues, and we may set $`U=U_m`$. ∎
We are now ready for the decisive step, analogous to \[24, Lemma 5.2.4\].
###### Proposition 4.4.6.
Let $`N_1,\mathrm{},N_m`$ be $`n\times n`$ matrices over $`\mathrm{\Gamma }_{\mathrm{con}}`$ with $`w(N_i)>0`$ for $`i=1,\mathrm{},m`$. Suppose that there exists an invertible $`n\times n`$ matrix $`A`$ over $`\mathrm{\Gamma }`$ such that $`w(AI_n)>w(p)/(p1)`$ and
$$N_iA+_i(A)=qAN_i^\sigma (i=1,\mathrm{},m).$$
Then there exists an invertible $`n\times n`$ matrix $`U`$ over $`\mathrm{\Gamma }`$ such that $`AU^\sigma =U`$.
###### Proof.
Suppose the contrary; then there exists some smallest integer $`h>w(p)/(p1)`$ such that the equation $`U^1AU^\sigma I_n(mod\pi ^{h+1})`$ cannot be solved for $`U`$ invertible over $`\mathrm{\Gamma }`$. Since $`\mathrm{\Gamma }_{\mathrm{con}}`$ is $`\pi `$-adically dense in $`\mathrm{\Gamma }`$, we may change basis over $`\mathrm{\Gamma }_{\mathrm{con}}`$ to reduce to the case where $`h=w(AI_n)`$ and we cannot write the reduction of $`\pi ^h(AI_n)`$ modulo $`\pi `$ in the form $`BB^\sigma `$.
Choose $`r_0>0`$ such that for $`i=1,\mathrm{},m`$, $`N_i`$ has entries in $`\mathrm{\Gamma }_{r_0}`$ and $`w_{r_0}(N_i)>0`$. Since $`h>w(p)/(p1)`$, we have $`hp/(h+w(p))>1`$; we can thus choose $`c`$ with $`1<c<hp/(h+w(p))`$. Write $`r_j=r_0p^jc^j`$.
Define a sequence $`U_0,U_1,\mathrm{}`$ of invertible matrices over $`\mathrm{\Gamma }_{\mathrm{con}}`$ as follows. Start with $`U_0=I_n`$. For $`j0`$, suppose that we have constructed an invertible matrix $`U_j`$ over $`\mathrm{\Gamma }_{r_j}`$ with the following properties:
1. $`w(U_jI_n)h`$ and $`w_{r_j}(U_jI_n)>0`$;
2. $`A_j=U_j^1AU_j^\sigma `$ and $`N_{i,j}=U_j^1N_iU_j+U_j^1_i(U_j)`$ are supported on $`p^jL`$ for $`i=1,\mathrm{},m`$;
3. $`w_{r_j}(N_{i,j})>jw(p)`$.
Write $`A_j^{}`$ and $`N_{i,j}^{}`$ for the matrices $`A_j`$ and $`p^jN_{i,j}`$ viewed in $`\mathrm{\Gamma }^{E_j}`$ for $`E_j=k((p^jL))_\lambda `$, put $`\mu _i^{}=p^j\mu _i`$, and let $`_i^{}`$ be the derivation on $`\mathrm{\Gamma }^{E_j}`$ corresponding to $`\mu _i^{}`$. Then $`w_{r_j}(N_{i,j}^{})>0`$, and
$$N_{i,j}^{}A_j^{}+_i^{}(A_j^{})=qA_j^{}(N_{i,j}^{})^\sigma .$$
Put $`s=r_j(h+w(p))c/(hp)<r_j`$, and apply Lemma 4.4.5 to produce $`U_{j+1}`$ over $`\mathrm{\Gamma }_s`$ supported on $`p^jL`$, with $`w(U_{j+1}I_n)h`$ and $`w_s(U_{j+1}I_n)>0`$, such that $`A_{j+1}=U_{j+1}^1A_jU_{j+1}^\sigma `$ is supported on $`p^{j+1}L`$.
Since $`r_{j+1}=sh/(h+w(p))<s`$, (a) is satisfied again. From the equation
$$N_{i,j+1}A_{j+1}+_i(A_{j+1})=qA_{j+1}N_{i,j+1}^\sigma $$
(4.4.6.1)
(a consequence of Remark 4.1.3), we see that each $`N_{i,j+1}`$ is also supported on $`p^{j+1}L`$ (the argument is as in the proof of Lemma 4.4.5). Hence (b) is satisfied again.
To check (c), note that on one hand, (4.4.6.1) and the fact that $`w(_i(A_{j+1}))h+(j+1)w(p)`$ imply that $`w(N_{i,j+1})h+(j+1)w(p)`$. On the other hand, the facts that $`w_s(N_{i,j})>w_{r_j}(N_{i,j})>jw(p)`$, $`U_{j+1}`$ is supported on $`p^jL`$, and $`w_s(U_{j+1}I_n)>0`$ imply that $`w_s(N_{i,j+1})>jw(p)`$, and so
$`w_{r_{j+1}}(N_{i,j+1})`$ $`=\underset{mh+(j+1)w(p)}{\mathrm{min}}\{r_{j+1}v_m(N_{i,j+1})+m\}`$
$`{\displaystyle \frac{r_{j+1}}{s}}w_s(N_{i,j+1})+(h+(j+1)w(p))\left(1{\displaystyle \frac{r_{j+1}}{s}}\right)`$
$`>{\displaystyle \frac{r_{j+1}}{s}}(jw(p))+(h+(j+1)w(p))\left(1{\displaystyle \frac{r_{j+1}}{s}}\right)`$
$`={\displaystyle \frac{h}{h+w(p)}}(jw(p))+(h+(j+1)w(p)){\displaystyle \frac{w(p)}{h+w(p)}}`$
$`=(j+1)w(p).`$
Hence (c) is satisfied again, and the iteration may continue.
Note that if we take $`X=U_jI_n`$, then $`AX+X^\sigma A_j(mod\pi ^h)`$. This means that on one hand, $`AX+X^\sigma `$ is congruent modulo $`\pi ^h`$ to a matrix supported on $`p^jL`$, and on the other hand,
$`v_h(AX+X^\sigma )`$ $`\mathrm{min}\{v_h(A),v_h(X),v_h(X^\sigma )\}`$
$`\mathrm{min}\{v_h(A),hr_0^1p^{j+1}c^j\}`$
since $`w_{r_j}(X)>0`$. However, since $`c>1`$, this last inequality contradicts Lemma 4.3.1 for $`j`$ large. This contradiction means that our original assumption was incorrect, i.e., the desired matrix $`U`$ does exist, as desired. ∎
###### Remark 4.4.7.
In his setting, Tsuzuki actually proves a stronger result \[24, Proposition 6.1.10\] that produces solutions of $`p`$-adic differential equations even without a Frobenius structure. As noted in Remark 3.1.6, one cannot hope to do likewise in our setting.
### 4.5 Trivialization
We now begin reaping the fruits of our labors, first in a restricted setting. Note that Convention 4.4.1 is no longer in force.
###### Proposition 4.5.1.
Suppose that the field $`k`$ is algebraically closed. Let $`M`$ be a $``$-module over $`\mathrm{\Gamma }_{\mathrm{con}}`$ which becomes an $`(F,)`$-module over $`\mathrm{\Gamma }`$. Then as a representation of $`G=\mathrm{Gal}(E^{\mathrm{sep}}/E)`$, $`D_\mathrm{\Gamma }(M)`$ has finite image; moreover, if $`D_\mathrm{\Gamma }(M)`$ is trivial modulo $`\pi ^m`$ for some integer $`m>w(p)/(p1)`$, then $`D_\mathrm{\Gamma }(M)`$ is trivial.
###### Proof.
We treat the second assertion first. Suppose that $`D_\mathrm{\Gamma }(M)`$ is trivial modulo $`\pi ^m`$ for some integer $`m>w(p)/(p1)`$. By Remark 4.2.6, we can change the choice of the Frobenius lift without affecting the fact that $``$ admits a compatible Frobenius structure over $`\mathrm{\Gamma }`$, or that $`D_\mathrm{\Gamma }(M)`$ is trivial modulo $`\pi ^m`$. In particular, we may assume that $`\sigma `$ is a standard Frobenius lift; we may then choose a coordinate system for $`\mathrm{\Gamma }`$ to drop back into the purview of Convention 4.4.1.
Given a basis $`𝐞_1,\mathrm{},𝐞_n`$ of $`M`$, define the $`n\times n`$ matrices $`A,N_1,\mathrm{},N_m`$ by
$`F(𝐞_l)`$ $`={\displaystyle \underset{j}{}}A_{jl}𝐞_j`$
$`\mathrm{\Delta }_i(𝐞_l)`$ $`={\displaystyle \underset{j}{}}(N_i)_{jl}𝐞_j;`$
as in Remark 4.1.2, we then have $`N_iA+_i(A)=qAN_i^\sigma `$ for $`i=1,\mathrm{},m`$. By hypothesis, we can arrange to have $`w(AI_n)>w(p)/(p1)`$; we may thus apply Proposition 4.4.6 to produce an invertible $`n\times n`$ matrix $`U`$ over $`\mathrm{\Gamma }`$ with $`AU^\sigma =U`$. Writing $`𝐯_l=_jU_{jl}𝐞_j`$, we then have $`F𝐯_j=𝐯_j`$ for $`j=1,\mathrm{},n`$; that is, $`M_{\mathrm{\Gamma }_{\mathrm{con}}}\mathrm{\Gamma }`$ admits an $`F`$-invariant basis, so $`D_\mathrm{\Gamma }(M)`$ is trivial.
We now proceed to the first assertion. Pick an integer $`m`$ with $`m>w(p)/(p1)`$. Let $`E^{}`$ be the fixed field of the kernel of the action of $`G`$ on $`D_\mathrm{\Gamma }(M)/\pi ^mD_\mathrm{\Gamma }(M)`$. Then by what we have just shown, the restriction of $`D_\mathrm{\Gamma }(M)`$ to $`\mathrm{Gal}(E^{\mathrm{sep}}/E^{})`$ is trivial; hence $`D_\mathrm{\Gamma }(M)`$, as a representation of $`G`$, has finite image. ∎
Finally, we give the analogue of Tsuzuki’s unit-root monodromy theorem \[24, Theorems 4.2.6 and 5.1.1\].
###### Theorem 4.5.2.
Let $`M`$ be a unit-root $`(F,)`$-module over $`\mathrm{\Gamma }_{\mathrm{con}}`$. Then there exists a finite separable extension $`E^{}`$ of $`E`$ such that $`M_{\mathrm{\Gamma }_{\mathrm{con}}}\mathrm{\Gamma }_{\mathrm{con}}^E^{}`$ admits a basis of elements which are horizontal, and also $`F`$-invariant in case $`k`$ is algebraically closed.
###### Proof.
Suppose $`k`$ is algebraically closed; by Proposition 4.5.1, for some finite separable extension $`E^{}`$ of $`E`$, there exists a basis of $`M_{\mathrm{\Gamma }_{\mathrm{con}}}\mathrm{\Gamma }^{}`$ (for $`\mathrm{\Gamma }^{}=\mathrm{\Gamma }^E^{}`$) consisting of horizontal $`F`$-invariant elements. Since $`M`$ is unit-root, we may apply \[14, Lemma 5.4.1\] to deduce that any $`F`$-invariant element of $`M_{\mathrm{\Gamma }_{\mathrm{con}}}\mathrm{\Gamma }^{}`$ actually belongs to $`M_{\mathrm{\Gamma }_{\mathrm{con}}}\mathrm{\Gamma }_{\mathrm{con}}^{}`$; this yields the claim.
For $`k`$ general, let $`E^{}`$ denote the completion of the compositum of $`E`$ and $`k^{\mathrm{alg}}`$ over $`k`$. Then the restriction of $`D_\mathrm{\Gamma }(M)`$ to $`\mathrm{Gal}((E^{})^{\mathrm{sep}}/E^{})`$ has finite image by Proposition 4.5.1. By a standard approximation argument, we can replace $`E`$ by a finite separable extension in such a way as to trivialize the action of the resulting $`\mathrm{Gal}((E^{})^{\mathrm{sep}}/E^{})`$; the resulting action of $`\mathrm{Gal}(k^{\mathrm{sep}}/k)`$ is trivial by Hilbert 90. This yields the claim. (Alternatively, one may proceed as in \[11, Proposition 6.11\] to reduce the case of $`k`$ general to the case of $`k`$ algebraically closed.) ∎
## 5 Monodromy of $`(F,)`$-modules
In this section, we recall the slope filtration theorem of (in the form presented in ), and combine it with the unit-root monodromy theorem (Theorem 4.5.2) to obtain the $`p`$-adic local monodromy theorem for fake annuli (Theorem 5.2.4).
Throughout this section, we retain Convention 4.0.1.
### 5.1 Isoclinicity
###### Definition 5.1.1.
Let $`M`$ be an $`F`$-module of rank 1 over $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$, and let $`𝐯`$ be a generator of $`M`$. Then $`F𝐯=r𝐯`$ for some $`r\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$ which is a unit, that is, $`r\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$. In particular, $`w(r)`$ is well-defined; it also does not depend on $`r`$, since changing the choice of generator multiplies $`r`$ by $`u^\sigma /u`$ for some $`u\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$, whereas $`w(u^\sigma /u)=0`$. We call the integer $`w(r)`$ the *degree* of $`M`$, and denote it by $`\mathrm{deg}(M)`$; if $`M`$ has rank $`n>1`$, we define the degree of $`M`$ as $`\mathrm{deg}(^nM)`$. We write $`\mu (M)=\mathrm{deg}(M)/\mathrm{rank}(M)`$ and call it the *slope* of $`M`$.
###### Definition 5.1.2.
An $`F`$-module $`M`$ over $`\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$ is *unit-root* (or *étale*) if it contains an $`F`$-stable $`\mathrm{\Gamma }_{\mathrm{con}}`$-lattice which forms a unit-root $`F`$-module over $`\mathrm{\Gamma }_{\mathrm{con}}`$. Note that if $`M`$ is a unit-root $`(F,)`$-module over $`\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$, then any unit-root $`\mathrm{\Gamma }_{\mathrm{con}}`$-lattice is stable under $``$ (as can be seen by applying Frobenius repeatedly).
###### Definition 5.1.3.
An $`F`$-module $`M`$ over $`\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$ is *isoclinic of slope $`s`$* if there exist integers $`c`$ and $`d`$ with $`c/d=s`$ such that $`([d]_{}M)(c)`$ is unit-root; note that necessarily $`s=\mu (M)`$. An $`F`$-module $`M`$ over $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$ is *isoclinic of slope $`s`$* if it is the base extension of an isoclinic $`F`$-module of slope $`s`$ over $`\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$; the base extension from isoclinic $`F`$-modules of a given slope over $`\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$ to isoclinic $`F`$-modules of that slope over $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$ is an equivalence of categories \[14, Theorem 6.3.3\].
###### Remark 5.1.4.
Note that this is not the definition of isoclinicity used in , but it is equivalent to it thanks to \[14, Proposition 6.3.5\].
The base extension functor mentioned above also behaves nicely with respect to connections; see \[14, Proposition 7.1.7\].
###### Proposition 5.1.5.
Let $`M`$ be an isoclinic $`F`$-module over $`\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$. Suppose that $`M\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$, with its given Frobenius, admits the structure of an $`(F,)`$-module. Then $`M`$, with its given Frobenius, already admits the structure of an $`(F,)`$-module.
### 5.2 Slope filtrations and local monodromy
The slope filtration theorem can be stated as follows; see Remark 5.2.5 for a precise citation.
###### Theorem 5.2.1.
Let $`M`$ be an $`F`$-module over $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$. Then there exists a unique filtration $`0=M_0M_1\mathrm{}M_l=M`$ by saturated $`F`$-submodules with the following properties.
1. For $`i=1,\mathrm{},l`$, the quotient $`M_i/M_{i1}`$ is isoclinic of some slope $`s_i`$.
2. $`s_1<\mathrm{}<s_l`$.
###### Definition 5.2.2.
In Theorem 5.2.1, we refer to the numbers $`s_1,\mathrm{},s_l`$ as the *Harder-Narasimhan slopes* (or *HN slopes* for short) of $`M`$, viewed as a multiset in which $`s_i`$ occurs with multiplicity $`\mathrm{rank}(M_i/M_{i1})`$. See \[14, § 4.6\] for more on the calculus of the HN slopes.
The relevance of the slope filtration theorem to $`(F,)`$-modules comes via the following fact \[14, Proposition 7.1.6\].
###### Proposition 5.2.3.
Let $`M`$ be an $`(F,)`$-module over $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$. Then each step of the filtration of Theorem 5.2.1 is an $`(F,)`$-submodule.
Using the slope filtration theorem, we easily obtain the $`p`$-adic local monodromy theorem.
###### Theorem 5.2.4 ($`p`$-adic local monodromy theorem).
Let $`M`$ be an $`(F,)`$-module over $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$. Then $`M`$ is quasi-unipotent; moreover, if $`M`$ is isoclinic, then $`M`$ is quasi-constant.
###### Proof.
Let $`0=M_0M_1\mathrm{}M_l=M`$ be the filtration of the underlying $`F`$-module of $`M`$ given by Theorem 5.2.1; by Proposition 5.2.3, this is also a filtration of $`(F,)`$-submodules. From the definition of isoclinicity plus Proposition 5.1.5, each successive quotient $`M_i/M_{i1}`$ can be written as $`N_i\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$, where $`N_i`$ is an $`(F,)`$-module over $`\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$ whose underlying $`F`$-module is isoclinic.
It suffices to check that each $`N_i`$ is quasi-constant; since that condition does not depend on the Frobenius structure, we may check after applying $`[d]_{}`$ and twisting. We may thus reduce to the case where $`N_i`$ is unit-root; in that case, Theorem 4.5.2 asserts that indeed $`N_i`$ is quasi-constant, as desired. ∎
###### Remark 5.2.5.
For true annuli, Theorem 5.2.4 is what is normally called the “$`p`$-adic (local) monodromy theorem”. The proof here, restricted to that case, is essentially the same as in \[11, Theorem 6.12\], except that the invocation of the slope filtration theorem \[11, Theorem 6.10\] is replaced with the more refined form \[14, Theorem 6.4.1\]. Proofs in the true annuli case have also been given by André \[1, Théorème 7.1.1\] and Mebkhout \[21, Corollaire 5.0-23\]; these rely not on a close analysis of Frobenius (as in the slope filtration theorem), but on the close analysis of connections on annuli given by the $`p`$-adic index theorem of Christol-Mebkhout . As per Remark 3.1.6, it seems unlikely that such an approach can be made to work in the fake annuli setting, at least without integrating Frobenius structures into the analysis.
###### Remark 5.2.6.
In case $`k`$ is algebraically closed, one can refine the conclusion of Theorem 5.2.4. Namely, given an constant $`(F,)`$-module, the $`K`$-span of the horizontal sections form an $`F`$-module over $`K`$, to which we may apply the classical Dieudonné-Manin theorem; the result is a decomposition of the given $`(F,)`$-module into pieces of the form $`[d]^{}\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}(c)`$. Such pieces are called *standard* in and .
## 6 Complements
In this section, we gather some consequences of the $`p`$-adic local monodromy theorem for fake annuli. These generalize known consequences of the ordinary $`p`$LMT: calculation of some extension groups, local duality, and full faithfulness of overconvergent-to-convergent restriction.
Throughout this section, retain Convention 4.0.1.
### 6.1 Kernels and cokernels
We calculate some Hom and Ext groups in the category of $`(F,)`$-modules.
###### Definition 6.1.1.
For $`M`$ an $`F`$-module over some ring, let $`H_F^0(M)`$ and $`H_F^1(M)`$ denote the kernel and cokernel of $`F1`$ on $`M`$. Note that $`\mathrm{Hom}_F(M_1,M_2)=H_F^0(M_1^{}M_2)`$ and $`\mathrm{Ext}_F^1(M_1,M_2)=H_F^1(M_1^{}M_2)`$.
###### Definition 6.1.2.
For $`M`$ an $`(F,)`$-module over some ring, let $`H_{F,}^0(M)`$ be the subgroup of $`𝐯M`$ with $`F(𝐯)=𝐯`$ and $`(𝐯)=0`$. Let $`H_{F,}^1(M)`$ be the set of pairs $`(𝐯,\omega )M\times (M\mathrm{\Omega }^1)`$ with
$$\omega +(𝐯)=(Fd\sigma )(\omega ),_1(\omega )=0,$$
(6.1.2.1)
modulo pairs of the form $`(F(𝐰)𝐰,(𝐰))`$ for some $`𝐰M`$. Note that $`\mathrm{Hom}_{F,}(M_1,M_2)=H_{F,}^0(M_1^{}M_2)`$ and $`\mathrm{Ext}_{F,}^1(M_1,M_2)=H_{F,}^1(M_1^{}M_2)`$. In particular, by Proposition 3.4.7, we can use any Frobenius lift $`\sigma `$ to compute $`H_{F,}^i(M)`$.
###### Lemma 6.1.3.
For $`d`$ an integer, we have
$$H_{F,}^0(\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}(d))=\{\begin{array}{cc}K_q\hfill & d=0\hfill \\ 0\hfill & d0.\hfill \end{array}$$
###### Proof.
Note that $`\mathrm{ker}(d:\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}\mathrm{\Omega }_{\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}/𝔬}^1)=K`$. Then note that for $`xK`$ nonzero, $`w(x^\sigma \pi ^d)=w(x)+d`$ can only equal $`w(x)`$ for $`d=0`$. ∎
###### Proposition 6.1.4.
Let $`M`$ be an $`(F,)`$-module over $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$ whose HN slopes are all positive. Then $`H_{F,}^1(M)=0`$.
###### Proof.
Consider a short exact sequence $`0MN\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}0`$; by \[14, Proposition 7.4.4\], the exact sequence splits if and only if $`N`$ has smallest HN slope zero. In particular, this may be checked after enlarging $`k`$, applying $`[a]_{}`$, and passing from $`k((L))_\lambda `$ to a finite separable extension. By Theorem 5.2.4, this allows us to reduce to the case where $`N`$ is a successive extension of twists of trivial $`(F,)`$-modules whose slopes are the HN slopes of $`N`$. If these slopes are all positive, then repeated application of Lemma 6.1.3 implies that the map $`N\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$ is zero, which it isn’t; hence $`N`$ has smallest HN slope zero. It follows that $`H_{F,}^1(M)=0`$, as desired. ∎
###### Proposition 6.1.5.
Assume that $`k`$ is algebraically closed. Put $`n=w(q)`$, let $`J`$ be the subgroup of $`xK`$ satisfying $`qx^\sigma =\pi ^nx`$, and let $`z_1,\mathrm{},z_m`$ be a basis of $`L`$. Then for $`d`$ an integer, we have
$$H_{F,}^1(\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}(d))=\{\begin{array}{cc}Jd\mathrm{log}\{z_1\}\mathrm{}Jd\mathrm{log}\{z_m\}\hfill & d=n\hfill \\ 0\hfill & dn.\hfill \end{array}$$
###### Proof.
For $`d>0`$, Proposition 6.1.4 implies that $`H_{F,}^1(\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}(d))=0`$; we may thus focus on $`d0`$. First suppose $`d=0`$. Let $`0\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}M\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}0`$ be a short exact sequence of $`(F,)`$-modules. Then by Theorem 5.2.1 and Lemma 6.1.3, $`M`$ cannot have any nonzero slopes, so $`M`$ is isoclinic of slope 0; as in Definition 5.1.3, we thus obtain a short exact sequence $`0\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]M_0\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]0`$ of $`(F,)`$-modules over $`\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$ from which the original sequence is obtained by tensoring up to $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$. Choose a basis $`𝐯,𝐰`$ of $`M_0`$ such that $`𝐯`$ is an $`F`$-stable element of $`\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$ within $`M_0`$, $`𝐰`$ maps to an $`F`$-stable element under the map $`M_0\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$, and $`F(𝐰)𝐰=c𝐯`$ with $`w(c)>w(p)/(p1)`$. By Proposition 4.5.1, for some finite separable extension $`E^{}`$ of $`E`$, $`M_0\mathrm{\Gamma }_{\mathrm{con}}^E^{}[\pi ^1]`$ admits a basis of $`F`$-stable elements $`𝐞_1,𝐞_2`$. By Lemma 6.1.3, $`𝐯`$ is a $`K_q`$-linear combination of $`𝐞_1,𝐞_2`$; we may thus assume that $`𝐞_1=𝐯`$. Similarly, we may assume that $`𝐞_2`$ and $`𝐰`$ have the same image under $`M\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$. Thus the original exact sequence splits, as desired.
Now suppose $`d<0`$; we may assume that $`\sigma `$ is a standard Frobenius lift. Write $`\mathrm{\Gamma }_{\mathrm{con}}^{}`$ for the subring of $`\mathrm{\Gamma }_{\mathrm{con}}`$ of series supported on the set $`\{zL:\lambda (z)0\}`$. We first check that given a pair $`(a,\omega )`$ representing an element of $`H_{F,}^1(\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}(d))`$, if $`a\mathrm{\Gamma }_{\mathrm{con}}^{}[\pi ^1]`$ and $`a`$ is supported on $`(LqL)\{0\}`$, then we must have $`aK`$. Put $`\omega =x_1d\mathrm{log}\{z_1\}+\mathrm{}+x_md\mathrm{log}\{z_m\}`$, so that
$$x_i+_i(a)=\pi ^dqx_i^\sigma (i=1,\mathrm{},m).$$
For $`zLqL`$, suppose that the coefficient of $`\{z\}`$ in $`a`$ is nonzero. Choose $`i`$ with $`\mu _i(z)0`$, so that the coefficient of $`\{z\}`$ in $`_i(a)`$ is nonzero, and for $`j=0,1,\mathrm{}`$, let $`c_j`$ be the coefficient of $`\{z\}^{q^j}`$ in $`x_i`$. Since $`\sigma `$ is standard, the coefficient of $`\{z\}`$ in $`qx_i^\sigma `$ is zero; hence $`c_00`$. Moreover, $`c_{j+1}=\pi ^dqc_j^\sigma `$ for $`j=0,1,\mathrm{}`$. It follows that $`w(c_j)=w(c_0)+j(w(q)+d)`$ for all $`j`$; however, by the definition of $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$, we must have $`lim\; inf_j(w(c_j)/q^j)>0`$, contradiction. Thus the coefficient of $`\{z\}`$ in $`a`$ is zero for each $`zLqL`$; since $`a`$ is supported on $`(LqL)\{0\}`$, we must have $`aK`$.
We next check that if $`a\mathrm{\Gamma }_{\mathrm{con}}^{}[\pi ^1]`$, then the pair $`(a,\omega )`$ represents the same class in $`H_{F,}^1(\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}(d))`$ as another pair $`(a^{},\omega ^{})`$ with $`a^{}\mathrm{\Gamma }_{\mathrm{con}}^{}[\pi ^1]`$ supported on $`(LqL)\{0\}`$, and hence $`a^{}K`$ as above. Write $`a=_{zL}a_z\{z\}`$, and for $`j=0,1,\mathrm{}`$, write $`f_j(a)`$ for the sum of $`a_z\{z\}`$ over all $`zq^jLq^{j+1}L`$. Then the sum
$$y=\underset{j=0}{\overset{\mathrm{}}{}}\underset{l=1}{\overset{j}{}}(\pi ^{\{l\}})^df_j(a)^{\sigma ^l},$$
where $`\pi ^{\{0\}}=1`$ and $`\pi ^{\{l+1\}}=(\pi ^{\{l\}})^\sigma \pi `$, converges in $`\mathrm{\Gamma }_{\mathrm{con}}^{}[\pi ^1]`$, so we can represent the same class in $`H_{F,}^1(\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}(d))`$ by a pair with first member $`a^{}=ay+\pi ^dy^\sigma `$. Since
$$a^{}=a_0+\underset{j=0}{\overset{\mathrm{}}{}}(\pi ^{\{j\}})^df_j(a)^{\sigma ^j},$$
$`a^{}`$ is supported on $`(LqL)\{0\}`$.
Next, we check that any pair $`(a,\omega )`$ represents the same class in $`H_{F,}^1(\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}(d))`$ as another pair $`(a^{},\omega ^{})`$ with $`a^{}\mathrm{\Gamma }_{\mathrm{con}}^{}[\pi ^1]`$. Write $`a=_{zL}a_z\{z\}`$, and let $`a_+,a_0,a_{}`$ be the sum of $`a_z\{z\}`$ over those $`zL`$ with $`\lambda (z)`$ positive, zero, negative, respectively. We can then represent the same class in $`H_{F,}^1(\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}(d))`$ by a pair with first member $`a^{}=ay+\pi ^dy^\sigma `$, for
$$y=\underset{i=0}{\overset{\mathrm{}}{}}(\pi ^{\{i\}})^da_+^{\sigma ^i}.$$
Then $`a^{}=a_0+a_{}\mathrm{\Gamma }_{\mathrm{con}}^{}[\pi ^1]`$.
Combining the previous paragraphs, we find that every element of $`H_{F,}^1(\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}(d))`$ is represented by a pair $`(a,\omega )`$ with $`aK`$, and consequently $`\pi ^dx_i=qx_i^\sigma `$. Since $`k`$ is algebraically closed, we can force $`a=0`$; moreover, the resulting class representative is in fact unique. This yields the desired result. ∎
###### Remark 6.1.6.
Note that in the notation of Proposition 6.1.5, $`J`$ is a one-dimensional vector space over $`K_q`$.
### 6.2 Duality and decompositions
###### Lemma 6.2.1.
Any irreducible $`(F,)`$-module over $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$ is isoclinic and quasi-constant.
###### Proof.
An irreducible $`(F,)`$-module admits a slope filtration on its underlying $`F`$-module by Theorem 5.2.1, and the slope filtration is $``$-stable by Proposition 5.1.5. Hence it must have a single step, i.e., the module is isoclinic. Since any isoclinic $`(F,)`$-module is quasi-constant (Theorem 5.2.1), the claims follow. ∎
###### Definition 6.2.2.
Let $`M,N`$ be $`(F,)`$-modules over a nearly admissible ring $`S`$. By the *cup product*, we will mean the natural bilinear map $`H_{F,}^0(M)\times H_{F,}^1(N)H_{F,}^1(MN)`$ sending $`(x,(𝐯,\omega ))`$ to $`(x𝐯,x\omega )`$. Define the *Poincaré pairing* on $`M`$ as the $`F`$-equivariant bilinear pairing obtained by composing the cup product map
$$H_{F,}^0(M)\times H_{F,}^1(M^{}(w(q)))H_{F,}^1(MM^{}(w(q)))$$
with the map
$$H_{F,}^1(MM^{}(w(q)))H_{F,}^1(S(w(q)))$$
given by the trace map $`MM^{}S`$.
###### Proposition 6.2.3.
Assume that $`k`$ is algebraically closed, and let $`M`$ be an $`(F,)`$-module over $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$. Then the Poincaré pairing $`H_{F,}^0(M)\times H_{F,}^1(M^{}(w(q)))H_{F,}^1(\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}(w(q)))`$ is perfect, i.e., it induces an isomorphism
$$H_{F,}^1(M^{}(w(q)))\mathrm{Hom}_K(H_{F,}^0(M),H_{F,}^1(\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}(w(q))).$$
###### Proof.
The argument consists of a series of reductions ending with an appeal to the calculation in Proposition 6.1.5. To begin with, by the snake and five lemmas, we may reduce to the case where $`M`$ is irreducible; then $`M`$ is isoclinic and quasi-constant by Lemma 6.2.1 (note that this relies on the full theory of slope filtrations). Let $`s=c/d`$ be the slope of $`M^{}`$ written in lowest terms. Since $`k`$ is algebraically closed, by Theorem 4.5.2, there exists a finite separable extension $`E^{}`$ of $`E`$ such that for $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}^{}=\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}^E^{}`$, $`([d]_{}M)(c)\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}^{}`$ admits a basis of horizontal vectors.
Put $`N=M\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}^{}`$ viewed as an $`(F,)`$-module over $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$; then the trace from $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}^{}`$ to $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$ induces a projector on $`N`$ with image $`M`$, and this map commutes with the Poincaré pairing. We may thus reduce to checking the perfectness of the Poincaré pairing for $`N`$ instead of $`M`$.
In other words, we have reduced to the case where $`([d]_{}M)(c)`$ admits a basis of horizontal vectors. At this point, we may apply Proposition 3.4.7 to reduce to considering a standard Frobenius. Also, we may replace $`K`$ by a Galois extension (since we can take traces down that extension); in particular, we can force $`K`$ to contain the $`p^d`$-th roots of unity.
Since $`K`$ contains the $`p^d`$-th roots of unity, we can form a trace for the morphism $`\sigma ^d:\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$, by averaging over automorphisms of $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$ of the form $`\{z_i\}\zeta _i\{z_i\}`$ for $`\zeta _1,\mathrm{},\zeta _m\mu _{p^d}`$; this gives a trace map from $`[d]^{}[d]_{}M`$ to $`M`$. This means that to check perfectness of the pairing for $`M`$, it is enough to do so for $`[d]^{}[d]_{}M`$. Since the formation of $`H^0`$ and $`H^1`$ is insensitive to application of $`[d]^{}`$, we are reduced to checking perfectness for $`[d]_{}M`$.
However, Theorem 4.5.2 actually asserts that $`([d]_{}M)(c)`$ admits a basis of horizontal vectors *stable under $`F^d`$*. That is, as a $`(F^d,)`$-module, $`[d]_{}M`$ splits up as a direct sum of copies of $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}(c)`$. Once more by the snake lemma, we now reduce perfectness of the Poincaré pairing for $`[d]_{}M`$ to perfectness for $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}(c)`$. Since the latter follows from Proposition 6.1.5, we are done. ∎
###### Proposition 6.2.4.
Assume that $`k`$ is algebraically closed. Let $`M_1,M_2`$ be irreducible $`(F,)`$-modules over $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$, neither of which is isomorphic to a twist of the other. Then
$$\mathrm{Hom}_{F,}(M_1,M_2)=\mathrm{Ext}_{F,}^1(M_1,M_2)=0.$$
###### Proof.
If $`M_1`$ and $`M_2`$ are irreducible, then $`M_1`$ and $`M_2`$ are isomorphic if and only if $`H_{F,}^0(M_1^{}M_2)0`$ if and only if $`H_{F,}^0(M_2^{}M_1)0`$. Now Proposition 6.2.3 yields the desired result. ∎
We may now refine the conclusion of Theorem 5.2.4 as follows.
###### Definition 6.2.5.
Let $`N`$ be an irreducible $`(F,)`$-module. We say that another $`(F,)`$-module $`M`$ is *$`N`$-typical* if $`M`$ admits an exhaustive filtration by saturated $`(F,)`$-submodules, whose successive quotients are isomorphic to twists of $`N`$. If $`N`$ is not to be specified, we say $`M`$ is *isotypical*.
###### Proposition 6.2.6.
Assume that $`k`$ is algebraically closed. Let $`M`$ be an $`(F,)`$-module over $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$. Then $`M`$ admits a unique direct sum decomposition $`M_1\mathrm{}M_l`$ into isotypical $`(F,)`$-submodules, such that each $`M_i`$ is $`N_i`$-typical for some $`N_i`$, and no two $`N_i`$ are twists of each other.
###### Proof.
The uniqueness follows from the fact that there are no nonzero morphisms between $`(F,)`$-modules which are isotypical for irreducible modules which are not twists of each other; this follows by repeated application of Proposition 6.2.4.
We prove existence by induction on $`M`$. If $`M`$ is irreducible, then $`M`$ itself is isotypical. Otherwise, choose a short exact sequence $`0M_0MN0`$ with $`M_0`$ irreducible. Decompose $`N=N_i`$ by the induction hypothesis, and let $`P_i`$ be the preimage of $`N_i`$ in $`M`$. For each $`i`$, if $`N_i`$ is not $`M_0`$-typical, then the exact sequence $`0M_0P_iN_i0`$ splits, again by repeated application of Proposition 6.2.4. This is true for all but possibly one $`i`$; we may thus split $`M`$ as a direct sum of those $`N_i`$ plus an $`M_0`$-typical factor. This completes the induction. ∎
###### Remark 6.2.7.
One can doubtless refine Proposition 6.2.6 with more work. For instance, it should be possible to drop the restriction that $`k`$ be algebraically closed. For another, it should be possible to show that an $`M`$-typical $`(F,)`$-module is isomorphic to the tensor product of $`M`$ with a unipotent $`(F,)`$-module; for true annuli, this amounts to a result of Matsuda \[19, Theorem 7.8\], which in turn mimics a result of Levelt in the context of classical differential equations. (However, one must keep Remark 3.1.6 in mind: while Matsuda’s result is purely about connections, one is compelled to use the Frobenius also when working on fake annuli.) This should in turn make it possible to construct a monodromy representation in this setting, as in \[13, Theorem 4.45\], and perhaps to relate it to some form of the Christol-Mebkhout construction, as in \[13, Theorem 5.23\]. The latter may be related to some conjectures of Matsuda; see .
### 6.3 Splitting exact sequences
###### Lemma 6.3.1.
Let
$$0M_1MM_20$$
(6.3.1.1)
be a short exact sequence of $`F`$-modules over $`\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$ or $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$, and put $`d=\mathrm{rank}(M_1)`$. Then the sequence splits if and only if the sequence
$$0^dM_1^dM(^dM)/(^dM_1)0$$
(6.3.1.2)
splits.
###### Proof.
If (6.3.1.1) splits, then (6.3.1.2) splits by the Künneth decomposition. Conversely, if (6.3.1.2) splits, let $`N`$ be the image of $`(^{d1}M_1)M`$ in $`^dM`$ under $``$; then the exact sequence
$$0^dM_1N(^{d1}M_1M_2)0$$
splits. Tensor with $`M_1`$ to obtain another split exact sequence
$$0(M_1^dM_1)(M_1N)(M_1^{d1}M_1M_2)0.$$
Twisting by $`(^dM_1)^{}`$, we obtain a split exact sequence
$$0M_1P(M_1M_1^{}M_2)0$$
for some $`P`$. Take the trace component within $`M_1M_1^{}`$, tensor with $`M_2`$, and let $`Q`$ be the inverse image in $`P`$; we then obtain yet another split exact sequence
$$0M_1QM_20.$$
By backtracking through the definitions, we see that this is none other than (6.3.1.1). ∎
###### Proposition 6.3.2.
Suppose that $`\sigma `$ is standard and that $`k`$ is algebraically closed. Let $`M`$ be an $`(F,)`$-module over $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$. Then any exact sequence $`0M_1MM_20`$ in the category of $`F`$-modules splits.
###### Proof.
By Lemma 6.3.1, we may assume that $`\mathrm{rank}(M_1)=1`$; by twisting, we may assume that $`M_1\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$. Also, since $`k`$ is algebraically closed, we may assume that $`\pi ^n=q`$ for some integer $`n`$.
Choose a coordinate system for $`\mathrm{\Gamma }`$. Let $`𝐯`$ be an $`F`$-stable element of $`M_1`$; then if $`𝐰=\mathrm{\Delta }_1^{i_1}\mathrm{}\mathrm{\Delta }_m^{i_m}(𝐯)`$, we have $`F(𝐰)=q^{i_1\mathrm{}i_m}𝐰`$. We thus obtain a nonzero map $`NM`$ for some unipotent $`(F,)`$-module $`N`$, whose image contains $`𝐯`$.
Since $`k`$ is algebraically closed, by Proposition 6.2.6, we can write $`M`$ as a direct sum of isotypical $`(F,)`$-submodules $`P_1\mathrm{}P_l`$. At most one of the $`P_i`$ is isotypical for the trivial $`(F,)`$-module $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$; if $`P_j`$ is one of the others, then the map $`NP_j`$ obtained by composing the map $`NM_1`$ and the projection $`MP_j`$ is zero by Proposition 6.2.4. We may thus reduce to the case where $`M`$ is unipotent.
In this case, by repeated application of Proposition 6.1.5, we deduce that $`M`$ is isomorphic as an $`F`$-module to a direct sum of twists of the trivial $`F`$-module. Then $`𝐯`$ must be a $`K_q`$-linear combination of $`F`$-stable generators of summands in this decomposition; from this observation, we may construct an $`F`$-stable complement of $`M_1`$, yielding the desired splitting. ∎
###### Remark 6.3.3.
The proof of Proposition 6.3.2 depends heavily on the hypothesis that $`\sigma `$ is a standard Frobenius lift. Whether the result should even hold otherwise is not entirely clear.
### 6.4 Descent of morphisms
Following the philosophy of (as imitated in ), we now parlay our splitting results into statements that let us descend morphisms of $`(F,)`$-modules from $`\mathrm{\Gamma }`$ to $`\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$.
###### Definition 6.4.1.
If $`M`$ is an $`F`$-module over $`\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$, we may associate to $`M`$ two sets of HN slopes, by passing to $`\mathrm{\Gamma }[\pi ^1]`$ and invoking Theorem 5.2.1 for the trivial valuation on $`E`$, or by passing to $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$; we refer to these as the *generic HN slopes* and *special HN slopes*, respectively. Note that the Newton polygon of the special slopes always lies on or above that of the generic slopes, with the same endpoint \[14, Proposition 5.5.1\].
###### Proposition 6.4.2.
Suppose that $`k`$ is algebraically closed. Let $`0M_1MM_20`$ be a short exact sequence of $`F`$-modules over $`\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$, for $`\sigma `$ a standard Frobenius lift, such that $`M`$ acquires the structure of an $`(F,)`$-module over $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$. Suppose that each generic HN slope of $`M_1`$ is greater than each generic HN slope of $`M_2`$. Then the exact sequence splits.
###### Proof.
By applying Lemma 6.3.1 and taking duals, we may assume that $`\mathrm{rank}(M_2)=1`$; by twisting, we may assume that $`M_2\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$. Then the slopes condition is that $`M_1`$ has all generic HN slopes positive. By Proposition 6.3.2, the exact sequence
$$0M_1\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}M\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}0$$
of $`F`$-modules splits; by \[14, Proposition 7.4.2\], the original sequence also splits. ∎
###### Theorem 6.4.3.
1. Let $`M`$ be an $`F`$-module over $`\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$, for $`\sigma `$ a standard Frobenius lift, which acquires the structure of an $`(F,)`$-module over $`\mathrm{\Gamma }_{\mathrm{an},\mathrm{con}}`$. Then the natural map
$$H_F^0(M)H_F^0(M\mathrm{\Gamma }[\pi ^1])$$
is a bijection.
2. Let $`M`$ be an $`(F,)`$-module over $`\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$. Then the natural map
$$H_{F,}^0(M)H_{F,}^0(M\mathrm{\Gamma }[\pi ^1])$$
is a bijection.
###### Proof.
In either case, there is no harm in assuming that $`k`$ is algebraically closed.
1. Given $`𝐯H_F^0(M\mathrm{\Gamma }[\pi ^1])`$, we obtain an $`F`$-equivariant dual map $`\varphi :M^{}\mathrm{\Gamma }[\pi ^1]`$. Let $`N_0`$ be the kernel of $`\varphi `$; by \[14, Proposition 7.5.1\] (applicable because any monomial field admits a valuation $`p`$-basis, namely any coordinate system), the preimage $`N_1=\varphi ^1(\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1])`$ has the property that $`N_1/N_0\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$, and $`M^{}/N_1`$ has all generic slopes negative.
By Proposition 6.3.2, $`N_1`$ admits an $`F`$-stable complement in $`M^{}`$, so $`N_1/N_0`$ admits an $`F`$-stable complement $`P`$ in $`M^{}/N_0`$. However, the generic slopes of $`P`$ are the same as those of $`M^{}/N_1`$, so they are all negative. Thus the map $`P\mathrm{\Gamma }[\pi ^1]`$ obtained by composition with $`\varphi `$ is forced to vanish by \[14, Proposition 7.5.1\], whereas $`\varphi :M^{}/N_0\mathrm{\Gamma }[\pi ^1]`$ is injective. We must then have $`P=0`$ and $`N_1=M^{}`$, so $`\varphi `$ maps into $`\mathrm{\Gamma }_{\mathrm{con}}[\pi ^1]`$ and $`𝐯M`$, as desired.
2. By Proposition 3.4.7, there is no harm in reducing to the case of $`\sigma `$ standard. The result now follows from (a).
###### Remark 6.4.4.
Note that even for true annuli, Theorem 6.4.3 makes an assertion (namely (a)) not covered by . |
warning/0507/nucl-th0507052.html | ar5iv | text | # Density distributions of superheavy nuclei
## I introduction
Proton and neutron density distributions, the associated root-mean-square (rms) radii $`R_\mathrm{p}`$ and $`R_\mathrm{n}`$, and the neutron skin thickness $`\mathrm{\Delta }R=R_\mathrm{n}R_\mathrm{p}`$ provide fundamental information on nuclear structure. For example, halo nuclei are charactered with long tails in density distributions. The density is a direct probe of the size of an atomic nucleus, and plays an important role in the cross sections of nuclear reactions. While charge densities can be measured from the elastic scattering of electrons, neutron densities are largely unknown. The parity-violating electron scattering has been suggested to measure neutron densities Donnelly89 . With the experimental method, the neutron densities of more nuclei can be expected to be measured. Theoretically, the calculated neutron skin thickness can be model dependent Horow01 .
Recently, Horowitz et al. Horowitz01 studied the relationship between the neutron skin of the spherical double magic nucleus <sup>208</sup>Pb and the properties of neutron-star crusts, showing the importance of the knowledge of nucleon densities in understanding the equation of state of neutron-rich matter and therefore the properties of neutron stars. The heavy nucleus <sup>208</sup>Pb has been measured to have a neutron skin thickness of about 0.15 fm trzcinska . With increasing mass number, the neutron excess becomes larger in general and it’s natural to think that superheavy nuclei provide the largest neutron excesses. The heaviest nuclei also have large proton numbers and thus large Coulomb repulsive forces that push the protons to larger radii and therefore change density distributions. Novel density distributions were predicted in extraordinary $`A>400`$ nuclei that have bubbles and showed coupling effects between density distributions and shell structures Dietrich98 ; Decharge99 ; Yu00 . The nucleon density in a bubble is reduced to be zero. Semi-bubble nuclei were also suggested with the considerable reduction of central densities for the $`Z120`$ nuclei Decharge99 located around the center of the predicted island of stability of superheavy nuclei. Bender et al. have investigated the density distributions of superheavy nuclei with the restriction of spherical shapes bender99 . Recent progress in experiments are motivating the structure study of superheavy nuclei Hofmann00 ; Oganessian99 . Many theoretical works have investigated the properties of superheavy nuclei cwiok96 ; cwiok ; xu ; smolanczuk97 ; nazarewicz ; bender00r ; ren01 ; ren ; meng ; wu , such as shell structure, $`\alpha `$ decay and spontaneous fission. Experiments have also provided the structure information of superheavy nuclei by the in-beam study of spectroscopy Reiter99 ; Herzberg01 ; Butler02 ; Herzberg04 . In the present work, we investigate the density distributions of superheavy nuclei and related structure properties, with deformation effects taken into account.
## II calculations
The deformed Skyrme-Hartree-Fock model (SHF) blum was used in the present investigation. Pairing correlations are treated in the BCS scheme using a $`\delta `$-pairing force, $`V_{\mathrm{pair}}=V_\mathrm{q}\delta (\stackrel{}{r}_1\stackrel{}{r}_2)`$ krieger ; bender00 . The pairing strength $`V_\mathrm{q}`$ (q=p, n for the protons and neutrons, respectively) has been parameterized throughout the chart of nuclei bender00 , but the actual values are dependent on Skyrme forces chosen. The detailed values of the pairing strengthes can be found in Ref. bender99 .
Calculations are performed in coordinate space with axially symmetric shape. The density distribution of protons or neutrons is given in the two-dimensional form as follows.
$$\rho (z,r)=\underset{k}{}2v_k^2(|\psi _k^+(z,r)|^2+|\psi _k^{}(z,r)|^2)$$
(1)
where, $`\psi _k^+`$ and $`\psi _k^{}`$ are the components of the wavefunctions with intrinsic spin s<sub>z</sub>=$`+\frac{\mathrm{}}{2}`$ and $`\frac{\mathrm{}}{2}`$, respectively, and $`v_k^2`$ is the pairing occupation probability of the $`k`$-th orbit. The ground states of most superheavy nuclei are expected to have axially symmetric or spherical shape smolanczuk97 . In this paper, we consider the most important axially symmetric deformations, $`\beta _2`$ and $`\beta _4`$.
In the present work, we investigated the densities and related structure problems of even-even superheavy nuclei with Z=104$``$120. In the SHF calculations, results are in general parameter dependent. For example, the SkI3 force predicts <sup>292</sup>120 for the next magic nucleus beyond <sup>208</sup>Pb, while SLy7 and SkI4 predict <sup>310</sup>126 and <sup>298</sup>114 for the magic nucleus, respectively cwiok96 ; bender99 . To make comparison, we used the different sets of parameters SLy4, SLy7 chabanat , SkI3 and SkI4 reinhard . These sets of parameters have been developed recently with good isospin properties, and we note that they have been recommended by Rikovska Stone et al. for their ability to describe realistic neutron stars and the properties of asymmetric nuclear matter Rik03 . In the superheavy mass region, Skyrme parameter sets can reproduce experimental binding energies within a few MeV burvenich and $`\alpha `$-decay energies within a few hundred keV cwiok . Table I lists the properties of the ground states calculated with SkI4 for experimentally known even-even superheavy nuclei and the predicted magic nucleus <sup>298</sup>114. The calculated $`\alpha `$-decay energies agree with experimental data within a few hundred keV (The largest difference with data is 640 keV in <sup>266</sup>Hs). Experimental binding energies can be reproduced within $`4`$ MeV for the nuclei listed in Table I with the SkI4 force. The obtained neutron-skin thicknesses are smaller than the calculations by the relativistic mean-field (RMF) ren . (It was pointed out that RMF calculations overestimate neutron-skin thicknesses Furnstahl02 .)
### II.1 Density distributions of spherical nuclei
In order to test potential parameters in the calculations of densities, we calculated the density distributions of the spherical doubly magic nucleus <sup>208</sup>Pb using the different sets of parameters. For <sup>208</sup>Pb, the charge distribution has been measured by electron scattering pb208 . Fig.1 shows the calculated density distributions with comparison with the experimental charge density. It can be seen that the proton density given by the SkI4 force is closest to the experimental measurement. The rms radii of <sup>208</sup>Pb are calculated with the SkI4 force to be 5.43 fm for the protons and 5.61 fm for the neutrons, leading to a neutron-skin thickness of 0.18 fm which is slightly larger than the values of 0.16 fm given by the SLy4 and SLy7 forces. These results agree with the 0.15$`\pm `$0.02 fm from the recent antiprotonic atom experiment trzcinska . The SkI3 force gives a larger neutron skin thickness of 0.23 fm compared to the other three forces, which may be due to the similar behavior of SkI3 force to the RMF model bender99 . It needs to be pointed out that nuclear ground-state correlations (see, e.g., Strayer ; Dang01 ) can have visible effects on nuclear properties, such as energies and densities. The oscillations observed in the calculated density distributions in Fig.1 could be reduced when the correlation is taken into account. Such correlations, which go beyond the mean-field approximation, are not included in the present work.
Fig.2 shows the calculated densities of the spherical nucleus <sup>298</sup>114. This nucleus was predicted to be the next doubly closed shell nucleus by a Macro-microscopic model smolanczuk97 and SHF with SkI4 force bender99 . For <sup>298</sup>114, the calculations with SkI4, SLy7 and SLy4 give similar density distributions. It can be seen that the charge density distribution of <sup>298</sup>114 has a central depression. The central density depression has been predicted to exist widely in the spherical superheavy nuclei bender99 . Fig.3 displays the calculated square wavefunctions of the proton 1i<sub>13/2</sub>, 1h<sub>9/2</sub> and 2f<sub>7/2</sub> orbits that locate the 82$`<`$Z$``$114 closed shell. It can be seen that the high-$`j`$ orbits have density contributions in the nuclear surface region. This is consistent with the classical picture in which orbits with large angular momentum locate at surface. In Fig.4, we show the proton densities of <sup>208</sup>Pb and <sup>298</sup>114 for comparison. The proton density of <sup>298</sup>114 is decomposed into two parts: i)the contribution from proton orbits below the Z=82 closed shell; and ii) from the orbits in the next closed shell with 82$`<`$Z$``$114. The contribution from the orbits below Z=82 has a similar behavior to the charge distribution of <sup>208</sup>Pb, without central depression. The 82$`<`$Z$``$114 orbits (that have high-$`j`$ values) have contributions in the surface region of the nucleus, leading to a central depression in the proton density of <sup>298</sup>114.
For the neutrons at spherical case, the high-$`j`$ orbits of 2g<sub>9/2</sub>, 1i<sub>11/2</sub>, 1j<sub>15/2</sub> and 2g<sub>7/2</sub> occur in the region of N=126$``$172. The low-$`j`$ orbits of 4s<sub>1/2</sub> and 3d<sub>3/2</sub> are at N=178$``$184. Fig.5 shows the density distributions for N=172$``$196 and Z=114 with assumed spherical shape. Indeed, these nuclei were predicted to be nearly spherical in their ground states smolanczuk97 . It can be seen that proton densities have central depressions and neutron densities become centrally depressed for N$``$178.
### II.2 Densities of deformed superheavy nuclei
The shell fillings of nucleons are sensitive to the deformations of nuclei. Therefore, density distributions should be expected to be shape dependent. Most superheavy nuclei known experimentally are believed to have deformed shapes. Fig.6 displays the calculated density distributions for the Z=110 isotopes with N=160, 170 and 180 nuclei. The equilibrium deformations are determined by minimizing calculated energies. The deformations determined with the SkI4 force are $`\beta _2`$=0.25, 0.19 and 0.05 for <sup>270,280,290</sup>110, respectively. To check the possible parameter dependence, we also used the SLy7 force to calculate the densities, shown in Fig.6. The deformations determined with SLy7 are $`\beta _2`$=0.25, 0.17 and 0.05 for <sup>270,280,290</sup>110, respectively. In order to see deformation effects, we calculated the densities assuming spherical shape. It can be seen for <sup>270,280</sup>110 that the densities become more centrally depressed in spherical cases. We also see some difference between the densities along the $`z`$\- and $`r`$-axes. Fig.7 shows the two-dimensional proton density distribution for the <sup>280</sup>110 nucleus. The density has two humps in the $`z`$-axis. The double-hump distribution has been suggested experimentally, e.g., in the deformed <sup>166</sup>Er and <sup>176</sup>Yb cooper . Our calculations with SkI4 and SLy7 show that such double-hump phenomenon is relatively pronounced for nuclei around <sup>280</sup>110. The macro-microscopic calculations show that the even-even nuclei around N=170 are particularly unstable against spontaneous fission (see Fig.9 in smolanczuk97 ). This could be related to the double-hump distributions in these deformed nuclei.
Fig.8 shows the distributions of the square wavefunctions of the $`j_z=1/2`$ orbits in <sup>280</sup>110, modified with pairing occupation probabilities (see Eq.(1)). These orbits are above the Z=82 and N=126 shells for the protons and neutrons, respectively. Shown as in Fig.8, the high-$`j`$ low-$`j_z`$ orbits have important contributions to densities in the surface region of the $`z`$-axis. High-$`j`$ low-$`j_z`$ orbits have strong prolate-driving effect. Hence, sufficient number of high-$`j`$ low-$`j_z`$ orbits occupied can result in prolate shapes and double-hump densities. In <sup>166</sup>Er and <sup>176</sup>Yb that were suggested experimentally to have double-hump densities, the low-$`j_z`$ orbits of the proton 1h<sub>11/2</sub>, 1g<sub>7/2</sub> subshells are occupied.
Density distributions given by different Skyrme forces can differ as shown in Fig.6. The SLy7 predicts larger central depression in neutron densities and less central depression in proton densities than the SkI4 force. This difference also occurs in the calculation of <sup>208</sup>Pb, see Fig.1. For different parameters, the opposite behavior of proton and neutron distributions reduce the difference in the total (proton+neutron) nuclear density distributions. The origin of the opposite behavior would be due to the self-consistent coupling between protons and neutrons in the SHF model, to approach the nuclear density saturation richter . The SkI4 force is better in reproducing the density of <sup>208</sup>Pb, compared to other Skyrme forces used in the present investigation. However, the SkI4 force was pointed out to have larger spin-orbit splittings in the calculations of single-particle level schemes for the superheavy region bender99 .
### II.3 The <sup>292</sup>120 nucleus
<sup>292</sup>120 is an interesting nucleus that was predicted to be a doubly magic bender99 and spherical semi-bubble nucleus Decharge99 . We calculated the energy curve with the SLy4, SLy7, SkI3 and SkI4 parameters, see Fig.9. The results of SLy4, SLy7 and SkI4 are close to the Hartree-Fock calculations with Gogny force Decharge99 . The SkI3 force gives a shallow spherical minimum. The SkI4 calculation predicts three minima at $`\beta _2=`$0.11, $`0.12`$ and 0.52 (Superdeformations in superheavy nuclei have been discussed by Ren et al. ren01 ; ren ). It needs to be mentioned that the present calculations are restricted to axially symmetric shapes without considering the possibility of triaxiality. The inclusion of the triaxial degree of freedom could alter the shallow minima. Fig.10 shows the SkI4 calculated density distributions at the prolate and oblate shapes. For comparison, the densities at the spherical shape are also displayed. It can be seen that significant central depressions or central semi-bubble appear in the spherical case. However, the situation is considerably altered even with a small shape change. Only weak central depressions are seen at the small deformations, see Fig.10. To have a further understanding, we calculated the corresponding single-particle potentials in the $`z`$-axis, shown inside Fig.10. It can be seen that the proton potential has a considerable change with changing the deformation. The spherical proton potential has a significant hump at the center of the nucleus. This implies that the Coulomb energy can be considerably reduced by forming the center semi-bubble at the spherical shape. The Coulomb energy can also be reduced by generating the deformation of the nucleus. In the reduction of the total energy of the nucleus, there is competition between forming the central semi-bubble and generating the deformation. The deformation can affect the shell structure and then density distributions, and vice versa. In a self-consistent model, such as the SHF approach, densities are fed back into the potential, which amplifies the coupling between deformations and densities.
## III summary
In summary, the density distributions of superheavy nuclei have been investigated with the Skyrme-Hartree-Fock model. To test the model and parameters, the $`\alpha `$-decay energies of even-even superheavy nuclei and the charge density of <sup>208</sup>Pb are calculated and compared to existing experimental data. For axially-symmetrically deformed nuclei, the density distributions in the symmetric $`z`$\- and the vertical $`r`$-axes or the two-dimension distributions were calculated. The distribution in different directions can be different. The high-$`j`$ low-$`j_z`$ orbits have important contributions to the densities at nuclear surface in the $`z`$-axis, while high-$`j`$ high-$`j_z`$ orbits have important contributions at surfaces in the $`r`$-axis. The deformation effect was found to be significant in the calculation of the density distribution in the <sup>292</sup>120 nucleus. Only a weak central depression was seen in the deformed case for <sup>292</sup>120, compared to the predicted semi-bubble at the spherical shape.
###### Acknowledgements.
We thank Prof. P.M. Walker for his valuable comments, and Prof. P.-G. Reinhard for the computer code. This work was supported by the Chinese Major State Basic Research Development Program No. G2000077400, the Natural Science Foundation of China (Grants No. 10175002 and No. 10475002), the Doctoral Foundation of Chinese Ministry of Education (20030001088), the U.K. Royal Society, and the U.K. Science and Engineering Research Council. |
warning/0507/hep-th0507047.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The conjectured duality between the type IIB superstring theory on the AdS$`{}_{5}{}^{}\times `$S<sup>5</sup> space (AdS superstring) and $`D=4,𝒩=4`$ Yang-Mills theory has been driven not only studies of variety of background theories but also studies of basic aspects such as integrability. The approach of the pp-wave background superstring theory was explored by Berenstein, Maldacena and Nastase and developed in, for example . For further development Mandal, Suryanarayan and Wadia pointed out the relevance with the integrability , and Bethe anzatz approach was explored by Minahan and Zarembo and in for example . The integrability is a powerful aspect expected in the large N QCD and shown to exist in the IIB superstring theory on the AdS$`{}_{5}{}^{}\times `$S<sup>5</sup> space by Bena, Polchinski and Roiban . The integrability provides hidden symmetry generated by an infinite number of conserved “non-local” charges as well as an infinite number of conserved “local” charges which are related by a spectral parameter at different points. Related aspects on the integrability of the AdS superstring were discussed in .
Recently the conformal symmetry of AdS superstrings was conjectured due to the $`\kappa `$ symmetry . The classical conformal symmetry of the AdS superstring theory also leads to an infinite number of conserved Virasoro operators. The naive questions are how the conformal generator is related to the infinite number of conserved “local” currents, and how many independent conserved currents exist. For principal chiral models the stress-energy tensor is written by trace of the square of the conserved flat current; for reviews see refs. . For the AdS superstring theory the Wess-Zumino term and the $`\kappa `$ symmetry make a difference. Recently issues related to the integrability and the conformal symmetry of the AdS superstring theory have been discussed . In this paper we will obtain the expression of the conformal generator, which is the stress-energy tensor relating to the lowest spin “local” current, and we calculate the higher spin “local” currents to clarify independent components.
The AdS space contains the Ramond/Ramond flux which causes difficulty of the standard Neveu-Schwarz-Ramond (NSR) formulation of the superstring theory. The AdS superstring was described in the Green-Schwarz (GS) formalism by Metsaev and Tseytlin based on the coset PSU(2,2$``$4)/\[SO(4,1)$`\times `$SO(5)\] . Later Roiban and Siegel reformulate it in terms of the unconstrained GL(4$``$4) supermatrix coordinate based on an alternative coset GL(4$``$4)/\[Sp(4)$`\times `$GL(1)\]<sup>2</sup> . In this formalism the local Lorentz is gauged, and it turns out that this treatment is essential for separation into $`+/`$ modes (right/left moving modes) easier. Furthermore the fermionic constraint including the first class and second class is necessary for separation of the fermionic modes into $`+/`$ modes. As the first step toward the CFT formulation of the AdS superstring, the affine Sugawara construction , the Virasoro algebra and the algebra of currents carrying the space-time indices are also listed.
The organization of this paper is the following; in the next section the notation is introduced. In section 3 we analyze the superparticle in the AdS$`{}_{5}{}^{}\times `$S<sup>5</sup> space, and the relation between the reparametrization constraint and the conserved right invariant (RI) current is given. In section 4 we analyze the superstring in the AdS$`{}_{5}{}^{}\times `$S<sup>5</sup> space, and the infinite number of conserved currents are presented both from the conformal point of view and from the integrability point of view. We show that the stress-energy tensor is written by the “supertrace” of the square of the RI current as the lowest spin “local” current. Then we calculate higher spin “local” currents to clarify independent components of the “local” currents.
## 2 GL(4$``$4) covariant coset
We review the Roiban-Siegel formulation of the AdS$`{}_{5}{}^{}\times `$S<sup>5</sup> coset and follow the notation in . The coset GL(4$``$4)/\[GL(1)$`\times `$Sp(4)\]<sup>2</sup> is used instead of PSU(2,2$``$4)/\[SO(4,1)$`\times `$SO(5)\] for the linear realization of the global symmetry after Wick rotations and introducing the auxiliary variables. A coset element $`Z_M^A`$ is an unconstrained matrix defined on a world-volume carrying indices $`M=(m,\overline{m}),A=(a,\overline{a})`$ with $`m,\overline{m},a,\overline{a}=1,\mathrm{},4`$. The left invariant (LI) current, $`L^L`$, is invariant under the left action $`Z_M{}_{}{}^{A}\mathrm{\Lambda }_M{}_{}{}^{N}Z_{N}^{}^A`$ with a global parameter GL(4$``$4)$`\mathrm{\Lambda }`$
$`(J^L)_A{}_{}{}^{B}=(Z^1dZ)_A{}_{}{}^{B}.`$ (2.1)
The LI current satisfies the flatness condition by definition
$`dJ^L=J^LJ^L.`$ (2.2)
The right invariant (RI) current, $`J^R`$, is invariant under the right action $`Z_M{}_{}{}^{A}Z_M{}_{}{}^{B}\lambda _{B}^{}^A`$ with a local parameter \[Sp(4)$``$GL(1)\]<sup>2</sup> $`\lambda `$
$`(J^R)_M{}_{}{}^{N}=(𝒟ZZ^1)_M{}_{}{}^{N},(𝒟Z)_M{}_{}{}^{A}dZ_M{}_{}{}^{A}+Z_M{}_{}{}^{B}A_{B}^{}^A`$ (2.3)
with
$`A\lambda A\lambda ^1+(d\lambda )\lambda ^1,`$ (2.4)
and
$`dJ^R=J^RJ^R+Z(dAAA)Z^1.`$ (2.5)
Originally $`A`$ is bosonic \[Sp(4)$``$GL(1)\]<sup>2</sup> $``$$`A`$, but we will show that the fermionic constraint i.e. $`\kappa `$ symmetry gives fermionic components of $`A`$.
The conjugate momenta are introduced
$`\{Z_M{}_{}{}^{A},\mathrm{\Pi }_B{}_{}{}^{N}\}=()^A\delta _B^A\delta _M^N`$ (2.6)
as the graded Poisson bracket and $`\{q,p\}=()^{qp}\{p,q\}`$. There are also two types of differential operators; the global symmetry generator (left action generator), $`G_M^N`$, and the supercovariant derivatives (right action generator), $`D_A^B`$,
$`G_M{}_{}{}^{N}=Z_M{}_{}{}^{A}\mathrm{\Pi }_{A}^{}{}_{}{}^{N},D_A{}_{}{}^{B}=\mathrm{\Pi }_A{}_{}{}^{M}Z_{M}^{}{}_{}{}^{B}.`$ (2.7)
In our coset approach $`8\times 8=64`$ variables for $`Z_M^A`$ are introduced and auxiliary variables are eliminated by the following constraints corresponding to the stability group \[Sp(4)$`\times `$GL(1)\]<sup>2</sup>,
$`(𝐃)_{(ab)}=(\overline{𝐃})_{(\overline{a}\overline{b})}=\mathrm{tr}𝐃=\mathrm{tr}\overline{𝐃}0,`$ (2.8)
where the bosonic components are denoted by boldfaced characters as $`𝐃_{ab}D_{ab}`$ and $`\overline{𝐃}_{\overline{a}\overline{b}}D_{\overline{a}\overline{b}}`$ of (2.7). The number of the coset constraints is $`10+10+1+1=22`$, so the number of the coset parameters is $`6422=42`$ where $`10`$ bosons and $`32`$ fermions. The $`[Sp(4)]^2`$ invariant metric is anti-symmetric and a matrix is decomposed into trace part, anti-symmetric-traceless part and the symmetric part, denoted by
$`𝐌_{ab}={\displaystyle \frac{1}{4}}\mathrm{\Omega }_{ab}𝐌^c{}_{c}{}^{}+𝐌_{ab}+𝐌_{(ab)}{\displaystyle \frac{1}{4}}\mathrm{\Omega }\mathrm{tr}𝐌+𝐌+(𝐌),`$ (2.9)
with $`M_{(ab)}=\frac{1}{2}(M_{ab}+M_{ba})`$, and similar notation for the barred sector.
Both $`G_M^N`$ and $`D_A^B`$ in (2.7) satisfy GL(4$``$4) algebra. If we focus on the AdS superalgebra part, the global symmetry generators $`G_M^N`$ satisfies the global AdS superalgebra
$`\{Q_{A\alpha },Q_{B,\beta }\}`$ $`=`$ $`2\left[\tau _3{}_{AB}{}^{}P_{\alpha \beta }^{}+ϵ_{AB}M_{\alpha \beta }\right]`$ (2.10)
$`Q_{1\alpha }`$ $`=`$ $`G_{m\overline{m}}+G_{\overline{m}m}`$
$`Q_{2\alpha }`$ $`=`$ $`G_{m\overline{m}}G_{\overline{m}m}`$
$`P_{\alpha \beta }`$ $`=`$ $`G_{mn}\mathrm{\Omega }_{\overline{m}\overline{n}}G_{\overline{m}\overline{n}}\mathrm{\Omega }_{mn}\mathrm{}\mathrm{total}\mathrm{momentum}`$
$`M_{\alpha \beta }`$ $`=`$ $`G_{(mn)}\mathrm{\Omega }_{\overline{m}\overline{n}}+G_{(\overline{m}\overline{n})}\mathrm{\Omega }_{mn}\mathrm{}\mathrm{total}\mathrm{Lorentz}.`$
The right hand side of (2.10) can not be diagonalized by the real SO(2) rotation of $`Q_A`$’s because of the total Lorentz charge term with $`ϵ_{AB}`$. On the other hand the local AdS supersymmetry algebra is given by
$`\{d_{A\alpha },d_{B,\beta }\}`$ $`=`$ $`2\left[\tau _3{}_{AB}{}^{}\stackrel{~}{p}_{\alpha \beta }^{}+ϵ_{AB}m_{\alpha \beta }\right]`$ (2.11)
$`d_{1\alpha }`$ $`=`$ $`D_{a\overline{a}}+\overline{D}_{\overline{a}a}`$
$`d_{2\alpha }`$ $`=`$ $`D_{a\overline{a}}\overline{D}_{\overline{a}a}`$
$`\stackrel{~}{p}_{\alpha \beta }`$ $`=`$ $`𝐃_{ab}\mathrm{\Omega }_{\overline{a}\overline{b}}\overline{𝐃}_{\overline{a}\overline{b}}\mathrm{\Omega }_{ab}\mathrm{}\mathrm{local}\mathrm{LI}\mathrm{momentum}`$
$`m_{\alpha \beta }`$ $`=`$ $`𝐃_{(ab)}\mathrm{\Omega }_{\overline{a}\overline{b}}+\overline{𝐃}_{(\overline{a}\overline{b})}\mathrm{\Omega }_{ab}\mathrm{}\mathrm{local}\mathrm{Lorentz}.`$
In our coset approach the local Lorentz generator is a constraint (2.8), so the local supercovariant derivative $`d_{A\alpha }`$’s can be separated into;
$`\{d_{1\alpha },d_{2\beta }\}=2m_{\alpha \beta }0,\{d_{1\alpha },d_{1\beta }\}=2\stackrel{~}{p}_{\alpha \beta },\{d_{2\alpha },d_{2\beta }\}=2\stackrel{~}{p}_{\alpha \beta }`$ (2.12)
Although the global superalgebra can not be separated into irreducible algebras in the AdS background, the local superalgebra can be separated into irreducible sets on the GL(4$``$4) covariant coset approach. This property allows simpler description of the AdS superstring as the flat case at least in the classical mechanics level.
## 3 AdS Superparticle
We begin with the action for a superparticle in the AdS$`{}_{5}{}^{}\times `$S<sup>5</sup>
$`S={\displaystyle 𝑑\tau \frac{1}{2e}\left\{𝐉_\tau ^{ab}𝐉_{\tau ,ab}+\overline{𝐉}_\tau ^{\overline{a}\overline{b}}\overline{𝐉}_{\tau ,\overline{a}\overline{b}}\right\}}`$ $`.`$ (3.1)
Here we omit the upper-subscript $`L`$ for the LI currents and their components are denoted as
$`(J_\mu ^L)_A{}_{}{}^{B}=\left(\begin{array}{cc}𝐉_{\mu ,}{}_{a}{}^{}^b& j_{\mu ,}{}_{a}{}^{}^{\overline{b}}\\ \overline{j}_{\mu ,}{}_{\overline{a}}{}^{}^b& \overline{𝐉}_{\mu ,}{}_{\overline{a}}{}^{}^{\overline{b}}\end{array}\right).`$ (3.4)
From the definition of the canonical conjugates, $`\mathrm{\Pi }_A{}_{}{}^{M}=\delta S/\delta _\tau Z_M{}_{}{}^{A}()_{}^{A}`$, we have the following primary constraints
$`𝒜_\mathrm{P}={\displaystyle \frac{1}{2}}\mathrm{tr}\left[𝐃^2\overline{𝐃}^2\right]=0,D_{a\overline{b}}=\overline{D}_{\overline{a}b}=0`$ (3.5)
with
$`D_A{}_{}{}^{B}=\left(\begin{array}{cc}𝐃_a^b& D_a^{\overline{b}}\\ \overline{D}_{\overline{a}}^b& \overline{𝐃}_{\overline{a}}^{\overline{b}}\end{array}\right).`$ (3.8)
The Hamiltonian is chosen as
$`=𝒜_\mathrm{P}={\displaystyle \frac{1}{2}}\mathrm{tr}\left[𝐃^2\overline{𝐃}^2\right]`$ (3.9)
and the $`\tau `$-derivative is determined by the Poisson bracket with $``$, $`_\tau 𝒪=\{𝒪,\}`$. The fact that a half of the fermionic constraints is second class requires the Dirac bracket in general. Fortunately the Dirac bracket with the Hamiltonian is equal to its Poisson bracket because the fermionic constrains are $``$ invariant.
The LI current is calculated as
$`J_\tau ^L=Z^1_\tau Z=\left(\begin{array}{cc}𝐃& 0\\ 0& \overline{𝐃}\end{array}\right),_\tau J^L=0.`$ (3.12)
The RI current, generating the global GL(4$``$4) symmetry, is given as
$`J_\tau ^RZ\mathrm{\Pi }=Z\left(J_\tau ^L+A_\tau \right)Z^1,A_\tau =\left(\begin{array}{cc}(𝐃)\frac{1}{4}\mathrm{\Omega }\mathrm{tr}𝐃& D\\ \overline{D}& (\overline{𝐃})\frac{1}{4}\mathrm{\Omega }\mathrm{tr}\overline{𝐃}\end{array}\right).`$ (3.15)
Although the stability group does not contain fermionic components originally, the fermionic components of the gauge connection $`A`$ in (3.15) is induced. It is noted that “$`A`$” is the gauge connection distinguishing from the reparametrization constraint “$`𝒜`$”. The RI current is conserved, since the Hamiltonian is written by LI currents which are manifestly global symmetry invariant
$`_\tau J^R=0.`$ (3.16)
The $`\kappa `$ symmetry generators are half of the fermionic constraints by projecting out with the null vector as
$`_\mathrm{P}{}_{a}{}^{}{}_{}{}^{\overline{b}}=𝐃_a{}_{}{}^{b}D_{b}^{}{}_{}{}^{\overline{b}}+D_a{}_{}{}^{\overline{a}}\overline{𝐃}_{\overline{a}}^{}{}_{}{}^{\overline{b}},\overline{}_\mathrm{P}{}_{\overline{a}}{}^{}{}_{}{}^{b}=\overline{𝐃}_{\overline{a}}{}_{}{}^{\overline{b}}\overline{D}_{\overline{b}}^{}{}_{}{}^{b}+\overline{D}_{\overline{a}}{}_{}{}^{a}𝐃_{a}^{}{}_{}{}^{b}.`$ (3.17)
If we construct the closed algebra including these $`\kappa `$ generators with keeping the bilinear of the fermionic constraints, the $`\tau `$-reparametrization constraint, $`𝒜_\mathrm{P}`$, is modified to
$`\stackrel{~}{𝒜}_\mathrm{P}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{tr}\left[𝐃^2\overline{𝐃}^2+2D\overline{D}\right].`$ (3.18)
This expression appears in the Poisson bracket of $``$ with $`\overline{}`$, when we keep the bilinear of fermionic constrains. The RR flux is responsible for the last term “$`D\overline{D}`$”. The term which is bilinear of the constraints does not change the Poisson bracket since its bracket with an arbitrary variable gives terms proportional to the constraints which are zero on the constrained surface. In another word $`𝒜_\mathrm{P}`$ has an ambiguity of bilinear of the constraints, and the $`\kappa `$ invariance fixes it. On the original coset constrained surface (2.8) it is also rewritten as
$`\stackrel{~}{𝒜}_\mathrm{P}={\displaystyle \frac{1}{2}}\mathrm{Str}[D_A{}_{}{}^{B}]^2={\displaystyle \frac{1}{2}}\mathrm{Str}[J_\tau ^R]^2.`$ (3.19)
This is zero-mode contribution of the classical Virasoro constraint for a superstring in the AdS$`{}_{5}{}^{}\times `$S<sup>5</sup> background.
## 4 AdS Superstring
### 4.1 Conserved currents
We take the action for a superstring in the AdS$`{}_{5}{}^{}\times `$S<sup>5</sup> given by
$`S={\displaystyle d^2\sigma \frac{1}{2}\left\{\sqrt{g}g^{\mu \nu }(𝐉_\mu ^{ab}𝐉_{\nu ,ab}\overline{𝐉}_\mu ^{\overline{a}\overline{b}}\overline{𝐉}_{\nu ,\overline{a}\overline{b}})+\frac{k}{2}ϵ^{\mu \nu }(E^{1/2}j_\mu ^{a\overline{b}}j_{\nu ,a\overline{b}}E^{1/2}\overline{j}_\mu ^{\overline{a}b}\overline{j}_{\nu ,\overline{a}b})\right\}}`$
(4.1)
where “$`k`$” represents the WZ term contribution with $`k=1`$ and $`E=\mathrm{sdet}Z_M^A`$. The consistent $`\tau `$ and $`\sigma `$ reparametrization generators are
$`𝒜_{}`$ $`=`$ $`𝒜_0+k\mathrm{tr}\left[E^{1/4}Fj_\sigma +E^{1/4}\overline{F}\overline{j}_\sigma \right]`$
$`𝒜_{}`$ $`=`$ $`𝒜_0+k\mathrm{tr}\left[E^{1/4}F\overline{j}_\sigma E^{1/4}\overline{F}j_\sigma \right]`$ (4.2)
with the following primary constraints
$`𝒜_0`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{tr}\left[(𝐃^2+𝐉_\sigma ^2)(\overline{𝐃}^2+\overline{𝐉}_\sigma ^2)\right]=0`$
$`𝒜_0`$ $`=`$ $`\mathrm{tr}\left[𝐃𝐉_\sigma \overline{𝐃}\overline{𝐉}_\sigma \right]=0`$ (4.3)
$`F_{a\overline{b}}`$ $`=`$ $`E^{1/4}D_{a\overline{b}}+{\displaystyle \frac{k}{2}}E^{1/4}(\overline{j}_\sigma )_{\overline{b}a}=0`$
$`\overline{F}_{\overline{a}b}`$ $`=`$ $`E^{1/4}\overline{D}_{\overline{a}b}+{\displaystyle \frac{k}{2}}E^{1/4}(j_\sigma )_{b\overline{a}}=0.`$ (4.4)
Their Poisson brackets are
$`\{𝒜_{}(\sigma ),𝒜_{}(\sigma ^{})\}`$ $`=`$ $`2𝒜_{}(\sigma )_\sigma \delta (\sigma \sigma ^{})+_\sigma 𝒜_{}(\sigma )\delta (\sigma \sigma ^{})`$
$`\{𝒜_{}(\sigma ),𝒜_{}(\sigma ^{})\}`$ $`=`$ $`2𝒜_{}(\sigma )_\sigma \delta (\sigma \sigma ^{})+_\sigma 𝒜_{}(\sigma )\delta (\sigma \sigma ^{})`$ (4.5)
$`\{𝒜_{}(\sigma ),𝒜_{}(\sigma ^{})\}`$ $`=`$ $`2𝒜_{}(\sigma )_\sigma \delta (\sigma \sigma ^{})+_\sigma 𝒜_{}(\sigma )\delta (\sigma \sigma ^{}).`$
The Hamiltonian is chosen as
$``$ $`=`$ $`{\displaystyle 𝑑\sigma 𝒜_{}}`$
$`=`$ $`{\displaystyle 𝑑\sigma \mathrm{tr}\left[\frac{1}{2}\left\{𝐃^2+𝐉_\sigma ^2\overline{𝐃}^2\overline{𝐉}_\sigma ^2\right\}+\left(kE^{1/2}\overline{D}\overline{j}_\sigma kE^{1/2}Dj_\sigma +j_\sigma \overline{j}_\sigma \right)\right]}.`$
From now on $`E=1`$ gauge is taken using the local GL(1) invariance. The global GL(1) symmetry is broken by the WZ term.
Using the Hamiltonian in (4.1) the $`\tau `$-derivative of $`𝒜_{}`$ and $`𝒜_{}`$ are given as
$`_\tau 𝒜_{}=_\sigma 𝒜_{},_\tau 𝒜_{}=_\sigma 𝒜_{}.`$ (4.7)
Although the coset parameter $`Z_M^A`$ does not satisfy the world-sheet free wave equation, it is essential to introduce the world-sheet lightcone coordinate
$`\sigma ^\pm =\tau \pm \sigma ,_\pm ={\displaystyle \frac{1}{2}}(_\tau \pm _\sigma ).`$ (4.8)
The differential equations (4.7) are rewritten as
$`_{}𝒜_+=0,_+𝒜_{}=0,𝒜_\pm =𝒜_{}\pm 𝒜_{},`$ (4.9)
so the infinite number of the conserved currents are
$`_{}\left[f(\sigma ^+)𝒜_+\right]=0,_+\left[f(\sigma ^{})𝒜_{}\right]=0`$ (4.10)
with an arbitrary function $`f`$. Then there exist infinite number of conserved charges
$`_{}\left[{\displaystyle 𝑑\sigma f(\sigma ^+)𝒜_+}\right]=0,_+\left[{\displaystyle 𝑑\sigma f(\sigma ^{})𝒜_{}}\right]=0.`$ (4.11)
On the other hand the integrability of the superstring will provide the infinite number of “local” charges as well as the “non-local” charges written down in . The LI currents is given by
$`\{\begin{array}{ccc}J_\tau ^L& =& \left(\begin{array}{cc}𝐃& k\overline{j}_\sigma \\ kj_\sigma & \overline{𝐃}\end{array}\right)=\left(\begin{array}{cc}𝐃& 2D2F\\ 2\overline{D}2\overline{F}& \overline{𝐃}\end{array}\right)\left(\begin{array}{cc}𝐃& 2D\\ 2\overline{D}& \overline{𝐃}\end{array}\right)\hfill \\ J_\sigma ^L& =& \left(\begin{array}{cc}𝐉_\sigma & j_\sigma \\ \overline{j}_\sigma & \overline{𝐉}_\sigma \end{array}\right)\hfill \end{array}`$ (4.22)
where the $`\tau `$ component is determined by (4.1). The LI currents satisfy the flatness condition but does not satisfy the conservation law. The RI currents are obtained in as
$`\{\begin{array}{ccc}J_\tau ^R& =& ZDZ^1=Z(J_\tau ^L+A_\tau )Z^1\hfill \\ J_\sigma ^R& =& Z(J_\sigma ^L+A_\sigma )Z^1,J_\sigma ^L+A_\sigma =\left(\begin{array}{cc}𝐉_\sigma & \overline{F}+\frac{1}{2}j_\sigma \\ F+\frac{1}{2}\overline{j}_\sigma & \overline{𝐉}_\sigma \end{array}\right)\hfill \end{array}`$ (4.27)
where the gauge connection $`A_\mu `$ is
$`\{\begin{array}{ccc}A_\tau & =& \left(\begin{array}{cc}(𝐃)\frac{1}{4}\mathrm{\Omega }\mathrm{tr}𝐃& D\\ \overline{D}& (\overline{𝐃})\frac{1}{4}\mathrm{\Omega }\mathrm{tr}\overline{𝐃}\end{array}\right)\hfill \\ A_\sigma & =& \left(\begin{array}{cc}(𝐉_\sigma )+\frac{1}{4}\mathrm{\Omega }\mathrm{tr}𝐉_\sigma & \overline{F}\frac{1}{2}j_\sigma \\ F\frac{1}{2}\overline{j}_\sigma & (\overline{𝐉}_\sigma )+\frac{1}{4}\mathrm{\Omega }\mathrm{tr}\overline{𝐉}_\sigma \end{array}\right)\hfill \end{array}.`$ (4.34)
The fermionic components of $`A_\mu `$ appear again. In this paper the fermionic constraints, $`F`$ and $`\overline{F}`$, in the fermionic components of $`A_\sigma `$ are kept while they were absent in our previous paper depending on the treatment of the constraint bilinear terms. Then the integrability of the superstring leads to the current conservation and the flatness condition for the RI current;
$`_\tau J_\tau ^R=_\sigma J_\sigma ^R,_\tau J_\sigma ^R_\sigma J_\tau ^R=2[J_\tau ^R,J_\sigma ^R].`$ (4.35)
They are rewritten as
$`_{}J_+^R=[J_{}^R,J_+^R],_+J_{}^R=[J_+^R,J_{}^R],J_\pm ^R=J_\tau ^R\pm J_\sigma ^R.`$ (4.36)
Taking the supertrace, denoting “Str”, leads to the infinite number of conserved “local” currents because $`J_\mu ^R`$ are supermatrices,
$`_{}\mathrm{Str}\left[(J_+^R)^n\right]=0,_+\mathrm{Str}\left[(J_{}^R)^n\right]=0,n=1,2,\mathrm{}.`$ (4.37)
It gives the infinite number of conserved “local” charges
$`_\tau \left[{\displaystyle 𝑑\sigma f(\sigma ^+)\mathrm{Str}(J_+^R)^n}\right]=0,_\tau \left[{\displaystyle 𝑑\sigma f(\sigma ^{})\mathrm{Str}(J_{}^R)^n}\right]=0.`$ (4.38)
In this way classical 2-dimensional conformal symmetry and integrability of AdS superstring lead to two infinite sets of conserved currents, (4.10) and (4.37). In next sections the relation between them is examined.
### 4.2 Stress-energy tensor ($`n=2`$)
The “$`+/`$” (right/left moving) modes of the RI currents on the original coset constrained space (2.8) are written as
$`J_\pm ^R=Z\left(\begin{array}{cc}𝐃_\pm & D\pm (\overline{F}+\frac{1}{2}j_\sigma )\\ \overline{D}\pm (F+\frac{1}{2}\overline{j}_\sigma )& \overline{𝐃}_\pm \end{array}\right)Z^1=Z\left(\begin{array}{cc}𝐃_\pm & d_\pm +\frac{1}{2}j_\pm \\ \pm (d_\pm \frac{1}{2}j_\pm )& \overline{𝐃}_\pm \end{array}\right)Z^1`$ (4.43)
with
$`𝐃_\pm =𝐃\pm 𝐉_\sigma ,\overline{𝐃}_\pm =\overline{𝐃}\pm \overline{𝐉}_\sigma ,d_\pm =F\pm \overline{F},j_\pm =j_\tau \pm j_\sigma =\overline{j}_\sigma \pm j_\sigma `$ (4.45)
carrying the LI currents indices, $`AB`$. This is supertraceless, Str$`J_\pm ^R=0`$, so $`n=1`$ case of (4.37) gives just trivial equation.
Let us look at the $`n=2`$ case of (4.37), $`\mathrm{Str}\left[(J_\pm ^R)^2\right]`$. Then the “+” sector is written as
$`{\displaystyle \frac{1}{2}}\mathrm{Str}\left[(J_+^R)^2\right]`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{Str}\left[\left(\begin{array}{cc}𝐃_+& d_++\frac{1}{2}j_+\\ d_+\frac{1}{2}j_+& \overline{𝐃}_+\end{array}\right)^2\right]`$
$`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{tr}\left[𝐃_+^2\overline{𝐃}_+^2+2(d_++{\displaystyle \frac{1}{2}}j_+)(d_+{\displaystyle \frac{1}{2}}j_+)\right]`$
$`=`$ $`\mathrm{tr}\left[{\displaystyle \frac{1}{2}}\left(𝐃_+^2\overline{𝐃}_+^2\right)+j_+d_+\right].`$ (4.49)
The “$``$” sector is
$`{\displaystyle \frac{1}{2}}\mathrm{Str}\left[(J_{}^R)^2\right]`$ $`=`$ $`\mathrm{tr}\left[{\displaystyle \frac{1}{2}}\left(𝐃_{}^2\overline{𝐃}_{}^2\right)j_{}d_{}\right].`$ (4.50)
On the other hand the conformal symmetry generator $`𝒜_\pm `$ is rewritten from the relation (4.2) and (4.45) as
$`𝒜_\pm `$ $`=`$ $`\mathrm{tr}\left[{\displaystyle \frac{1}{2}}\left(𝐃_\pm ^2\overline{𝐃}_\pm ^2\right)\pm j_\pm d_\pm \right]={\displaystyle \frac{1}{2}}\mathrm{Str}\left[(J_\pm ^R)^2\right].`$ (4.51)
If we take care of the square of the fermionic constraints, the closure of the first class constraint set including the $`\kappa `$ symmetry generators,
$`_\pm `$ $`=`$ $`𝐃_\pm d_\pm +d_\pm \overline{𝐃}_\pm `$ (4.52)
determines the ambiguity of bilinear of the constraints as
$`\stackrel{~}{𝒜}_\pm =\mathrm{tr}\left[{\displaystyle \frac{1}{2}}\left(𝐃_\pm ^2\overline{𝐃}_\pm ^2\right)\pm ({\displaystyle \frac{1}{2}}d_{}+j_\pm )d_\pm \right]=𝒜_\pm +\mathrm{tr}F\overline{F}`$ (4.53)
obtained in as a generator of the $`𝒜𝒞𝒟`$ constraint set known to exist for a superstring in a flat space . Then the stress-energy tensor is
$`T_{\pm \pm }\stackrel{~}{𝒜}_\pm 𝒜_\pm =\mathrm{Str}J_\pm ^RJ_\pm ^R.`$ (4.54)
This is $`\kappa `$ symmetric stress-energy tensor in a supersymmetric generalization of Sugawara form.
### 4.3 Supercovariant derivative algebra
Existence of the conformal invariance should present the irreducible coset components of supercovariant derivatives ;
$`𝐃_\pm `$ $`=`$ $`𝐃\pm 𝐉_\sigma ,\overline{𝐃}_\pm =\overline{𝐃}\pm \overline{𝐉}_\sigma `$
$`d_\pm `$ $`=`$ $`F\pm \overline{F}=(D\pm {\displaystyle \frac{1}{2}}j_\sigma )\pm (\overline{D}\pm {\displaystyle \frac{1}{2}}\overline{j}_\sigma ).`$
On the constraint surface (2.8) and (4.4) the $`+/`$ sector supercovariant derivatives are separated as
$`\{𝐃_+_{ab}(\sigma ),𝐃_{}_{cd}(\sigma ^{})\}`$ $`=`$ $`2\mathrm{\Omega }_{c|b}(𝐃)_{a|d}\delta (\sigma \sigma ^{})0`$
$`\{𝐃_+_{ab}(\sigma ),d_{,c\overline{d}}(\sigma ^{})\}`$ $`=`$ $`\mathrm{\Omega }_{cb}d_{+,a\overline{d}}\delta (\sigma \sigma ^{})=\mathrm{\Omega }_{cb}(F+\overline{F})_{a\overline{d}}\delta (\sigma \sigma ^{})0`$
$`\{d_{+,a\overline{b}}(\sigma ),d_{,c\overline{d}}(\sigma ^{})\}`$ $`=`$ $`2\left[\mathrm{\Omega }_{ac}(\overline{𝐃})_{\overline{b}\overline{d}}+\mathrm{\Omega }_{\overline{b}\overline{d}}(𝐃)_{ac}\right]\delta (\sigma \sigma ^{})0`$
with analogous relation for the barred sector, $`\overline{𝐃}_\pm `$.
The “+” sector supercovariant derivative algebra is
$`\{𝐃_+_{ab}(\sigma ),𝐃_+_{cd}(\sigma ^{})\}`$ $`=`$ $`2\mathrm{\Omega }_{c|b}\mathrm{\Omega }_{a|d}\delta ^{}(\sigma \sigma ^{})+4\mathrm{\Omega }_{c|b}(𝐉_\sigma )_{a|d}\delta (\sigma \sigma ^{})`$
$``$ $`2\mathrm{\Omega }_{c|b}_{a|d}\delta (\sigma \sigma ^{})`$
$`\{d_{+,a\overline{b}}(\sigma ),d_{+,c\overline{d}}(\sigma ^{})\}`$ $`=`$ $`2\left[\mathrm{\Omega }_{\overline{b}\overline{d}}𝐃_+_{ac}\mathrm{\Omega }_{ac}\overline{𝐃}_+_{\overline{b}\overline{d}}\right]\delta (\sigma \sigma ^{})`$
$`\{𝐃_+_{ab}(\sigma ),d_{+,c\overline{d}}(\sigma ^{})\}`$ $`=`$ $`\mathrm{\Omega }_{cb}(d_{}+2j_+)_{a\overline{d}}\delta (\sigma \sigma ^{})2\mathrm{\Omega }_{cb}\omega _{+,a\overline{d}}\delta (\sigma \sigma ^{})`$
$`\{d_{+,a\overline{b}}(\sigma ),\omega _{+,c\overline{d}}(\sigma ^{})\}`$ $`=`$ $`2\mathrm{\Omega }_{\overline{b}\overline{d}}\mathrm{\Omega }_{ac}\delta ^{}(\sigma \sigma ^{})2\left[\mathrm{\Omega }_{\overline{b}\overline{d}}(𝐉_\sigma )_{ac}\mathrm{\Omega }_{ac}(\overline{𝐉}_\sigma )_{\overline{b}\overline{d}}\right]\delta (\sigma \sigma ^{})`$
$``$ $`2_{\overline{b}\overline{d};ac}\delta (\sigma \sigma ^{})`$
$`\{𝐃_+_{ab}(\sigma ),\omega _{+,c\overline{d}}(\sigma ^{})\}`$ $`=`$ $`\mathrm{\Omega }_{cb}\omega _{,a\overline{d}}\delta (\sigma \sigma ^{})`$
$`\{\omega _{+,a\overline{b}}(\sigma ),\omega _{+,c\overline{d}}(\sigma ^{})\}`$ $`=`$ $`0`$
where
$`\omega _\pm `$ $`=`$ $`j_\pm =\overline{j}_\sigma \pm j_\sigma .`$ (4.56)
This is comparable with the flat case where the non-local term, $`_\sigma \delta (\sigma \sigma ^{})`$, is replaced by the local Lorentz ( \[Sp(4)\]<sup>2</sup> ) covariant non-local term, $`_\sigma \delta (\sigma \sigma ^{})`$. For the fifth Poisson bracket, $`\{𝐃_+,\omega \}`$, it is zero for the flat case but it is not for the AdS case. For a superstring in a flat space the consistency of the $`\kappa `$ symmetry constraint requires the first class constraint set, namely “$`𝒜𝒞𝒟`$” constraint, which are bilinear of the supercovariant derivatives . For the AdS case the situation is completely the same, despite of this anomalous term .
### 4.4 “Local” currents ($`n3`$)
Next let us look at $`n3`$ cases of the infinite number of conserved “local” current (4.37). For simplicity we focus on the “+” sector and replace $`\mathrm{`}\mathrm{`}+\mathrm{"}`$ by $`\mathrm{`}\mathrm{`}\widehat{}\mathrm{"}`$, as $`J_+\widehat{J}`$. The first three powers of the RI current, $`(J^R)^n`$ with $`n=1,2,3`$, are listed as below:
$`\left[Z^1\widehat{J}^RZ\right]_{AB}`$ $`=`$ $`\left(\begin{array}{cc}\widehat{𝐃}_{ab}& (\widehat{d}+\frac{1}{2}\widehat{j})_{a\overline{b}}\\ \pm (\widehat{d}\frac{1}{2}\widehat{j})_{b\overline{a}}& \widehat{\overline{𝐃}}_{\overline{a}\overline{b}}\end{array}\right)`$ (4.59)
$`\left[Z^1(\widehat{J}^R)^2Z\right]_{AB}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left(\begin{array}{cc}\mathrm{\Omega }_{ab}\mathrm{tr}(\widehat{𝐃}^2+\widehat{j}\widehat{d})& \\ & \mathrm{\Omega }_{\overline{a}\overline{b}}\mathrm{tr}(\widehat{\overline{𝐃}}^2\widehat{j}\widehat{d})\end{array}\right)`$ (4.65)
$`+\left(\begin{array}{cc}(\widehat{d}^2\frac{1}{4}\widehat{j}^2)_{(ab)}+\widehat{j}\widehat{d}_{ab}& \widehat{}_{a\overline{b}}+\frac{1}{2}(\widehat{𝐃}\widehat{j}+\widehat{j}\widehat{\overline{𝐃}})_{a\overline{b}}\\ \widehat{}_{b\overline{a}}\frac{1}{2}(\widehat{𝐃}\widehat{j}+\widehat{j}\widehat{\overline{𝐃}})_{b\overline{a}}& (\widehat{d}^2\frac{1}{4}\widehat{j}^2)_{(\overline{a}\overline{b})}\widehat{j}\widehat{d}_{\overline{a}\overline{b}}\end{array}\right)`$
$`\left[Z^1(\widehat{J}^R)^3Z\right]_{AB}={\displaystyle \frac{1}{4}}\left(\begin{array}{cc}\mathrm{\Omega }_{ab}\mathrm{tr}\left[\widehat{}\widehat{j}(\widehat{𝐃}\widehat{d})\widehat{j}\right]& \\ & \mathrm{\Omega }_{\overline{a}\overline{b}}\mathrm{tr}\left[\widehat{}\widehat{j}(\widehat{d}\widehat{\overline{𝐃}})\widehat{j}\right]\end{array}\right)`$ (4.68)
$`\left(\begin{array}{cc}\left[\frac{1}{4}\mathrm{tr}(\widehat{𝐃}^2+\widehat{j}\widehat{d})\widehat{𝐃}\widehat{𝐃}(\widehat{j}\widehat{d})+\widehat{}\widehat{j}\right]_{ab}& \frac{1}{4}\mathrm{tr}(\widehat{𝐃}^2\widehat{\overline{𝐃}}^2)(\widehat{d}+\frac{1}{2}\widehat{j})_{a\overline{b}}\\ \frac{1}{4}\mathrm{tr}(\widehat{𝐃}^2\widehat{\overline{𝐃}}^2)(\widehat{d}\frac{1}{2}\widehat{j})_{b\overline{a}}& \left[\frac{1}{4}\mathrm{tr}(\widehat{\overline{𝐃}}^2+\widehat{j}\widehat{d})\widehat{\overline{𝐃}}(\widehat{j}\widehat{d})\widehat{\overline{𝐃}}\widehat{}\widehat{j}\right]_{\overline{a}\overline{b}}\end{array}\right)`$ (4.71)
$`+\left(\begin{array}{cc}\left[2(\widehat{d}^2\frac{1}{4}\widehat{j}^2)\widehat{𝐃}+\widehat{d}\widehat{\overline{𝐃}}\widehat{d}\frac{1}{4}\widehat{j}\widehat{\overline{𝐃}}\widehat{j}\right]_{(ab)}& \frac{1}{4}\mathrm{tr}(\widehat{j}\widehat{d})(\widehat{d}+\frac{1}{2}\widehat{j})_{a\overline{b}}+\left[\widehat{𝐃}(\widehat{d}+\frac{1}{2}\widehat{j})\widehat{\overline{𝐃}}\right]_{a\overline{b}}\\ \frac{1}{4}\mathrm{tr}(\widehat{j}\widehat{d})(\widehat{d}\frac{1}{2}\widehat{j})_{b\overline{a}}+\left[\widehat{𝐃}(\widehat{d}\frac{1}{2}\widehat{j})\widehat{\overline{𝐃}}\right]_{b\overline{a}}& \left[2(\widehat{d}^2\frac{1}{4}\widehat{j}^2)\widehat{\overline{𝐃}}+\widehat{d}\widehat{𝐃}\widehat{d}\frac{1}{4}\widehat{j}\widehat{𝐃}\widehat{j}\right]_{(\overline{a}\overline{b})}\end{array}\right)`$ (4.74)
$`+\left(\begin{array}{cc}& \left[\left\{(\widehat{d}^2\frac{1}{4}\widehat{j}^2)+\widehat{j}\widehat{d}\right\}(\widehat{d}+\frac{1}{2}\widehat{j})\right]_{a\overline{b}}\\ \left[\left\{(\widehat{d}^2\frac{1}{4}\widehat{j}^2)\widehat{j}\widehat{d}\right\}(\widehat{d}\frac{1}{2}\widehat{j})\right]_{\overline{a}b}& \end{array}\right)`$ (4.77)
In this computation 5-dimensional $`\gamma `$-matrix relations are used, for example $`𝐕^{ab}𝐔_{bc}+𝐔^{ab}𝐕_{bc}=\frac{1}{2}\delta _c^a\mathrm{tr}\mathrm{𝐕𝐔}`$ for bosonic vectors $`𝐕,𝐔`$.
The conserved “local” current with $`n=3`$ becomes
$`\mathrm{Str}(\widehat{J}^R)^3`$ $`=`$ $`\mathrm{tr}\left[2\widehat{}\widehat{j}(\widehat{𝐃}\widehat{d})\widehat{j}(\widehat{d}\widehat{\overline{𝐃}})\widehat{j}\right]=\mathrm{tr}(\widehat{}\widehat{j})`$ (4.78)
where $`\widehat{}`$ is the $`\kappa `$ generating constraint (4.52). The conserved “local” current with $`n=4`$ becomes
$`\mathrm{Str}(\widehat{J}^R)^4`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{tr}\left(\widehat{𝐃}^2+\widehat{\overline{𝐃}}^2\right)\widehat{𝒜}+\left(\mathrm{}\right)\mathrm{tr}(\widehat{}\widehat{j}).`$ (4.79)
The conserved “local” current with $`n=5,6`$ are given as; Str$`(\widehat{J}^R)^5=`$( $`\widehat{}`$ dependent terms), Str$`(\widehat{J}^R)^6=`$( $`\widehat{𝒜}`$ and $`\widehat{}`$ dependent terms). In general for even $`n=2m`$ its bosonic part is given as
$`\mathrm{Str}(\widehat{J}^R)^{2m}_{\mathrm{bosonic}}`$ $`=`$ $`\left(\mathrm{tr}\widehat{𝐃}^2\right)^m\left(\mathrm{tr}\widehat{\overline{𝐃}}^2\right)^m`$ (4.80)
$`=`$ $`\mathrm{tr}\left(\widehat{𝐃}^2\widehat{\overline{𝐃}}^2\right)\left\{\left(\mathrm{tr}\widehat{𝐃}^2\right)^{m1}+\mathrm{}+\left(\mathrm{tr}\widehat{\overline{𝐃}}^2\right)^{m1}\right\}`$
$``$ $`(\mathrm{})\widehat{𝒜}+\left(\mathrm{}\right)\mathrm{tr}(\widehat{}\widehat{j})`$
where the last equality is guaranteed by the $`\kappa `$ invariance. It is also pointed out that the conserved supertraces of multilinears in the currents factorize in traces of lower number of currents and that for an even number of currents one of the factors is the stress tensor in . For odd $`n=2m+1`$ its bosonic part is given as
$`\mathrm{Str}(\widehat{J}^R)^{2m+1}_{\mathrm{bosonic}}=0\left(\mathrm{}\right)\mathrm{tr}(\widehat{}\widehat{j})`$ (4.81)
where the possible fermionic variable dependence is a term proportional to $`\widehat{}`$ guaranteed by the $`\kappa `$ invariance.
In this way, after taking supertrace the even $`n`$-th power of $`J^R`$ reduces terms proportional to $`𝒜`$ and $``$, and the odd $`n`$-th power of $`J^R`$ reduces a term proportional to $``$ only. In this paper $`𝒞𝒟`$ constraints in the $`𝒜𝒞𝒟`$ first class constraint set are not introduced for simpler argument, and set to zero because they are bilinears of constraints.
## 5 Conclusion and discussions
We obtained the expression of the conserved “local” currents derived from the integrability of a superstring in the AdS$`{}_{5}{}^{}\times `$S<sup>5</sup> background. The infinite number of the conserved “local” currents are written by the supertrace of the $`n`$-th power of the RI currents. The lowest nontrivial case, $`n=2`$, is nothing but the stress-energy tensor which is also Virasoro constraint, Str$`(J^R)_\pm ^2`$ in (4.49) and (4.50). For even $`n`$ the “local” current reduces to terms proportional to the Virasoro constraint and the $`\kappa `$ symmetry constraint. For odd $`n`$ it reduces to a term proportional to the $`\kappa `$ symmetry constraint. In another word the integrability reduces to the $`𝒜(𝒞𝒟)`$ first class constraint set where $`𝒜`$ is the Virasoro generator and $``$ is the $`\kappa `$ symmetry generator. The $`𝒜𝒞𝒟`$ first class constraint set is the local symmetry generator of superstrings both on the flat space and on the AdS space. It is natural that the physical degrees of freedom of a superstring is common locally, independently of flat or AdS backgrounds.
It seems that the combination of the $`_\pm j_\pm `$ in (4.78) plays the role of the world-sheet supersymmetry operator in a sense of the grading of the conformal generator. However it is not straightforward to construct the worldsheet supersymmetry operator. As in the flat case where the lightcone gauge makes the relation between the GS fermion and the NSR fermion more transparent, the $`\kappa `$ gauge fixing will be a clue to make a connection to the world-sheet supersymmetry. We leave this problem in addition to the quantization problem for future investigations.
Acknowledgments
The author thanks to K. Kamimura, S. Mizoguchi and K. Yoshida for fruitful discussions. |
warning/0507/physics0507055.html | ar5iv | text | # Introduction to thermodynamics of spin models in the Hamiltonian limit
## 1 Introduction
A rather interesting method for studying thermal fluctuations of $`(d+1)`$-dimensional interacting classical degrees of freedom is provided through the corresponding analysis of quantum fluctuations in a $`d`$-dimensional system of interacting quantum degrees of freedom. The correspondence is detailed in the remarkable review of Kogut . It can be understood as a re-wording of the transfer matrix formalism of classical systems in a special limit where the transfer matrix takes a simplified form, the so-called Hamiltonian limit . At the origin of the correspondence, there is the transition amplitude between quantum states in the path integral formulation,
$$\mathrm{𝙰𝚖𝚙𝚕𝚒𝚝𝚞𝚍𝚎}=\underset{\mathrm{𝚙𝚊𝚝𝚑𝚜}}{}\mathrm{exp}\frac{i}{\mathrm{}}S[x(t)]$$
(1)
where the classical action is a functional of $`x(t)`$, the time integral of a Lagrangian $`S=_{t_a}^{t_b}L(x,\dot{x})𝑑t`$. Here we consider the transition amplitude of a point particle (space dimensionality of the quantum system $`d=0`$) for the sake of simplicity. Equation (1) looks a bit similar to the partition function of something to be specified,
$$𝙿.𝙵.=\underset{\mathrm{𝚌𝚘𝚗𝚏𝚒𝚐𝚞𝚛𝚊𝚝𝚒𝚘𝚗𝚜}}{}\mathrm{exp}(\beta E\{𝚍.𝚘.𝚏.\}),$$
(2)
provided that we change the imaginary argument of the exponential to a real one and we give sense to the degrees of freedom (d.o.f.) and to the sum over their configurations. This transformation is achieved through the definition of an imaginary time, also called Euclidean time,
$$t=i\tau $$
(3)
in terms of which $`L(x,\dot{x})=H(x,p)`$, $`p=mx^{}`$, $`\dot{x}=dx/dt`$ and $`x^{}=dx/d\tau `$. The phase factor of equation (1) becomes, with equation (3),
$$\frac{i}{\mathrm{}}S[x(t)]=\frac{1}{\mathrm{}}_{\tau _a}^{\tau _b}H(x,p)𝑑\tau $$
(4)
where $`H(x,p)=\frac{p^2}{2m}+V(x)`$. Equation (1) has a mathematical sense through discretization of the time axis, $`t_{i+1}t_i=\delta =iϵ`$, $`ϵ0`$, $`t_bt_a=N\delta `$, $`x_i=x(t_i)`$. The action becomes <sup>1</sup><sup>1</sup>1We denote by $`S[x(\tau )]`$ the functional of a continuous function and $`_iH\{x_i\}`$ the corresponding sum over a set of discrete variables.
$$S[x(\tau )]ϵ\underset{i=1}{\overset{N1}{}}\left(\frac{1}{2}m\left(\frac{x_{i+1}x_i}{ϵ}\right)^2+V(x_i)\right)$$
(5)
and the transition amplitude eventually reads as the partition function of a set of $`N`$ classical d.o.f. $`\{x_i\}`$ on a 1-dimensional lattice
$`\mathrm{𝙰𝚖𝚙𝚕𝚒𝚝𝚞𝚍𝚎}`$ $`=`$ $`{\displaystyle \left(\underset{i}{}dx_i\right)\mathrm{exp}(\mathrm{}^1\underset{i}{}H\{x_i\})}`$
$`𝙿.𝙵.`$ $`=`$ $`{\displaystyle \underset{\{x_i\}}{}}\mathrm{exp}(\beta {\displaystyle \underset{i}{}}H\{x_i\}).`$ (6)
$`H\{x_i\}`$ is a classical energy density depending on the set of values taken by the d.o.f., ($`\mathrm{}<x_i<+\mathrm{}`$ if there is no restriction specified in the original quantum problem) and leads after summation to the energy of a configuration appearing in the Boltzmann weight. The role of $`\beta `$ (thermal fluctuations) is played by $`\mathrm{}^1`$ (quantum fluctuations). The classical limit $`\mathrm{}0`$ where quantum fluctuations are suppressed corresponds to $`\beta \mathrm{}`$ in the classical system, i.e. suppression of thermal fluctuations. The correspondence (6) is more generally valid than in the simple case of a point particle. To the statistical physics problem of classical degrees of freedom living in $`d+1`$ space dimensions, there corresponds a quantum problem in $`d`$ dimensions where the fluctuations result from competition between non commutating variables. The partition function of the former problem involves quantum transitions between multi-particle states in the latter formulation and the Hamiltonian limit is the simplest formulation of the quantum problem when only survive transitions between single-particle quantum states.
This approach deserves some attention and might be taught in graduate statistical physics courses. In the following, we remind how the thermodynamic properties of a classical system might be obtained from their quantum counterpart, then we apply the technique to an approximate determination of the critical properties of two-dimensional classical spin models, namely the Ising model and the $`3`$ and $`4`$state Potts models.
## 2 Thermodynamics
In the quantum version of the problem, we may define a time evolution operator in terms of which the Feynman kernel is
$`K(x_b,t_b|x_a,t_a)`$ $`=`$ $`x_b|\mathrm{exp}({\displaystyle \frac{i}{\mathrm{}}}\widehat{\text{H}}(t_bt_a))|x_a`$ (7)
$`=`$ $`{\displaystyle x_b|\widehat{\text{T}}|x_{N1}𝑑x_{N1}x_{N1}|\widehat{\text{T}}|x_{N2}𝑑x_{N2}x_{N2}|\mathrm{}}`$
$`\mathrm{}dx_2x_2|\widehat{\text{T}}|x_1dx_1x_1|\widehat{\text{T}}|x_a`$
$`=`$ $`x_b|\widehat{\text{T}}^N|x_a,`$
where $`\widehat{\text{T}}`$ is the infinitesimal time evolution operator which, in the Euclidean time, becomes the transfer matrix $`\widehat{\text{T}}=\mathrm{exp}(ϵ\widehat{\text{H}}/\mathrm{})`$, $`\delta =iϵ`$. $`Nϵ`$ is the length of the system in the supplementary time direction (see Fig. 1). Summing over initial states, when periodic boundary conditions in the time direction are imposed, $`x_b=x_a=x_0`$, we get the partition function from
$$Z=𝑑x_0K(x_0,iNϵ|x_0,0)=Tr\widehat{\text{T}}^N.$$
(8)
The thermodynamic limit $`N\mathrm{}`$ ensures projection onto the ground state,
$`\widehat{\text{T}}^N`$ $`=`$ $`t_0^N\left[|00|+{\displaystyle \underset{\alpha 0}{}}|\alpha (t_\alpha /t_0)^N\alpha |\right]|0t_0^N0|t_0>t_1>\mathrm{}`$ (9)
The partition function is thus determined by the largest eigenvalue $`t_0`$ of the transfer matrix,
$$Zt_0^N=\mathrm{}^{NϵE_0/\mathrm{}},$$
(10)
where $`E_0`$ is the ground state energy of the quantum Hamiltonian $`\widehat{\text{H}}`$. The free energy density (i.e. per time slice) follows
$$f=\underset{N\mathrm{}}{lim}(Nϵ)^1\mathrm{}\mathrm{ln}ZE_0.$$
(11)
The (time) correlation function of some local quantities $`\varphi _x`$ depending on the classical d.o.f. $`x`$, $`\varphi _{x_i}\varphi _{x_{i+j}}`$, is expressed in terms of the transfer matrix through
$`\varphi _{x_i}\varphi _{x_{i+j}}`$ $`=`$ $`\underset{N\mathrm{}}{lim}Z^1{\displaystyle \underset{\{x_i\}}{}}\varphi _{x_i}\varphi _{x_{i+j}}\mathrm{}^{\beta E\{x_i\}}`$ (12)
$`=`$ $`\underset{N\mathrm{}}{lim}Z^1{\displaystyle x_0|\widehat{\text{T}}|x_{N1}𝑑x_{N1}x_{N1}|\mathrm{}}`$
$`\mathrm{}\widehat{\text{T}}|x_{i+j}\varphi _{x_{i+j}}dx_{i+j}x_{i+j}|\mathrm{}\widehat{\text{T}}|x_i\varphi _{x_i}dx_ix_i|\mathrm{}`$
$`\mathrm{}\widehat{\text{T}}|x_1dx_1x_1|\widehat{\text{T}}|x_0`$
$`=`$ $`\underset{N\mathrm{}}{lim}{\displaystyle \frac{Tr(\widehat{\text{T}}^i\widehat{\mathit{\varphi }}\widehat{\text{T}}^j\widehat{\mathit{\varphi }}\widehat{\text{T}}^{Nij})}{Tr\widehat{\text{T}}^N}},`$
where diagonal operators in the $`\{|x_k\}`$ basis have been introduced,
$$\widehat{\mathit{\varphi }}=𝑑x_k|x_k\varphi _{x_k}x_k|.$$
(13)
Using the eigenstates of the transfer matrix, the correlation function (12) becomes
$$\varphi _{x_i}\varphi _{x_{i+j}}=|0|\widehat{\mathit{\varphi }}|0|^2+\underset{\alpha 0}{}|0|\widehat{\mathit{\varphi }}|\alpha |^2\left(\frac{t_\alpha }{t_0}\right)^j.$$
(14)
The connected part of the correlation function, $`G_\varphi (j)`$, is obtained after subtraction of the ground state expectation value, and in the thermodynamic limit it is dominated by the first eigenstate $`|\beta _\varphi `$ such that the matrix element of $`\widehat{\mathit{\varphi }}`$ between this state and the ground state does not vanish,
$$G_\varphi (j)|0|\widehat{\mathit{\varphi }}|\beta _\varphi |^2\left(\frac{t_\beta }{t_0}\right)^j=|0|\widehat{\mathit{\varphi }}|\beta _\varphi |^2\mathrm{}^{jϵ\mathrm{}^1(E_\beta E_0)}.$$
(15)
The matrix element in prefactor measures the average of the field squared,
$$\mathrm{𝚊𝚟𝚎𝚛𝚊𝚐𝚎}\mathrm{𝚏𝚒𝚎𝚕𝚍}=|0|\widehat{\mathit{\varphi }}|\beta _\varphi |,$$
(16)
and the exponential decay along the time direction allows to define the correlation length in terms of an inverse gap ($`\mathrm{}=1`$),
$$\frac{1}{\xi _\varphi }=E_\beta E_0=\mathrm{𝚐𝚊𝚙}_\varphi .$$
(17)
A challenging problem in critical phenomena is the identification of the universality class of a given model. Due to scaling relations among critical exponents, the knowledge of two of them determines the whole set of exponents (see table 1). From finite size-scaling (FSS) behaviour of the critical densities , we expect power law behaviours of local quantities in terms of the finite size $`L`$ in the space direction,
$`\mathrm{𝙴𝚗𝚎𝚛𝚐𝚢}\mathrm{𝚍𝚎𝚗𝚜𝚒𝚝𝚢}L^{x_ϵ},`$ (18)
$`\mathrm{𝙾𝚛𝚍𝚎𝚛}\mathrm{𝚙𝚊𝚛𝚊𝚖𝚎𝚝𝚎𝚛}L^{x_\sigma }.`$ (19)
These relations may be taken as definitions of the scaling dimensions of the energy and order parameter density, $`x_ϵ=(1\alpha )/\nu `$ and $`x_\sigma =\beta /\nu `$. The commonly accepted notation for the critical exponents is reminded in table 1.
Imagine that we consider a more general quantum problem in $`d1`$-space dimensions. The classical counterpart is defined in $`d`$-dimensions (space and time) and the susceptibility for example follows from the integration of the correlation function over these $`d`$ dimensions. Let us write schematically what happens in the case of a translation invariant (in $`d1`$ dimensions) matrix element,
$$G_\varphi (\mathrm{𝚝𝚒𝚖𝚎})=|0|\widehat{\mathit{\varphi }}|\beta _\varphi |^2\mathrm{}^{\mathrm{𝚐𝚊𝚙}_\varphi \times \mathrm{𝚝𝚒𝚖𝚎}},$$
(20)
$`\chi `$ $`=`$ $`{\displaystyle \underset{\mathrm{𝚜𝚙𝚊𝚌𝚎}}{}}{\displaystyle G_\varphi (\mathrm{𝚝𝚒𝚖𝚎})𝑑\mathrm{𝚝𝚒𝚖𝚎}}`$ (21)
$``$ $`L^{d1}|0|\widehat{\mathit{\varphi }}|\beta _\varphi |^2{\displaystyle \frac{1}{\mathrm{𝚐𝚊𝚙}_\varphi }}.`$
The matrix element $`|0|\widehat{\mathit{\varphi }}|\beta _\varphi |`$ should scale like $`L^{x_\varphi }`$ (see e.g. equations (18) and (19)), hence relation (21) requires that the gap scales according to $`\mathrm{𝚐𝚊𝚙}_\varphi L^1`$ in order to restore the usual scaling of the susceptibility $`\chi L^{\gamma /\nu }`$ with $`\gamma /\nu =2\eta =d2x_\varphi `$ (see table 1). Note that the gap being an inverse correlation length, its scaling inversely proportional to the typical linear size of the system logically means that the correlation length at criticality is locked at that size.
The scaling of the gap is more constrained in $`2d`$ where rather powerful techniques apply. Conformal invariance indeed provides quite efficient methods for the determination of critical exponents of two-dimensional critical systems . The cylinder geometry is relevant in the study of quantum chains, since such a chain with periodic boundary conditions in the space direction just corresponds to an infinitely long classical cylinder ($`1+1`$ dimensions) of complex coordinates $`w=\mathrm{𝚝𝚒𝚖𝚎}+i\times \mathrm{𝚜𝚙𝚊𝚌𝚎}`$ (here, Euclidean time is assumed). This former geometry follows from the infinite two-dimensional plane $`z=r\mathrm{}^{i\theta }`$ through the standard logarithmic mapping $`w(z)=\frac{L}{2\pi }\mathrm{ln}z`$ and, at the critical point, the correlation functions transform according to conformal covariance, leading along the cylinder to
$$G_\varphi (\mathrm{𝚝𝚒𝚖𝚎})=\left(\frac{2\pi }{L}\right)^{2x_\varphi }\mathrm{}^{\frac{2\pi }{L}x_\varphi \times \mathrm{𝚝𝚒𝚖𝚎}}.$$
(22)
From comparison with equation (20), the critical correlation length amplitude appears universal and its value related to the corresponding critical exponent through the simple relation
$$\mathrm{𝚐𝚊𝚙}_\varphi =\frac{2\pi }{L}x_\varphi .$$
(23)
The matrix elements are also predicted by conformal invariance and follow from equation (22),
$$|0|\widehat{\mathit{\varphi }}|\beta |=\left(\frac{2\pi }{L}\right)^{x_\varphi }.$$
(24)
A prescription before using any conformal invariance result concerns the scaling of the whole spectrum. Multiplying the Hamiltonian by an arbitrary number of course changes the scale of the spectrum. Gap scaling from equation (23) requires to fix the normalization in such a way that the sound velocity is unity . The sound velocity might be defined in the long-wavelength limit, $`\mathrm{𝚜𝚘𝚞𝚗𝚍}\mathrm{𝚟𝚎𝚕𝚘𝚌𝚒𝚝𝚢}=\mathrm{\Delta }E/\mathrm{\Delta }k`$, where $`\mathrm{\Delta }E`$ is for instance measured by gaps in the bottom of the spectrum and $`\mathrm{\Delta }k=2\pi /L`$ is given by the quantization step of wave vectors.
## 3 Quantum Ising chain
The two-dimensional Ising model is often used as the paradigmatic illustration of second-order phase transitions. It is one of the most simple non trivial models and in the following we will show how its quantum counterpart, the quantum Ising chain in a transverse field, is built.
### 3.1 The Hamiltonian limit
Let us first consider a ladder of classical Ising spins $`s_{i,j}=\pm 1`$, $`\mathrm{}<i<+\mathrm{}`$, and $`j=1,2`$ (see figure 2). The nearest neighbour interactions are denoted as $`K_s=\beta J_s`$ in the space ($`j`$) direction and $`K_t=\beta J_t`$ in the time ($`i`$) direction.
The energy of a configuration is a sum over time slices,
$$\beta E\{s_{i,j}\}=\underset{i}{}(K_ss_{i,1}s_{i,2}+K_t(s_{i,1}s_{i+1,1}+s_{i,2}s_{i+1,2})),$$
(25)
and the partition function reads as
$`Z`$ $`=`$ $`\underset{\{s_{i,j}\}}{Tr}{\displaystyle \underset{i}{}}\mathrm{}^{K_ss_{i,1}s_{i,2}+K_t(s_{i,1}s_{i+1,1}+s_{i,2}s_{i+1,2})}=\underset{\{s_{i,j}\}}{Tr}{\displaystyle \underset{i}{}}T_{i,i+1},`$ (26)
where $`T_{i,i+1}`$ is an element of some transfer matrix between row states $`|\{\sigma _j\}=|s_{i,1},s_{i,2}`$ and $`|\{\sigma _{}^{}{}_{j}{}^{}\}=|s_{i+1,1},s_{i+1,2}`$ at rows $`i`$ and $`i+1`$ (see figure 2),
$$T_{i,i+1}=\{\sigma _j^{}\}|\widehat{\text{T}}|\{\sigma _j\}.$$
(27)
$`\widehat{\text{T}}`$ is a $`4\times 4`$ matrix, since the two spins $`\sigma _1,\sigma _2`$ have four different configurations. The space coupling term is diagonal, $`\delta _{\{\sigma _j\},\{\sigma _j^{}\}}\times \mathrm{exp}(K_s\sigma _1\sigma _2)`$,
$$\begin{array}{cc}\text{}& \begin{array}{cccc}|,& |,& |,& |,\end{array}\\ \multicolumn{2}{c}{}\\ & \\ \multicolumn{2}{c}{}\\ \begin{array}{c},|\\ ,|\\ ,|\\ ,|\end{array}& \left(\begin{array}{cccc}\mathrm{}^{K_s}& 0& 0& 0\\ 0& \mathrm{}^{K_s}& 0& 0\\ 0& 0& \mathrm{}^{K_s}& 0\\ 0& 0& 0& \mathrm{}^{K_s}\end{array}\right),\end{array}$$
(28)
while a non-diagonal contribution due to possible spin flips in the time direction, $`\mathrm{exp}[K_t(\sigma _1\sigma _1^{}+\sigma _2\sigma _2^{})]`$, leads to the following matrix
$$\begin{array}{cc}\text{}& \begin{array}{cccc}|,& |,& |,& |,\end{array}\\ \multicolumn{2}{c}{}\\ & \\ \multicolumn{2}{c}{}\\ \begin{array}{c},|\\ ,|\\ ,|\\ ,|\end{array}& \left(\begin{array}{cccc}\mathrm{}^{2K_t}& 1& 1& \mathrm{}^{2K_t}\\ 1& \mathrm{}^{2K_t}& \mathrm{}^{2K_t}& 1\\ 1& \mathrm{}^{2K_t}& \mathrm{}^{2K_t}& 1\\ \mathrm{}^{2K_t}& 1& 1& \mathrm{}^{2K_t}\end{array}\right).\end{array}$$
(29)
The transfer matrix $`\widehat{\text{T}}`$ should be identified through its matrix elements. For the diagonal term, we introduce diagonal operators in the $`|\{\sigma _j\}`$ basis, namely
$$\widehat{𝝈}_z(1)=\underset{\sigma _1,\sigma _2}{}|\sigma _1,\sigma _2\sigma _1\sigma _1,\sigma _2|,[\widehat{𝝈}_z(1)]=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)_1\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)_2,$$
(30)
$$\widehat{𝝈}_z(2)=\underset{\sigma _1,\sigma _2}{}|\sigma _1,\sigma _2\sigma _2\sigma _1,\sigma _2|,[\widehat{𝝈}_z(2)]=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)_1\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)_2,$$
(31)
where $`(\widehat{\text{1}})_j`$ and $`(\widehat{𝝈}_z)_j`$ represent the $`2\times 2`$ identity or Pauli matrix acting only on variables $`\sigma _j`$ at site $`j`$. The matrix in equation (28) is identified to that of the operator
$$\mathrm{exp}(K_s\widehat{𝝈}_z(1)\widehat{𝝈}_z(2)).$$
(32)
For the flipping term, we introduce transition operators
$$\widehat{𝝈}_x(1)|\sigma _1,\sigma _2=|\sigma _1,\sigma _2,[\widehat{𝝈}_x(1)]=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)_1\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)_2,$$
(33)
$$\widehat{𝝈}_x(2)|\sigma _1,\sigma _2=|\sigma _1,\sigma _2,[\widehat{𝝈}_x(2)]=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)_1\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)_2,$$
(34)
and since $`\widehat{𝝈}_x^2(j)=\widehat{\text{1}}`$, we have, for any value of $`K_t^{}`$, the useful identity
$$\mathrm{exp}[K_t^{}\widehat{𝝈}_x(j)]=\widehat{\text{1}}\mathrm{cosh}K_t^{}+\widehat{𝝈}_x(j)\mathrm{sinh}K_t^{},$$
(35)
so that the matrix in equation (29) is identified to that of the operator $`\mathrm{}^{K_t^{}(\widehat{𝝈}_x(1)+\widehat{𝝈}_x(2))}`$,
$$\left(\begin{array}{cccc}\mathrm{cosh}^2K_t^{}& \mathrm{cosh}K_t^{}\mathrm{sinh}K_t^{}& \mathrm{cosh}K_t^{}\mathrm{sinh}K_t^{}& \mathrm{sinh}^2K_t^{}\\ \mathrm{cosh}K_t^{}\mathrm{sinh}K_t^{}& \mathrm{cosh}^2K_t^{}& \mathrm{sinh}^2K_t^{}& \mathrm{cosh}K_t^{}\mathrm{sinh}K_t^{}\\ \mathrm{cosh}K_t^{}\mathrm{sinh}K_t^{}& \mathrm{sinh}^2K_t^{}& \mathrm{cosh}^2K_t^{}& \mathrm{cosh}K_t^{}\mathrm{sinh}K_t^{}\\ \mathrm{sinh}^2K_t^{}& \mathrm{cosh}K_t^{}\mathrm{sinh}K_t^{}& \mathrm{cosh}K_t^{}\mathrm{sinh}K_t^{}& \mathrm{cosh}^2K_t^{}\end{array}\right),$$
(36)
(up to a prefactor which only shifts the free energy by a constant) provided that we demand
$$\mathrm{tanh}K_t^{}=\mathrm{}^{2K_t}.$$
(37)
This relation is equivalent to the usual duality relation
$$\mathrm{sinh}2K_t\mathrm{sinh}2K_t^{}=1.$$
(38)
The transfer matrix eventually follows (we use $`\mathrm{}^{𝐀+𝐁}=\mathrm{}^𝐀\mathrm{}^𝐁\mathrm{}^{[𝐀,𝐁]/2}\mathrm{}`$ and already take into account the simplification due to the extreme anisotropic limit which eliminates all correction terms)
$$\widehat{\text{T}}=\mathrm{exp}(K_s\widehat{𝝈}_z(1)\widehat{𝝈}_z(2))\times \mathrm{exp}(K_t^{}(\widehat{𝝈}_x(1)+\widehat{𝝈}_x(2)).$$
(39)
The multiple spin flip terms appear in the expansion of the exponential through product of terms like in equation (35). The generalization to a lattice of width $`L`$ in the space direction (we choose now for the rest of the paper periodic boundary conditions $`\widehat{𝝈}_{x,z}(L+1)=\widehat{𝝈}_{x,z}(1)`$ in space direction) is straightforward, summing over $`j=1,L`$,
$$\widehat{\text{T}}=\mathrm{exp}(K_s\underset{j=1}{\overset{L}{}}\widehat{𝝈}_z(j)\widehat{𝝈}_z(j+1))\times \mathrm{exp}(K_t^{}\underset{j=1}{\overset{L}{}}\widehat{𝝈}_x(j)).$$
(40)
The transfer matrix has a complicated structure, since arbitrary large numbers of single spin transitions may simultaneously occur and the matrix representation of $`\widehat{\text{T}}`$ is expected to be dense. It may be considerably simplified in the extreme anisotropic limit also called Hamiltonian limit, since then only single spin transitions survive and $`\widehat{\text{H}}`$ is a sparser matrix. Remember that the relation between the transfer matrix and the Hamiltonian involves the lattice spacing $`ϵ`$ in the time direction which tends to zero, $`\widehat{\text{T}}=\mathrm{exp}(ϵ\widehat{\text{H}})=\widehat{\text{1}}ϵ\widehat{\text{H}}+O(ϵ^2)`$. A simple limit of $`\widehat{\text{T}}`$ in equation (40) is obtained when $`K_s0`$, $`K_t\mathrm{}`$ (or $`K_t^{}0`$), with $`K_s/K_t^{}=\lambda =O(1)`$, and in order to avoid unnecessary constants, we may set $`ϵ=2K_t^{}`$ <sup>2</sup><sup>2</sup>2The normalization with a factor 2 ensures a sound velocity equal to 1 (i.e. a linear dispersion relation at long wavelength with slope unity) and is necessary in order to compare later with conformal invariance predictions (see discussion). to get
$$\widehat{\text{H}}=\frac{1}{2}\lambda \underset{j=1}{\overset{L}{}}\widehat{𝝈}_z(j)\widehat{𝝈}_z(j+1))\frac{1}{2}_{j=1}^L\widehat{𝝈}_x(j).$$
(41)
The critical line of the two-dimensional anisotropic classical system, obtained through duality i.e. when the relation $`\mathrm{sinh}2K_s\mathrm{sinh}2K_t=1`$ is fulfilled, is equivalent in the extreme anisotropic limit to $`2K_s(2K_t^{})^1=\lambda _c=1`$ when equation (38) is also used. Fluctuations in the ground state structure have their origin in the competition between two non-commuting terms in the Hamiltonian. In expression (41), the term $`\lambda \widehat{𝝈}_z(j)\widehat{𝝈}_z(j+1)`$ acts like a ferromagnetic interaction which reinforces order in the $`z`$-direction (order parameter) in the chain, while $`\widehat{𝝈}_x(j)`$ appears as a disordering flipping term. When $`\lambda >1`$, the ordering term dominates and we expect ferromagnetic $`z`$-order in the ground state, while the disordered phase corresponds to a dominant role of the flipping term when $`\lambda <1`$.
### 3.2 The symmetries of the model
The classical model exhibits the $`Z_2`$ symmetry (invariance under global change $`+11`$ on each site). In the quantum case we expect a similar property <sup>3</sup><sup>3</sup>3For simplicity, we use a short notation $`|,,,\mathrm{},`$ for $`|_1|_2|_3\mathrm{}|_L`$. $`|,,,\mathrm{},_z|,,,\mathrm{},_z`$. This property has an algebraic manifestation through the commutator
$$[\widehat{\text{H}},\widehat{\text{P}}]=0,\widehat{\text{P}}=\underset{j=1}{\overset{L}{}}\widehat{𝝈}_x(j).$$
(42)
The eigenstates of $`\widehat{\text{H}}`$ may then be classified according to their parity, $`P=\pm 1`$, $`\widehat{\text{P}}|\mathrm{𝚎𝚟𝚎𝚗}\mathrm{𝚜𝚝𝚊𝚝𝚎}=+1|\mathrm{𝚎𝚟𝚎𝚗}\mathrm{𝚜𝚝𝚊𝚝𝚎}`$ and $`\widehat{\text{P}}|\mathrm{𝚘𝚍𝚍}\mathrm{𝚜𝚝𝚊𝚝𝚎}=1|\mathrm{𝚘𝚍𝚍}\mathrm{𝚜𝚝𝚊𝚝𝚎}`$ and an obvious consequence is that $`\widehat{\text{H}}`$ has vanishing matrix elements between states of different parities. More generally, for any even operator $`\widehat{\text{E}}`$ such that $`\widehat{\text{E}}\widehat{\text{P}}=\widehat{\text{P}}\widehat{\text{E}}`$, the only possibly non vanishing matrix elements of $`\widehat{\text{E}}`$ are between states of the same parity while odd operators satisfying $`\widehat{\text{O}}\widehat{\text{P}}=\widehat{\text{P}}\widehat{\text{O}}`$ have surviving matrix elements between states of opposite parities. The ground state of the system is expected to be even (it is of the symmetry of the Hamiltonian itself), so that a measure of local order in the system is given in agreement with equation (16) by
$$\mathrm{𝙾𝚛𝚍𝚎𝚛}\mathrm{𝚙𝚊𝚛𝚊𝚖𝚎𝚝𝚎𝚛}=|\mathrm{𝙻𝚘𝚠𝚎𝚜𝚝}\mathrm{𝚘𝚍𝚍}\mathrm{𝚜𝚝𝚊𝚝𝚎}|\widehat{𝝈}_z(j)|\mathrm{𝙶𝚗𝚍}|.$$
(43)
A similar definition of the energy density may be given by inserting terms appearing into the Hamiltonian inside states of the same parity, e.g.
$$\mathrm{𝙴𝚗𝚎𝚛𝚐𝚢}\mathrm{𝚍𝚎𝚗𝚜𝚒𝚝𝚢}=|\mathrm{𝟷}𝚜𝚝\mathrm{𝚎𝚡𝚌𝚒𝚝𝚎𝚍}\mathrm{𝚎𝚟𝚎𝚗}\mathrm{𝚜𝚝𝚊𝚝𝚎}|\widehat{𝝈}_x(j)|\mathrm{𝙶𝚗𝚍}|.$$
(44)
Matrix elements of $`\widehat{𝝈}_z(j)\widehat{𝝈}_z(j+1)`$ might have been chosen as well. The reason for being mainly interested in the bottom of the spectrum (see figure 3) lies in the fact that the quantum phase transition takes place at zero temperature. Together with the eigenstates already mentioned, we define the state $`|\mathrm{𝚜𝚘𝚞𝚗𝚍}\mathrm{𝚟𝚎𝚕𝚘𝚌𝚒𝚝𝚢}`$ and the corresponding relation which fixes the value of the sound velocity,
$$E_{v_s}E_{\mathrm{𝙾𝚛𝚍𝚎𝚛}}=\frac{2\pi }{L}v_s.$$
(45)
If the normalization is not properly chosen, or if the sound velocity is not known, equation (23) should be replaced by
$$\mathrm{𝚐𝚊𝚙}_\varphi =\frac{2\pi }{L}v_sx_\varphi $$
(46)
and $`v_s`$ obtained by the equation above. In the case of the quantum Ising model, the Hamiltonian (41) is conveniently normalized and $`v_s=1`$ .
### 3.3 Small Chains
A quantum chain of length $`L`$ corresponds to an infinitely long ladder or strip of classical spins $`s_{i,j}`$. In the case of a chain of 2 spins with periodic boundary conditions (the classical counterpart is thus a cylinder), the Hamiltonian may be written
$$2\widehat{\text{H}}=\lambda \widehat{𝝈}_z(1)\widehat{𝝈}_z(2)+\lambda \widehat{𝝈}_z(2)\widehat{𝝈}_z(1)+\widehat{𝝈}_x(1)+\widehat{𝝈}_x(2).$$
(47)
The construction of the matrix of $`\widehat{\text{H}}`$ is easy,
$$[2\widehat{\text{H}}]=\left(\begin{array}{cccc}2\lambda & 1& 1& 0\\ 1& 2\lambda & 0& 1\\ 1& 0& 2\lambda & 1\\ 0& 1& 1& 2\lambda \end{array}\right)\begin{array}{c}|,_z\\ |,_z\\ |,_z\\ |,_z\end{array}$$
(48)
but it is not written here in the simplest way. Due to the parity property of $`\widehat{\text{H}}`$, it is easier to write the $`4\times 4`$ matrix in the basis of $`\widehat{𝝈}_x`$-eigenstates, where it is block diagonal. In the basis $`\{|,_x,|,_x,|,_x,|,_x\}`$, we have
$$[2\widehat{\text{H}}]=\left(\begin{array}{cc}\begin{array}{cc}2& 2\lambda \\ 2\lambda & 2\end{array}& \text{0}\\ & \\ \text{0}& \begin{array}{cc}0& 2\lambda \\ 2\lambda & 0\end{array}\end{array}\right)\begin{array}{c}|,_x\\ |,_x\\ |,_x\\ |,_x\end{array}$$
(49)
In the even sector at the critical coupling $`\lambda _c=1`$, the eigenvalues of $`\widehat{\text{H}}`$ are $`E_{\mathrm{𝙶𝚗𝚍}}=\sqrt{2}`$ and $`E_{\mathrm{𝙴𝚗𝚎𝚛𝚐𝚢}}=+\sqrt{2}`$. The ground state $`|\mathrm{𝙶𝚗𝚍}=a_{\mathrm{𝙶𝚗𝚍}}^{}|,_x+a_{\mathrm{𝙶𝚗𝚍}}^{}|,_x`$ has normalized components $`a_{\mathrm{𝙶𝚗𝚍}}^{}=0.924`$ and $`a_{\mathrm{𝙶𝚗𝚍}}^{}=0.383`$ while the energy excited state is given by $`|\mathrm{𝙴𝚗𝚎𝚛𝚐𝚢}=a_{\mathrm{𝙴𝚗𝚎𝚛𝚐𝚢}}^{}|,_x+a_{\mathrm{𝙴𝚗𝚎𝚛𝚐𝚢}}^{}|,_x`$ with $`a_{\mathrm{𝙴𝚗𝚎𝚛𝚐𝚢}}^{}=0.383`$ and $`a_{\mathrm{𝙴𝚗𝚎𝚛𝚐𝚢}}^{}=0.924`$. The energy matrix element follows,
$`\mathrm{𝙴𝚗𝚎𝚛𝚐𝚢}\mathrm{𝚍𝚎𝚗𝚜𝚒𝚝𝚢}`$ $`=`$ $`|\mathrm{𝙴𝚗𝚎𝚛𝚐𝚢}|\widehat{𝝈}_x(1)|\mathrm{𝙶𝚗𝚍}|`$ (50)
$`=`$ $`|a_{\mathrm{𝙴𝚗𝚎𝚛𝚐𝚢}}^{}a_{\mathrm{𝙶𝚗𝚍}}^{}a_{\mathrm{𝙴𝚗𝚎𝚛𝚐𝚢}}^{}a_{\mathrm{𝙶𝚗𝚍}}^{}|`$
$`=`$ $`0.707.`$
In the odd sector, the lowest eigenvalue is $`E_{\mathrm{𝙾𝚛𝚍𝚎𝚛}}=1`$ and the corresponding eigenvector $`|\mathrm{𝙾𝚛𝚍𝚎𝚛}=a_{\mathrm{𝙾𝚛𝚍𝚎𝚛}}^{}|,_x+a_{\mathrm{𝙾𝚛𝚍𝚎𝚛}}^{}|,_x`$ has components $`a_{\mathrm{𝙾𝚛𝚍𝚎𝚛}}^{}=a_{\mathrm{𝙾𝚛𝚍𝚎𝚛}}^{}=0.707`$. The order parameter matrix element follows,
$`\mathrm{𝙾𝚛𝚍𝚎𝚛}\mathrm{𝚙𝚊𝚛𝚊𝚖𝚎𝚝𝚎𝚛}`$ $`=`$ $`|\mathrm{𝙾𝚛𝚍𝚎𝚛}|\widehat{𝝈}_z(1)|\mathrm{𝙶𝚗𝚍}|`$ (51)
$`=`$ $`|a_{\mathrm{𝙾𝚛𝚍𝚎𝚛}}^{}a_{\mathrm{𝙶𝚗𝚍}}^{}+a_{\mathrm{𝙾𝚛𝚍𝚎𝚛}}^{}a_{\mathrm{𝙶𝚗𝚍}}^{}|`$
$`=`$ $`0.924.`$
In the case of a chain of 3 spins with periodic boundary conditions, we have in the basis of $`\sigma _x`$ eigenstates
$$[2\widehat{\text{H}}]=\left(\begin{array}{ccccccccc}3& \lambda & \lambda & \lambda & & & & & \\ \lambda & 1& \lambda & \lambda & & & & & \\ \lambda & \lambda & 1& \lambda & & & & & \\ \lambda & \lambda & \lambda & 1& & & & & \\ & & & & & & & & \\ & & & & 3& \lambda & \lambda & \lambda & \\ & & & & \lambda & 1& \lambda & \lambda & \\ & & & & \lambda & \lambda & 1& \lambda & \\ & & & & \lambda & \lambda & \lambda & 1& \end{array}\right)\begin{array}{c}|,,_x\\ |,,_x\\ |,,_x\\ |,,_x\\ |,,_x\\ |,,_x\\ |,,_x\\ |,,_x\end{array}$$
(52)
The relevant energy levels and corresponding matrix elements for three small sizes are collected in table 2. Use will be made of these results to get approximate values of the critical exponents in the discussion in section 5.
## 4 Quantum Potts chain
The Potts model generalizes the Ising model. The sites are occupied by Potts variables (abusively called spins) with $`q`$ different states, $`n_{i,j}=0,1,\mathrm{}q1`$, and bonds between nearest neighbour sites may have two different energy levels, depending on the relative states of the site variables, e.g. $`J(q\delta _{n,n^{}}1)`$.
### 4.1 Hamiltonian limit
We follow the same steps as in the case of the Ising model . The energy of a configuration is written as a sum over time slices,
$$\beta E\{n_{i,j}\}=\underset{i}{}\left[K_s\underset{j=1}{\overset{L}{}}(q\delta _{n_{i,j},n_{i,j+1}}1)+K_t\underset{j=1}{\overset{L}{}}(q\delta _{n_{i,j},n_{i+1,j}}1)\right],$$
(53)
where the fact that interactions are limited to nearest neighbours allows a factorized partition function $`Z=Tr_{\{n_{i,j}\}}_iT_{i,i+1}`$ with matrix elements between row states $`|\{\nu _j\}=|n_{i,1},n_{i,2},\mathrm{}n_{i,L}`$ and $`|\{\nu _{}^{}{}_{j}{}^{}\}=|n_{i+1,1},n_{i+1,2},\mathrm{}n_{i+1,L}`$ at time indexes $`i`$ and $`i+1`$,
$$T_{i,i+1}=\{\nu _{}^{}{}_{j}{}^{}\}|\widehat{\text{T}}|\{\nu _j\}.$$
(54)
Now $`\widehat{\text{T}}`$ is a $`q^L\times q^L`$ matrix (each of the $`L`$ spins $`n_{i,j}`$ have $`q`$ different possible states). As in the case of the Ising model, the space coupling term is diagonal, $`\delta _{\{\nu _j\},\{\nu _{}^{}{}_{j}{}^{}\}}\times \mathrm{exp}(qK_s_k(\delta _{n_{i,k},n_{i,k+1}}1))`$. Its description in $`\widehat{\text{T}}`$ requires the introduction of a combination of diagonal operators in the basis $`|\{\nu _j\}`$, the matrix element of which reproduces the Kronecker delta. For that purpose, we define state vector $`|\nu _1|\nu _2\mathrm{}|\nu _L|\nu _1,\nu _2,\mathrm{},\nu _L`$ and the corresponding diagonal operators $`\widehat{\text{C}}_k`$ and $`\widehat{\text{C}}_k^{}`$ such that (in this section, we stay close to the notations of Ref. )
$`\widehat{\text{C}}_k|\nu _1,\nu _2,\mathrm{},\nu _k,\mathrm{},\nu _L`$ $``$ $`\mathrm{}^{2i\pi \nu _k/q}|\nu _1,\nu _2,\mathrm{},\nu _k,\mathrm{},\nu _L,`$
$`\widehat{\text{C}}_k^{}|\nu _1,\nu _2,\mathrm{},\nu _k,\mathrm{},\nu _L`$ $``$ $`\mathrm{}^{2i\pi \nu _k/q}|\nu _1,\nu _2,\mathrm{},\nu _k,\mathrm{},\nu _L.`$ (55)
Due to the property
$$\frac{1}{q}\underset{p=0}{\overset{q1}{}}\mathrm{}^{2ip\pi (\nu _k\nu _l)/q}=\delta _{\nu _k,\nu _l}$$
we may write the Kronecker delta as the diagonal matrix element
$$\delta _{\nu _k,\nu _l}=\{\nu _j\}|\frac{1}{q}\underset{p=0}{\overset{q1}{}}(\widehat{\text{C}}_k^{}\widehat{\text{C}}_l)^p|\{\nu _j\}.$$
(56)
and thus the space contribution to the transfer matrix follows
$$\mathrm{exp}\left(K_s\underset{j=1}{\overset{L}{}}\underset{p=1}{\overset{q1}{}}(\widehat{\text{C}}_j^{}\widehat{\text{C}}_{j+1})^p\right)$$
(57)
where the term $`p=0`$ is subtracted, since it compensates the $`1`$ in the definition of the pair energy in equation (53). For the time contribution, it is necessary to introduce flipping operators and we define
$`\widehat{\text{R}}_k|\nu _1,\nu _2,\mathrm{},\nu _k,\mathrm{},\nu _L`$ $``$ $`|\nu _1,\nu _2,\mathrm{},\nu _k+1,\mathrm{},\nu _L,`$
$`\widehat{\text{R}}_k^{}|\nu _1,\nu _2,\mathrm{},\nu _k,\mathrm{},\nu _L`$ $``$ $`|\nu _1,\nu _2,\mathrm{},\nu _k1,\mathrm{},\nu _L,`$ (58)
where periodicity in space state is assumed, $`|\nu +q=|\nu `$. The matrix representation of these operators is the following
$$[\widehat{\text{C}}_j]=\widehat{\text{1}}\widehat{\text{1}}\mathrm{}\left(\begin{array}{ccccc}1& 0& 0& \mathrm{}& 0\\ 0& \mathrm{}^{2i\pi /q}& 0& \mathrm{}& 0\\ \mathrm{}& 0& \mathrm{}^{4i\pi /q}& \mathrm{}& \mathrm{}\\ 0& \mathrm{}& \mathrm{}& \mathrm{}& 0\\ 0& 0& \mathrm{}& 0& \mathrm{}^{2i\pi (q1)/q}\end{array}\right)_j\mathrm{}\widehat{\text{1}},$$
(59)
$$[\widehat{\text{R}}_j]=\widehat{\text{1}}\widehat{\text{1}}\mathrm{}\left(\begin{array}{ccccc}0& 0& \mathrm{}& 0& 1\\ 1& 0& \mathrm{}& 0& 0\\ 0& 1& \mathrm{}& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& 1& 0& 0\\ 0& \mathrm{}& 0& 1& 0\end{array}\right)_j\mathrm{}\widehat{\text{1}}.$$
(60)
They commute on different sites and obey the following algebra on a given site $`\widehat{\text{R}}_j\widehat{\text{C}}_j=\mathrm{}^{2i\pi /q}\widehat{\text{C}}_j\widehat{\text{R}}_j`$, $`\widehat{\text{R}}_j^{}\widehat{\text{C}}_j=\mathrm{}^{2i\pi /q}\widehat{\text{C}}_j\widehat{\text{R}}_j^{}`$, $`\widehat{\text{R}}_j\widehat{\text{C}}_j^{}=\mathrm{}^{2i\pi /q}\widehat{\text{C}}_j^{}\widehat{\text{R}}_j`$, $`\widehat{\text{R}}_j^{}\widehat{\text{C}}_j^{}=\mathrm{}^{2i\pi /q}\widehat{\text{C}}_j^{}\widehat{\text{R}}_j^{}`$, $`\widehat{\text{C}}_j^q=\widehat{\text{R}}_j^q=\widehat{\text{1}}`$. Temporarily forgetting about the site index, diagonalization of $`\widehat{\text{R}}`$ leads to
$$\widehat{\text{R}}|r=\mathrm{}^{2ir\pi /q}|r,|r=q^{1/2}\underset{\nu =0}{\overset{q1}{}}\mathrm{}^{(2ir\pi /q)\nu }|\nu .$$
It follows that the operator $`\widehat{𝚯}\frac{1}{q}_{p=0}^{q1}\widehat{\text{R}}^p`$ has the property $`\widehat{𝚯}|r=\delta _{r,0}|r`$ (it is equal the projector on the “zero eigenstate” of $`\widehat{\text{R}}`$, $`\widehat{𝚯}=|r=0r=0|`$, and so $`\{\nu _{}^{}{}_{j}{}^{}\}|\widehat{𝚯}|\{\nu _j\}=1/q`$ for any pair of states $`|\{\nu _j\}`$ and $`|\{\nu _{}^{}{}_{j}{}^{}\}`$) which enables to write
$$\mathrm{exp}(qK_t^{}\widehat{𝚯}_j)=\widehat{\text{1}}+\widehat{𝚯}_j(\mathrm{}^{qK_t^{}}1)$$
(61)
for arbitrary value of $`K_t^{}`$. The time contributions, $`\mathrm{}^{K_t(q\delta _{\nu _j,\nu _{}^{}{}_{j}{}^{}}1)}`$, which can take values $`\mathrm{}^{K_t(q1)}`$ if $`\nu _{}^{}{}_{j}{}^{}=\nu _j`$ and $`\mathrm{}^{K_t}`$ otherwise, are obtained from matrix elements of terms $`\mathrm{exp}(qK_t^{}\widehat{𝚯}_j)`$ provided that $`K_t^{}`$ satisfies the duality relation
$$(\mathrm{}^{qK_t}1)(\mathrm{}^{qK_t^{}}1)=q.$$
(62)
The transfer matrix eventually reads as
$$\widehat{\text{T}}=\mathrm{exp}\left(K_s\underset{j=1}{\overset{L}{}}\underset{p=1}{\overset{q1}{}}(\widehat{\text{C}}_j^{}\widehat{\text{C}}_{j+1})^p\right)\times \mathrm{exp}\left(K_t^{}\underset{j=1}{\overset{L}{}}\underset{p=1}{\overset{q1}{}}\widehat{\text{R}}_j^p\right).$$
(63)
We have shifted the sum in the last term, starting from $`p=1`$ which only changes the transfer matrix by a constant prefactor $`\mathrm{}^{K_t^{}}`$. This modification affects the free energy density (or the Hamiltonian) by a constant only and does not change the thermodynamic properties. In the Hamiltonian limit $`K_s0`$, $`K_t\mathrm{}`$ with fixed $`\lambda =K_s/K_t^{}`$, we obtain the Hamiltonian of the quantum Potts chain
$$\widehat{\text{H}}=\frac{1}{2}\lambda \underset{j=1}{\overset{L}{}}\underset{p=1}{\overset{q1}{}}(\widehat{\text{C}}_j^{}\widehat{\text{C}}_{j+1})^p\frac{1}{2}\underset{j=1}{\overset{L}{}}\underset{p=1}{\overset{q1}{}}\widehat{\text{R}}_j^p.$$
(64)
It is easy to check that in the case $`q=2`$, we recover the Hamiltonian (41) of the quantum Ising chain with $`\widehat{\text{C}}`$ playing the role of $`\widehat{𝝈}_z`$ and $`\widehat{\text{R}}`$ that of $`\widehat{𝝈}_x`$. The limit $`\lambda \mathrm{}`$ leads to $`q`$ degenerate ordered ground states $`|n,n,n,\mathrm{},n`$, $`n=0,1,\mathrm{},q1`$, while in the other limit $`\lambda 0`$, the term in $`\widehat{\text{R}}_j`$ introduces disorder in the ground state through local rotations between $`\widehat{\text{C}}_j`$-eigenstates. The critical point of the classical two-dimensional model is given by duality, $`(\mathrm{}^{qK_s}1)(\mathrm{}^{qK_t}1)=q`$, i.e. $`K_s=K_t^{}`$ or $`\lambda _c=1`$.
A different, but simpler route in order to get the time contribution to the Hamiltonian matrix is the following: in the “row states basis” $`|\{\nu _j\}`$, the time contribution $`_k\mathrm{}^{K_t(q\delta _{\nu _k,\nu _{}^{}{}_{k}{}^{}}1)}`$ is a product over “single bond time transfer operators”
$$\mathrm{𝚝𝚒𝚖𝚎}\mathrm{𝚌𝚘𝚗𝚝𝚛𝚒𝚋𝚞𝚝𝚒𝚘𝚗}\mathrm{𝚝𝚘}\widehat{\text{T}}\underset{k}{}\underset{\{\nu _j\},\{\nu _{}^{}{}_{j}{}^{}\}}{}|\{\nu _{}^{}{}_{j}{}^{}\}\mathrm{}^{K_t(q\delta _{\nu _k,\nu _{}^{}{}_{k}{}^{}}1)}\{\nu _j\}|$$
(65)
which take values $`\mathrm{}^{K_t(q1)}`$ on the diagonal and the same value $`\mathrm{}^{K_t}`$ for all single spin flipping terms, $`\mathrm{}^{2K_t}`$ for simultaneous double spin flipping terms and so on. Anticipating further simplifications which occur while taking the Hamiltonian limit $`K_t\mathrm{}`$, we restrict ourselves to single flipping terms. We introduce the symbol $`_{\{\nu _j\},\{\nu _{}^{}{}_{j}{}^{}\}}^{}`$ with the meaning of a double sum over row states $`|\{\nu _j\}`$ and $`|\{\nu _{}^{}{}_{j}{}^{}\}`$ such that $`jk`$, $`\nu _{}^{}{}_{j}{}^{}=\nu _j`$ (single spin flipping terms, since only $`\nu _k`$ is likely to be flipped between the two row states). The contribution to the Hamiltonian follows from taking the logarithm of this expression which leads to $`ϵ\widehat{\text{H}}`$.
$`\mathrm{ln}(\mathrm{𝚝𝚒𝚖𝚎}\mathrm{𝚌𝚘𝚗𝚝𝚛𝚒𝚋𝚞𝚝𝚒𝚘𝚗}\mathrm{𝚝𝚘}\widehat{\text{T}})`$ $`=`$ $`{\displaystyle \underset{k}{}}\mathrm{ln}[\mathrm{}^{K_t(q1)}(\widehat{\text{1}}+`$ (66)
$`\mathrm{}^{qK_t}{\displaystyle \underset{\genfrac{}{}{0pt}{}{\{\nu _j\},\{\nu _{}^{}{}_{j}{}^{}\}}{\nu _{}^{}{}_{k}{}^{}\nu _k}}{\overset{}{}}}|\{\nu _{}^{}{}_{j}{}^{}\}\{\nu _j\}|)+O(\mathrm{}^{(q+1)K_t})]`$
$``$ $`LK_t(q1)+\mathrm{}^{qK_t}{\displaystyle \underset{\genfrac{}{}{0pt}{}{\{\nu _j\},\{\nu _{}^{}{}_{j}{}^{}\}}{\nu _{}^{}{}_{k}{}^{}\nu _k}}{\overset{}{}}}|\{\nu _{}^{}{}_{j}{}^{}\}\{\nu _j\}|`$
in the limit of a strong coupling in the time direction, $`K_t\mathrm{}`$.
### 4.2 Symmetries
We proceed as in the Ising case to get the symmetries of the Hamiltonian. The classical Potts model is obviously globally unchanged by the cyclic transformation $`j`$, $`0_j1_j`$, $`1_j2_j`$, …$`(q1)_j0_j`$. Such a rotation in spin states is realized by the operator $`\widehat{\text{R}}_j`$. We may thus define a charge operator which simultaneously rotates Potts variables at all sites,
$$\widehat{\text{Q}}=\underset{j}{}\widehat{\text{R}}_j.$$
(67)
Using the commutation relations between the $`\widehat{\text{R}}_j`$’s and the $`\widehat{\text{C}}_j`$’s, it is easy to prove that the charge operator commutes with the Hamiltonian. Since the eigenvalues of $`\widehat{\text{R}}_j`$ are the $`q`$ distinct complex numbers $`\omega _j=\mathrm{}^{2i\pi r_j/q}`$, those of the charge operator are also given by $`q`$ numbers $`\mathrm{exp}(_j2i\pi r_j/q)`$ which enable to write the Hamiltonian matrix under a block-diagonal structure with $`q`$ different sectors, depending on the value of $`_jr_j\mathrm{mod}(q)`$ (the sectors can be referred to as sector # $`0`$, $`1`$, …, $`q1`$).
Any operator $`\widehat{\text{O}}_p`$ with the commutation property $`\widehat{\text{O}}_p\widehat{\text{Q}}=\mathrm{}^{2i\pi p/q}\widehat{\text{Q}}\widehat{\text{O}}_p`$ ($`\widehat{\text{R}}`$ is such an example with $`p=0`$, $`\widehat{\text{C}}^{}`$ is an example with $`p=1`$, and $`\widehat{\text{C}}`$ with $`p=q1`$) has non vanishing matrix elements between states of defined symmetry provided that the charge sectors obey a simple relation $`\nu _\psi \nu _\varphi =p`$:
$`|\varphi ,|\psi \mathrm{such}\mathrm{that}\widehat{\text{Q}}|\varphi =\mathrm{}^{2i\pi \nu _\varphi /q}|\varphi ,\widehat{\text{Q}}|\psi =\mathrm{}^{2i\pi \nu _\psi /q}|\psi ,`$
$`\nu _\psi \nu _\varphi p\varphi |\widehat{\text{O}}_p|\psi =0`$ (68)
As an energy density, we may choose any single term appearing inside the Hamiltonian, e.g. $`\mathrm{𝙴𝚗𝚎𝚛𝚐𝚢}\mathrm{𝚍𝚎𝚗𝚜𝚒𝚝𝚢}=\widehat{\text{R}}_j`$. A measure of the order parameter would be provided by $`\mathrm{𝙻𝚘𝚌𝚊𝚕}\mathrm{𝚘𝚛𝚍𝚎𝚛}\mathrm{𝚙𝚊𝚛𝚊𝚖𝚎𝚝𝚎𝚛}=\widehat{\text{C}}_j`$, or $`\widehat{\text{C}}_j^{}`$ as well. The energy density gives access to the critical exponent $`x_ϵ`$. For the order parameter, although the two choices correspond to different values of $`p`$, since two critical exponents are sufficient in order to determine the whole universality class, the same value of $`x_\sigma `$ is expected from both definitions. From gap scaling in particular, we expect $`(q1)`$-fold degeneracy of the sectors $`p0`$.
### 4.3 Diagonalization of small chains
In the case of a quantum $`3`$state Potts chain of length $`L=2`$ with periodic boundary conditions, the action of the Hamiltonian in the $`\widehat{\text{C}}`$operators eigenbasis is the following,
$``$ $`2\widehat{\text{H}}|\nu _1,\nu _2=2\lambda (\omega _1\omega _2^1+\omega _1^2\omega _2^2)|\nu _1,\nu _2`$ (70)
$`+|\nu _1+1,\nu _2+|\nu _1+2,\nu _2+|\nu _1,\nu _2+1+|\nu _1,\nu _2+2,`$
$`\omega _j=\mathrm{exp}(2i\pi \nu _j/3).`$
We obtain the following $`3^2\times 3^2`$ matrix,
$$[2\widehat{\text{H}}]=\left(\begin{array}{ccccccccc}4\lambda & 1& 1& 1& 0& 0& 1& 0& 0\\ 1& 2\lambda & 1& 0& 1& 0& 0& 1& 0\\ 1& 1& 2\lambda & 0& 0& 1& 0& 0& 1\\ 1& 0& 0& 2\lambda & 1& 1& 1& 0& 0\\ 0& 1& 0& 1& 4\lambda & 1& 0& 1& 0\\ 0& 0& 1& 1& 1& 2\lambda & 0& 0& 1\\ 1& 0& 0& 1& 0& 0& 2\lambda & 1& 1\\ 0& 1& 0& 0& 1& 0& 1& 2\lambda & 1\\ 0& 0& 1& 0& 0& 1& 1& 1& 4\lambda \end{array}\right)\begin{array}{c}|00_C\\ |01_C\\ |02_C\\ |10_C\\ |11_C\\ |12_C\\ |20_C\\ |21_C\\ |22_C\end{array}$$
(71)
the eigenvalues of which are (eigenvalues of $`\widehat{\text{H}}`$ at $`\lambda _c=1`$), $`2.732`$, $`2.137`$, $`2.137`$, $`+0.500`$, $`+0.500`$, $`+0.732`$, $`+1.637`$, $`+1.637`$, and $`+2.000`$. It is of course more efficient to exploit the symmetries and to work in the $`\widehat{\text{R}}`$operators eigenbasis where,
$``$ $`2\widehat{\text{H}}|r_1,r_2=2\lambda (|r_1+1,r_21+|r_1+2,r_22)`$ (73)
$`+(\omega _1+\omega _1^2+\omega _2+\omega _2^2)|r_1,r_2,`$
$`\omega _j=\mathrm{exp}(2i\pi r_j/3).`$
The corresponding matrix is now
$$[2\widehat{\text{H}}]=\left(\begin{array}{ccc}\begin{array}{ccc}4& 2\lambda & 2\lambda \\ 2\lambda & 2& 2\lambda \\ 2\lambda & 2\lambda & 2\end{array}& & \\ & \begin{array}{ccc}& & \\ 2& 2\lambda & 2\lambda \\ 2\lambda & 1& 2\lambda \\ 2\lambda & 2\lambda & 1\end{array}& \\ & & \begin{array}{c}& & \multicolumn{-1}{c}{}\\ 2& 2\lambda & 2\lambda \\ 2\lambda & 1& 2\lambda \\ 2\lambda & 2\lambda & 1\end{array}\end{array}\right)\begin{array}{c}|00_R\\ |12_R\\ |21_R\\ |11_R\\ |20_R\\ |02_R\\ |22_R\\ |01_R\\ |10_R\end{array}$$
(74)
It is block diagonal, with eigenvalues of the first block ($`0`$ sector) $`2.732`$, $`+0.732`$ and $`+2.000`$, and the two-fold degenerate eigenvalues of the two remaining identical blocks (sectors $`1`$ and $`2`$), $`2.137`$, $`+0.500`$ and $`+1.637`$. The ground state and energy eigenstates are given by $`|\mathrm{𝙶𝚗𝚍}=a_{\mathrm{𝙶𝚗𝚍}}^{00}|00+a_{\mathrm{𝙶𝚗𝚍}}^{12}|12+a_{\mathrm{𝙶𝚗𝚍}}^{21}|21`$ and $`|\mathrm{𝙴𝚗𝚎𝚛𝚐𝚢}=a_{\mathrm{𝙴𝚗𝚎𝚛𝚐𝚢}}^{00}|00+a_{\mathrm{𝙴𝚗𝚎𝚛𝚐𝚢}}^{12}|12+a_{\mathrm{𝙴𝚗𝚎𝚛𝚐𝚢}}^{21}|21`$ and the energy density follows
$`\mathrm{𝙴𝚗𝚎𝚛𝚐𝚢}\mathrm{𝚍𝚎𝚗𝚜𝚒𝚝𝚢}`$ $`=`$ $`|\mathrm{𝙴𝚗𝚎𝚛𝚐𝚢}|\widehat{\text{R}}_1|\mathrm{𝙶𝚗𝚍}|`$ (75)
$`=`$ $`|a_{\mathrm{𝙴𝚗𝚎𝚛𝚐𝚢}}^{00}a_{\mathrm{𝙶𝚗𝚍}}^{00}+a_{\mathrm{𝙴𝚗𝚎𝚛𝚐𝚢}}^{12}a_{\mathrm{𝙶𝚗𝚍}}^{12}\mathrm{}^{2i\pi /3}+a_{\mathrm{𝙴𝚗𝚎𝚛𝚐𝚢}}^{21}a_{\mathrm{𝙶𝚗𝚍}}^{21}\mathrm{}^{4i\pi /3}|`$
$`=`$ $`0.613.`$
In the $`p=1`$ sector, the order parameter density is given by (the notation is obvious)
$`\mathrm{𝙾𝚛𝚍𝚎𝚛}\mathrm{𝚙𝚊𝚛𝚊𝚖𝚎𝚝𝚎𝚛}`$ $`=`$ $`|\mathrm{𝙾𝚛𝚍𝚎𝚛}p=1|\widehat{\text{C}}_1^{}|\mathrm{𝙶𝚗𝚍}|`$ (76)
$`=`$ $`|a_{\mathrm{𝙾𝚛𝚍𝚎𝚛}}^{10}a_{\mathrm{𝙶𝚗𝚍}}^{00}+a_{\mathrm{𝙾𝚛𝚍𝚎𝚛}}^{22}a_{\mathrm{𝙶𝚗𝚍}}^{12}+a_{\mathrm{𝙾𝚛𝚍𝚎𝚛}}^{01}a_{\mathrm{𝙶𝚗𝚍}}^{21}|`$
$`=`$ $`0.916,`$
and, due to the complete degeneracy of the two sectors $`p=1`$ and $`p=2`$, we get the same result in the remaining sector
$`\mathrm{𝙾𝚛𝚍𝚎𝚛}\mathrm{𝚙𝚊𝚛𝚊𝚖𝚎𝚝𝚎𝚛}`$ $`=`$ $`|\mathrm{𝙾𝚛𝚍𝚎𝚛}p=2|\widehat{\text{C}}_1|\mathrm{𝙶𝚗𝚍}|`$ (77)
$`=`$ $`|a_{\mathrm{𝙾𝚛𝚍𝚎𝚛}}^{20}a_{\mathrm{𝙶𝚗𝚍}}^{00}+a_{\mathrm{𝙾𝚛𝚍𝚎𝚛}}^{02}a_{\mathrm{𝙶𝚗𝚍}}^{12}+a_{\mathrm{𝙾𝚛𝚍𝚎𝚛}}^{11}a_{\mathrm{𝙶𝚗𝚍}}^{21}|`$
$`=`$ $`0.916,`$
The relevant eigenvalues and matrix elements for systems of sizes $`L=2`$ and 3 are collected in table 3.
Interested readers might find helpful to have also the matrix representing the $`4`$state Potts chain. When $`L=2`$ (the matrix is now $`4^L\times 4^L`$) at the critical coupling $`\lambda _c=1`$ in its block diagonal form in the $`\widehat{\text{R}}`$eigenbasis $`\{|00,|13,|22,|31,|01,|10,|23,|32,|02,|11,|20,|33,|03,|12,|21,|30\}`$, this matrix reads as
$$\left[2\widehat{\text{H}}\right]=\left(\begin{array}{cccc}\begin{array}{cccc}\hfill 6& \hfill 2& \hfill 2& \hfill 2\\ \hfill 2& \hfill 2& \hfill 2& \hfill 2\\ \hfill 2& \hfill 2& \hfill 2& \hfill 2\\ \hfill 2& \hfill 2& \hfill 2& \hfill 2\end{array}& & & \\ & \begin{array}{cccc}& & & \\ \hfill 2& \hfill 2& \hfill 2& \hfill 2\\ \hfill 2& \hfill 2& \hfill 2& \hfill 2\\ \hfill 2& \hfill 2& \hfill 2& \hfill 2\\ \hfill 2& \hfill 2& \hfill 2& \hfill 2\end{array}& & \\ & & \begin{array}{cc}& & \multicolumn{-1}{c}{}& \\ \hfill 2& \hfill 2& \hfill 2& \hfill 2\\ \hfill 2& \hfill 2& \hfill 2& \hfill 2\\ \hfill 2& \hfill 2& \hfill 2& \hfill 2\\ \hfill 2& \hfill 2& \hfill 2& \hfill 2\end{array}& \\ & & & \begin{array}{c}& & & \multicolumn{-2}{c}{}\\ \hfill 2& \hfill 2& \hfill 2& \hfill 2\\ \hfill 2& \hfill 2& \hfill 2& \hfill 2\\ \hfill 2& \hfill 2& \hfill 2& \hfill 2\\ \hfill 2& \hfill 2& \hfill 2& \hfill 2\end{array}\end{array}\right)$$
(78)
The following section, where an approximate determination of the critical exponents of the $`3`$ and $`4`$state Potts models will be proposed, will make use of the numerical values listed in table 3.
## 5 Discussion
Finite-size estimators of the critical exponents follow from equations (18) and (19),
$$x_\varphi =\frac{\mathrm{ln}\varphi (L)\mathrm{ln}\varphi (L^{})}{\mathrm{ln}L\mathrm{ln}L^{}}.$$
(79)
The results, collected in table 4 in the case of the Ising, $`3`$ and $`4`$state Potts models, show as expected a very weak convergence. A strip of width $`2`$ to $`4`$ is obviously a poor approximation of the thermodynamic limit. In the case of the Ising model, exact diagonalization through Jordan-Wigner transformation into fermion operators, then Bogoljubov-Valatin canonical transformation into free fermions leads to an exact expression for the energy density
$$\mathrm{𝙴𝚗𝚎𝚛𝚐𝚢}\mathrm{𝚍𝚎𝚗𝚜𝚒𝚝𝚢}=\frac{2}{L}\mathrm{cos}\left(\frac{\pi }{2L}\right).$$
(80)
Unfortunately, there is no closed expression for the order parameter matrix element which couples two different sectors of the Hamiltonian. Hence, the presence of a boundary term breaks the quadratic expression necessary for the diagonalization.
The resort to gap scaling (23) as predicted from conformal invariance is thus desirable. Indeed, as emphasized in Section 2, the cylinder geometry is the geometrical shape of the classical problem corresponding to the one-dimensional quantum chain. Since such an infinitely long cylinder follows from the infinite plane geometry through conformal mapping, we may argue that the thermodynamic limit is somehow encoded in the results following from conformal rescaling and a better convergence for the critical exponents is expected. With the normalization of equation (41) the fermion excitations take the form $`\epsilon _k=|2\mathrm{sin}\frac{k}{2}|`$. The sound velocity is thus fixed to unity as previously announced. The results following from gap scaling,
$$x_\varphi =\frac{L\times \mathrm{𝚐𝚊𝚙}_\varphi }{2\pi v_s},$$
(81)
are collected in table 5. The quality of the results is indisputably better than through FSS.
In the case of the Potts model, the sound velocity is not known and it is worth referring to the literature. Gehlen et al. have studied numerically quantum Potts chains in refs. where tables of numerical results are reported for $`q=3`$ up to a size $`L=13`$ and for $`q=4`$ up to $`L=11`$ (they denote $`R`$ ($`P`$) the energy gap (magnetic gap) multiplied by the strip size) . These authors used a different normalization and their Hamiltonian for $`q=3`$ in is related to ours through $`\frac{2}{3}\widehat{\text{H}}_{\mathrm{𝚑𝚎𝚛𝚎}}(q=3)+\frac{4}{3}\widehat{\text{1}}=\widehat{\text{H}}_{\mathrm{𝙶𝚎𝚑𝚕𝚎𝚗}}(q=3)`$, while in it is related to ours through $`\widehat{\text{H}}_{\mathrm{𝚑𝚎𝚛𝚎}}(q=4)=2\widehat{\text{H}}_{\mathrm{𝙶𝚎𝚑𝚕𝚎𝚗}}(q=4)`$. Using the values of the sound velocity quoted in , we deduce that in our case the sound velocity is $`3/2`$ larger for $`q=3`$ and $`2`$ times larger for $`q=4`$, i.e.
$$v_s(q=3)=1.299,v_s(q=4)=1.578.$$
(82)
These values are used to calculate the exponents reported in table 5. Further numerical values denoted with an asterisk are taken from refs. to complete the table with results at larger sizes. The expected result in the thermodynamic is also mentioned in the table .
To summarize, we note that the quality of the results is acceptable for the relatively small effort of diagonalization of small matrices. The quest for critical exponents is an important step in the characterization of the nature of a phase transition, since these quantities are universal. The study of models of statistical physics which display second-order phase transitions usually requires sophisticated methods, and approximate determinations are often desirable. Considering quite small systems, we reach a few percent accuracy in the determination of the critical exponents in the case of the Ising and $`3`$state Potts models. The convergence is poor in the $`4`$state Potts case. This is essentially due to the logarithmic corrections present in this model.
Eventually, regarding the relatively small efforts, we believe that an introduction to the study of quantum chains in courses on statistical physics or many-body problems might be of interest.
## References |
warning/0507/cond-mat0507621.html | ar5iv | text | # Effects of 𝜏₁ Scattering on Fourier-Transformed Inelastic Tunneling Spectra in High-𝑇_𝑐 Cuprates with Bosonic Modes
(June 21, 2005)
## Abstract
We study the $`\tau _1`$-impurity induced $`𝐪`$-space pattern of the energy derivative local density of states (LDOS) in a $`d`$-wave superconductor. We are motivated in part by the recent scanning tunneling microscopy (STM) observation of strong gap inhomogeneity with weak charge density variation in Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+δ</sub> (BSCCO). The hypothesis is that the gap inhomogeneity might be triggered by the disorder in pair potential. We focus on the effects of electron coupling to various bosonic modes, at the mode energy shifted by the $`d`$-wave superconducting gap. The pattern due to a highly anisotropic coupling of electrons to the $`B_{1g}`$ phonon mode is similar to preliminary results from the Fourier transformed inelastic electron tunneling spectroscopy (FT-IETS) STM experiment in BSCCO. We discuss the implications of our results in the context of band renormalization effects seen in the ARPES experiments, and suggest means to further explore the electron-boson coupling in the high-$`T_c`$ cuprates.
The extent to which collective excitations of high-$`T_c`$ cuprates are manifested in their single particle spectra is a long standing issue. The band renormalization effects, seen in the ARPES Damascelli03 (as well as in the break junction tunneling experiments Zasadzinski01 ), have the characteristics of an electron-bosonic mode coupling. The “41 mev” spin resonance mode, a prominent feature in the spin excitation spectrum, is a natural candidate for electrons to couple to Norman98 ; Abanov99 ; Eschrig00 . However, there are also phonon modes of similar energies, and they may instead have the strongest influence on the single electron spectra Lanzara01 ; Cuk04 ; Sandvik04 ; Devereaux04 . At present, ARPES alone appears inadequate to differentiate the two scenarios. Sometime ago, several of us proposed Zhu04 an FT-IETS STM technique as a complimentary means to study this issue. The technique takes advantage of the pioneering work of the Fourier transformed STM Hoffman02b ; McElroy03 , and combines it with the vintage IETS Jakievic66 ; Scalapino67 ; Balatsky03 .
Central to this technique is the Fourier transform of the energy derivative of tunneling conductance map in real space $`d^2I/dV^2(𝐫,eV)d^2I/dV^2(𝐪,eV)`$. This $`𝐪`$-space map, which can also be called Fourier map, contains information about inelastic scattering in the system. Theoretically, one finds peaks in $`𝐪`$ space and energy $`eV`$ in this Fourier map of IETS that are related to the inelastic scattering off some collective excitations in the system. In the case of electron-spin mode coupling, the FT-LDOS near an ordinary potential scattering center (a $`\tau _3`$ impurity in the Nambu space) at the energy of $`E=\pm (\mathrm{\Delta }_0+\mathrm{\Omega }_0)`$ was shown to have sharp features at momenta close to $`(\pi ,\pi )`$. (Here $`\mathrm{\Delta }_0`$ is the maximum of the $`d`$-wave superconducting energy gap and $`\mathrm{\Omega }_0`$ the mode energy.)
Recently preliminary results from the first FT-IETS STM experiment has been reported in BSCCO Lee05 . While features are observed in the Fourier-transformed $`d^2I/dV^2`$ at the expected energy range, observed intensity near $`(\pi ,\pi )`$ is low. Instead, the strongest intensity appears at the wavevectors parallel to the Cu-O bond directions. These experimental results have in turn motivated us to compare the FT-IETS spectra near a potential scatterer in the cases of electrons coupled to the spin resonance and various phonon modes Zhu05 . It was shown that all cases contain sharp features near $`(\pi ,\pi )`$, in disagreement with the experimental spectrum.
This raises the question of whether the $`\tau _3`$ scatterer correctly describes the impurities in Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+δ</sub> (BSCCO). The origin of ubiquitously observed nanoscale inhomogeneity in BSCCO has been investigated by studying the correlation between the inhomogeneity and the position of oxygen dopants McElroy05 . It is shown that the local electronic states are not associated with charge density variations. To account for those features, Nunner et al. Nunner05 proposed to look at the effect of random pairing potential fluctuations, the so-called $`\tau _1`$ disorder. In what follows, we will address the effect of $`\tau _1`$ disorder on FT-IETS signatures.
In this Letter, we take the disorder in pair potential as the scattering center, which is of $`\tau _1`$ character in the Nambu space, and study the Fourier component of the energy derivative local density of states, $`d^2I/dV^2`$, at $`E\pm (\mathrm{\Omega }_0+\mathrm{\Delta }_0)`$ around such a scatterer. We considered a few typical bosonic modes as a possible scattering modes that produce IETS fingerprints. We find that the results for $`\tau _1`$ are qualitatively different from the case of potential disorder: 1) there are no strong signatures near in the $`𝐪`$-space near $`(\pi ,\pi )`$ in any of the electron-boson couplings; 2) the highly anisotropic coupling of electrons to the out-of-plane out-of-phase oxygen buckling $`B_{1g}`$ phonon mode, gives rise to a Fourier pattern similar to the IETS-STM experiment in BSCCO Lee05 . Our results are also consistent with the in-plane breathing mode although the agreement with the data is not as good.
We start with a model Hamiltonian for a two-dimensional $`d`$-wave superconductor with the coupling of electrons to bosonic modes:
$$=_{BCS}+_{elboson}+_{imp}.$$
(1)
Here the bare BCS Hamiltonian, $`_{BCS}=_{𝐤,\sigma }\xi _𝐤c_{𝐤\sigma }^{}c_{𝐤\sigma }+_𝐤(\mathrm{\Delta }_𝐤c_𝐤^{}c_𝐤^{}+\mathrm{\Delta }_𝐤^{}c_𝐤c_𝐤)`$, where the normal state energy dispersion is given by Norman95 , $`\xi _𝐤=2t_1(\mathrm{cos}k_x+\mathrm{cos}k_y)4t_2\mathrm{cos}k_x\mathrm{cos}k_y2t_3(\mathrm{cos}2k_x+\mathrm{cos}2k_y)4t_4(\mathrm{cos}2k_x\mathrm{cos}k_y+\mathrm{cos}k_x\mathrm{cos}2k_y)4t_5\mathrm{cos}2k_x\mathrm{cos}2k_y\mu `$, with $`t_1=1`$, $`t_2=0.2749`$, $`t_3=0.0872`$, $`t_4=0.0938`$, $`t_5=0.0857`$, and $`\mu =0.8772`$, and the $`d`$-wave gap dispersion $`\mathrm{\Delta }_𝐤=\frac{\mathrm{\Delta }_0}{2}(\mathrm{cos}k_x\mathrm{cos}k_y)`$. Unless specified explicitly, the energy is measured in units of $`t_1`$ hereafter. The coupling of the electrons to bosonic modes is modeled by the Hamiltonian $`_{elboson}=\frac{1}{\sqrt{N_L}}_{\genfrac{}{}{0pt}{}{𝐤,𝐪}{\sigma }}g_\nu (𝐤,𝐪)c_{𝐤+𝐪,\sigma }^{}c_{𝐤\sigma }(b_{\nu 𝐪}+b_{\nu ,𝐪}^{})`$ for the buckling $`B_{1g}`$ ($`\nu =1`$) and the in-plane half breathing ($`\nu =2`$) modes, while $`_{elboson}=\frac{g_0}{2N_L}_{\genfrac{}{}{0pt}{}{𝐤,𝐪}{\sigma ,\sigma ^{}}}c_{𝐤+𝐪,\sigma }^{}(𝐒_𝐪𝝈_{\sigma \sigma ^{}})c_{𝐤,\sigma ^{}}`$ for the spin resonance mode. For the phonon modes, we consider the cases where the coupling matrix element is either highly anisotropic, dependent on both $`𝐤`$ and $`𝐪`$, or is only $`𝐪`$ dependent. In the following we use the notation $`B_{1g}`$-I and $`br`$-I for the former type of phonon modes while $`B_{1g}`$-II and $`br`$-II for the latter type. Detailed expression of the coupling matrix elements for these types of phonon modes can be found in Ref. Zhu05 . The third term describes the quasiparticles scattered off a $`\tau _1`$ impurity due to the inhomogeneity in pair potential rather than off a conventional $`\tau _3`$ impurity. In the following, we consider a single $`\tau _1`$ impurity in a $`dwave`$ superconductor. The resulting Fourier pattern should survive a white-noise random distribution of such $`\tau _1`$ impurities in a realistic system. The impurity part of Hamiltonian can then be written as:
$$_{imp}=\delta \mathrm{\Delta }\underset{\delta }{}\eta _\delta [c_0^{}c_\delta ^{}+c_\delta ^{}c_0^{}+H.c.],$$
(2)
where $`\eta _\delta =1(1)`$ for $`\delta =\widehat{x}(\widehat{y})`$.
To be relevant to recent experimental realization, where no impurity-induced resonance state was observed, we assume the $`\tau _1`$ impurity to have a weak scattering potential $`\delta \mathrm{\Delta }`$. In this limit, we employ the Born approximation and arrive at the correction to the LDOS at the $`i`$-th site, summed over two spin components:
$$\delta \rho (𝐫_i,E)=\frac{2\delta \mathrm{\Delta }}{\pi }\underset{\delta }{}\eta _\delta \text{Im}[\widehat{𝒢}(i,0;E+i\gamma )\widehat{\tau }_1\widehat{𝒢}(\delta ,i;E+i\gamma )]_{11},$$
(3)
where $`\widehat{𝒢}`$ is the Green’s function dressed with the bosonic renormalization effect and defined in the Nambu space Zhu05 . From the perspective of the IETS, the energy derivative of the LDOS, $`\frac{d\delta \rho (𝐫_i,E)}{dE}`$, is the quantity we are most interested in. It corresponds to the derivative of the local differential tunneling conductance (i.e., $`d^2I/dV^2`$). The Fourier component of the differential LDOS is then given by $`\frac{d\delta \rho (𝐪,E)}{dE}=_i\frac{d\delta \rho (𝐫_i,E)}{dE}e^{i𝐪𝐫_i}`$ with the spectral weight defined as $`P(𝐪,E)=\left|\frac{d\delta \rho (𝐪,E)}{dE}\right|`$.
We consider here for comparison the coupling of electrons to spin resonance mode, $`B_{1g}`$ and breathing phonon modes. For the numerical calculation, we take the superconducting energy gap $`\mathrm{\Delta }_0=0.1`$, the frequency of all bosonic modes $`\mathrm{\Omega }_0=0.15`$. The $`\tau _1`$ impurity scattering strength $`\delta \mathrm{\Delta }`$ is taken to be 50% of the superconducting energy gap. The coupling strength for all types of bosonic modes is calibrated to give at the Fermi energy $`E=0`$ an identical frequency renormalization factor in the self energy. The same procedure as in Ref. Zhu05 is followed to obtain the Fourier spectral weight $`P(𝐪,E)`$.
In Figs. 1, we present the results of the Fourier spectrum, $`P(𝐪,E)`$, at the energy $`E=(\mathrm{\Delta }_0+\mathrm{\Omega }_0)`$ for a $`d`$-wave superconductor with the electronic coupling to the bosonic modes. For comparison, the same spectrum is also shown (last panel) for the case of no mode coupling. Note that the case without the mode coupling, the energy $`\mathrm{\Omega }_0`$ has no special meaning in the context of the electronic properties, and the energy $`E=(\mathrm{\Delta }_0+\mathrm{\Omega }_0)`$ is chosen merely for comparison to the case of mode coupling. The energy $`E=(\mathrm{\Delta }_0+\mathrm{\Omega }_0)`$ corresponds to the position where the bosonic modes are excited, signaling a peak in the IETS $`d^2I/dV^2`$-$`V`$ tunneling spectrum Zhu05 . First of all, the Fourier maps for all cases does not display any peak structure at large $`𝐪`$ near $`(\pm \pi ,\pm \pi )`$ and $`(\pm \pi ,\pi )`$, which appears persistently with the $`\tau _3`$ scattering Zhu05 . Instead the Fourier spectral weight is minimal in intensity (dark blue area in the figure) in these regions. The map for the case of electronic coupling to the $`B_{1g}`$-I phonon mode shows locally strongest intensity (red spots) at $`𝐪`$ about $`(\pm \frac{2\pi }{4},0)`$ and $`(0,\pm \frac{2\pi }{4})`$. The peak intensity at $`𝐪`$ near $`(\pm \frac{2\pi }{4},\pm \frac{2\pi }{4})`$ and $`(\pm \frac{2\pi }{4},\frac{2\pi }{4}`$) is much weaker than those along the bond directions. The map for the case of the coupling to the $`br`$-I phonon mode exhibits locally strongest intensity (red spots) at $`𝐪`$ near $`(\pm \frac{3\pi }{10},0)`$ and $`(0,\pm \frac{3\pi }{10})`$. In addition, each of these red spot has a double-tail structure, which is absent in the case of $`B_{1g}`$-I mode coupling. The maps for the cases of the $`B_{1g}`$-II and spin resonance mode coupling exhibit similar $`𝐪`$ structure. The finite intensity is uniformly distributed on a circular strip near $`|𝐪|=\frac{2\pi }{4}`$ and becomes stronger as $`𝐪`$ approaches the zero point. No locally distinguishable strongest intensity peak can be identified at $`𝐪=(\pm \frac{2\pi }{4},0)`$ and $`(0,\pm \frac{2\pi }{4})`$. The map for the coupling to the $`br`$-II phonon mode exhibits locally the highest intensity (red spots) at $`𝐪`$ near $`(\pm \frac{2\pi }{4},\pm \frac{2\pi }{4})`$ and $`(\pm \frac{2\pi }{4},\frac{2\pi }{4})`$. No peaks are found at $`𝐪`$ near $`(\pm \frac{2\pi }{4},0)`$ and $`(0,\pm \frac{2\pi }{4})`$. The map for the case of no mode coupling shows an eight-tail star shape at $`𝐪=(0,0)`$, which consists of the head-on overlap of four red spots, such as those appearing in the case of the $`br`$-I phonon mode coupling each with two tails. As we have already emphasized, experimentally, the Fourier map of $`d^2I/dV^2`$ shows intensity peaks only at $`𝐪=(\pm \frac{2\pi }{5},0)\pm 15\%`$ and $`(0,\pm \frac{2\pi }{5})\pm 15\%`$ Lee05 . Therefore, by comparison with the experimental data, our new FT-IETS STM analysis also supports the notion Zasadzinski01 ; Norman98 ; Abanov99 ; Eschrig00 ; Lanzara01 ; Cuk04 ; Sandvik04 ; Devereaux04 that the electronic band must be renormalized by its coupling to the bosonic modes. In particular, the results based on the scenario of highly anisotropic coupling of electrons to the $`B_{1g}`$ phonon mode are in best agreement with the IETS-STM data in BSCCO.
In Fig. 2, we present the energy evolution Fourier pattern for the electronic coupling to the $`B_{1g}`$-I mode. It shows that the intensity peak structure along the bond direction is robust against the energy change. The characteristic $`q`$ vector, at which the locally highest intensity is located, decreases slightly with the increased energy. This result is also not inconsistent with the preliminary experiment.
Our results demonstrate the important role the character of the scattering center plays in the Fourier spectrum. To explore this issue further, we note that, in the case of $`\tau _1`$ scattering considered here, the Fourier spectrum can be expressed as follows,
$`\delta \rho (𝐪;E)={\displaystyle \frac{2\delta \mathrm{\Delta }}{N_L}}{\displaystyle \underset{𝐤}{}}(\mathrm{cos}k_x\mathrm{cos}k_y)`$
$`\times \{[A(𝐤;E)K(𝐤+𝐪;E)+K(𝐤;E)A(𝐤+𝐪;E)]`$
$`+[J(𝐤;E)B(𝐤+𝐪;E)+B(𝐤;E)J(𝐤+𝐪;E)]\},`$ (4)
where
$`A(𝐤;E)`$ $`=`$ $`{\displaystyle \frac{2}{\pi }}\text{Im}[𝒢_{11}(𝐤;E+i\gamma )],`$ (5)
$`B(𝐤;E)`$ $`=`$ $`\text{Re}[𝒢_{11}(𝐤;E+i\gamma )],`$ (6)
$`J(𝐤;E)`$ $`=`$ $`{\displaystyle \frac{2}{\pi }}\text{Im}[𝒢_{12}(𝐤;E+i\gamma )],`$ (7)
$`K(𝐤;E)`$ $`=`$ $`\text{Re}[𝒢_{12}(𝐤;E+i\gamma )].`$ (8)
This expression shows that, for the $`\tau _1`$ scattering, the Fourier spectrum is determined by the $`𝐤`$-summation of the product terms constituting the imaginary (real) parts of the single-particle $`(𝒢_{11})`$ with the real (imaginary) parts of anomalous $`(𝒢_{12})`$ Green’s function in the superconducting state, weighted by a $`d`$-wave type form factor $`\mathrm{cos}k_x\mathrm{cos}k_y`$. This scattering process with the $`\tau _1`$ impurity is significantly different than the case of a $`\tau _3`$ impurity scattering Zhu05 , where the convolution takes place between the real and imaginary parts of the same Green’s function without the form factor. This difference of the scattering process matters significantly in the resulting Fourier map. As shown in Eq. (4), the form factor $`\mathrm{cos}k_x\mathrm{cos}k_y`$ appearing in the $`\tau _1`$ scattering case is identically zero along the diagonals in the first Brillouin, but reaches a maximum at the $`M`$ points \[$`𝐤=(\pm \pi ,0)`$ and $`(0,\pm \pi )`$\] on the zone boundary . It then follows that any stronger intensity from the product of $`AK`$, $`BJ`$ connected by a $`𝐪`$ oriented close to the diagonals is strongly suppressed, while the intensity connected by $`𝐪`$ oriented parallel to the bond direction is enhanced. For the electronic coupling to the $`B_{1g}`$-I phonon mode, it has been found Zhu05 that there are moderately strong intensity on the two split beams oriented perpendicular to the zone boundary at M points in the function $`A`$, $`B`$, $`J`$, and $`K`$. The form factor $`\mathrm{cos}k_x\mathrm{cos}k_y`$ then tips the relative contribution from the product $`AK`$ and $`BJ`$, giving rise to locally highest intensity at $`𝐪=(\pm \frac{2\pi }{4},0)`$ and $`(0,\pm \frac{2\pi }{4})`$ in the Fourier map (see the first panel of Fig. 1). These split beams of intensity are absent for the electronic coupling to other modes.
Our results naturally suggest additional means to further explore the electron-bosonic mode coupling experimentally. For instance, a Zn impurity acts as a nonmangetic potential center – a $`\tau _3`$ scatterer. In this case, the sharp features near $`(\pi ,\pi )`$ should be observed in the FT-IETS spectrum. In the case of low-energy elastic scattering interference of quasiparticles, related effects have in fact been demonstrated. Indeed, strong signatures near $`(\pm \pi ,\pm \pi )`$ and $`(\pm \pi ,\pi )`$ appear in the theoretical spectra near a $`\tau _3`$ scatterer Wang03 . These features are not observed experimentally in the stoichiometrical BSCCO Hoffman02b ; McElroy03 , but are seen around a nomagnetic Zn impurity in the doped BSCCO.
In conclusion, we have studied, for the first time, the $`\tau _1`$-impurity induced Fourier pattern of the energy derivative local density of states in a $`d`$-wave superconductor with the coupling of electrons to various bosonic modes. We consider $`B_{1g}`$, half-breathing, and spin $`(\pi ,\pi )`$ modes. Our results show that, at the mode energy shifted by the $`d`$-wave superconducting gap energy $`\mathrm{\Delta }_0`$, i.e., $`E=\pm (\mathrm{\Delta }_0+\mathrm{\Omega }_0)`$, the coupling of electrons to the $`B_{1g}`$ or breathing phonon modes, gives rise to a Fourier pattern similar to the preliminary Fourier transformed IETS-STM experiment in BSCCO Lee05 . The coupling of electrons to the spin resonance mode, on the other hand, yields a Fourier spectrum that is inconsistent with the experiment. These results do not necessarily rule out the role of the spin-spin interactions as being relevant for superconductivity in BSCCO, instead they imply that electron-phonon coupling has a strong impact on the superconducting electronic structure. These results have important implications for our understanding of the electronic properties of the cuprates. They also demonstrate the potential of the FT-IETS STM technique, and highlight the importance of $`\tau _1`$ scattering in the impurity-free BSCCO Nunner05 .
We thank D.-H. Lee, N. Nagaosa, M. R. Norman, D. J. Scalapino, and Z. X. Shen for very useful discussions. This work was supported by the US DOE (J.X.Z. and A.V.B.), the NSERC, the Office of Naval Research under Grant No. N00014-05-1-0127, and the A. von Humboldt Foundation (T.P.D.), the NSF under Grant No. DMR-0424125 and the Robert A. Welch Foundation (Q.S.), the Office of Naval Research under grant N00014-03-1-0674, the NSF under Grant No. DMR-9971502, the NSF-ITR FDP-0205641, and the Army Research Office under Grant No. DAAD19-02-1-0043 (K.M., J.L., and J.C.D). |
warning/0507/hep-th0507067.html | ar5iv | text | # Flow of Geometries and Instantons on the Null Orbifold
## 1 Introduction
The treatment of time-dependent backgrounds and spacelike singularities remains among the main puzzles of string theory. The importance of the problem has led to a considerable amount of work, and some progress has been made using perturbative and non-perturbative techniques, such as exact CFT’s, AdS/CFT etc. (-). Examples of the goals of this program include understanding the quantum state at the big-bang singularity, or at the singularity of the black hole .
In particular, the role of $`\alpha ^{}`$ corrections vs. perturbative $`g_s`$ corrections vs. non-perturbative stringy corrections, at such singularities, is unclear. It is likely that all are needed in order to understand the singularity in detail. This is the case in more familiar stringy singularities, such as the conifold, for which key effects are understood by quantitative non-perturbative effects, such as wrapped D-branes becoming light , but a more detailed understanding uses LST or double-scaled LST . The latter focuses on a smaller set of degrees of freedom localized near the singularity (and is a solvable CFT with a varying $`g_s`$), which clarifies where and how non-perturbative effects set in, as well as being computationally useful.
Other effects which are likely to be important are effects beyond 2nd quantized string theory. Already at the level of the Einstein-Hilbert action generic singularities exhibit a strong mixing property - the BKL dynamics (for a review of its recent appearance in supergravity see ). There are no proposals for a tractable stringy formalism to deal with such mixing (beyond the low energy effective action, which breaks down close to the singularity).
We list below some of the motivation for exploring $`\alpha ^{}`$ effects around the singularity:
1. In , a search was carried out for remnants of the BH singularity at large N but weak ’t-Hooft coupling. No indications of the singularity were found. This might suggest that the singularity is resolved by $`\alpha ^{}`$ corrections.
2. $`\alpha ^{}`$ effects already exhibit interesting and unexpected behavior near the singularity. In it was shown that already at sphere level a non-commutative geometry like structure appears. The latter delocalized the twisted sector states over large distances in spacetime. In particular in it was shown that twisted sector pair creation occurs near the singularity.
3. For the non-rotating extremal BTZ in $`AdS_3`$, one can construct candidates for the microstates of the BH (for a recent review see ). The proposed microstates of the BH are solutions of Einstein-Hilbert action, and in particular do not require higher $`g_s`$ corrections. However, they do require $`\alpha ^{}`$ corrections in order, for example, to understand the $`_N`$ singularities which occur in some of these configurations.
4. It was recently shown that the singularity and horizon structure of a class of supersymmetric black holes changes significantly by tree level or 1-loop higher order curvature corrections to the effective action
5. Although not directly related yet, a problem of singularities in causal structure also appears in purely 2D field theory context. The Coulomb branch of $`\mathrm{SU}(2)`$ (4,4) gauge theory in two dimensions with a single quark flavor has a metric which is not positive definite in the IR<sup>1</sup><sup>1</sup>1The metric is believed to receive only a 1-loop corrections in perturbation theory.. Since the model is the IR of a perfectly well defined unitary supersymmetric field theory, this problem has to be resolved within the 2D field theory. One usually says that the correct degrees of freedom were not identified properly in the IR (already in terms of the 2D field theory).
We are interested in mapping what string theory considers to be small deformations of a singularity, in our case the null orbifold singularity , as a step towards understanding its large deformation, which might be relevant for its resolution. For the null orbifold, we will explore the relation between it and familiar $`_N`$ orbifolds that posses a mild (and well understood in string theory) time like singularity.
One can also consider a more detailed role that “near by geometries” can play. Consider for example the relation between the fuzzball states of and the BTZ geometry. One way to reconcile the validity of the two descriptions is examine the amount of mixing the micro-states undergo under any attempt to probe them. Since they differ from each other only over small length scales<sup>2</sup><sup>2</sup>2We would like to thank S. Ross for a discussion of this point., they clearly mix under any such perturbation. This implies that the effective geometry may not be that of the microstates but something else - perhaps for some purposes it is the original BTZ black hole geometry. The situation might be analogous to that of some field theories which at zero temperature exhibit spontaneous symmetry breaking, but exhibit symmetry restoration at finite temperature (above a threshold). The microstates are analogous, in this very rough analogy, to the true vacua, and the finite temperature minimum in the origin is analogous to the original black hole singularity, which dominates the dynamics once mixing is taken into account. In both cases a complicated enough process, with enough energy, will be dominated by the symmetric phase (=the background with the black hole singularity) although most pure states (and in particular the ground state) are in the broken symmetry phase (which is like the micro-state description). To go from the symmetric phase to the broken phase one usually condenses a tachyon (at zero temperature), and hence we would like to explore the analogues of these tachyons in the case of the spacelike (or null) singularity.
The null-orbifold singularity has a very concrete relation to the $`_N`$ singularity. Already in this relation was touched upon, and we develop it further in the current paper. In section 2 we show how precisely the two-cone null-orbifold is a large N limit of the better understood single cone $`/_N`$ orbifold (as well as the subtleties associated with this limit). In section 3 we discuss the action of D(-1) branes in this background. In section 4 we explore the transition from the null orbifold towards the $`_N`$ orbifolds, after condensing an $`N`$-twisted sector state in the null-orbifold.
This situation can be summarized in the following diagram
$$\begin{array}{ccc}^1\times /_n& \underset{\mathrm{}\text{boosted}}{\overset{n\mathrm{}}{}}& ^{1,2}/\text{Null}\\ \text{tachyon}\text{sector }k& & \text{tachyon}\text{sector }m& & \\ ^1\times /_{\frac{n}{k}=m}& \underset{\mathrm{}\text{boosted}}{\overset{n,k\mathrm{},\frac{n}{k}=m}{}}& ^1\times /_m\end{array}$$
(1)
The left downward point arrow is the flow of . The upper rightward pointing arrow is section 2. The right downward pointing arrow is section 4 and is the main conclusion of the paper, which presents evidence that upon condensation of an $`N`$ twisted sector mode of the null-orbifold a $`_N`$ singularity appears<sup>3</sup><sup>3</sup>3Although in the extreme boost limit, as we will discuss later..
Section 5 contains a summary and conclusions.
Further details on the $`_Nnullorbifold`$ are provided in appendix A. Appendix B quantizes the RNS string on the null orbifold and shows the emergence of a logarithmic CFT.
As we completed this, two papers appeared which discuss related models from a different perspective
## 2 The Null-orbifold and $`\times /_N`$ Orbifold
String theory in orbifolds of the form $`(^{1,2}/\mathrm{\Gamma })\times 𝒞^{}`$ with $`\mathrm{\Gamma }`$ generating a group isomorphic to $``$ or $`_N`$ were extensively studied. The $`_N`$ orbifolds (where $`\mathrm{\Gamma }`$ is in the elliptic class of $`\mathrm{SO}(1,2)`$) are well understood (for a review see ), and the geometry is a cone perpendicular to the time direction. The null-orbifold studied in is generated by $`\mathrm{\Gamma }`$ in the parabolic class of $`\mathrm{SO}(1,2)`$, and the geometry consists of two three dimensional cones with a common tip and a singular plane crossing the tip. Unlike the $`_N`$ orbifolds the null-orbifolds is a singular time-dependent background<sup>4</sup><sup>4</sup>4Although it possesses a null isometry.. The singularity at the origin of the cones is still not completely understood.
In this section we review the quantization of the $`_N`$ and the null orbifold. We introduce the construction of the latter from the former by a infinite boost. The result we obtain is that the $`_N`$ orbifold converges to the two-cone null orbifold. This is shown using both the classical geometry and the 1st quantized string.
### 2.1 Classical Results (Geometry)
To describe orbifolds of $`^{1,2}`$ we take the coordinates $`x^0,x^1,x^2`$ on $`^{1,2}`$, with the flat metric $`ds^2=d(x^0)^2+d(x^1)^2+d(x^2)^2`$ and consider the Killing vector:
$$J(a,b)=bJ^{02}+aJ^{12}$$
(2)
where $`J^{02}=x^0_2+x^2_0`$ is a boost and $`J^{12}=x^1_2x^2_1`$ is a rotation.
For $`b<a`$,$`J(a,b)`$ is in the elliptic class, and conjugate to $`J(\sqrt{a^2b^2},0)`$ using some Lorentz transformation M. The null orbifold is given by $`a=b`$. Choosing $`\sqrt{a^2b^2}=1/N`$, the generator of the null orbifold, $`a=b`$, is identified as the limit $`N\mathrm{}`$ of a sequence $`M_NJ(1/N,0)M_N^1`$, where $`M_N`$ is an N-dependent Lorentz transformation. This Lorentz transformation is singular when $`N\mathrm{}`$. It is useful to write the explicit form of $`M_N`$.
$$M_N=\left(\begin{array}{ccc}\frac{a}{\sqrt{a^2b^2}}& 0& \frac{b}{\sqrt{a^2b^2}}\\ \frac{b}{\sqrt{a^2b^2}}& 0& \frac{a}{\sqrt{a^2b^2}}\\ 0& 1& 0\end{array}\right)$$
(3)
$$J(a,b)=M_NJ(1/N,0)M_N^1$$
(4)
For the pure rotation case ($`b=0`$) we will use the familiar convention:
$$(x^0,z,\overline{z})(x^0,e^{\frac{2\pi i}{N}}z,e^{\frac{2\pi i}{N}}\overline{z})\text{with}z=\frac{x^1+ix^2}{\sqrt{2}}.$$
(5)
The generator of the null Orbifold twisting operator acts on the light cone coordinates
$$x^+=\frac{x^0x^1}{\sqrt{2}}x=x^2x^{}=\frac{x^0+x^1}{\sqrt{2}}$$
(6)
by
$$J^{null}=a\left(\begin{array}{ccc}0& 0& 0\\ \sqrt{2}& 0& 0\\ 0& \sqrt{2}& 0\end{array}\right)$$
(7)
We shall follow the convention of in defining the null orbifold by choosing $`a=\frac{\nu }{2\pi \sqrt{2}}`$ and the identification takes the following form
$`x^+x^+`$ (8)
$`xx+\nu x^+`$ (9)
$`x^{}x^{}+\nu x+{\displaystyle \frac{\nu ^2}{2}}x^+`$ (10)
A fundamental domain of the space looks like two 3 dimensional cones emanating from $`x^+=x=0`$, one towards $`x^+>0`$ and one towards $`x^+<0`$, and two codimension 1 cones in the plane $`x^+=0`$ pinching at $`x=x^{}=0`$. Any two null orbifolds (differing by the value of $`\nu `$) are related by boosts in the $`x^1`$ direction.
We show how by the procedure described above, of boosting by $`M_N`$, we may actually understand geometrically that there is a singular limit which relates the $`_N`$ space with the null orbifold space. This isn’t straightforward as the $`_N`$ is an one cone space and the null orbifold has two cones (that are not contained in the singular plane). Indeed for any finite $`N`$, we remain with an one cone fundamental domain. Some of the orbits in the $`x^+,x,x^{}`$ coordinates for large N are plotted in figure 1. One confirms the impression from the figure, that the orbits become localized around a fixed $`x^+`$ with a spread in $`x^+`$ which is $`\frac{1}{N}`$ that of $`x^{}`$. At $`N\mathrm{}`$ we observe that the slopes go to infinity. In this limit the orbits are contained in the $`xx^{}`$ plane, at fixed $`x^+`$. The infinity slopes orbits are exactly the parabolas of the null orbifold. Hence the single cone orbifold maps onto the two cone null-orbifold geometry.
### 2.2 Hilbert space: Untwisted sector
As is usual in orbifolds, the untwisted wave functions are projections of the wave functions of the covering space. The latter are plane waves on $`^{1,2}`$, and the untwisted sector of the orbifold (focusing on scalar functions) is
$$\mathrm{\Psi }_{k,s}^{orb.}=_{\mathrm{}}^{\mathrm{}}dse^{2\pi s\left(il+\widehat{J}\right)}\psi _k(x),l$$
(11)
where $`\widehat{J}`$ is the action of the null boost generator on the function $`\psi _k`$ (the wave function in flat space).
The formula for the $`\times /_N`$ orbifold is similar. The invariant wave functions of the elliptic orbifold are
$$\mathrm{\Psi }_{k,l}^N=\frac{N}{2\pi }e^{ik^0x^0}_0^{2\pi }d\theta e^{ik\overline{z}e^{i\theta }+i\overline{k}ze^{i\theta }+iN\theta l},l$$
(12)
$$\text{with,}k=\frac{k^1+ik^2}{\sqrt{2}}z=\frac{x^1+ix^2}{\sqrt{2}}$$
(13)
These integrals can be evaluated in terms of Bessel functions:
$$\mathrm{\Psi }_{k,l}^N=Ne^{ik^0x^0+iN\varphi l}J_{Nl}(2u),ue^{i\varphi }i\overline{k}z$$
(14)
It is easy to verify the completeness of this basis. We also choose $`k`$ to be real.
For the null-orbifold the integration in (11) is Gaussian (The operator $`\widehat{J}^{null}`$ is nilpotent of order 3), and the wave-function matches the results of (after fixing the phase and taking $`\nu =2\pi `$ until the end of the section):
$$\mathrm{\Psi }_{k,l}^{null}=\frac{\mathrm{exp}\left[i\left(k^+x^{}k^{}x^++\frac{(lk^+x^2)^2}{2k^+x^+}\right)\right]}{\sqrt{2\pi }\sqrt{ik^+x^+}},l$$
(15)
$$\text{where}k^\pm =\frac{1}{\sqrt{2}}(k^0k^1)x^\pm =\frac{1}{\sqrt{2}}(x^0x^1)$$
(16)
This is the wave function on the three dimensional cones. On the singular co-dimension 1 cones, it is a distribution.
We have shown that one can take the limit of the geometry. However, since we are interested in a CFT statement, it is more meaningful to show that the limit of the set of wave functions on the single cone $`R^2/Z_N`$ is the set of all wave functions on the null orbifold. Under an $`M_N`$ boost, the wave-function transforms as:
$`\mathrm{\Psi }_{k,l}^{(a,b)}(x)`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑se^{2\pi s\left(il+\widehat{J}(a,b)\right)}\psi _k(x)=`$ (17)
$``$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑s\widehat{M_N}e^{2\pi s\left(il+\widehat{J}(\frac{1}{N},0)\right)}\widehat{M_N^1}\psi _k(x)\mathrm{\Psi }_{\stackrel{~}{k},l}^N(M_N^1x)`$ (18)
where $`\stackrel{~}{k}_\mu =k_\nu (M_N)_{}^{\nu }{}_{\mu }{}^{}`$. In the limit $`N\mathrm{}`$, after normalizing the wave functions, we expect to find that:
$$\mathrm{\Psi }_{k,l}^{Null}(x)=\underset{N\mathrm{}}{lim}\mathrm{\Psi }_{kM_N,l}^N(M_N^1x).$$
(19)
To demonstrate that this limit is well defined (and not just formal) we will show it explicitly using the wave-function (14) and (15). The limit is defined so that the parameters $`a,b`$ in $`M_N`$ approach their final (common) value symmetrically<sup>5</sup><sup>5</sup>5This is a necessary requirement. For other prescriptions we have not been able to show that one obtains the null orbifold wave functions. keeping $`a>b`$
$`a={\displaystyle \frac{1}{\sqrt{2}}}+{\displaystyle \frac{1}{2\sqrt{2}}}N^2+O(N^4)b={\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle \frac{1}{2\sqrt{2}}}N^2+O(N^4)`$ (20)
Applying this transformation to the wave-functions (14) in the large N limit we find (further details are provided in appendix A):
$`\mathrm{\Psi }_{k;l}^{(a,b)}(x)=`$ $`{\displaystyle \frac{(\text{sign}(k^+x^+))^{Nl}}{\sqrt{2\pi ix^+k^+}}}\{`$ (21)
$`\mathrm{exp}\left[i{\displaystyle \frac{\left(lx^2k^++x^+k^2\right)^2}{2x^+k^+}}+ik^2x^2ik^+x^{}ik^{}x^++O(N^1)\right]+`$ (22)
$`+()^{Nl+1}\mathrm{exp}[i{\displaystyle \frac{\left(l+x^2k^+x^+k^2\right)^2}{2x^+k^+}}ik^2x^2i2N^2k^+x^++O(N^1)]\}`$ (23)
$`k^2`$ may be set to zero using a Lorentz transformation. Taking the limit $`N\mathrm{}`$ the second term is zero<sup>6</sup><sup>6</sup>6To see that, it should be considered as a distribution: any integral with a well behaved function vanishes in the limit. The exact statement is explained in appendix A. and the first term reduces exactly to the null-orbifold wave function<sup>7</sup><sup>7</sup>7The sign in front is irrelevant as $`Nl`$ can be chosen even throughout. (15), thus completing our proof.
### 2.3 Hilbert space: Twisted sectors
The authors of showed that the zero mode wave functions in the twisted sectors of the null-orbifold are (in our conventions of normalizing the wave functions):
$$\mathrm{\Psi }_{m,p^+,J}^w=\sqrt{\frac{1}{2\pi ip^+x^+}}\mathrm{exp}\left[ip^+x^{}i\frac{m^2}{2p^+}x^++i\frac{p^+}{2x^+}\left(x+\frac{J}{p^+}\right)^2i\frac{w^2(x^+)^3}{6(\alpha ^{})^2p^+}\right]$$
(25)
with $`J`$ and $`m`$ is the three dimensional mass given by the on shell condition:
$$2p^+p^{}=m^2=\frac{4}{\alpha ^{}}+\stackrel{}{p}_{}^{\mathrm{\hspace{0.33em}2}}$$
One observes that this wave function is the invariant combination of wave functions describing particle in the presence of an extremal configuration of electric-magnetic fields $`E=\pm B`$ (in analogy to ). The absolute value of this field is proportional to the twist parameter.
We now write the twisted sector wave functions of the $`\times /_N`$ orbifold in the same manner by considering a particle in the presence of a magnetic field<sup>8</sup><sup>8</sup>8For completeness, we derive this idea on general orbifold of $`^{1,2}`$ in the next subsection.. The most convenient covering space wave functions are radial strips.
$$\mathrm{\Psi }_{p^0,j,n_r}^{k/N}=e^{ij\varphi }e^{ip^0x^0}R_{\left|j\right|,n_r}\left(\frac{k}{N\alpha ^{}}r^2\right)$$
(26)
Where $`r,\varphi `$ are the usual polar coordinates on the plane, $`j`$ and the radial function is given by
$$R_{\left|j\right|,n_r}(\xi )=e^{\xi /2}\xi ^{\left|j\right|/2}F(n_r,\left|j\right|+1,\xi )$$
(27)
$`F`$ is the degenerate (confluent) hypergeometric function. The orbifold invariance condition is simple in these coordinates and is given by $`jN`$. The wave function describes a particle with mean distance of $`n_r\sqrt{\alpha ^{}}`$ from the origin and the wave function has $`j`$ oscillations. The conformal weight (energy) of this state is
$$E=\frac{\alpha ^{}}{4}(p^0)^2+\frac{k}{N}\left(n_r+\frac{1}{2}(\left|j\right|j+1)\right)$$
(28)
The relation to the usual states of the CFT, which are given by acting with creation quasi zero modes, is the following :
$$L_0=E\frac{k}{N}j=L_0\stackrel{~}{L}_0$$
(29)
Using these relations one may obtain the CFT meaning of $`n_r`$
$$\frac{k}{N}n_r=\frac{\alpha ^{}}{4}(p_0)^2+L_0\frac{1}{2}\left(\left|L_0\stackrel{~}{L}_0\right|(L_0\stackrel{~}{L}_0)+1\right)$$
(30)
In order to boost the wave function we boost the coordinates as in (19) but the quantum numbers are more subtle. In order to match the quantum numbers of the boosted wave function with the quantum numbers of the null orbifold it is very suggestive to use the geometric interpretation of $`n_r`$. Indeed we expect it to transform as $`(p^1)^2+(p^2)^2`$. The exact way is inferred from (28) and we simply replace $`n_r`$ everywhere by
$$n_r\frac{\alpha ^{}N}{4k}((p^1)^2+(p^2)^2).$$
(31)
The quantum number $`j/N`$ (which is integer) will be mapped exactly to $`l`$ as can already be seen from the fact the azimuthal part of the wave function (26) coincides with the untwisted wave function azimuthal part (14).
To recapitulate, the wave function of the twisted $`\times /_N`$ which is ready to be boosted with appropriately chosen quantum numbers
$$\mathrm{\Psi }_{p^0,\left|p\right|,l}^{k,N}=e^{iNl\varphi }e^{ip^0x^0}R_{\left|Nl\right|,\frac{\alpha ^{}N}{4k}\left|p\right|^2}\left(\frac{k}{N\alpha ^{}}r^2\right)$$
(32)
In doing the boost it is convenient to move to Whittaker functions which are related to $`R`$ according to
$$M_{\frac{\left|j\right|}{2}+n_r+\frac{1}{2},\frac{\left|j\right|}{2}}(z)=z^{1/2}R_{\left|j\right|,n_r}$$
(33)
Using the $`M_N`$ above, and the asymptotics of the Whittaker function one can show that (32) converges to (25).
### 2.4 First Quantization of the string
Applying the quantization scheme introduced in we quantize the bosonic string on $`(^{1,2}/\mathrm{\Gamma })\times _{}^{d3}`$. A more detailed discussion including quantization of the superstring is postponed to appendix B.
The worldsheet action and monodromies are:
$`S={\displaystyle \frac{1}{4\pi \alpha ^{}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\tau {\displaystyle _0^{2\pi }}𝑑\sigma \eta _{\mu \nu }\left(_\tau X^\mu _\tau X^\nu _\sigma X^\mu _\sigma X^\nu \right)`$ (34)
$`X(\sigma +2\pi ,\tau )=e^{2\pi w𝒥}X(\sigma ,\tau )`$ (35)
where $`w`$ is the twisted sector number and $`𝒥`$ is matrix defined from the differential realization of (2):
$$\widehat{J}(a,b)=X^\mu 𝒥_{\mu }^{}{}_{}{}^{\nu }\frac{}{X^\nu }$$
The mode expansion in a twisted sector $`w0`$:
$$\begin{array}{c}X^\mu (\sigma ,\tau )=\left[e^{w𝒥\sigma }\right]_{\nu }^{}{}_{}{}^{\mu }X_{z}^{}{}_{}{}^{\nu }(x,p;\tau )+i\sqrt{\frac{\alpha ^{}}{2}}\underset{n0}{}\left[\frac{e^{i(n+iw𝒥)(\sigma +\tau )}}{n+iw𝒥}\right]_{\nu }^{}{}_{}{}^{\mu }\alpha _n^\nu +\hfill \\ \hfill +i\sqrt{\frac{\alpha ^{}}{2}}\underset{n0}{}\left[\frac{e^{i(niw𝒥)(\sigma \tau )}}{niw𝒥}\right]_{\nu }^{}{}_{}{}^{\mu }\stackrel{~}{\alpha }_n^\nu \end{array}$$
(36)
With the zero-mode
$$\text{Null:}X_{z}^{}{}_{}{}^{\mu }(x,p;\tau )=\mathrm{cosh}\left(w\tau 𝒥\right)_{}^{\mu }{}_{\nu }{}^{}x^\mu +[11+\frac{1}{2}\mathrm{cosh}\left(w\tau 𝒥\right)]_{\nu }^{}{}_{}{}^{\mu }\frac{2\alpha ^{}\tau }{3}p^\nu $$
(37)
$$\text{Other:}X_{z}^{}{}_{}{}^{\mu }(x,p;\tau )=\mathrm{cosh}\left(w\tau 𝒥\right)_{}^{\mu }{}_{\nu }{}^{}x^\mu +\left[(w𝒥)^1\mathrm{sinh}(w𝒥\tau )\right]_{\nu }^{}{}_{}{}^{\mu }\alpha ^{}p^\nu $$
(38)
Where $`x^\mu `$ and $`p^\mu `$ satisfy the canonical commutation relations $`[x^\mu ,p^\nu ]=i\eta ^{\mu \nu }`$.
Introduce the operators:
$$\alpha _0=\sqrt{\frac{1}{2\alpha ^{}}}(w𝒥x+\alpha ^{}p)\stackrel{~}{\alpha }_0=\sqrt{\frac{1}{2\alpha ^{}}}(w𝒥x\alpha ^{}p)$$
(39)
Although assigned subscript $`0`$, these aren’t real zero modes in general and may possess nonzero conformal weight as we discuss below. Hence, we will properly name them quasi zero modes. The commutation relations between the modes:
$$[\alpha _n,\alpha _m]=\delta _{n+m}(n+iw𝒥)\eta ,[\stackrel{~}{\alpha }_n,\stackrel{~}{\alpha }_m]=\delta _{n+m}(niw𝒥)\eta ,[\alpha _n,\stackrel{~}{\alpha }_m]=0$$
(40)
The Virasoro generators are
$$L_n=\frac{1}{\alpha ^{}}\frac{dz}{2\pi i}z^{n+1}\text{:}XX(z)\text{:}=\frac{1}{2}\underset{m}{}{\scriptscriptstyle \genfrac{}{}{0.0pt}{}{}{}}\alpha _m\eta \alpha _{nm}{\scriptscriptstyle \genfrac{}{}{0.0pt}{}{}{}}+\delta _{n,0}a^X(w)$$
(41)
For any choice of normal ordering scheme of the quasi zero modes the zero point energy in the bosonic sector is
$$a^X(w)=\frac{w^2}{4}\mathrm{Tr}(𝒥^2)+vac\left|\frac{1}{2}\eta _{\mu \rho }\alpha _0^\mu \alpha _0^\rho \right|vac_w$$
(42)
The correct choice of normal-ordering depends on the $`\mathrm{SO}(1,2)`$ class of the orbifold identification generator:
* Elliptic class: the quasi zero modes have positive and negative conformal weights. Therefore, we naturally choose the positive conformal weight to annihilate the twisted vacuum state. The collection of states generated by acting with the creation operators is the Hilbert space of a (1st quantized) particle in a uniform magnetic field.
* Hyperbolic class: the quasi zero modes have pure imaginary conformal weight, as discussed in we should choose the reality conditions to be consistent with the commutation relations. This way, $`L_0`$ and $`\stackrel{~}{L}_0`$ are manifestly hermitian. The result is the Hilbert space of a particle in uniform electric field.
* Parabolic class: the quasi zero modes have conformal weight zero. The Hilbert space is of a particle in electric-magnetic fields which are equal in magnitude. One choice of quantization is that of (25).
## 3 D(-1) Instantons probes
### 3.1 The World-Volume Theory
We wish to examine the null orbifold deformed by some tachyon condensate. We shall probe the space with instantons (D(-1) branes), following the discussion of . In this section we develop the technology needed, beginning with the non deformed null orbifold. The open string theory on the instantons is a matrix theory of the collective coordinates (fermionic and bosonic). We focus on the bosonic degrees of freedom, they are parameterized by the covering space coordinates:
$$(X^+,X^{},X)^{2,1},(Y^3,Y^4,\mathrm{}Y^9)_{}^7$$
Under the projection each D(-1) instanton has infinitely many images, we use Chan-Paton indices that span the adjoint of $`\mathrm{U}(\mathrm{})`$ in the covering space. The null-orbifold projection should break the $`U(\mathrm{})`$ to $`U(1)^{\mathrm{}}`$, even at the singularity. We use the orbifold projection:
$`Y_{i,j}^a=Y_{i1,j1}^a,a=3\mathrm{}9`$ (43)
$`X_{i,j}^+=X_{i1,j1}^+`$ (44)
$`X_{i,j}=X_{i1,j1}+\nu X_{i1,j1}^+`$ (45)
$`X_{i,j}^{}=X_{i1,j1}^{}+\nu X_{i1,j1}+{\displaystyle \frac{1}{2}}\nu ^2X_{i1,j1}^+`$ (46)
These recursive equations are linear and can be easily solved. The solution can be neatly written using the following matrices (which also define an infinite closed algebra):
$$(\beta _l^m)_{i,j}\frac{(i+j)^l}{2^l}\delta _{i,jm}$$
(47)
$$[\beta _l^m,\beta _l^{}^m^{}]=2\underset{p=0}{\overset{l}{}}\underset{p^{}=0}{\overset{l^{}}{}}\left(\frac{m}{2}\right)^{l^{}p^{}}\left(\frac{m^{}}{2}\right)^{lp}\beta _{p+p^{}}^{m+m^{}}\delta _{p^{}p2+l^{}l+1}\left(\genfrac{}{}{0pt}{}{l}{p}\right)\left(\genfrac{}{}{0pt}{}{l^{}}{p^{}}\right)$$
(48)
The collective coordinates (bosonic) fields which solve (43):
$`Y^a={\displaystyle \underset{m}{}}y_m^a\beta _0^m`$ $`X={\displaystyle \underset{m}{}}x_m\beta _0^m+\nu x_m^+\beta _1^m`$ (49)
$`X^+={\displaystyle \underset{m}{}}x_m^+\beta _0^m`$ $`X^{}={\displaystyle \underset{m}{}}x_m^{}\beta _0^m+\nu x_m\beta _1^m+{\displaystyle \frac{\nu ^2}{2}}x_m^+\beta _2^m`$ (50)
along with the reality conditions
$$(y_m^a)^{}=y_m^a(x_m^+)^{}=x_m^+(x_m^{})^{}=x_m^{}(x_m)^{}=x_m.$$
The world volume low energy effective Lagrangian of the instantons may be obtained by dimensionally reducing the 10d Super Yang Mills. The action on the world-volume is<sup>9</sup><sup>9</sup>9The normalization factor $`Z_0`$ (which is simply $`_{\mathrm{}}^{\mathrm{}}1`$) is set to remove an overall infinite factor coming from the traces.:
$`\mathrm{S}=`$ $`{\displaystyle \frac{1}{2Z_0}}{\displaystyle \underset{\mu ,\nu =0}{\overset{9}{}}}\mathrm{Tr}\left([X^\mu ,X^\nu ][X_\mu ,X_\nu ]\right)=`$ (51)
$`=`$ $`{\displaystyle \underset{m+n+m^{}+n^{}=0}{}}[y_m^ay_m^{}^ax_n^+x_n^{}^+mm^{}+x_nx_m^+x_n^{}x_m^{}^+(mm^{}2nm^{})+`$ (52)
$`+2x_m^+x_n^{}x_m^{}^+x_n^{}^+nm^{}+\nu ^2x_m^+x_n^+x_m^{}^+x_n^{}^+{\displaystyle \frac{mm^{}n^2}{4}}]`$ (53)
It is useful to represent the above action using real bosonic fields living on $`S^1`$, identifying the modes as Fourier coefficients for a real field on $`S^1`$ as in :
$$u_m=_0^{2\pi }\frac{d\sigma }{\sqrt{2\pi }}U(\sigma )e^{im\sigma },U=(X^+,X^{},X,Y^3,\mathrm{}Y^9)$$
(54)
with the action:
$$\mathrm{S}=_0^{2\pi }\frac{d\sigma }{2\pi }\left[(X^+\dot{Y^a})^2(X\dot{X^+})^2+2\dot{X}\dot{X}^+XX^+2(X^+)^2\dot{X}^+\dot{X}^{}\frac{\nu ^2}{12}(\dot{X}^+)^4\right]$$
(55)
One may think of this action as the (analog of) T dual action, for D0 branes wrapping a space-like cycles. We show in the next subsection how we may identify the gauge fields and gauge invariant operators of the theory.
### 3.2 Symmetries
The action of the $`U(1)_m`$ factor in the gauge group on the fields (49) is generated by $`\beta _0^m`$, inducing the following transformation law for the modes:
$`e^{\alpha Q_m}y_n^a=y_n^a`$ $`e^{\alpha Q_m}x_n=x_n+\alpha \nu mx_{nm}^+`$ (56)
$`e^{\alpha Q_m}x_n^+=x_n^+`$ $`e^{\alpha Q_m}x_n^{}=x_n^{}+\alpha \nu mx_{nm}+{\displaystyle \frac{1}{2}}(\alpha \nu m)^2x_{n2m}^+`$ (57)
Although it is a $`0+0`$ model, one still considers the above transformations, which are symmetries of the Lagrangian, as gauge transformation. Using the $`\sigma `$ representation we can combine the transformation into a single gauged $`U(1)`$ and an associated arbitrary $`\mathrm{\Lambda }(\sigma )`$ living on the $`S^1`$ which is the gauge transformation $`0`$-form:
$`Y^a{}_{}{}^{}=Y^a`$ $`X{}_{}{}^{}=X+\nu _\sigma \mathrm{\Lambda }X^+`$ (58)
$`X^+{}_{}{}^{}=X^+`$ $`X^{}{}_{}{}^{}=X^{}+\nu _\sigma \mathrm{\Lambda }X+{\displaystyle \frac{\nu ^2}{2}}\left(_\sigma \mathrm{\Lambda }\right)^2X^+`$ (59)
The gauge transformations above are all connected to the identity of $`U(1)^{\mathrm{}}`$, however there is another element in the group which is disconnected from the identity (we will refer to it as ’large gauge transformation’). It’s action on the fields (49):
$$U_{ij}^{}=\underset{kl}{}(\beta _1^0)_{ik}^1U_{kl}(\beta _1^0)_{lj}$$
The transformations of this element on $`\sigma `$-representation fields (54):
$`Y^a{}_{}{}^{}=Y^a`$ $`X{}_{}{}^{}=X+\nu X^+`$ (60)
$`X^+{}_{}{}^{}=X^+`$ $`X^{}{}_{}{}^{}=X^{}+\nu X+{\displaystyle \frac{\nu ^2}{2}}X^+`$ (61)
These transformations (and all successive transformations generated by it) can be viewed as a modification on (58) by allowing a specific non-periodic boundary condition for the $`0`$-form $`\mathrm{\Lambda }(\sigma )`$.
The action (55) is also invariant under translations in $`\sigma `$, this is the quantum $`\mathrm{U}(1)`$ symmetry which insures the conservation of winding number of the string and the momentum in the T dual picture.
### 3.3 Moduli space
As a check of the formalism we will verify that the classical moduli space becomes the position of a single instanton in the fundamental domain of the null orbifold. The equations of motion derived are:
$`{\displaystyle \frac{d}{d\sigma }}\left[(X^+)^2{\displaystyle \frac{dY^a}{d\sigma }}\right]=0`$ (62a)
$`{\displaystyle \frac{d}{d\sigma }}\left[(X^+)^2{\displaystyle \frac{dX^+}{d\sigma }}\right]=0`$ (62b)
$`\left[{\displaystyle \frac{d^2(X^+)^2}{d\sigma ^2}}\left({\displaystyle \frac{d}{d\sigma }}X^+\right)^2\right]X=0`$ (62c)
$`X^+\left[X^+{\displaystyle \frac{d^2X^{}}{d\sigma ^2}}{\displaystyle \frac{1}{2}}{\displaystyle \frac{d^2(X^2)}{d\sigma ^2}}+{\displaystyle \underset{a}{}}(\dot{Y}^a)^2\right]+{\displaystyle \frac{d}{d\sigma }}\left[X^2{\displaystyle \frac{dX^+}{d\sigma }}+{\displaystyle \frac{\nu ^2}{6}}\left({\displaystyle \frac{dX^+}{d\sigma }}\right)^3\right]=0`$ (62d)
Solving the second and third equations we see that $`X^+=const`$. By gauge invariance we can make $`X`$ constant (for a nonzero $`X^+`$). Then the equations of motion and periodicity constraint for $`X^{}`$ and $`Y`$ to be constants. The large gauge transformation (60) is still not fixed, so the moduli space is:
$$=\{X^+,X^{},X,Y^3,\mathrm{}Y^9\}/\left\{\begin{array}{c}XX+\nu X^+\hfill \\ X^{}X^{}+\nu X+\frac{\nu ^2}{2}X^+\hfill \end{array}\right\}$$
(63)
The above solution is the Higgs branch and as expected it is the null-orbifold.
The analog of fractional branes (Coulomb branch) is more difficult to understand. We set $`X^+=0`$ and combine the small and large gauge transformations in
$$X^{}(\sigma )X^{}(\sigma )+\nu \stackrel{~}{\mathrm{\Lambda }}(\sigma )X(\sigma ),_0^{2\pi }\frac{d\sigma }{2\pi }\stackrel{~}{\mathrm{\Lambda }}$$
(64)
Where $`\stackrel{~}{\mathrm{\Lambda }}`$ is a periodic function.
For $`X(\sigma )0`$ we choose a gauge fixed solution by setting $`X^{}=constant`$ identified by the shifts:
$$X^{}=X_0^{},X^{}X^{}+\nu (\frac{1}{2\pi }\frac{d\sigma }{X(\sigma )})^1.$$
(65)
In addition to this $`X^{}`$, the Coulomb branch is parameterized by $`X(\sigma )`$ and $`Y(\sigma )`$ arbitrary functions of $`\sigma `$. Therefore the solutions allow a separation of the fractional branes both in the singular plane and in the transverse directions. Note, however, that the gauge symmetry is completely broken for $`X0`$. This is different from the fractional branes of the $`\times /_N`$ orbifold, which have a $`U(1)`$ gauge symmetry for each fractional brane. We identify this peculiarity as arising from the existence of a 2-dim space of orbifold fixed points.
Pursuing the analogy of T duality, we expect that there is a field redefinition which brings the variables to the form of some gauge fields. Indeed for $`X^+0`$ we define the fields
$$A\frac{X}{X^+}Y^+=X^+Y=X^{}\frac{1}{2}\frac{X}{X^+}$$
(66)
$`Y^+`$,$`Y`$ are gauge invariant and $`AA+\nu _\sigma \mathrm{\Lambda }`$. Of course, in the action expressed using these variables the gauge field is decoupled from the gauge invariant operators. In addition, it confirms the geometric picture we have. We know of space-like cycles the null orbifold possesses. The gauge field is exactly the coordinate which parameterizes the space-like cycles. On the singular plane, we encounter the same phenomenon, the field
$$A^{\text{sing}}=\frac{X^{}}{X}$$
(67)
is again the natural gauge field in agreement with the existence of null cycles.
## 4 Twisted Closed String Condensation in the Null-Orbifold
This section contains the main result of the paper. We show that after condensing a closed twisted sector state, the D(-1) action flows to that of a $`_N`$ orbifold (at the limit of infinite boost which we elaborate below).
### 4.1 The action after closed string condensation
We denote by $`T_k`$ the modulus squared of the k twisted sector closed string state, $`k0`$. This choice follows . We repackage the twisted condensate field using the $`\sigma `$ variable as
$$T(\sigma )=\underset{k0}{}T_ke^{ik\sigma }$$
(68)
The quantum $`U(1)`$ symmetry of the $``$ orbifold is given by shifts in $`\sigma `$, and $`T(\sigma )`$ breaks it with the appropriate charges. In order to flow to a $`_N`$ orbifold, we turn on only $`T_k`$ for $`k=N`$, i.e, $`T(\sigma )=T(\sigma +2\pi /N)`$.
We also define the function $`U(\sigma )`$ to be the double integral of $`T`$ which is periodic and integrates to zero
$$U(\sigma )=^\sigma T(\sigma )=\underset{k0}{}\frac{T_k}{k^2}e^{ik\sigma }$$
(69)
The couplings of the twisted condensate to open string fields are determined uniquely from three properties. The first is locality in the $`\sigma `$ variable. The second is that $`T`$ couples to a quadratic form of the $`X`$’s. By gauge invariance the possibilities are $`X^+(\sigma )X^+(\sigma )`$ and $`X^22X^+X^{}`$. The third requirement is that since the background is invariant under translations in $`X^{}`$, we can choose the twisted condensate field to have a similar property. This determines the coupling to be
$$𝑑\sigma T(\sigma )X^+(\sigma )X^+(\sigma )$$
(70)
This is also the result obtained from taking the limit of the $`_N`$ orbifold. As in , it means that the twisted sector closed string state couples to the open string field which measures the effective distance from the singularity.
It is convenient to change variables to the following gauge invariant ones by
$$(X^+,X^{},X)(X^+,L),L(\sigma )=2X^+(\sigma )X^{}(\sigma )X^2(\sigma )$$
(71)
Using these variables the action becomes
$$\mathrm{S}=_0^{2\pi }\left[(X^+)^2(\dot{Y}^a)^2\frac{1}{3}\frac{d(X^+)^3}{d\sigma }\left(\frac{d}{d\sigma }\left(\frac{L}{X^+}\right)+\frac{2}{X^+}\dot{U}(\sigma )\right)\frac{\nu ^2}{12}\left(\dot{X}^+\right)^4\right]$$
(72)
The equations of motions, in terms of the gauge invariant variables are:
$`{\displaystyle \frac{d}{d\sigma }}\left[(X^+)^2{\displaystyle \frac{dY^a}{d\sigma }}\right]=0`$ (73a)
$`{\displaystyle \frac{d}{d\sigma }}\left(X^+{\displaystyle \frac{dX^+}{d\sigma }}\right)+\left({\displaystyle \frac{dX^+}{d\sigma }}\right)^2=0`$ (73b)
$`X^+\left[{\displaystyle \underset{a}{}}\left({\displaystyle \frac{dY^a}{d\sigma }}\right)^2+{\displaystyle \frac{1}{2}}{\displaystyle \frac{d^2L}{d\sigma ^2}}+T(\sigma )\right]{\displaystyle \frac{d}{d\sigma }}\left[\left(L{\displaystyle \frac{dX^+}{d\sigma }}\right){\displaystyle \frac{\nu ^2}{6}}\left({\displaystyle \frac{dX^+}{d\sigma }}\right)^3\right]=0`$ (73c)
Using the fact that $`X^+`$ and $`Y^a`$ have to be constant by the first two equations (and periodicity conditions), we can reduce the last equation to:
$$X^+\left[\frac{1}{2}\frac{d^2L}{d\sigma ^2}+T(\sigma )\right]=0$$
(74)
and the general solution is:
$$X^+(\sigma )=X_0^+,L(\sigma )=L_02U(\sigma )$$
(75)
In the next subsection we will deal carefully with the IR action around this solution, but we can already see the glimpses of $`_N`$ orbifold. Given a value of the gauge invariant $`X_0^+`$ and $`L_0`$ we have an entire gauge orbit of the gauge symmetry (58). For generic values of $`X_0^+`$ and $`L_0`$, the gauge group is completely broken on this orbit. The maximal unbroken gauge group occur on the orbit
$$X_0^+=0,L_0=\overline{L}_02\underset{\sigma [0,2\pi ]}{\mathrm{min}}U(\sigma )$$
(76)
where it is $`\mathrm{U}(1)^{N1}`$. This point corresponds to bringing as many D(-1) instantons as possible close to the singularity, and the symmetry suggests $`D(1)`$ instantons near a $`_N`$ singularity.
To show this we note first that for $`X^+0`$ the gauge symmetry is completely broken due to the transformation $`XX+\nu \mathrm{\Lambda }X^+`$. For $`X^+=0`$, $`X`$ is gauge invariant and the transformation of $`X^{}`$ (64):
$$X^{}X^{}+\nu \stackrel{~}{\mathrm{\Lambda }}X=X^{}+\nu \stackrel{~}{\mathrm{\Lambda }}\sqrt{L}$$
Note that $`L0`$ everywhere for $`X^+=0`$. We see that if $`L_0`$ attains its value from 76, there are $`N`$ points<sup>10</sup><sup>10</sup>10$`N`$ is the number of points where $`^\sigma T(\sigma ^{})`$ reaches it’s maximal value. where $`L=0`$, and hence $`X=0`$, and the symmetry is restored. We identify the gauge symmetry as
$$\mathrm{\Lambda }\delta (\sigma \sigma _i),\text{for each }\sigma _i\text{ such that }L(\sigma _i)=X(\sigma _i)=0$$
The constraint $`\stackrel{~}{\mathrm{\Lambda }}2\pi `$ removes one transformation to obtain the gauge symmetry $`U(1)^{N1}`$ of D(-1) instantons in $`\times /_N`$ like orbifold.
Although this is an indication, it can not be taken to be the complete story. This is so because the localized transformations as we have written them are difficult to extend to the entire Higgs branch (where the gauge symmetry is generically broken). In this case we have to smear the localized gauge transformations (to avoid $`\delta ^2(\sigma \sigma _0)`$ terms), and there seems to be considerable arbitrariness in how to do so.
### 4.2 The ”IR” theory
Since the action is not positive definite, we first need to clarify the notion of “IR” physics. By this we mean that we separate the degrees of freedom in the action into slow and fast variables, where for the latter we carry out a stationary phase approximation. To obtain a clear separation of scales we write the twisted condensate field as
$$T(\sigma )=M\widehat{T}(\sigma )$$
and work in the scaling of $`\widehat{T}`$ fixed and $`M\mathrm{}`$.
We are interested in working slightly off-shell - i.e, we do not impose the equations of motion, but restrict our attention to fluctuations which have finite action (does not scale with M). We also would like to work near the singularity. To do so, we need to work close to the special solution (76) for which the $`U(1)^{N1}`$ is unbroken. This means that we should take
$$L(\sigma )+\delta L(\sigma )=\left(\overline{L}_02U(\sigma )\right)+\delta L(\sigma )$$
(77)
For this class $`L`$ scales linearly in $`M`$ in all of the interval (due to the term in the parenthesis) except in the $`N`$ points $`\sigma _i`$, where its value is held fixed as $`M\mathrm{}`$ (and determined by $`\delta L`$). Expanding the action to quadratic order in $`\delta L`$, $`\delta X^+`$ and $`\delta Y`$ around this solution, we obtain
$$\mathrm{S}_{fluc}=\frac{d\sigma }{2\pi }\left[(X_0^+)^2(\dot{\delta Y^a})^2(X_0^+)\frac{d(\delta X^+)}{d\sigma }\frac{d(\delta L)}{d\sigma }+L\left(\frac{d(\delta X^+)}{d\sigma }\right)^2\right]$$
(78)
In the intervals between the points $`\sigma _i`$, the value of $`L`$ scales with $`M`$, and the variations $`\delta X^+`$ are fast in these intervals. We therefore “integrate out” the variations of $`X^+`$ between the $`\sigma _i`$’s. These give the constrains that the function $`X^+(\sigma )`$ is piecewise constant, and may jump at the points $`\sigma _i`$. The function is therefore:
$$X^+(\sigma )=X_0^++\delta X_i^+,\sigma _i<\sigma <\sigma _{i+1}$$
(79)
However (78) is not convenient at the points $`\sigma _i`$ since the last term is $`0\times \delta ^2(\sigma \sigma _i)`$, rather we go back to (72). Using the derivative of $`X^+`$ we see that the action localizes at the points $`\sigma _i`$.
Next we evaluate the term $`_\sigma (L/X^+)2\dot{U}/X^+`$ at the points $`\sigma _i`$. The first step is to evaluate the behavior of $`L`$ near the points $`\sigma _i`$. In order to obtain a finite action we need to impose that $`L/X^+`$ be continuous at the points $`\sigma _i`$. Otherwise we will have a $`\delta ^2(\sigma \sigma _i)`$ divergence. We will use the notation
$$L_i^\pm =\underset{\sigma \sigma _i\pm }{lim}L(\sigma )$$
(80)
and hence
$$\frac{L_i^+}{X_i^+}=\frac{L_i^{}}{X_{i1}^+}$$
(81)
The degrees of freedom of $`L`$ in the intervals between the points $`\sigma _i`$ do not appear in the action and can be integrated out. To obtain some physical intuition for the remaining degrees of freedom in $`L`$, we impose on L the equations of motion piecewise in these intervals. The solution we obtain in each interval is
$$L(\sigma )=(L_i^++2U(\sigma _i))+\frac{L_{i+1}^{}L_i^+}{\sigma _{i+1}\sigma _i}(\sigma \sigma _i)2U(\sigma ),\sigma (\sigma _i,\sigma _{i+1})$$
(82)
Note that the value of $`U`$ is the same in all the points $`\sigma _i`$.
Applying the expansion (82) to the action (72), (we rename the fluctuation fields by omitting the $`\delta `$’s in front to unclutter the equations) we find the leading IR action <sup>11</sup><sup>11</sup>11In the calculation of (83) we used a ”non-symmetrized” value to the integration of a delta function over a discontinues functions
$$𝑑\sigma \delta (\sigma \sigma _i)F(\sigma )\underset{ϵ0}{lim}F(\sigma _iϵ)$$
Any other definition of the integral will result in an equivalent action up to a linear combination of the $`L_i^\pm `$’s. :
$$\begin{array}{c}\mathrm{S}_{IR}=(X_0^+)^2_0^{2\pi }\frac{d\sigma }{2\pi }(\dot{Y}^a)^2(X_0^+)^2\underset{\sigma _i}{}\left(X_{i+1}^+X_i^+\right)\frac{1}{(\sigma _{i+1}\sigma _i)}\left(\frac{L_{i+1}^+}{X_{i+1}^+}\frac{L_i^+}{X_i^+}\right)\hfill \end{array}$$
(83)
In the next subsection we will match the $`X^+L`$ action to that of the boosted $`_N`$ singularity.
The situation of the $`Y`$’s is less clear. To the order that we are working in the open string fields and in the twisted condensate field we do not get a clear separation into $`N`$ degrees of freedom below a gap. Perhaps this happens at higher orders.
#### 4.2.1 The boosted $`\times /_N`$
We would like to clarify the term ”infinitely boosted $`_N`$” used in the beginning of the section. We start from the action of a fractional $`D`$-brane in the $`\times /_N`$ orbifold as computed in , apply infinite boost similar to the one used in section 2 but keeping the parameter N constant<sup>12</sup><sup>12</sup>12In other words we redefine coordinate but do not change the identification of the orbifold. The action for a $`D(1)`$ instanton reads<sup>13</sup><sup>13</sup>13We are using a slightly different normalization for the $`Z`$’s then .:
$$\begin{array}{c}\mathrm{S}^_N=\underset{j=0}{\overset{N1}{}}\left(X_{j+1,j+1}^0X_{j,j}^0\right)^2\left|Z_{j,j+1}\right|^2+\underset{a=3}{\overset{9}{}}\underset{j=0}{\overset{N1}{}}\left(Y_{j+1,j+1}^aY_{j,j}^a\right)^2\left|Z_{j,j+1}\right|^2+\hfill \\ \hfill +\frac{1}{4}\underset{j=0}{\overset{N1}{}}\left(\left|Z_{j,j+1}\right|^2\left|Z_{j1,j}\right|^2\right)^2\end{array}$$
(84)
Expanding the action around a general classical solution in the Higgs branch:
$`X_{j,j}^0=X^0+\delta X_{j,j}^0`$ (85)
$`Y_{j,j}^a=Y^a+\delta Y_{j,j}^a`$ (86)
$`\left|Z_{j,j+1}\right|=\left|Z\right|+\delta \left|Z_{j,j+1}\right|`$ (87)
and dropping the $`\delta `$’s we find:
$$\begin{array}{c}\mathrm{S}_{fluc}^_N=\left|Z\right|^2\underset{j=0}{\overset{N1}{}}\left[\left(X_{j+1,j+1}^0X_{j,j}^0\right)^2\left(\left|Z_{j,j+1}\right|\left|Z_{j1,j}\right|\right)^2\right]+\hfill \\ \hfill +\left|Z\right|^2\underset{a=3}{\overset{9}{}}\underset{j=0}{\overset{N1}{}}\left(Y_{j+1,j+1}^aY_{j,j}^a\right)^2\end{array}$$
(88)
We boost and rotate the coordinates by:
$$\frac{1}{\sqrt{2}}\left(\begin{array}{c}X^++X^{}\\ X^{}X^+\\ \sqrt{2}X\end{array}\right)=\left(\begin{array}{ccc}\frac{\alpha }{\sqrt{\beta }}& \sqrt{\beta }\alpha & 0\\ \sqrt{\beta }\alpha & \frac{\alpha }{\sqrt{\beta }}& 0\\ 0& 0& 1\end{array}\right)\left(\begin{array}{c}X^0\\ \mathrm{Re}Z\\ \mathrm{Im}Z\end{array}\right),\alpha =\frac{\sqrt{\beta }}{\sqrt{1\beta ^2}}$$
(89)
Taking the limit<sup>14</sup><sup>14</sup>14The explicit parametrization of the transformation was chosen to produce a finite limit. $`\alpha \mathrm{}`$ we find:
$$X^0=\sqrt{2}\alpha X^++\frac{X^{}}{2\sqrt{2}\alpha }+O(\alpha ^2)$$
(90)
$$\left|Z\right|^2=2\alpha ^2(X^+)^2+\left((X)^2X^+X^{}\right)+O(\alpha ^1)$$
(91)
Plugging the transformation into (88) and taking the leading order in $`\alpha `$ (using only the first index of each field).
$$\begin{array}{c}\mathrm{S}_{fluc}^{\text{boosted-}_N}=2\alpha ^2(X_0^+)^2\underset{j=0}{\overset{N1}{}}\left(X_{j+1}^+X_j^+\right)\left(\frac{L_{j+1}}{X_{j+1}^+}\frac{L_j}{X_j^+}\right)+\hfill \\ \hfill +2\alpha ^2\left|X_0^+\right|^2\underset{a=3}{\overset{9}{}}\underset{j=0}{\overset{N1}{}}\left(Y_{j+1,j+1}^aY_{j,j}^a\right)^2+O(\alpha ^0)\end{array}$$
(92)
The factor $`\alpha `$ can be absorbed into a rescaling of the coordinates.
#### 4.2.2 Relating the IR physics
We found the action of the IR physics around a Higgs branch VEV both in the boosted-$`_N`$ (92) and the null-orbifold (83). Focusing on the orbifold directions:
$`\mathrm{S}_{IR}^{\text{boosted-}_N}`$ $`(X_0^+)^2{\displaystyle \underset{j=0}{\overset{N1}{}}}\left(X_{i+1}^+X_i^+\right)\left({\displaystyle \frac{L_{i+1}}{X_{i+1}^+}}{\displaystyle \frac{L_i}{X_i^+}}\right)`$
$`\mathrm{S}_{IR}`$ $`(X_0^+)^2{\displaystyle \underset{\sigma _i}{}}{\displaystyle \frac{1}{(\sigma _{i+1}\sigma _i)}}\left(X_{i+1}^+X_i^+\right)\left({\displaystyle \frac{L_{i+1}^+}{X_{i+1}^+}}{\displaystyle \frac{L_i^+}{X_i^+}}\right)`$
The above actions are exactly the same up to rescaling of fields. The action is ill-defined at the singularity, (i.e the VEV of the $`X^+`$ field vanish) which indicates the emergence of new degrees of freedom, which are the fractional branes (Coulomb branch) together with the expected $`\mathrm{U}(1)^{N1}`$ gauge symmetry.
Note, however, that after the rescaling of the fields, the issue of whether one is dealing in a finite boost or an infinite boost is subleading in the boost parameter. In order to see this effect in the flow from the null orbifold to the $`\times /_N`$ case, we need to look at subleading corrections to the action - in this case subleading in $`M`$. We do not know how to do this precisely, and hence can not answer the detailed question of whether the deformed null singularity is the infinitely boosted $`\times /_N`$ or boosted by some parameter proportional to a power of $`M`$.
Note also that we cannot make a clear study of the IR physics in the Coulomb branch or the $`Y`$ variables with our probes as already pointed out.
## 5 Summary and Conclusions
The purpose of this paper is to study the geometry of small deformations of the null orbifold, which occur after condensation of twisted sector states. This might be a first step towards understanding both the situation in which twisted sector states are condensed with large VEV’s, or the situation of pair creation of twisted sector state, which might be ways in in which the singularity might be tamed.
We focused on the case of the null-orbifold since it has a clear relationship to the $`\times /_N`$ orbifold - the latter is the large N limit of the former. We have exhibited, using D(-1) brane probes, some evidence that indeed there is a transition from the null orbifold to the $`\times /_N`$ orbifold.
In the case of flows between $`\times /_N`$ theories, one can eventually flow to flat space, smoothing out the singularity completely. We expect that a similar situation occurs here - hence we conclude that the null orbifold may be smoothed out by the condensation of twisted sector states.
The intermediate step - of flowing to a $`\times /_N`$ \- might be interesting by itself. For example, it might be interesting to study the condensation of twisted sector states of the BTZ black hole and examine their relation of the deformed geometry to the microstates of .
###### Acknowledgments.
The authors are happy to thank O.Aharony, D.Kutasov, B.Pioline, S.Ross, M.Rozali, J.Simon and S.Shenker for useful discussions. The work is supported in part by the Israel Science Foundation, by the Braun-Roger-Siegl foundation, by the European network HPRN-CT-2000-00122, by a grant from the G.I.F. (the German-Israeli Foundation for Scientific Research and Development), by the Minerva Foundation, by the Einstein Center for Theoretical Physics and the by Blumenstein foundation.
## Appendix A Calculations of the Wave-Functions Limit
In this appendix we demonstrate (in the untwisted case) the limiting procedure between wave-functions on the $`\times /_N`$ orbifold and wave-functions on the null-orbifold. We start with the wave-functions of the $`\times /_N`$ orbifold in the static coordinates (14)
$$\mathrm{\Psi }_{k,l}^N=Ne^{ik^0x^0+iN\varphi l}J_{Nl}(2u),ue^{i\varphi }i\overline{k}z$$
(93)
The limit procedure is defined in (19) and (20), the boost matrix is
$$M_N=\left(\begin{array}{ccc}\frac{a}{\sqrt{a^2b^2}}& 0& \frac{b}{\sqrt{a^2b^2}}\\ \frac{b}{\sqrt{a^2b^2}}& 0& \frac{a}{\sqrt{a^2b^2}}\\ 0& 1& 0\end{array}\right),\text{with}a,b=\frac{\sqrt{1\pm 1/N^2}}{\sqrt{2}}$$
(94)
We take care of each term of (93) separately. First the phase factors:
$`e^{ik^0x^0}\stackrel{\text{boost}}{}`$ $`e^{iN^2(ak_0+bk_1)(ax^0bx^1)}=`$ (95)
$`=\mathrm{exp}\left[{\displaystyle \frac{i}{2}}N^2(k_0+k_1)(x^0x^1)+{\displaystyle \frac{i}{2}}k_0x^0+{\displaystyle \frac{i}{2}}k_1x^1+O(1/N)\right]`$ (96)
$`e^{iN\varphi l}=`$ $`\left({\displaystyle \frac{k\overline{z}}{\overline{k}z}}\right)^{\frac{Nl}{2}}\stackrel{\text{boost}}{}\left[{\displaystyle \frac{\left(\frac{i(ak_1+bk_0)}{\sqrt{a^2b^2}}k_2\right)\left(\frac{i(bx^0ax^1)}{\sqrt{a^2b^2}}x^2\right)}{c.c}}\right]^{\frac{Nl}{2}}=`$ (97)
$`=(1)^{\frac{Nl}{2}}\left[{\displaystyle \frac{(k_0+k_1)(x^0x^1)i\sqrt{2}\frac{x^2(k_0+k_1)+k_2(x^0x^1)}{N}+O(1/N^2)}{c.c}}\right]^{\frac{Nl}{2}}`$ (98)
The large N limit is easily calculated via $`(1+x/N)^Ne^x`$ and reduces to
$$e^{iN\varphi l}\stackrel{\text{boost}}{}(1)^{\frac{Nl}{2}}\mathrm{exp}\left[il\sqrt{2}\frac{x^2(k_0+k_1)+k_2(x^0x^1)}{(x^0x^1)(k_1+k_0)}+O(1/N)\right].$$
(100)
The remaining part of the story is the Bessel function. We quote the asymptotic expansion of the Bessel function from which shall play a major role
$`J_\nu (z)=`$ $`\sqrt{{\displaystyle \frac{2}{\pi z}}}cos(z{\displaystyle \frac{\pi }{2}}\nu {\displaystyle \frac{\pi }{4}})\left[{\displaystyle \underset{k=0}{\overset{n1}{}}}{\displaystyle \frac{(1)^k}{(2z)^{2k}}}{\displaystyle \frac{\mathrm{\Gamma }(\nu +2k+\frac{1}{2})}{(2k)!\mathrm{\Gamma }(\nu 2k+\frac{1}{2})}}+R_1\right]`$ (101)
$`\sqrt{{\displaystyle \frac{2}{\pi z}}}sin(z{\displaystyle \frac{\pi }{2}}\nu {\displaystyle \frac{\pi }{4}})\left[{\displaystyle \underset{k=0}{\overset{n1}{}}}{\displaystyle \frac{(1)^k}{(2z)^{2k+1}}}{\displaystyle \frac{\mathrm{\Gamma }(\nu +2k+\frac{3}{2})}{(2k+1)!\mathrm{\Gamma }(\nu 2k\frac{1}{2})}}+R_2\right]`$ (102)
Where for $`n>\nu /21/4`$ the remainders satisfy
$$\left|R_1\right|<\left|\frac{\mathrm{\Gamma }(\nu +2n+1/2)}{(2z)^{2n}(2n)!\mathrm{\Gamma }(\nu 2n+1/2)}\right|$$
(104)
$$\left|R_2\right|<\left|\frac{\mathrm{\Gamma }(\nu +2n+3/2)}{(2z)^{2n+1}(2n+1)!\mathrm{\Gamma }(\nu 2n1/2)}\right|$$
(105)
The Bessel functions in (93) transforms under the boost:
$`J_{Nl}`$ $`\left[\sqrt{\left((k^1)^2+(k^2)^2\right)\left((x^1)^2+(x^2)^2\right)}\right]\stackrel{\text{boost}}{}`$ (106)
$`=`$ $`J_{Nl}\left[N^2\left|x^+k^+\right|+{\displaystyle \frac{(k^2x^+)^2k^+k^{}(x^+)^2(k^+)^2x^+x^{}+(k^+x^2)^2}{2\left|x^+k^+\right|}}+O(1/N)\right]`$ (107)
Using the useful limit
$$\underset{z\mathrm{}}{lim}\frac{\mathrm{\Gamma }(z+a)}{\mathrm{\Gamma }(z)}z^a=1$$
We observe that the terms in the Bessel function expansion don’t scale as powers of N but only of n and hence no terms can be dropped. Fortunately, we are able to re-sum the series (101) in the $`N\mathrm{}`$ limit. As emphasized in (101) we should carefully estimate the remainder. Let $`\nu =Nl`$ then we take $`n`$ such that (at least) $`\nu =2n`$ and also $`z=\frac{\alpha }{2}n^2`$ for some constant $`\alpha `$. Estimating the remainder we obtain
$$\left|\frac{\mathrm{\Gamma }(\nu +2n+1/2)}{(2z)^{2n}(2n)!\mathrm{\Gamma }(\nu 2n+1/2)}\right|=\left|\frac{\mathrm{\Gamma }(4n+1/2)}{(\alpha n^2)^{2n}(2n)!\mathrm{\Gamma }(1/2)}\right|<\frac{1}{\alpha ^{2n}}\frac{(4n)^{2n}}{(n^2)^{2n}}=(\frac{4}{\alpha })^{2n}(\frac{1}{n})^{2n}$$
These are preferable circumstances and it holds that
$$\underset{k\mathrm{}}{lim}\underset{l=0}{\overset{k}{}}P(l,k)=\underset{l=0}{\overset{\mathrm{}}{}}P(l,\mathrm{})$$
Which means that
$$\begin{array}{c}\underset{N\mathrm{}}{lim}J_{Nl}(N^2\left|x^+k^+\right|+A)=\underset{N\mathrm{}}{lim}\sqrt{\frac{2}{\pi N^2\left|x^+k^+\right|}}\mathrm{cos}\left(N^2\left|x^+k^+\right|+A\frac{\pi }{2}Nl\frac{\pi }{4}\right)\hfill \\ \hfill \left[\underset{k=0}{\overset{Nl/21}{}}\frac{(1)^k}{(2N^2x^+k^+)^{2k}}\frac{\mathrm{\Gamma }(Nl+2k+\frac{1}{2})}{(2k)!\mathrm{\Gamma }(Nl2k+\frac{1}{2})}+R_1\right]2^{nd}\text{ term}\end{array}$$
(109)
where
$$A\frac{(k^2x^+)^2k^+k^{}(x^+)^2(k^+)^2x^+x^{}+(k^+x^2)^2}{2\left|x^+k^+\right|}+O(1/N)$$
Taking the limit inside the square brackets gives
$`=`$ $`\left[\underset{N\mathrm{}}{lim}\sqrt{{\displaystyle \frac{2}{\pi N^2\left|x^+k^+\right|}}}\mathrm{cos}\left(N^2\left|x^+k^+\right|+A{\displaystyle \frac{\pi }{2}}Nl{\displaystyle \frac{\pi }{4}}\right)\right]\left[{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^kl^{4k}}{(2k)!(2x^+k^+)^{2k}}}\right]2^{nd}\text{ term}`$ (110)
$`=`$ $`\left[\underset{N\mathrm{}}{lim}\sqrt{{\displaystyle \frac{2}{\pi N^2x^+k^+}}}\mathrm{cos}(N^2x^+k^++A{\displaystyle \frac{\pi }{2}}Nl{\displaystyle \frac{\pi }{4}})\right]\mathrm{cos}\left({\displaystyle \frac{l^2}{2\left|x^+k^+\right|}}\right)`$ (111)
$`\left[\underset{N\mathrm{}}{lim}\sqrt{{\displaystyle \frac{2}{\pi N^2\left|x^+k^+\right|}}}\mathrm{sin}\left(N^2\left|x^+k^+\right|+A{\displaystyle \frac{\pi }{2}}Nl{\displaystyle \frac{\pi }{4}}\right)\right]\mathrm{sin}\left({\displaystyle \frac{l^2}{2\left|x^+k^+\right|}}\right)`$ (112)
$`=`$ $`\underset{N\mathrm{}}{lim}\sqrt{{\displaystyle \frac{2}{\pi N^2|x^+k^+|}}}\mathrm{cos}\left(N^2\left|x^+k^+\right|+A{\displaystyle \frac{\pi }{2}}Nl{\displaystyle \frac{\pi }{4}}+{\displaystyle \frac{l^2}{2\left|x^+k^+\right|}}\right)`$ (113)
Combining the pieces of the full wave function together
$`\mathrm{\Psi }_{k;l}^{\text{boosted}}`$ $`(x^+,x^{},x^2)=\sqrt{{\displaystyle \frac{2}{\pi \left|x^+k^+\right|}}}e^{\pi i\frac{Nl}{2}}e^{il\frac{x^2k^+x^+k_2}{x^+k^+}}e^{\left[iN^2k^+x^+\frac{i}{2}\left(k^+x^{}+k^{}x^+\right)+O(1/N)\right]}`$ (114)
$`\mathrm{cos}\left[N^2\right|x^+k^+|{\displaystyle \frac{\pi }{2}}Nl{\displaystyle \frac{\pi }{4}}+{\displaystyle \frac{l^2}{2\left|x^+k^+\right|}}+`$ (115)
$`+{\displaystyle \frac{(k^2x^+)^2k^+k^{}(x^+)^2(k^+)^2x^+x^{}+(k^+x^2)^2}{2\left|x^+k^+\right|}}+O(1/N)]=`$ (116)
$`=`$ $`{\displaystyle \frac{(\text{sign}(k^+x^+))^{Nl}}{\sqrt{2\pi ix^+k^+}}}[e^{\left[i\frac{\left(lx^2k^++x^+k^2\right)^2}{2x^+k^+}+ik^2x^2ik^+x^{}ik^{}x^++O(N^1)\right]}+`$ (117)
$`+()^{Nl+1}e^{\left[i\frac{\left(l+x^2k^+x^+k^2\right)^2}{2x^+k^+}ik^2x^2i2N^2k^+x^++O(N^1)\right]}]`$ (118)
This is the expression (up to trivial algebra) quoted in the text (21). The First exponential is the wave function of the null-orbifold and the second produces an expression which is interpreted as a vanishing distribution. In order to prove the last statement we study the integration of the second exponential in the boosted wave-function with a test function $`g(x^+,x)`$ at the large $`N`$ limit:
$$𝑑x^+𝑑x^2g(x^+,x^2)\frac{1}{\sqrt{2\pi ix^+k^+}}\mathrm{exp}\left[i\frac{\left(l+x^2k^+\right)^2}{2x^+k^+}i2N^2k^+x^+\right]$$
Without loss of generality we substituted $`k^2=0`$. The $`x^+`$ integration can now be evaluated using a saddle point method:
$`{\displaystyle 𝑑x^+𝑑x^2}`$ $`g(x^+,x^2){\displaystyle \frac{1}{\sqrt{2\pi ix^+k^+}}}\mathrm{exp}\left[i{\displaystyle \frac{\left(l+x^2k^+\right)^2}{2x^+k^+}}i2N^2k^+x^+\right]=`$ (119)
$``$ $`{\displaystyle \underset{\pm }{}}{\displaystyle 𝑑x^2\frac{1}{Nk^+}g(\pm \frac{(l+x^2k^+)}{2Nk^+},x^2)e^{2iN(l+x^2k^+)}}`$ (120)
$``$ $`{\displaystyle \underset{\pm }{}}{\displaystyle \frac{\stackrel{~}{G}_\pm ^N(\pm 2Nk^+)}{N}}\stackrel{N\mathrm{}}{}0`$ (121)
Where $`\stackrel{~}{G}_\pm ^N`$ is the Fourier transform (with respect to $`x^2`$) of:
$$G_\pm ^Ng(\pm \frac{(l+x^2k^+)}{2Nk^+},x^2)$$
For a large class of functions <sup>15</sup><sup>15</sup>15We didn’t carry a full classification of the functions $`g(x^+,x^2)`$ that have the above property. A large enough set of examples are polynomials in $`x^+,x^2`$ multiplied by decaying exponentials and Gaussian $`g(x^+,x^2)`$ (which are in particular $`𝕃_2`$) the limit (119) is indeed zero which proves our claim.
## Appendix B First Quantization of the String on $`^{1,2}/\mathrm{\Gamma }`$
The orbifolds in mind are defined by a flat space CFT with a quotient by a twist:
$$ds^2=d(x^0)^2+d(x^1)^2+d(x^2)^2+dx_{}^2$$
(122)
$$Xe^{2\pi \widehat{J}}X=\left[e^{2\pi 𝒥}\right]_{}^{\mu }{}_{\nu }{}^{}X^\nu ,𝒥\mathrm{SO}(1,2)$$
(123)
We discuss 3 cases according the 3 classes of SO(1,2):
* The Milne orbifold (J is hyperbolic): $`J_\mathrm{\Delta }=2\pi \mathrm{\Delta }\widehat{J}_{02}`$
* The Null-orbifold (J is parabolic): $`J_v=\frac{2\pi v}{\sqrt{2}}\left(J_{02}+J_{12}\right)`$
* The $`_N`$ orbifold (J is elliptic): $`J_N=\frac{2\pi }{N}J_{12}`$
The operators above can be defined using the matrices in $`(x^0,x^1,x^2)`$ basis:
$$𝒥_\mathrm{\Delta }=\mathrm{\Delta }\left(\begin{array}{ccc}0& 0& 1\\ 0& 0& 0\\ 1& 0& 0\end{array}\right)𝒥_v=\frac{v}{\sqrt{2}}\left(\begin{array}{ccc}0& 0& 1\\ 0& 0& 1\\ 1& 1& 0\end{array}\right)𝒥_N=\frac{1}{N}\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& 1\\ 0& 1& 0\end{array}\right)$$
The worldsheet action is:
$$S=\frac{1}{4\pi }d^2z\frac{2}{\alpha ^{}}\eta _{\mu \nu }\left(X^\mu \overline{}X^\nu +\eta _{\mu \nu }\psi ^\mu \overline{}\psi ^\nu +\eta _{\mu \nu }\stackrel{~}{\psi }^\mu \stackrel{~}{\psi }^\nu \right)$$
(124)
And the monodromies <sup>16</sup><sup>16</sup>16Remember $`z=e^{i(\sigma +\tau )},\overline{z}=e^{i(\sigma \tau )}`$.:
$$X(\sigma +2\pi ,\tau )=e^{2\pi w𝒥}X(\sigma ,\tau )$$
(125)
$$\psi \left(ze^{2\pi i}\right)=e^{2\pi i\left((\nu \frac{1}{2})11+iw𝒥\right)}\psi (z)$$
(126)
$$\stackrel{~}{\psi }\left(\overline{z}e^{2\pi i}\right)=e^{2\pi i\left((\stackrel{~}{\nu }\frac{1}{2})11iw𝒥\right)}\stackrel{~}{\psi }(\overline{z})$$
(127)
where $`\nu `$ and $`\stackrel{~}{\nu }`$ take the values of $`0`$ for R-sector and $`\frac{1}{2}`$ for NS-sector and $`w`$ is the twisted sector number.
The bosonic part mode expansion in analogy to ($`w0`$):
$$\begin{array}{c}\sqrt{\frac{2}{\alpha ^{}}}X^\mu (\sigma ,\tau )=\sqrt{\frac{2}{\alpha ^{}}}\left[e^{w𝒥\sigma }\right]_\rho ^\mu X_z^\rho (\alpha _0,\stackrel{~}{\alpha }_0;\tau )+\hfill \\ \hfill +i\underset{n0}{}\left[\frac{e^{i(n+iw𝒥)(\sigma +\tau )}}{n+iw𝒥}\right]_\rho ^\mu \alpha _n^\rho +i\underset{n0}{}\left[\frac{e^{i(niw𝒥)(\sigma \tau )}}{niw𝒥}\right]_\rho ^\mu \stackrel{~}{\alpha }_n^\rho \end{array}$$
(128)
The zero-mode part is
$`\text{Null},X_{z}^{}{}_{}{}^{\mu }(\tau )=\mathrm{cosh}\left(w\tau 𝒥\right)_{}^{\mu }{}_{\nu }{}^{}x^\mu +[11+{\displaystyle \frac{1}{2}}\mathrm{cosh}\left(w\tau 𝒥\right)]_{}^{\mu }{}_{\nu }{}^{}{\displaystyle \frac{2\alpha ^{}\tau }{3}}p^\nu `$ (129a)
$`\text{Other},X_{z}^{}{}_{}{}^{\mu }(\tau )=\mathrm{cosh}\left(w\tau 𝒥\right)_{}^{\mu }{}_{\nu }{}^{}x^\mu +\left[(w𝒥)^1\mathrm{sinh}(w𝒥\tau )\right]_\nu ^\mu \alpha ^{}p^\nu `$ (129b)
Using the Euclidean world-sheet, the left/right moving parts can be expanded:
$$X(z)=iz^{(1+iw𝒥)}(w𝒥x+\alpha ^{}p)i\sqrt{\frac{\alpha ^{}}{2}}\underset{n0}{}z^{(n+1+iw𝒥)}\alpha _n$$
(130a)
$$\overline{}X(\overline{z})=i\overline{z}^{(1+iw𝒥)}(w𝒥x\alpha ^{}p)+i\sqrt{\frac{\alpha ^{}}{2}}\underset{n0}{}\overline{z}^{(n+1iw𝒥)}\stackrel{~}{\alpha }_n$$
(130b)
$$\psi ^\mu (z)=\underset{r+\nu }{}\left[\frac{1}{z^{r+\frac{1}{2}+iw𝒥}}\right]_\rho ^\mu \psi _r^\rho $$
(130c)
$$\stackrel{~}{\psi }^\mu (\overline{z})=\underset{s+\stackrel{~}{\nu }}{}\left[\frac{1}{\overline{z}^{s+\frac{1}{2}iw𝒥}}\right]_{}^{\mu }{}_{\rho }{}^{}\stackrel{~}{\psi }_s^\rho $$
(130d)
The commutation relations are calculated form the above expansions using the OPE and canonical quantization relations (needed for the quasi-zero modes).
$`[\stackrel{~}{\alpha }_n^\mu ,\stackrel{~}{\alpha }_m^\nu ]=\delta _{n+m}(n\eta iw𝒥\eta )^{\mu \nu }`$ $`\{\psi _r^\mu ,\psi _r^{}^\nu \}=\delta _{r+r^{},0}\eta ^{\mu \nu }`$ (131)
$`[\alpha _n^\mu ,\alpha _m^\nu ]=\delta _{n+m}(n\eta +iw𝒥\eta )^{\mu \nu }`$ $`\{\stackrel{~}{\psi }_s^\mu ,\stackrel{~}{\psi }_s^{}^\nu \}=\delta _{s+s^{},0}\eta ^{\mu \nu }`$ (132)
Where we adopted the following definition of quasi-zero modes<sup>17</sup><sup>17</sup>17Note that the above definition do not consist of a full set of zero-mode operators. This is similar to flat space where $`\alpha _0^\mu =\stackrel{~}{\alpha }_0^\mu =\sqrt{\frac{\alpha ^{}}{2}}p^\mu `$.:
$$\alpha _0=\frac{\alpha ^{}p+w𝒥x}{\sqrt{2\alpha ^{}}}\stackrel{~}{\alpha }_0=\frac{\alpha ^{}pw𝒥x}{\sqrt{2\alpha ^{}}}[\alpha _0,\stackrel{~}{\alpha }_0]=0$$
(133)
The Virasoro generators (matter part) are
$`L_n^\psi ={\displaystyle \frac{1}{4}}{\displaystyle \underset{r+\nu }{}}\left[(2rn)\eta +iw\left(\eta 𝒥𝒥^T\eta \right)\right]_{\mu \nu }{\scriptscriptstyle \genfrac{}{}{0.0pt}{}{}{}}\psi _{nr}^\mu \psi _r^\nu {\scriptscriptstyle \genfrac{}{}{0.0pt}{}{}{}}+a^\psi (w,\nu )\delta _{n,0}`$ (134a)
$`L_n^x={\displaystyle \frac{1}{2}}{\displaystyle \underset{m}{}}{\scriptscriptstyle \genfrac{}{}{0.0pt}{}{}{}}\alpha _m\eta \alpha _{nm}{\scriptscriptstyle \genfrac{}{}{0.0pt}{}{}{}}+\delta _{n,0}a^x(w)`$ (134b)
Using $`[L_1,L_1]=2L_0`$ we find:
$`a^x={\displaystyle \frac{w^2}{4}}\mathrm{Tr}(𝒥^2)+vac\left|{\displaystyle \frac{1}{2}}\eta _{\mu \rho }\alpha _0^\mu \alpha _0^\rho \right|vac_w`$ (135a)
$`a_{NS}^\psi ={\displaystyle \frac{w^2}{4}}\mathrm{Tr}(𝒥^2)`$ (135b)
$`a_R^\psi ={\displaystyle \frac{D}{16}}{\displaystyle \frac{w^2}{4}}\mathrm{Tr}(𝒥^2)+vac\left|\left({\displaystyle \frac{iw}{2}}\eta 𝒥\right)_{\mu \sigma }\psi _0^\mu \psi _0^\sigma \right|vac_{w,R}`$ (135c)
In order not to break worldsheet supersymmetry the quantization scheme for the bosonic zero modes and the R-sector fermionic zero modes must obey
$$a_R^\psi +a^x=\frac{D}{16}$$
Applying this constraint we can fix the quantization scheme for the different classes of orbifolds (using the already discussed scheme of the bosonic part). The zero-point energies can then be calculated according to:
* Elliptic Orbifolds Class: Find d+1 vectors that diagonalize the matrix $`iw𝒥\eta `$. Vectors corresponding to positive eigenvalue $`\lambda `$ annihilate the vacuum. Vectors corresponding to zero eigenvalue are ’standard’ zero-modes which annihilate the vacuum (these are momentum operators). The bosonic zero point energy have a contribution from negative eigenvalue modes, $`\frac{1}{2}\lambda _i`$ for each negative eigenvalue (the R-sector will have the opposite contribution).
* Hyperbolic Orbifolds Class: Find d+1 vectors that diagonalize the matrix $`iw𝒥\eta `$, their eigenvalues are pure imaginary. Vectors corresponding to zero eigenvalue are ’standard’ zero-modes which annihilate the vacuum (actually these are momentum operators). Vector corresponding to (non-vanishing) pure imaginary zero mode should be treated as in to eliminate imaginary contribution to the zero point energies.
* Parabolic Orbifolds Class: Either by taking the limit from the hyperbolic or elliptic cases we set the zero-point energies to zero (the orbifold part).
The different zero-point energies are calculated for the representatives of the classes ($`_N`$ , Null and Milne orbifolds):
| | $`a^X`$ | $`a_{NS}^\psi `$ | $`a_R^\psi `$ | $`a_{NS}^g`$ | $`a_R^g`$ |
| --- | --- | --- | --- | --- | --- |
| Flat | $`0`$ | $`0`$ | $`D/16`$ | $`1/2`$ | $`5/8`$ |
| $`_n`$ | $`w/2N(1w/N)`$ | $`w^2/2N^2`$ | $`D/16w/2N(1w/N)`$ | $`1/2`$ | $`5/8`$ |
| Milne | $`\frac{w^2\mathrm{\Delta }^2}{2}`$ | $`\frac{w^2\mathrm{\Delta }^2}{2}`$ | $`D/16\frac{w^2\mathrm{\Delta }^2}{2}`$ | $`1/2`$ | $`5/8`$ |
| Null | $`0`$ | $`0`$ | $`D/16`$ | $`1/2`$ | $`5/8`$ |
### B.1 The Vacuum Structure
By a simple manipulation of commutation relations we can rewrite $`L_0`$ as:
$$L_0=L_0^{(\text{diag})}+\frac{w}{2}(\eta 𝒥)_{\mu \nu }\mathrm{\Sigma }^{\nu \mu }A$$
(136)
With the Lorentz generators defined as
$$\mathrm{\Sigma }^{\mu \nu }\frac{i}{2}\underset{r+\nu }{}[\psi _r^\mu ,\psi _r^\nu ],\frac{w}{2}(\eta 𝒥)_{\mu \nu }[\mathrm{\Sigma }^{\nu \mu },\psi _r^\sigma ]=iw𝒥_{}^{\sigma }{}_{\mu }{}^{}\psi _r^\mu $$
The diagonal part of $`L_0`$:
$$L_0^{(\text{diag})}=\frac{1}{2}\underset{m}{}{\scriptscriptstyle \genfrac{}{}{0.0pt}{}{}{}}\alpha _m^\mu \alpha _{m}^{}{}_{\mu }{}^{}{\scriptscriptstyle \genfrac{}{}{0.0pt}{}{}{}}+\frac{1}{2}\underset{r+\nu }{}r{\scriptscriptstyle \genfrac{}{}{0.0pt}{}{}{}}\psi _r^\mu \psi _{r}^{}{}_{\mu }{}^{}{\scriptscriptstyle \genfrac{}{}{0.0pt}{}{}{}}+L_0^{\text{ghost}}+a(w,\nu )$$
And the constant A defined as:
$$A(w,\nu )=vac\left|\left(\frac{iw}{2}\eta 𝒥\right)_{\mu \sigma }\psi _\nu ^\mu \psi _\nu ^\sigma \right|vac_{w,\nu }$$
The constant A rises from ”undoing” the normal ordering of the $`\psi `$’s. It takes a non vanishing value (equal to $`\frac{w}{2N}`$) only in the R-sector of the elliptic orbifold where there exist fermion zero modes <sup>18</sup><sup>18</sup>18In the Elliptic case in the sector $`w=\frac{N}{2}`$ there are fermion zero modes in the NS sector, we ignore that subtlety which is not important for the out discussion. (and not only quasi-zero modes).
Remembering that a physical state in the CFT must be a zero eigenstate of $`L_0`$ we study the effect of the operator $`\frac{w}{2}(\eta 𝒥)_{\mu \nu }\mathrm{\Sigma }^{\mu \nu }`$ on the spectrum.
* Elliptic Orbifolds Class: The operator $`\frac{w}{2}(\eta 𝒥)_{\mu \nu }\mathrm{\Sigma }^{\mu \nu }`$ is a generator of a rotations group $`\mathrm{SO}(2)`$ in the direction of the orbifold and has real eigenvalues. By changing the charges of the left moving and right moving sides it is possible to generate physical states obeying $`L_0=\stackrel{~}{L}_0=0`$ of ”mixed” type $`(\text{NS},\text{R})`$.
As is well known from the literature , in these theories if $`N`$ is an even integer, untwisted fermions cannot be introduced. By taking $`N`$ odd one is able to construct type II theory where the untwisted tachyon is projected out. The physical condition $`L_0=\stackrel{~}{L}_0=0`$ forces us to consider $`(\text{NS-},\text{R})`$ sectors for odd $`\omega `$ and there are tachyons in all twisted sectors.
* Hyperbolic Orbifolds Class: The operator $`\frac{w}{2}(\eta 𝒥)_{\mu \nu }\mathrm{\Sigma }^{\mu \nu }`$ is a generator of a boost group $`\mathrm{SO}(1,1)`$ in the direction of the orbifold it has pure imaginary eigenvalues. These eigenvalues can be compensated by a suitable bosonic wave-function.
* Parabolic Orbifolds Class: The operator $`\frac{w}{2}(\eta 𝒥)_{\mu \nu }\mathrm{\Sigma }^{\mu \nu }`$ is a generator of a null-boost group in the direction of the orbifold. The operator is nilpotent such that $`L_0`$ can be written in a Jordan form. Thus the CFT is logarithmic, at this point we cannot determine whether it is possible to find a consistent set of constraints on the spectrum (i.e BRST + GSO + orbifold projection) such that the resulting theory will have no branch cuts (mutual locality between operators) and will be modular invariant. |
warning/0507/cond-mat0507268.html | ar5iv | text | # Normal state Nernst effect, semiconducting-like resistivity and diamagnetism of underdoped cuprates
## Abstract
Semiconducting-like low-temperature in-plane resistivity indicates that there are no remnants of superconductivity above the resistive phase transition at $`T>T_c`$ in underdoped cuprates. The model with the chemical potential pinned near the mobility edge inside the charge-transfer optical gap describes quantitatively the Nernst effect, thermopower, diamagnetism and the unusual low-temperature resistivity of underdoped cuprates as normal state phenomena above $`T_c`$.
I
In the framework of the weak-coupling BCS theory the superconducting state is described by a nonzero Gor’kov anomalous average $`(𝐫,𝐫^{})=\psi _{}(𝐫)\psi _{}(𝐫^{})`$, which is zero above the resistive phase transition temperature $`T_c`$. When the BCS theory is extended to the strong-coupling regime, electrons are paired into lattice bipolarons, which are real-space pairs dressed by phonons, *both* below and above $`T_c`$ alebook . The state above $`T_c`$ is a normal charged Bose-liquid and below $`T_c`$ phase coherence of the preformed bosons sets in. In this regime $`(𝐫,𝐫^{})`$ describes bosons in the Bose-Einstein condensate similar to the Bogoliubov anomalous average of the annihilation operator in the Bose-gas. As in the BCS theory the state above $`T_c`$ is perfectly ”normal” in the sense that the off-diagonal order parameter $`(𝐫,𝐫^{})`$ is zero at $`T>T_c`$.
In disagreement with the weak-coupling BCS and the strong-coupling bipolaron theories a significant fraction of research in the field of superconducting cuprates claims that the superconducting transition is only a phase ordering while the superconducting order parameter $`(𝐫,𝐫^{})`$ remains nonzero above the resistive $`T_c`$. One of the key experiments supporting this viewpoint is the large Nernst signal observed in the normal state of cuprates (see xu ; cap ; cap2 and references therein). Refs xu ; ong propose a ”vortex scenario”, where the long-range phase coherence is destroyed by mobile vortices, but the amplitude of the off-diagonal order parameter remains finite and the Cooper pairing with a large binding energy exists well above $`T_c`$ supporting the so-called ”preformed Cooper-pairs” or ”the phase fluctuation” model kiv . The model is based on the assumption that superfluid density is small compared with the normal carrier density in cuprates. These claims seriously undermine many theoretical and experimental works on superconducting cuprates, which consider the state above $`T_c`$ as perfectly normal with no off-diagonal order.
However, the vortex scenario is unreconcilable with the extremely sharp resistive transitions at $`T_c`$ in high-quality samples of cuprates. For example, the in-plane and out-of-plane resistivity of $`Bi2212`$, where the anomalous Nernst signal has been measured xu , is perfectly normal above $`T_c`$, showing only a few percent positive or negative magnetoresistance zavale . The preformed Cooper-pairs model kiv is clearly incompatible with a great number of thermodynamic, magnetic, kinetic and optical measurements, which show that only holes (density $`x`$), doped into a parent insulator are carriers *both* in the normal and the superconducting states of cuprates. The assumption kiv that the superfluid density $`x`$ is small compared with the normal-state carrier density $`1x`$ is also inconsistent with the theorem leg , which proves that the number of supercarriers at $`T=0`$K should be the same as the number of normal-state carriers in any clean superfluid.
Faced with these inconsistences we have recently described the unusual Nernst signal in overdoped $`La_{1.8}`$Sr<sub>0.2</sub>CuO<sub>4</sub> in a different manner as the normal state phenomenon alezav . Here we extend our description to cuprates with low doping level accounting not only for their anomalous Nernst signal, but also for the thermopower, normal state diamagnetism and a semiconducting-like in-plane low-temperature resistivity as observed in recent xu ; cap ; cap2 ; ong and more earlier experiments.
In underdoped cuprates strong on-site repulsive correlations (Hubbard $`U`$) are essential in shaping the insulating state of parent compounds. The Mott-Hubbard insulator arises from a potentially metallic half-filled band as a result of the Coulomb blockade of electron tunnelling to neighboring sites mott . The first band to be doped in cuprates is the oxygen band inside the Hubbard gap. The strong electron-phonon interaction (see for experimental facts Ref. alebook ) creates oxygen hole polarons and inter-site bipolarons. Hence the chemical potential remains inside the optical charge-transfer gap, as clearly observed in the tunnelling experiments by Bozovic et al. boz0 . Disorder, inevitable with doping, creates localised impurity states for holes separated by a mobility edge from their extended states like in conventional amorphous semiconductors mott ; ell . Then the chemical potential should be found at or near the mobility edge in slightly doped cuprates, if they superconduct.
Naturally carriers, localised below the mobility edge, contribute to the normal-state longitudinal transport together with the itinerant carriers in extended states. On the other hand, the contribution of localised carriers of any statistics to the *transverse* transport is usually small as in many amorphous semiconductors ell . Importantly, if the localised-carrier contribution is not negligible, it *adds* to the contribution of itinerant carriers to produce a large Nernst signal, $`e_y(T,B)E_y/_xT`$, while it *reduces* the thermopower $`S`$ and the Hall angle $`\mathrm{\Theta }`$. This unusual ”symmetry breaking” is at variance with ordinary metals where the familiar ”Sondheimer” cancelation sond makes $`e_y`$ much smaller than $`S\mathrm{tan}\mathrm{\Theta }`$ because of the electron-hole symmetry near the Fermi level. Such behavior originates in the ”sign” (or ”$`pn`$”) anomaly of the Hall conductivity of localised carriers. The sign of their Hall effect is often $`opposite`$ to that of the thermopower as observed in many amorphous semiconductors ell and described theoretically fri .
The Nernst signal can be expressed in terms of the kinetic coefficients $`\sigma _{ij}`$ and $`\alpha _{ij}`$ as
$$e_y=\frac{\sigma _{xx}\alpha _{yx}\sigma _{yx}\alpha _{xx}}{\sigma _{xx}^2+\sigma _{xy}^2},$$
(1)
where the current density is given by $`j_i=\sigma _{ij}E_j+\alpha _{ij}_jT`$. When the chemical potential $`\mu `$ is at the mobility edge, the localised carriers contribute to the transport, so $`\sigma _{ij}`$ and $`\alpha _{ij}`$ in Eq.(1) can be expressed as $`\sigma _{ij}^{ext}+\sigma _{ij}^l`$ and $`\alpha _{ij}^{ext}+\alpha ^lij`$, respectively alezav . Since the Hall mobility of carriers localised below $`\mu `$, $`\sigma _{yx}^l`$, has the sign opposite to that of carries in the extended states above $`\mu `$, $`\sigma _{yx}^{ext}`$, the sign of the off-diagonal Peltier conductivity $`\alpha _{yx}^l`$ should be the same as the sign of $`\alpha _{yx}^{ext}`$. Then neglecting the magneto-orbital effects in the resistivity (since $`\mathrm{\Theta }1`$ xu ) we obtain
$$S\mathrm{tan}\mathrm{\Theta }\frac{\sigma _{yx}\alpha _{xx}}{\sigma _{xx}^2+\sigma _{xy}^2}\rho (\alpha _{xx}^{ext}|\alpha _{xx}^l|)(\mathrm{\Theta }^{ext}|\mathrm{\Theta }^l|)$$
(2)
and
$$e_y\rho (\alpha _{yx}^{ext}+|\alpha _{yx}^l|)S\mathrm{tan}\mathrm{\Theta },$$
(3)
where $`\mathrm{\Theta }^{ext}\sigma _{yx}^{ext}/\sigma _{xx}`$, $`\mathrm{\Theta }^l\sigma _{yx}^l/\sigma _{xx}`$, and $`\rho =1/\sigma _{xx}`$ is the resistivity.
Clearly the model, Eqs.(2,3) can account for a low value of $`S\mathrm{tan}\mathrm{\Theta }`$ compared with a large value of $`e_y`$ in underdoped cuprates xu ; cap2 because of the sign anomaly. Even in the case when localised carriers contribute little to the conductivity their contribution to the thermopower, $`S=\rho (\alpha _{xx}^{ext}|\alpha _{xx}^l|))`$, could almost cancel the opposite sign contribution of itinerant carriers. Indeed, if the carriers are bosons, their longitudinal conductivity in two-dimensions, $`\sigma ^{ext}_0𝑑EE𝑑f(E)/𝑑E`$ diverges logarithmically when $`\mu `$ in the Bose-Einstein distribution function $`f(E)=[\mathrm{exp}((E\mu )/T)1]^1`$ goes to zero and the relaxation time $`\tau `$ is a constant (here and further we take $`\mathrm{}=c=k_B=1`$). At the same time $`\alpha _{xx}^{ext}_0𝑑EE(E\mu )𝑑f(E)/𝑑E`$ remains finite, and it could have a magnitude comparable with $`\alpha _{xx}^l`$. Statistics of bipolarons effectively changes from Bose to Fermi-like statistics with lowering energy below the mobility edge because of the Coulomb repulsion of bosons in localised states alegile . Hence one can use the same expansion near the mobility edge as in ordinary amorphous semiconductors to obtain the familiar textbook result $`S=S_0T`$ with a constant $`S_0`$ at low temperatures mott3 .
The model becomes particularly simple, if we neglect the localised carrier contribution to $`\rho `$, $`\mathrm{\Theta }`$ and $`\alpha _{xy}`$, and take into account that $`\alpha _{xy}^{ext}B/\rho ^2`$ and $`\mathrm{\Theta }^{ext}B/\rho `$ in the Boltzmann theory. Then Eqs.(2,3) yield
$$S\mathrm{tan}\mathrm{\Theta }T/\rho $$
(4)
and
$$e_y(T,B)(1T/T_1)/\rho .$$
(5)
According to our earlier suggestion alelog the semiconducting-like dependence of $`\rho (T)`$ in underdoped cuprates (cap ; cap2 and references therein) at low temperatures originates from the elastic scattering of non-degenerate itinerant carriers by charged impurities, different from scenarios based on any kind of metal-insulator transitions. The relaxation time of *non-degenerate* carriers depends on temperature as $`\tau T^{1/2}`$ for scattering by short-range deep potential wells, and as $`T^{1/2}`$ for scattering by very shallow wells as discussed in Ref. alelog . Combining both scattering rates yields
$$\rho =\rho _0[(T/T_2)^{1/2}+(T_2/T)^{1/2}].$$
(6)
Eq.(6) with $`\rho _0=0.236`$ m$`\mathrm{\Omega }`$cm and $`T_2=44.6`$K fits well the experimental semiconducting-like normal state resistivity of underdoped La<sub>1.94</sub> Sr<sub>0.06</sub>CuO<sub>4</sub> in the whole low-temperature range from 2K up to 50K, Fig.1, as revealed in the field $`B=12`$ Tesla cap ; cap2 . Another high quality fit can be obtained combining the Brooks-Herring formula for the 3D scattering off screened charged impurities, as proposed in Ref.kast for almost undoped $`LSCO`$, or the Coulomb scattering in 2D ($`\tau T`$) and a temperature independent scattering rate off neutral impurities with the carrier exchange erg similar to the scattering of slow electrons by hydrogen atoms in three dimensions. Hence the scale $`T_2`$, which determines the crossover toward an insulating behavior, depends on the relative strength of two scattering mechanisms. Importantly the expressions (4,5) for $`S\mathrm{tan}\mathrm{\Theta }`$ and $`e_y`$ do not depend on particular scattering mechanisms, but only on the experimental $`\rho (T)`$. Taking into account the excellent fit of Eq.(6) to the experiment, these expressions can be parameterized as
$$S\mathrm{tan}\mathrm{\Theta }=e_0\frac{(T/T_2)^{3/2}}{1+T/T_2},$$
(7)
and
$$e_y(T,B)=e_0\frac{(T_1T)(T/T_2)^{1/2}}{T_2+T},$$
(8)
where $`T_1`$ and $`e_0`$ are temperature independent.
In spite of many simplifications, the model describes remarkably well both $`S\mathrm{tan}\mathrm{\Theta }`$ and $`e_y`$ measured in La<sub>1.94</sub> Sr<sub>0.06</sub>CuO<sub>4</sub> with a $`single`$ fitting parameter, $`T_1=50`$K using the experimental $`\rho (T)`$. The constant $`e_0=2.95`$ $`\mu `$V/K scales the magnitudes of $`S\mathrm{tan}\mathrm{\Theta }`$ and $`e_y`$. The magnetic field $`B=12`$ Tesla destroys the superconducting state of the low-doped La<sub>1.94</sub> Sr<sub>0.06</sub>CuO<sub>4</sub> down to $`2`$K, Fig.1, so any residual superconducting order above $`2`$K is clearly ruled out. At the same time the Nernst signal, Fig.2, is remarkably large. The coexistence of the large Nernst signal and a nonmetallic resistivity is in sharp disagreement with the vortex scenario, but is in agreement with our model. Taking into account the field dependence of the conductivity of localised carriers, their contribution to the transverse magnetotransport and the phonon-drug effect (at elevated temperatures) can well describe the magnetic field dependence of the Nernst signal alezav and improve the fit in Fig.2 but at the expense of the increasing number of fitting parameters.
Another experimental observation, which has been linked with the Nernst signal and mobile vortexes above $`T_c`$ ong , is enhanced diamagnetism. A number of experiments (see, for example, mac ; jun ; hof ; nau ; igu ; ong and references therein), including torque magnetometries, showed enhanced diamagnetism near and above $`T_c`$, which has been explained as fluctuation diamagnetism in quasi-2D superconducting cuprates (see, for example Ref. hof ). The data taken at relatively low magnetic fields (typically below 5 Tesla) revealed a crossing point in the magnetization $`M(T,B)`$ of most anisotropic cuprates (e.g. $`Bi2212`$), or in $`M(T,B)/B^{1/2}`$ of less anisotropic $`YBCO`$ jun . The dependence of magnetization (or $`M/B^{1/2}`$) on the magnetic field has been shown to vanish at some characteristic temperature below $`T_c`$ in agreement with conventional fluctuations. However the data taken in high magnetic fields (up to 30 Tesla) have shown that the crossing point, anticipated for low-dimensional superconductors and associated with superconducting fluctuations, does not explicitly exist in magnetic fields above 5 Tesla nau . Most surprisingly the torque magnetometery mac ; nau uncovered a diamagnetic signal somewhat above $`T_c`$ which *increases* in magnitude with applied magnetic field.
Here we argue that such behaviors are incompatible with the vortex scenario but can be understood with bipolarons. Accepting the vortex scenario and fitting the magnetization data for $`Bi2212`$ with the conventional logarithmic field dependence ong , one obtains surprisingly high upper critical fields $`H_{c2}>120`$ Tesla even at temperatures close to $`T_c`$, and a very large Ginzburg-Landau parameter, $`\kappa =\lambda /\xi >450`$ . The in-plane low-temperature magnetic field penetration depth is $`\lambda 200`$ nm in optimally doped $`Bi2212`$ (see, for example tal ). Hence the zero temperature coherence length $`\xi `$ turns out to be about the lattice constant, $`\xi =0.45`$nm, or even smaller. Such a small coherence length rules out the ”preformed Cooper pairs” kiv , since the pairs are virtually not overlapped at any size of the Fermi surface in $`Bi2212`$. Moreover the magnetic field dependence of $`M(T,B)`$ at and above $`T_c`$ is entirely inconsistent with what one expects of a vortex liquid. While $`M(B)`$ decreases logarithmically at temperatures well below $`T_c`$, the experimental curves mac ; nau ; ong clearly show that $`M(B)`$ increases with the field at and above $`T_c`$ , just the opposite of what one could expect in a vortex liquid. This significant departure from the London liquid behavior clearly indicates that the vortex liquid does not appear above the resistive phase transition mac .
Some time ago den we proposed that anomalous diamagnetism $`M(T,B)`$ in cuprates could be the Landau normal-state diamagnetism of preformed bosons. When the strong magnetic field is applied perpendicular to the copper-oxygen plains the quasi-2D bipolaron energy spectrum is quantized, $`E=\omega (n+1/2)+2t_c[1\mathrm{cos}(k_zd)]`$, where $`\omega =2eB/m_b`$, $`n=0,1,2,\mathrm{}`$, and $`t_c`$, $`k_z`$, $`d`$ are the hopping integral, the momentum and the lattice period perpendicular to the planes. Differentiating the thermodynamic potential one can readily obtain $`M(0,B)=n_b\mu _b`$ at low temperatures, $`TT_c`$, which is the familiar Schafroth’s result sha . Here $`n_b`$ is the bipolaron density, $`\mu _b=e/m_b`$ is the ”bipolaron” Bohr magneton, and $`m_b`$ is the bipolaron in-plane mass. The magnetization of charged bosons is field-independent at low temperatures. At high temperatures, $`TT_c`$ the bipolaron gas is almost classical. The experimental conditions are such that $`T\omega `$, when $`T`$ is of the order of $`T_c`$ or higher, so $`M(T,B)n_b\mu _b\omega /6T`$. It is the familiar Landau orbital diamagnetism of non-degenerate carriers. The bipolaron in-plane mass in cuprates is about $`m_b10m_e`$ as follows from a number of independent experiments and numerical (QMC) simulations alebook . Using this mass yields $`M(0,B)2000`$ A/m with the bipolaron density $`n_b=10^{21}`$ cm<sup>-3</sup>. Then the magnitude and the field/temperature dependence of $`M(T,B)`$ at and above $`T_c`$ are about the same as experimentally observed in Refs nau ; ong .
The pseudogap temperature $`T^{}`$, which is half of the bipolaron binding energy in the model, depends on the magnetic field because of spin-splitting of the single-polaron band by the magnetic-field. Also the singlet-triplet exchange energy of inter-site bipolarons depends on the field for the same reason. As a result the number of singlet bipolarons and thermally excited triplet pairs and single polarons depend on the field and on the temperature. When the depletion of the bipolaron density with temperature and magnetic field is taken into account, the crossing point in $`M(T,B)`$ disappears at high magnetic fields as observed, and the normal state magnetization of singlet bipolarons fits experimental $`M(T,B)`$ curves ong in the whole normal state and critical regions alemag .
In summary, we have described the normal state Nernst effect, the thermopower, the diamagnetism and the semiconducting-like in-plane resistivity of underdoped cuprates at low temperatures as the normal-state properties of non-degenerate oxygen holes doped into the Mott-Hubbard charge-transfer insulator with the chemical potential close to the mobility edge. The familiar ”sign” (or ”$`pn`$”) anomaly of the Hall conductivity of localised carriers accounts for a small value of $`S\mathrm{tan}\mathrm{\Theta }_H`$ compared with a large value of $`e_y`$. The semiconducting-like temperature dependence of the in-plane resistivity at low temperatures originates from the elastic scattering of non-degenerate itinerant carriers by charged impurities, rather than from any localisation. The enhanced diamagnetism at $`T>T_c`$ is the normal state orbital diamagnetism of bipolarons.
I thank V.V. Kabanov and V.N. Zavaritsky for valuable discussions. The work was supported by EPSRC (UK) (grant EP/C518365/1). |
warning/0507/cond-mat0507531.html | ar5iv | text | # Photo-induced ordering and anchoring properties of azo-dye films
## I Introduction
When a nematic liquid crystal (NLC) is brought into contact with an anisotropic substrate, the energy of the NLC molecules in the interfacial layer and thus the surface tension \[the excess free energy per unit area\] will be orientationally dependent. The anisotropic part of the surface tension — the so-called anchoring energy — gives rise to the phenomenon known as anchoring, i.e., surface induced alignment of the nematic director along the vector of preferential orientation referred to as the easy axis.
Over the past few decades anchoring properties of NLCs have been the subject of intense studies for both technological and more fundamental reasons. There are a number of surface ordering and anchoring transitions that were observed experimentally and were studied using different theoretical approaches \[see, e.g., Refs. Sluckin and Poniewierski (1986); Jérôme (1991); Barbero and Durand (1996) for reviews\].
Technologically, producing substrates with anisotropic anchoring properties is of vital importance in the fabrication of liquid crystal electrooptic devices. The traditional technique widely used to align liquid crystal display cells involves mechanical rubbing of aligning layers. This method, however, has the well known difficulties related to physical damage, impurities, dust contamination and generation of electrostatic charge Chigrinov (1999).
An alternative photoalignment technique avoiding the drawbacks of the mechanical surface treatment was suggested in Refs. Gibbons et al. (1991); Schadt et al. (1992); Dyadyusha et al. (1992). It uses linearly polarized ultraviolet (UV) light to induce anisotropy of the angular distribution of molecules in a photosensitive film O’Neill and Kelly (2000).
The photoalignment has been extensively studied in a number of different polymer systems such as dye doped polymer layers Gibbons et al. (1991); Furumi et al. (1999), cinnamate polymer derivatives Schadt et al. (1992); Dyadyusha et al. (1992); Galabova et al. (1996); Perny et al. (2000); Yaroshchuk et al. (2002) and side chain azopolymers Petry et al. (1993); Holme et al. (1996); Blinov et al. (1998); Wu et al. (2000); Yaroshchuk et al. (2001a, b); Yaroshchuk et al. (2003). Light induced ordering in the photosensitive materials, though not being understood very well, can occur by a variety of photochemically induced processes. These typically may involve such transformations as photoisomerization, crosslinking, photodimerization and photodecomposition (a recent review can be found in Ref. Chigrinov et al. (2003a)).
In this paper we examine anchoring properties of the films containing photochemically stable azo dye structures that were recently studied as new photoaligning materials for NLC cells Chigrinov et al. (2002); Chigrinov et al. (2003b). Dependence of the surface anchoring strengths on the photoinduced anisotropy will be of our primary interest.
More specifically, we are aimed to study the effects of the photoinduced ordering in azo-dye films on the polar and azimuthal anchoring energies. The key point is that the photoalignment technique provides a means for controlling the photoinduced ordering that affects anchoring properties of photoaligning layers by changing ordering of azo-dye molecules at the surface and, thus, the surface anchoring strengths.
Recently, the anchoring properties of aligning photopolymer layers in relation to the photoinduced ordering were studied experimentally in Ref. Thieghi et al. (2003). The relationship between the rubbing strength and the azimuthal anchoring energy was discussed in Ref. Oka et al. (2004)
The photopolymer-NLC interface was also described theoretically in Refs. Yaroshchuk et al. (2001c); Alexe-Ionescu et al. (2001) using a modified version of the variational mean field approach which is also known as the Maier-Saupe theory. By contrast, the azo-dye films have not yet received a proper attention and we intend to fill in the gap.
The paper is organized as follows. In Sec. II we apply the mean field theoretical approach Sen and Sullivan (1987); Teixeira and Sluckin (1992a, b) to express the surface anchoring energy in terms of the tensorial order parameters which characterize angular distribution of the azo-dye and NLC molecules at the interfacial boundary surface. The general result is then used to derive the expressions for the azimuthal and polar anchoring strengths that, in addition to the order parameters, depend on the harmonics of the intermolecular potentials.
Experimental details are given in Sec. III. The polymerizable azo-dye monomer SDA-2 was used to prepare the photoaligning layers. Absorption dichroism spectra were measured in the films irradiated with linearly polarized UV light at various irradiation doses. Anchoring energy measurements were performed in NLC cells where NLC is sandwiched between the glass plates coated with the azo-dye film.
In Sec. IV we present the experimental results and apply the theory of Sec. II to interpret the data. Discussion and concluding remarks are given in Sec. V. Details on some technical results are relegated to appendices A-B.
## II Theory
In this section we begin with introducing general notations and apply the mean-field approach to express the Landau-de Gennes surface free energy in terms of both azo-dye and NLC order parameters. Expressions for the azimuthal and polar anchoring strengths, $`W_\varphi `$ and $`W_\theta `$, are then derived from the orientationally dependent part of the surface energy in Sec. II.3. In the concluding part of this section we consider effects of spatial variations of the azo-dye order parameter using a simple model formulated in Sec. II.4.
### II.1 Order parameter and dichroic ratio
Assuming that the unit vector, $`\widehat{𝐮}=(\mathrm{sin}\theta \mathrm{cos}\varphi ,\mathrm{sin}\theta \mathrm{sin}\varphi ,\mathrm{cos}\theta )`$, directed along the long molecular axis defines orientation of a molecule in both azo-dye film and NLC cell, quadrupolar orientational ordering of the molecules can be characterized using the traceless symmetric second-rank tensor de Gennes and Prost (1993)
$$𝐐(\widehat{𝐮})=(3\widehat{𝐮}\widehat{𝐮}𝐈)/2,$$
(1)
where $`𝐈`$ is the identity matrix. The dyadic (1) averaged over orientation of molecules with the one-particle distribution function $`\rho _\alpha (𝐫,\widehat{𝐮})`$, describing the orientation-density profile of azo-dye ($`\alpha =\mathrm{A}`$) and NLC ($`\alpha =\mathrm{N}`$) molecules, is proportional to the *order parameter tensor* $`𝐒_\alpha (𝐫)`$
$`{\displaystyle \rho _\alpha (𝐫,\widehat{𝐮})𝐐(\widehat{𝐮})d\widehat{𝐮}}=\rho _\alpha (𝐫)𝐒_\alpha (𝐫),`$ (2)
where $`\mathrm{d}\widehat{𝐧}\mathrm{sin}\theta \mathrm{d}\theta \mathrm{d}\varphi `$, $`\rho _\alpha (𝐫,\widehat{𝐮})=\rho _\alpha (𝐫)f_\alpha (𝐫,\widehat{𝐮})`$, $`\rho _\alpha (𝐫)=\rho _\alpha (𝐫,\widehat{𝐮})d\widehat{𝐮}`$ is the density profile and $`f_\alpha (𝐫,\widehat{𝐮})`$ is the normalized angular distribution. The general expression for the order parameter is given in appendix A \[see Eq. (60)\] along with technical details on the technique of irreducible tensors.
Now we dwell briefly on the relation between the order parameter $`𝐒_\mathrm{A}`$ characterizing orientational distribution of azo-dye molecules $`f_\mathrm{A}(\widehat{𝐮})`$ and the absorption dichroic ratio
$$R=\frac{D_{}D_{}}{D_{}+2D_{}},$$
(3)
where $`D_{}`$ \[$`D_{}`$\] is the absorption coefficient measured for a testing beam linearly polarized parallel \[perpendicular\] to the polarization vector of the activating UV light which is directed along the $`x`$ axis, $`𝐄_{\mathrm{ex}}=E_{\mathrm{ex}}\widehat{𝐱}`$. We shall also assume that the testing and the pumping waves are both propagating along the $`z`$ axis which is normal to the film substrate.
When the absorption tensor of an azo-dye molecule is uniaxially anisotropic with $`\sigma _{ij}(\widehat{𝐮})=\sigma _{}\delta _{ij}+(\sigma _{}\sigma _{})u_iu_j`$, its orientational average takes the following matrix form
$`𝝈=\left(\sigma _{\mathrm{av}}𝐈+2\mathrm{\Delta }\sigma 𝐒_\mathrm{A}\right)/3,`$ (4)
$`\sigma _{\mathrm{av}}=\sigma _{}+2\sigma _{},\mathrm{\Delta }\sigma =\sigma _{}\sigma _{},`$ (5)
where the angular brackets $`\mathrm{}`$ denote orientational averaging.
In the low concentration approximation, the optical densities $`D_{}`$ and $`D_{}`$ are proportional to the corresponding components of the tensor (4)
$`D_{}\rho _\mathrm{A}\left(\sigma _{\mathrm{av}}+2\mathrm{\Delta }\sigma S_{xx}^{(\mathrm{A})}\right)/3,`$ (6)
$`D_{}\rho _\mathrm{A}\left(\sigma _{\mathrm{av}}+2\mathrm{\Delta }\sigma S_{yy}^{(\mathrm{A})}\right)/3,`$ (7)
so that the average absorption coefficient $`D_{\mathrm{av}}`$ is given by
$$D_{\mathrm{av}}=D_{}+2D_{}\rho _\mathrm{A}\left(\sigma _{\mathrm{av}}+2/3\mathrm{\Delta }\sigma \left[S_{yy}^{(\mathrm{A})}S_{zz}^{(\mathrm{A})}\right]\right).$$
(8)
When the absorption coefficient $`D_{\mathrm{av}}`$ does not depend on irradiation dose (and, thus, on the order parameter), from the expression (8) we may conclude that anisotropy of the azo-dye film is uniaxial and $`S_{yy}^{(\mathrm{A})}=S_{zz}^{(\mathrm{A})}=S_{xx}^{(\mathrm{A})}/2S_\mathrm{A}/2`$. In this case we have
$$𝐒_\mathrm{A}=S_\mathrm{A}(3\widehat{𝐱}\widehat{𝐱}𝐈)/2,R=\frac{\mathrm{\Delta }\sigma }{\sigma _{\mathrm{av}}}S_\mathrm{A}.$$
(9)
As is seen from Eq. (9), the dichroic ratio equals the order parameter only in the limiting case where absorption of waves propagating along the long molecular axis is negligibly small and $`\sigma _{}0`$.
### II.2 Anisotropic part of surface energy in the mean-field approximation
In the previous section it was shown that the light induced ordering of azo-dye molecules can be described by the order parameter (9) which is expected to affect the surface free energy at the nematic-substrate interface. So, in this section, the order parameter dependent part of the surface energy will be of our primary concern.
In the case of a flat structureless substrate, the expression for the surface energy was originally obtained by Sen and Sullivan in Ref. Sen and Sullivan (1987). Subsequently, similar results have been derived by using the mean-field approximation Teixeira and Sluckin (1992a) and the density functional theory Osipov and Hess (1993); Osipov and Sluckin (1993); Osipov et al. (1997).
Similarly to Ref. Teixeira and Sluckin (1992a) , we adopt the mean-field approach and use the Fowler approximation for the one-particle distribution functions
$`\rho _\mathrm{N}(𝐫,\widehat{𝐮})=H(z)\rho _\mathrm{N}(z,\widehat{𝐮}),\rho _\mathrm{A}(𝐫,\widehat{𝐮})=H(z)\rho _\mathrm{A}(z,\widehat{𝐮}),`$ (10)
where $`H(z)`$ is the Heaviside step function which equals unity when $`z`$ is positive and vanishes otherwise.
Applying the mean-field theory Teixeira and Sluckin (1992a) gives the Landau-de Gennes surface free energy as an excess Helmholtz free energy per unit area that depends on two pair intermolecular potentials: (a) the potential of interaction between NLC molecules, $`U_{\mathrm{N}\mathrm{N}}(𝐫_{12},\widehat{𝐮}_1,\widehat{𝐮}_2)`$; and (b) the potential of interaction between NLC and azo-dye molecules, $`U_{\mathrm{A}\mathrm{N}}(𝐫_{12},\widehat{𝐮}_1,\widehat{𝐮}_2)`$, where $`𝐫_{12}=𝐫_1𝐫_2`$ is the vector of intermolecular separation and $`\widehat{𝐮}_i`$ is the orientation coordinates of the interacting molecules. The resulting expression is given by
$`\mathrm{\Delta }F/A={\displaystyle _{\mathrm{}}^0}dz_1{\displaystyle _0^{\mathrm{}}}dz_2{\displaystyle d\widehat{𝐮}_1d\widehat{𝐮}_2}`$
$`\times [\rho _\mathrm{A}(z_1,\widehat{𝐮}_1)V_\mathrm{A}(z_{12},\widehat{𝐮}_1,\widehat{𝐮}_2)\rho _\mathrm{N}(z_2,\widehat{𝐮}_2)`$
$`{\displaystyle \frac{1}{2}}\rho _\mathrm{N}(z_1,\widehat{𝐮}_1)V_\mathrm{N}(z_{12},\widehat{𝐮}_1,\widehat{𝐮}_2)\rho _\mathrm{N}(z_2,\widehat{𝐮}_2)],`$ (11)
$`V_\alpha (z_{12},\widehat{𝐮}_1,\widehat{𝐮}_2)={\displaystyle _A}U_{\alpha \mathrm{N}}(𝐫_{12},\widehat{𝐮}_1,\widehat{𝐮}_2)dx_{12}dy_{12},`$ (12)
where $`A`$ is the area of the substrate and $`V_\alpha `$ is the potential averaged over in-plane coordinates.
It should be noted that the potentials $`U_{\mathrm{N}\mathrm{N}}`$ and $`U_{\mathrm{A}\mathrm{N}}`$ actually represent the perturbative part of interaction that can be treated in the mean-field approximation. They can be written in the form of expansion over spherical harmonics given in Eq. (B) of appendix B. For our purposes, however, it is more convenient to use the tensorial representation for the averaged potentials $`V_\alpha `$ that was introduced in Ref. Ronis and Rosenblatt (1980). In appendix B, the coefficients that enter this representation \[see Eq. (B)\] are related to the coefficients, $`v_{j_1j_2j}(z)`$ with $`j_i<4`$, in the spherical harmonics expansion (B). This relation is given by Eqs. (78)–(81).
Substituting the representation (B) into Eq. (II.2) and assuming homogeneity of the contacting phases, we obtain the Landau-de Gennes expression for the surface free energy in the final form:
$`f_S(𝐒_\mathrm{N},𝐒_\mathrm{A})=f_\mathrm{N}(𝐒_\mathrm{N})++f_\mathrm{A}(𝐒_\mathrm{N},𝐒_\mathrm{A}),`$ (13)
$`f_\mathrm{N}(𝐒_\mathrm{N})`$ $`=c_0\widehat{𝐳}𝐒_\mathrm{N}\widehat{𝐳}+c_\mathrm{N}^{(1)}Tr(𝐒_\mathrm{N}^2)`$
$`+c_\mathrm{N}^{(2)}\widehat{𝐳}𝐒_\mathrm{N}^2\widehat{𝐳}+c_\mathrm{N}^{(3)}\left[\widehat{𝐳}𝐒_\mathrm{N}\widehat{𝐳}\right]^2,`$ (14)
$`f_\mathrm{A}(𝐒_\mathrm{N},𝐒_\mathrm{A})=c_\mathrm{A}^{(1)}Tr(𝐒_\mathrm{N}𝐒_\mathrm{A})`$
$`+c_\mathrm{A}^{(2)}\widehat{𝐳}𝐒_\mathrm{N}𝐒_\mathrm{A}\widehat{𝐳}+c_\mathrm{A}^{(3)}\left[\widehat{𝐳}𝐒_\mathrm{N}\widehat{𝐳}\right]\left[\widehat{𝐳}𝐒_\mathrm{A}\widehat{𝐳}\right],`$ (15)
where the coefficients are given by
$`c_0=b_\mathrm{A}^{(0)}b_\mathrm{N}^{(0)},`$ (16)
$`c_\mathrm{A}^{(i)}=b_\mathrm{A}^{(i)},c_\mathrm{N}^{(i)}=b_\mathrm{N}^{(i)}/2,`$ (17)
$`b_\alpha ^{(i)}=\rho _\alpha \rho _\mathrm{N}{\displaystyle _0^{\mathrm{}}}z\beta _\alpha ^{(i)}(z)dz,`$ (18)
$`\beta _\alpha ^{(i)}(z)`$ denote the coefficients in the representation (B) for the potential (12).
Eqs. (13)–(II.2) can be viewed as a generalization of the expression by Sen and Sullivan Sen and Sullivan (1987) supplemented with the term $`f_\mathrm{A}(𝐒_\mathrm{N},𝐒_\mathrm{A})`$ resulting from the interaction between NLC and azo-dye molecules. Note that this result can also be derived by constructing invariants from the order parameter tensors $`𝐒_\alpha `$ and the normal to the substrate $`\widehat{𝐳}`$. In this case, the surface of the azo-dye aligning film is treated phenomenologically as a bounding surface which, in addition to the normal $`\widehat{𝐳}`$, is characterized by the order parameter $`𝐒_\mathrm{A}`$.
### II.3 Bare anchoring energy
Separating out the director dependent part of the surface free energy requires the order parameters of azo-dye and NLC molecules be substituted into Eqs. (13)–(II.2). Since the order parameter at the surface may differ from its value in the bulk, we generalize the expression for the azo-dye order parameter (9) as follows
$`2𝐒_\mathrm{A}|_{z=0}=S_\mathrm{A}(3\widehat{𝐱}\widehat{𝐱}𝐈)+P_\mathrm{A}(\widehat{𝐳}\widehat{𝐳}\widehat{𝐲}\widehat{𝐲}).`$ (19)
Similarly, for NLC order parameter tensor at $`z=0`$, from Eq. (60) we have
$`2𝐒_\mathrm{N}|_{z=0}=S(3\widehat{𝐧}\widehat{𝐧}𝐈)+P(\widehat{𝐦}\widehat{𝐦}\widehat{𝐥}\widehat{𝐥})`$
$`=(3S+P)\widehat{𝐧}\widehat{𝐧}+2P\widehat{𝐦}\widehat{𝐦}(S+P)𝐈,`$ (20)
where $`\widehat{𝐧}=(\mathrm{sin}\mathrm{\Theta }\mathrm{cos}\mathrm{\Phi },\mathrm{sin}\mathrm{\Theta }\mathrm{sin}\mathrm{\Phi },\mathrm{cos}\mathrm{\Theta })`$ is the NLC director, $`\widehat{𝐧}\widehat{𝐦}=\mathrm{cos}\gamma 𝐞_x(\widehat{𝐧})\mathrm{sin}\gamma 𝐞_y(\widehat{𝐧})`$, $`𝐞_x(\widehat{𝐧})=(\mathrm{cos}\mathrm{\Theta }\mathrm{cos}\mathrm{\Phi },\mathrm{cos}\mathrm{\Theta }\mathrm{sin}\mathrm{\Phi },\mathrm{sin}\mathrm{\Theta })`$, $`𝐞_y(\widehat{𝐧})=(\mathrm{sin}\mathrm{\Phi },\mathrm{cos}\mathrm{\Phi },0)`$ and $`\widehat{𝐥}=\widehat{𝐧}\times \widehat{𝐦}`$.
Eqs. (19) and (II.3) suggest that the order parameters of azo-dye and NLC molecules though being both uniaxial in the bulk can be biaxial at the surface. In addition, the scalar order parameters $`S_\mathrm{A}`$ and $`S`$ at the surface may also deviate from their values in the bulk.
The surface free energy can now be expressed as a sum of two contributions
$`f_S(𝐒,𝐒_\mathrm{A})=W(\widehat{𝐧},\widehat{𝐦})+f_{\mathrm{scal}},`$ (21)
where $`W(\widehat{𝐧},\widehat{𝐦})`$ is the orientationally dependent part of the surface energy. This part can be calculated by substituting the order parameters (19) and (II.3) into the surface energy (13) to yield the expression for the bare anchoring energy
$`W(\widehat{𝐧},\widehat{𝐦})=N_z\left(\widehat{𝐧}\widehat{𝐳}\right)^2+M_z\left(\widehat{𝐦}\widehat{𝐳}\right)^2`$
$`+N_x\left(\widehat{𝐧}\widehat{𝐱}\right)^2+M_x\left(\widehat{𝐦}\widehat{𝐱}\right)^2`$
$`+c_\mathrm{N}^{(3)}/4\left[(3S+P)\left(\widehat{𝐧}\widehat{𝐳}\right)^2+2P\left(\widehat{𝐦}\widehat{𝐳}\right)^2\right]^2`$ (22)
with the coefficients defined by the relations:
$`4N_z=(3S+P)\left[q+c_\mathrm{N}^{(2)}(SP)\right],`$ (23)
$`4N_x=c_\mathrm{A}^{(1)}(3S_\mathrm{A}+P_\mathrm{A})(3S+P),`$ (24)
$`2M_z=P\left[q2c_\mathrm{N}^{(2)}S\right],2M_x=c_\mathrm{A}^{(1)}(3S_\mathrm{A}+P_\mathrm{A})P,`$ (25)
$`q2c_0\left(c_\mathrm{A}^{(2)}+c_\mathrm{A}^{(3)}\right)(S_\mathrm{A}P_\mathrm{A})+2c_\mathrm{A}^{(1)}P_\mathrm{A}2c_\mathrm{N}^{(3)}(S+P).`$ (26)
The second term on the right hand side of Eq. (21)
$`4f_{\mathrm{scal}}=\left[2c_0(S_\mathrm{A}P_\mathrm{A})(c_\mathrm{A}^{(2)}+c_\mathrm{A}^{(3)})+3c_\mathrm{A}^{(1)}(S_\mathrm{A}+P_\mathrm{A})\right]`$
$`\times (S+P)+(c_\mathrm{N}^{(2)}+c_\mathrm{N}^{(3)})(S+P)^2+2c_\mathrm{N}^{(1)}(3S^2+P^2).`$ (27)
is a quadratic function of NLC scalar order parameter $`S`$ and the biaxiality $`P`$. From Eq. (II.3) it, similarly to the anchoring energy (II.3), depends linearly on the azo-dye parameters $`S_\mathrm{A}`$ and $`P_\mathrm{A}`$.
For the anchoring energy (II.3), we consider the simplest case which occurs when the surface induced NLC biaxiality $`P`$ is negligibly small and the quadrupolar term $`v_{224}`$ in the expansion of the intermolecular potential $`V_\mathrm{N}`$ can be ignored. Under these circumstances, setting $`P=c_\mathrm{N}^{(3)}=0`$ and $`M_z=M_x=0`$, we arrive at the simplified formula for the anchoring energy
$$W(\widehat{𝐧})=N_z\left(\widehat{𝐧}\widehat{𝐳}\right)^2+N_x\left(\widehat{𝐧}\widehat{𝐱}\right)^2$$
(28)
which agrees with the expression for the anchoring energy recently proposed in Refs. Zao et al. (2000); Zhao et al. (2002); Guo-Chen et al. (2004).
From Eq. (28) it is clear that the easy axis is directed along the $`y`$ axis, $`𝐞_s=\widehat{𝐲}`$, only if the coefficients $`N_z`$ and $`N_z`$ are both positive. In this case the polar and the azimuthal anchoring strengths, $`W_\theta `$ and $`W_\varphi `$, are given by
$$W_\theta =2N_z,W_\varphi =2N_x.$$
(29)
From Eq. (24) we immediately deduce a more explicit expression for the azimuthal anchoring strength
$`W_\varphi =w_\varphi \left[S_{yy}^{(\mathrm{A})}|_{z=0}S_{xx}^{(\mathrm{A})}|_{z=0}\right],`$ (30)
$`2w_\varphi =3c_\mathrm{A}^{(1)}S|_{z=0},`$ (31)
where notations indicate the plane $`z=0`$ as a surface separating the phases.
Similarly, Eq. (23) gives the polar anchoring strength in the explicit form
$`W_\theta =w_\theta ^{(0)}+w_\theta ^{(1)}S_{zz}^{(\mathrm{A})}|_{z=0}w_\varphi \left[S_{zz}^{(\mathrm{A})}|_{z=0}S_{yy}^{(\mathrm{A})}|_{z=0}\right],`$ (32)
$`2w_\theta ^{(1)}=3\left(c_\mathrm{A}^{(2)}+c_\mathrm{A}^{(3)}\right)S|_{z=0}.`$ (33)
The formulas (30)-(33) will be subsequently used in Sec. IV to interpret the experimental data. At this stage, it is worth noting that, for the order parameter (9), the relations (31) and (33) provide the inequalities
$$c_\mathrm{A}^{(1)}<0,c_\mathrm{A}^{(2)}+c_\mathrm{A}^{(3)}>0$$
(34)
as conditions for the anchoring strengths to increase linearly as the scalar order parameter $`S_\mathrm{A}`$ decreases.
### II.4 Model of spatially varying order parameter
As it was pointed out at the beginning of the previous section, the surface order parameter tensor of an azo-dye film (19) may differ from the bulk order parameter of the film (9). The latter, according to the experimental results presented in the subsequent section III, is uniaxially anisotropic with the in-plane anisotropy axis that is normal to the polarization vector of the activating UV light. In addition, the light induced scalar order parameter, which is proportional to the dichroic ratio (3), turns out to be negative, $`S_\mathrm{A}^{(b)}<0`$.
From the other hand, assuming that the anisotropic part of the surface energy can be taken in the general form by Sen and Sullivan Sen and Sullivan (1987), the boundary conditions may favor either homeotropic or planar alignment of the azo-dye molecules, thus, counteracting the action of light. So, it can be expected that the effects caused by interplay between the light induced and the surface ordering are of importance in explaining the order parameter dependencies of the polar and azimuthal anchoring energies.
In this section we discuss these effects on the basis of a simple phenomenological model formulated by using the polar representation (66) for the azo-dye order parameter. The latter can be conveniently rewritten in the form
$$S_\mathrm{A}=s_\mathrm{A}\mathrm{cos}\psi ,P_\mathrm{A}=\sqrt{3}s_\mathrm{A}\mathrm{sin}\psi ,$$
(35)
where the angle $`\psi `$ is shifted by $`\pi `$ so as to have the angle $`\psi `$ vanishing in the bulk.
In what follows we shall assume that, similarly to nematic liquid crystals Lyuksyutov (1978); Penzenstadler and Trebin (1989); Rosso and Virga (1996); Kralj and Virga (2001), the amplitude $`s_\mathrm{A}`$ varies in space much slower than the angle $`\psi `$. So, in our model, the amplitude will be fixed at its bulk value, $`s_\mathrm{A}=|S_\mathrm{A}^{(b)}|`$, and we consider the limiting case of thick films in which the characteristic length of spatial variations of the angle $`\psi `$ is much shorter than the film thickness. In this case the film can be regarded as a semi-infinite sample filling the upper half space, $`z0`$.
Technically, our task will be to find the spatially varying angle $`\psi `$ as a function of $`z`$ that minimizes the excess free energy per unit area, $`F_\mathrm{A}`$, taken in the following nematic-like form
$`F_\mathrm{A}={\displaystyle _0^{\mathrm{}}}\left[Ls_\mathrm{A}^2\left(_z\psi \right)^2+Bs_\mathrm{A}^3\left(1\mathrm{cos}(3\psi )\right)\right]dz`$
$`+G_1s_\mathrm{A}\mathrm{cos}(\psi _0+\pi /3)+G_2s_\mathrm{A}^2\mathrm{cos}^2(\psi _0+\pi /3),`$ (36)
where $`\psi _0\psi |_{z=0}`$ and $`s_\mathrm{A}\mathrm{cos}(\psi _0+\pi /3)=\widehat{𝐳}𝐒_\mathrm{A}\widehat{𝐳}|_{z=0}`$. The first part of the excess free energy (II.4) is of integral form with the integrand describing the energy costs for deviations of $`\psi `$ from the equilibrium value, $`\psi =0`$. The gradient term of the energy density is taken to be proportional to $`(_z𝐒_\mathrm{A})^2`$, whereas the other term gives an increase in energy caused by spatially uniform changes in $`\psi `$. This term is written as a linear function of the angle dependent invariant (64), $`4Tr[𝐒_\mathrm{A}^3]=3s_\mathrm{A}^3\mathrm{cos}(3\psi )`$.
For the order parameter (19), the surface part of the energy (II.4) can be represented by a quadratic polynomial of $`\widehat{𝐳}𝐒_\mathrm{A}\widehat{𝐳}`$. By contrast to the elastic constant $`L`$ and the coefficient $`B`$, the surface coupling constants, $`G_1`$ and $`G_2`$, can generally be negative leading to different boundary conditions. For example, if $`G_2=0`$, minimizing the surface term requires the $`z`$-component of the order parameter, $`\widehat{𝐳}𝐒_\mathrm{A}\widehat{𝐳}=S_{zz}^{(\mathrm{A})}`$, to attain its maximal (minimal) value at the surface provided the coefficient $`G_1`$ is negative (positive). These can be referred to as the homeotropic (planar) boundary conditions.
The Euler-Lagrange equation for the free energy functional (II.4) can be readily solved to yield the relation
$`\mathrm{tan}(3\psi /4)=\mathrm{tan}(3\psi _0/4)\mathrm{exp}(z/\xi ),`$ (37)
where $`9\xi ^2=2L(Bs_\mathrm{A})^1`$. This relation can now be substituted into Eq. (II.4) to derive the free energy as a function of the angle $`\psi _0`$. The result is
$`F_\mathrm{A}/h\stackrel{~}{f}_\mathrm{A}(\psi _0)`$ $`=s_\mathrm{A}^{5/2}\mathrm{sin}^2(3\psi _0/4)+g_1s_\mathrm{A}\mathrm{cos}(\psi _0+\pi /3)`$
$`+g_2s_\mathrm{A}^2\mathrm{cos}^2(\psi _0+\pi /3),`$ (38)
where $`h=4(2BL)^{1/2}/3`$ and $`g_i=G_i/h`$. The angle $`\psi `$ at the surface then can be found as the value of $`\psi _0`$ that minimizes the function (II.4) on the interval ranged from $`2\pi /3`$ to $`2\pi /3`$.
Qualitatively, dependence of $`\psi _0`$ on the coupling constant $`g_1`$ can be analyzed using elementary methods. For $`g_20`$, we find that the angle $`\psi _0`$ is localized within different intervals depending on the value of $`g_1`$. These are given by
$$\{\begin{array}{cc}\pi /3<\psi _00,\hfill & g_1g_c^{(1)}=g_2s_\mathrm{A},\hfill \\ 0<\psi _0\pi /3,\hfill & g_c^{(1)}<g_1g_c^{(2)}=g_2s_\mathrm{A}+\sqrt{3}/2s_\mathrm{A}^{3/2},\hfill \\ \pi /3<\psi _02\pi /3,\hfill & g_c^{(2)}<g_1g_c^{(3)}=2g_2s_\mathrm{A}+9/8s_\mathrm{A}^{3/2},\hfill \\ \psi _0=2\pi /3,\hfill & g_1>g_c^{(3)}.\hfill \end{array}$$
(39)
The end points of the intervals in Eq. (39), $`\psi _0=k\pi /3`$ with $`1k2`$, represent the uniaxially anisotropic structures at the surface
$`𝐒_\mathrm{A}(0)=s_\mathrm{A}(3\widehat{𝐳}\widehat{𝐳}𝐈)/2,`$ $`\psi _0=\pi /3,`$ (40)
$`𝐒_\mathrm{A}(0)=s_\mathrm{A}(3\widehat{𝐱}\widehat{𝐱}𝐈)/2,`$ $`\psi _0=0,`$ (41)
$`𝐒_\mathrm{A}(0)=s_\mathrm{A}(3\widehat{𝐲}\widehat{𝐲}𝐈)/2,`$ $`\psi _0=\pi /3,`$ (42)
$`𝐒_\mathrm{A}(0)=s_\mathrm{A}(3\widehat{𝐳}\widehat{𝐳}𝐈)/2,`$ $`\psi _0=2\pi /3.`$ (43)
From Eqs. (40) and (42) it is seen that, for the angles $`\psi _0=\pi /3`$ and $`\psi _0=\pi /3`$, surface alignment will be homeotropic and homogeneous (monostable planar), respectively. The structure (41) coincides with uniaxial ordering in the bulk (9) and the surface order parameter tensor (43) corresponds to planar (random in-plane) alignment.
According to Eq. (39), the case of planar alignment occurs only if the coupling constant $`g_1`$ is positive and the bulk order parameter $`s_\mathrm{A}`$ is below its critical value $`s_c`$ defined by the relation
$`g_1=2g_2s_c+9/8s_c^{3/2}.`$ (44)
Fig. 1 shows that the critical order parameter $`s_c`$ is an increasing function of $`g_1`$.
In Fig. 2(a) we have plotted the curves representing the components of the order parameter tensor at the surface in relation to the order parameter in the bulk to illustrate that destruction of the planar alignment takes place in a second order transition manner.
However, it should be stressed that our model becomes inapplicable in the immediate vicinity of $`s_c`$ where $`z`$-dependence of the biaxiality parameter $`P_\mathrm{A}`$ critically slows down. Actually, as it can be inferred from Fig. 3, the characteristic length of spatially varying biaxiality diverges logarithmically as $`s_\mathrm{A}`$ approaches $`s_c`$ from above. Under these circumstances, the assumption that scale of the spatial variations is much shorter than the film thickness is no more justified.
By contrast to the boundary conditions with $`g_1>0`$, there are no second order transitions provided the coupling constant $`g_1`$ is negative. At sufficiently large values of $`|g_1|`$, the surface ordering remains nearly homeotropic. Otherwise, the surface order parameter changes smoothly with $`s_\mathrm{A}`$ towards the bulk order parameter (41). From Eq. (39), for $`g_1=g_2s_\mathrm{A}`$, difference between the order parameters vanishes. The curves presented in Fig. 2(b) illustrate this point.
Leaving aside a detailed discussion of what happen when the coupling constant $`g_2`$ is negative, we just note that in this case the above discussed transition will generally be first order leading to jump-like behavior of the order parameter at the surface.
## III Experiment
Now we pass on to describing the experimental procedure employed to obtain the data linking the anchoring energy strengths and the dichroic ratio as a measure of the photo-induced ordering. To this end we carried out the absorption spectra and the anchoring energy measurements for the azo-dye films irradiated at varying exposure time. Thus, we used the samples prepared at different irradiation doses to measure the anchoring strengths and the dichroic ratio in relation to the dose. The data then can be recalculated to obtain the required anchoring energy vs dichroic ratio dependence.
### III.1 Sample preparation
Following the method described in Ref. Chigrinov et al. (2002), the azobenzene sulfuric dye SD-1 was synthesized from corresponding benzidinedisulfonic acid using azo coupling. The azo-dye compound SD-1 was mixed with the polymerizable azo-dye SDA-2 in the ratio 40% to 60%. The mixture was dissolved in N,N-dimethylformamide (DMF) and a heat initiator V-65 (from Wako pure chemical industries, Ltd.) that was added in relation of 1:50 to SDA-2.
The solution was spin-coated onto glass substrates with indium-tin-oxide (ITO) electrodes at 800 rpm for 5 seconds and, subsequently, at 3000 rpm for 30 seconds. The solvent was evaporated on a hot plate at 100C for 10 minutes.
The surface of the coated film was illuminated with linearly polarized UV light using super-high pressure Hg lamp through an interference filter at the wavelength 365 nm. The intensity of light irradiated on the film surface at varying time exposure was 2.7 mW/cm<sup>2</sup>. After the photoaligning procedure, the SDA-2 films were polymerized by heating at 150C for 1 hour in vacuum. In order to recover quality of the photoalignment degraded after the polymerization, the films were exposed to the UV light for 1 minute regardless of the initial time exposure.
Two glass substrates with the photoaligned films were assembled to form liquid crystal cells to measure the azimuthal and polar anchoring energy strengths. The cell thickness was 5 $`\mu `$m and 18 $`\mu `$m, respectively. Liquid crystal mixtures MLC-6080 (from Merck) in an isotropic phase were injected into the cell by capillary action.
### III.2 Absorption spectra
The UV-visible absorption spectra of the films were measured in the spectral range from 250 nm to 600 nm for the normally incident probing light which is linearly polarized parallel (along the $`x`$ axis) and perpendicular (along the $`y`$ axis) to the polarization vector of the activating light.
For nonirradiated films, the curve shown in Fig. 4 as a solid line demonstrates that the absorption coefficient does not depend on the polarization state of the testing beam. By contrast, as it is illustrated in Fig. 4, the absorption coefficients $`D_{}`$ and $`D_{}`$ differ in the irradiated films, thus, revealing the light induced absorption dichroism. This dichroism is mainly caused by the photo-induced angular redistribution of the azo-dye molecules.
By varying the exposure time the films were prepared at different irradiation doses and the optical density components $`D_{}`$ and $`D_{}`$ at the absorption maximum of azo-dyes ($`\lambda _m350`$ nm) were estimated from the measured absorption spectra. The dichroic ratio then can be computed from the formula (3). The results for the absorption coefficients and the absorption order parameters, which are proportional to the dichroic ratio, are presented in Figs. 5(a) and 5(b), respectively.
### III.3 Anchoring energy strengths
The azimuthal anchoring strength, $`W_\varphi `$, was measured in a twisted nematic cell using the torque balance method Iimura et al. (1995); Vorflusev et al. (1995). The azo-dye aligning film and a rubbed polyimide layer were used as confining substrates. The twist angle was 90 degrees.
Measurements of the polar anchoring strength, $`W_\theta `$, in anti-parallel aligned cells were carried out using the high-voltage technique Yokoyama and van Sprang (1985); Nastishin et al. (1999); Yokoyama and Sun (2000); Chigrinov et al. (2003b). The experimental data for the azimuthal and polar anchoring strengths are plotted against the irradiation dose in Fig. 8(a) and Fig. 9(a), respectively. At low irradiation doses, when the exposure energy is below 1 J/cm<sup>2</sup>, our experimental technique failed to provide accurate estimates for the anchoring strengths because of poor quality of NLC alignment in the cells.
The experimentally measured dependence of the dichroic ratio on the irradiation dose can now be combined with the results for $`W_\varphi `$ and $`W_\theta `$ so as to recalculate the anchoring energy strengths as functions of the dichroic ratio, $`R`$. The resulting data are shown in Figs. 8(b) and 9(b).
## IV Results
As is shown in Fig. 5(a) representing the absorption coefficients, $`D_{}`$ and $`D_{}`$, measured in the film irradiated at various irradiation doses, within the limits of experimental error, the average absorption coefficient, $`D_{\mathrm{av}}`$, defined by the relation (8), remains unchanged at irradiation doses higher than 1 J/cm<sup>2</sup>. So, from the discussion given at the end of Sec. II.1 we conclude that the azo-dye order parameter in the bulk of the film is uniaxial and is of the form given by Eq. (9). In Fig. 5(b), it is indicated that the components of the order parameter tensor are proportional to the dichroic ratio (3). Clearly, for $`D_{}>D_{}`$ and $`\sigma _{}>\sigma _{}`$, the azo-dye order parameter $`S_\mathrm{A}`$ and $`R`$ are both negative.
The experimental results for the azimuthal and polar anchoring strengths measured in NLC cells with photoaligned azo-dye films used as aligning substrates are presented in Figs. 6(a) and 7(a), respectively. The films differ in an amount of photoinduced anisotropy which is controlled by varying exposure time and the anchoring strengths are plotted in relation to the irradiation dose.
However, the fundamentally important characteristic describing degree of the photoinduced anisotropy is the azo-dye order parameter. So, in order to compare the experimental data and the theory, we need to relate the anchoring strengths and the dichroic ratio. Combining the anchoring energy data and the curve depicted in Fig. 5(b) gives the result shown in Figs. 6(b) and 7(b).
From the relations (30)-(33) the anchoring strengths depend linearly on the dichroic ratio provided the order parameters at the surface do not differ from their values in the bulk. The results of linear fitting of the experimental data are shown as solid straight lines in Figs. 6(b) and 7(b).
Referring to Fig. 6(b), the linear approximation for $`W_\varphi `$ predicts that the azimuthal anchoring strength vanish at certain non-zero value of the dichroic ratio, $`R0.16`$. By contrast, from the formulas (30) and (31) the anchoring strength $`W_\varphi `$ is proportional to $`R`$ and, thus, disappear only in the limit of weak photoinduced anisotropy where $`R0`$. Assuming that this discrepancy can be attributed to the difference between the bulk and surface order parameters of azo-dye, we can apply the phenomenological model described in Sec. II.4 to interpret the experimental data.
In the angle-amplitude representation (35), the surface order parameters that enter the expression (30) are given by
$`S_{xx}^{(\mathrm{A})}|_{z=0}=s_\mathrm{A}\mathrm{cos}\psi _0,S_{yy}^{(\mathrm{A})}|_{z=0}=s_\mathrm{A}\mathrm{cos}(\psi _0\pi /3),`$ (45)
where $`s_\mathrm{A}=S_\mathrm{A}=(\mathrm{\Delta }\sigma /\sigma _{\mathrm{av}})R`$ is the scalar order parameter in the bulk of the azo-dye film \[see Eq. (9)\]. According to our model, $`\psi _0`$ is the angle that minimizes the energy (II.4).
So, the computational procedure involves two steps: (a) minimization of the energy (II.4) to find the angle $`\psi _0`$; and (b) using the relations (45) to compute the azimuthal anchoring energy (30). Following this procedure, we may calculate dependence of the anchoring strength $`W_\varphi `$ on the photoinduced order parameter $`s_\mathrm{A}`$.
As it is discussed in Sec. II.4, the result crucially depend on the boundary conditions that are determined by two coupling constants, $`g_1`$ and $`g_2`$. At $`g_10`$, the surface favors planar (random in-plane) alignment of the azo-dye molecules. In the opposite case of negative coupling constant $`g_1`$, the alignment is homeotropic.
In figure 8, we show the theoretical curves calculated for both planar and homeotropic boundary conditions. The corresponding numerical results for the polar anchoring strength are presented in Fig. 9.
The curve plotted in Fig. 8(a) indicates that, for planar boundary conditions, the azimuthal anchoring strength $`W_\varphi `$ takes non-zero values and starts growing only if the azo-dye order parameter $`s_\mathrm{A}`$ exceeds its critical value $`s_c`$. Such threshold behavior is a consequence of the second order transition discussed in Sec. II.4. Contrastingly, as is shown in Fig. 8(b), $`W_\varphi `$ is a smoothly increasing function of $`s_\mathrm{A}`$ when the boundary conditions favor the homeotropic alignment at the surface.
At first glance, the curves representing the polar anchoring strength $`W_\theta `$ plotted against $`s_\mathrm{A}`$ in Figs. 9(a)-(b) do not show any noticeable differences. It, however, should be stressed that the planar boundary conditions prevent the polar anchoring energy from decaying to zero as the order parameter $`s_\mathrm{A}`$ decreases. From Fig. 9(b) it can be seen that this is no longer the case when the coupling constant $`g_1`$ becomes negative.
These results show that interplay between photoinduced ordering in the bulk of the azo-dye films and the preferred alignment of the azo-dye molecules at the surface may have a profound effect on the order parameter dependence of the anchoring energies. For the planar alignment with $`g_1>0`$, our model predicts that the surface ordering may change through the second order transition as the photoinduced anisotropy increases. This transition bears close resemblance to the second order transitions in nematic liquid crystals previously studied in Refs. Sheng et al. (1992); Qian and Sheng (1997).
Quantitatively, the polar anchoring energy $`W_\theta `$ appears to be an order of magnitude higher than the azimuthal energy $`W_\varphi `$. From our estimates, the ratio of the coefficients $`w_\theta ^{(1)}`$ and $`w_\varphi `$, $`w_\theta ^{(1)}/w_\varphi =(c_\mathrm{A}^{(2)}+c_\mathrm{A}^{(3)})/c_\mathrm{A}^{(1)}`$, is likely to be well above 10.
The coefficient $`c_\mathrm{A}^{(1)}`$ is negative and its sign is determined by the dominating contribution from the Maier-Saupe term $`v_{220}`$ of the spherical harmonics expansion for the azo-dye – NLC intermolecular potential $`U_{\mathrm{A}\mathrm{N}}`$. From Eq. (B) the absolute value of $`c_\mathrm{A}^{(2)}`$ can be significantly reduced when the quadrupole term $`v_{224}`$ and the harmonics $`v_{222}`$ are predominately positive, so that $`_0^{\mathrm{}}zv_{224}(z)dz>0`$ and $`_0^{\mathrm{}}zv_{222}(z)dz>0`$.
Under these circumstances, the condition (34), that requires the sum of $`c_\mathrm{A}^{(2)}`$ and $`c_\mathrm{A}^{(3)}`$ to be positive, can be satisfied only if the contribution of the quadrupole term to the sum $`c_\mathrm{A}^{(2)}+c_\mathrm{A}^{(3)}`$ is dominating. Thus, we may conclude that the quadrupole term is of vital importance for the understanding the reasons behind the significant difference in magnitude between the photoinduced parts of the polar and azimuthal anchoring strengths.
## V Discussion and Conclusions
In this paper we have studied both theoretically and experimentally effects of photo-induced ordering in the azo-dye aligning films on the anchoring energy strengths. These effects are governed by dependence of the strengths on the azo-dye order parameter.
Our theoretical approach relies on the mean field theory Teixeira and Sluckin (1992a) and provides general expressions for the Landau-de Gennes surface free energy (13) and the anchoring energy (II.3). The theoretical results for the azimuthal and polar anchoring energy strengths are obtained under certain simplifying assumptions and used to interpret the experimental data relating the anchoring strengths and the dichroic ratio.
We found that linear fitting of the data for the azimuthal anchoring strength, though giving good results, predicts, in contradiction to the bare anchoring theory, the effect of vanishing anchoring which occurs at certain non-zero value of the dichroic ratio (and, thus, the irradiation dose). By using a simple phenomenological model we have shown that this effect can be attributed to an interplay between the light induced ordering in the bulk and the boundary conditions at the surface of the film which may counteract the action of light. Thus, the bulk and surface values of the azo-dye order parameter are generally different.
For planar boundary conditions that favor the random in-plane alignment of the azo-dye molecules, our model predicts threshold behavior of the azimuthal anchoring strength that starts growing provided the bulk order parameter, $`s_\mathrm{A}`$, exceeds its critical value. When the boundary conditions are homeotropic, this is no longer the case and the azimuthal strength smoothly increases with the dichroic ratio.
The results presented in Figs. 8 and 9 demonstrate that the theoretical curves calculated for both types of the boundary conditions can fit the experimental data equally well. There are, however, differences concerning the polar anchoring strength. By contrast to the homeotropic boundary conditions, it never equals zero at the planar conditions. From the previously published results Chigrinov et al. (2003b); Murauski et al. (2005) it can be concluded that the homeotropic boundary conditions is unlikely to occur in the azo-dye films under consideration.
Our final remark concerns some of the simplifying assumptions taken in our theoretical analysis. A more sophisticated theory that goes beyond the Fowler approximation (10) is required to take into account surface adsorption phenomena. A self-consistent treatment of two order parameter tensors in the interfacial layer also remains a challenge. We hope that our results will stimulate further progress in the field.
###### Acknowledgements.
This research was partially supported by RGC grants HKUST6102/03E and HKUST6149/04E. A.D.K. is indebted to Professor T.J. Sluckin for numerous useful conversations and acknowledges partial support from INTAS under grant No. 03-51-5448. We are also grateful to Professor H.-S. Kwok for stimulating comments.
## Appendix A Irreducible tensors, order parameter and invariants
In this appendix we introduce notations and definitions used throughout the paper. In addition, we express the order parameter in terms of irreducible tensors and deduce a number of algebraic relations simplifying the derivation of the tensorial form of the intermolecular potential given in Appendix B.
The irreducible tensors, $`𝐓_m`$, with the azimuthal number $`m`$ ranged from $`2`$ to 2 can be defined as linear combinations of the following form Biedenharn and Louck (1981):
$$𝐓_m=\underset{\mu ,\nu =1}{\overset{1}{}}C_{\mu \nu m}^{\mathrm{1\hspace{0.17em}1\hspace{0.17em}2}}𝐞_\mu 𝐞_\nu ,$$
(46)
where $`C_{m_1m_2m}^{j_1j_2j}`$ is the Clebsch-Gordon (Wigner) coefficient and $`𝐞_{\pm 1}=(\widehat{𝐱}\pm i\widehat{𝐳})/\sqrt{2}`$, $`𝐞_0=\widehat{𝐳}`$ are the vectors of spherical basis ($`\widehat{𝐱}`$, $`\widehat{𝐲}`$ and $`\widehat{𝐳}`$ are the unit vectors directed along the corresponding coordinate axes). Substituting the values of the Wigner coefficient into Eq. (46) gives the expressions for $`𝐓_m`$
$`𝐓_0=(3𝐞_0𝐞_0𝐈)/\sqrt{6},`$ (47)
$`𝐓_{\pm 1}=(𝐞_0𝐞_{\pm 1}+𝐞_{\pm 1}𝐞_0)/\sqrt{2},`$ (48)
$`𝐓_{\pm 2}=𝐞_{\pm 1}𝐞_{\pm 1},`$ (49)
so that it is not difficult to verify the validity of the orthogonality relation
$$Tr[𝐓_m𝐓_n]=(1)^m\delta _{mn}$$
(50)
and the algebraic identities
$$\widehat{𝐳}𝐓_m\widehat{𝐳}=\sqrt{2/3}\delta _{m0},\widehat{𝐳}𝐓_m𝐓_n\widehat{𝐳}=c_{|m|}\delta _{mn},$$
(51)
where $`c_0=2/3`$, $`c_1=1/2`$ and $`c_2=0`$.
Under the action of rotation the vectors of spherical basis transform as follows
$`𝐞_\mu 𝐞_\mu (\widehat{𝐮})={\displaystyle \underset{\nu =1}{\overset{1}{}}}D_{\nu \mu }^1(\widehat{𝐮})𝐞_\nu ,`$ (52)
$`𝐞_0(\widehat{𝐮})=\widehat{𝐮},`$ (53)
$`𝐞_{\pm 1}(\widehat{𝐮})=(𝐞_x(\widehat{𝐮})\pm i𝐞_y(\widehat{𝐮}))/\sqrt{2},`$ (54)
where $`D_{nm}^j(\widehat{𝐮})D_{nm}^j(\theta ,\varphi )`$ is the Wigner $`D`$ function Biedenharn and Louck (1981); Gelfand et al. (1963); $`\theta `$ and $`\varphi `$ are Euler angles of the unit vector $`\widehat{𝐮}`$; $`𝐞_x(\widehat{𝐮})=(\mathrm{cos}\theta \mathrm{cos}\varphi ,\mathrm{cos}\theta \mathrm{sin}\varphi ,\mathrm{sin}\theta )`$, $`𝐞_y(\widehat{𝐮})=(\mathrm{sin}\varphi ,\mathrm{cos}\varphi ,0)`$.
The definition (46) implies that transformation properties of the tensors $`𝐓_m`$ under rotations are determined by the irreducible representation of the rotation group with $`j=2`$, where $`j`$ is the angular momentum number. So, we have
$$𝐓_m𝐓_m(\widehat{𝐮})=\underset{\mu ,\nu =1}{\overset{1}{}}C_{\mu \nu m}^{\mathrm{1\hspace{0.17em}1\hspace{0.17em}2}}𝐞_\mu (\widehat{𝐮})𝐞_\nu (\widehat{𝐮})=\underset{k=2}{\overset{2}{}}D_{km}^2(\widehat{𝐮})𝐓_k.$$
(55)
Eq. (55) can now be combined with the relations (53) and (47) to yield the expression for the order parameter tensor $`𝐐(\widehat{𝐮})`$:
$$𝐐(\widehat{𝐮})=\sqrt{3/2}𝐓_0(\widehat{𝐮})=(3\widehat{𝐮}\widehat{𝐮}𝐈)/2,$$
(56)
where the unit vector $`\widehat{𝐮}`$ is directed along the long molecular axis.
The director $`\widehat{𝐧}`$ is defined as an eigenvector of the orientationally averaged order parameter tensor
$$𝐐(\widehat{𝐮})_{\widehat{𝐮}}=\sqrt{3/2}\underset{k=2}{\overset{2}{}}D_{k0}^2(\varphi ^{},\theta ^{})𝐓_k(\widehat{𝐧}),$$
(57)
where $`\theta ^{}`$ and $`\varphi ^{}`$ are Euler angles of the vector $`\widehat{𝐮}`$ related to the basis vectors $`𝐞_i(\widehat{𝐧})`$.
Since $`\widehat{𝐧}`$ is the director, the averages $`D_{\pm 10}^2(\varphi ^{},\theta ^{})_{\varphi ^{},\theta ^{}}`$ vanish. Other averages
$`D_{00}^2(\varphi ^{},\theta ^{})=\sqrt{2/3}S,`$ (58)
$`D_{\pm 20}^2(\varphi ^{},\theta ^{})=P\mathrm{exp}(\pm 2i\gamma )/\sqrt{6}`$ (59)
are proportional to the scalar order parameter $`S`$ and the biaxiality parameter $`P`$.
By using the orthogonality conditions (50) and Eqs. (57)-(59) we recover the relations in the traditional form de Gennes and Prost (1993):
$`𝐐(\widehat{𝐮})_{\widehat{𝐮}}𝐒(\widehat{𝐧})=S𝐐(\widehat{𝐧})+P(\widehat{𝐦}\widehat{𝐦}\widehat{𝐥}\widehat{𝐥})/2,`$ (60)
$`S=3\left(\widehat{𝐮}\widehat{𝐧}\right)^21/2,`$ (61)
$`P=3\left(\widehat{𝐮}\widehat{𝐦}\right)^2\left(\widehat{𝐮}\widehat{𝐥}\right)^2/2,`$ (62)
where $`\widehat{𝐦}=\mathrm{cos}\gamma 𝐞_x(\widehat{𝐧})\mathrm{sin}\gamma 𝐞_y(\widehat{𝐧})`$ and $`\widehat{𝐥}=\mathrm{sin}\gamma 𝐞_x(\widehat{𝐧})+\mathrm{cos}\gamma 𝐞_y(\widehat{𝐧})`$.
The order parameter (60) is a traceless symmetric tensor. Therefore, there are two non-vanishing independent invariants
$`I_2=Tr[𝐒^2(\widehat{𝐧})]=\left(3S^2+P^2\right)/2,`$ (63)
$`I_3=Tr[𝐒^3(\widehat{𝐧})]=3S\left(S^2P^2\right)/4,`$ (64)
which enter the non-elastic part of the well known phenomenological expression for the Landau-de Gennes free energy density
$`f_{\mathrm{LG}}={\displaystyle \frac{2a}{3}}(TT^{})I_2{\displaystyle \frac{4B}{3}}I_3+{\displaystyle \frac{4C}{9}}I_2^2,`$ (65)
where $`T`$ is the temperature and $`T^{}`$ is the supercooling temperature.
For the scalar order parameters (61) and (62) combined into a pair $`(S,P)`$, it is convenient to introduce what might be called the “polar” (or amplitude-angle) representation
$`S=s\mathrm{cos}\psi ,P=\sqrt{3}s\mathrm{sin}\psi ,`$ (66)
where $`s^2=2I_2/3`$. Using the representation (66) the free energy density (65) can be recast into the form
$`f_{\mathrm{LG}}(s,\psi )=a(TT^{})s^2Bs^3\mathrm{cos}(3\psi )+Cs^4B^2C^1U_{\mathrm{LG}},`$ (67)
$`U_{\mathrm{LG}}(\eta ,\psi )={\displaystyle \frac{8+t}{32}}\eta ^2\eta ^3\mathrm{cos}(3\psi )+\eta ^4,\eta {\displaystyle \frac{C}{B}}s,`$ (68)
where $`t=32aCB^2(TT_c)`$ is the dimensionless temperature parameter and $`T_c=T^{}+B^2/(4aC)`$ is the temperature of the bulk nematic-isotropic transition. The rescaled density (68) is a generalized version of the dimensionless free energy density previously used in Refs. Sheng et al. (1992); Qian and Sheng (1997).
Finally, we write down the components of the order parameter tensor (61) in the polar representation
$`S_{ij}=s\left[n_in_j\mathrm{cos}\psi +m_im_j\mathrm{cos}(\psi 2\pi /3)+l_il_j\mathrm{cos}(\psi +2\pi /3)\right]`$ (69)
and notice that the stationary points of the free energy density (67) where the angle $`\psi `$ is a multiple of $`\pi /3`$ represent uniaxially anisotropic states. The latter immediately recovers the well known result about uniaxial anisotropy of NLC equilibrium states Patashinskii and Pokrovskii (1982).
## Appendix B Tensorial form of intermolecular potential
We begin with the intermolecular potential between two rigid, axially symmetric molecules expanded in a series of spherical harmonics as follows Pople (1954)
$`U(𝐫,\widehat{𝐮}_1,\widehat{𝐮}_2)={\displaystyle \underset{j_1,j_2,j}{}}u_{j_1j_2j}(r){\displaystyle \underset{m_1,m_2,m}{}}C_{m_1m_2m}^{j_1j_2j}`$
$`\times Y_{j_1m_1}(\widehat{𝐮}_1)Y_{j_2m_2}(\widehat{𝐮}_2)Y_{jm}^{}(\widehat{𝐫})`$ (70)
where $`𝐫𝐫_{12}=𝐫_1𝐫_2`$, $`𝐫_i`$ and $`\widehat{𝐮}_i`$ are the position and orientation (equivalently, Euler angles of the long molecular axis) coordinates of the interacting molecules, respectively; $`Y_{jm}(\widehat{𝐮})=\sqrt{(2j+1)/(4\pi )}D_{m0}^j(\widehat{𝐮})`$ is the spherical function Biedenharn and Louck (1981); Abramowitz and Stegun (1972). The form of the expansion (B) implies that the potential is invariant under translations, $`𝐫_i𝐫_i+\mathrm{\Delta }𝐫`$, and rotations, $`\{𝐫_i,\widehat{𝐮}_i\}\{R𝐫_i,R\widehat{𝐮}_i\}`$. In addition, we shall assume the head-tail symmetry
$$U(𝐫,\widehat{𝐮}_1,\widehat{𝐮}_2)=U(𝐫,\widehat{𝐮}_1,\widehat{𝐮}_2)=U(𝐫,\widehat{𝐮}_1,\widehat{𝐮}_2),$$
(71)
so that the outer sum in Eq. (B) is restricted to run over even values of $`j_1`$ and $`j_2`$.
It is now our task to link the pairwise potential integrated over in-plane coordinates
$`V(z,\widehat{𝐮}_1,\widehat{𝐮}_2)={\displaystyle _S}U(𝐫,\widehat{𝐮}_1,\widehat{𝐮}_2)dxdy`$
$`={\displaystyle \underset{j_1,j_2,j}{}}[(2j_1+1)(2j_2+1)(2j+1)/(4\pi )^3]^{1/2}`$
$`\times v_{j_1j_2j}(z){\displaystyle \underset{m}{}}C_{mm\mathrm{\hspace{0.17em}0}}^{j_1j_2j}D_{m0}^{j_1}(\widehat{𝐮}_1)D_{m0}^{j_2}(\widehat{𝐮}_2)D_{00}^j(\widehat{𝐳})`$ (72)
and the tensorial representation that was originally suggested by Ronis and Rosenblatt in Ref. Ronis and Rosenblatt (1980)
$`V(z,\widehat{𝐮}_1,\widehat{𝐮}_2)V(z,𝐐_1,𝐐_2)=V_{\mathrm{iso}}(z)+\beta ^{(0)}(z)\widehat{𝐳}(𝐐_1+𝐐_2)\widehat{𝐳}`$
$`+\beta ^{(1)}(z)Tr(𝐐_1𝐐_2)+\beta ^{(2)}(z)\widehat{𝐳}𝐐_1𝐐_2\widehat{𝐳}`$
$`+\beta ^{(3)}(z)[\widehat{𝐳}𝐐_1\widehat{𝐳}][\widehat{𝐳}𝐐_2\widehat{𝐳}],`$ (73)
where $`𝐐_i𝐐(\widehat{𝐮}_i)`$ is defined by Eq. (56). In other words, the problem is to express the coefficients $`\beta ^{(i)}(z)`$ in terms of the harmonics $`v_{j_1j_2j}(z)`$. To this end we restrict ourselves to the lowest order harmonics of the expansion (B) with $`j_i<4`$ and consider the equation
$`x_1[\widehat{𝐳}𝐓_0(\widehat{𝐮}_1)\widehat{𝐳}][\widehat{𝐳}𝐓_0(\widehat{𝐮}_2)\widehat{𝐳}]+x_2Tr(𝐓_0(\widehat{𝐮}_1)𝐓_0(\widehat{𝐮}_2))`$
$`+x_3\widehat{𝐳}𝐓_0(\widehat{𝐮}_1)𝐓_0(\widehat{𝐮}_2)\widehat{𝐳}={\displaystyle \underset{m=0}{\overset{2}{}}}\alpha _mD_{m0}^2(\widehat{𝐮}_1)D_{m0}^2(\widehat{𝐮}_2),`$ (74)
that need to be solved for $`x_1`$, $`x_2`$ and $`x_3`$. The sum on the right hand side of Eq. (B) represents sum of the harmonics with $`j_1=j_2=2`$. The case of the harmonics with $`j_1j_2=0`$ is much easier to treat as we only have to use the relation
$$\widehat{𝐳}𝐓_0(\widehat{𝐮})\widehat{𝐳}=D_{00}^2(\widehat{𝐮}).$$
(75)
By using Eq. (55) combined with the relations (50) and (51) it is straightforward to transform Eq. (B) into a system of linear equations. The solution of the system is given by
$`x_1=(3\alpha _0+4\alpha _1+\alpha _2)/2,`$
$`x_2=2(\alpha _1+\alpha _2),`$
$`x_3=\alpha _2.`$ (76)
Given the values of the Wigner coefficients, we can now use the relations (75) and (B) to derive the final result in the following form:
$$V_{\mathrm{iso}}(z)=v_{000}(z)/(4\pi )^{3/2},$$
(77)
$$[(4\pi )^{3/2}/5]\beta ^{(0)}(z)=v_{202}(z)=v_{022}(z),$$
(78)
$`[(4\pi )^{3/2}/5]\beta ^{(1)}(z)=2/(3\sqrt{5})\{v_{220}(z)`$
$`+10/\sqrt{14}v_{222}(z)+3/\sqrt{14}v_{224}(z)\},`$ (79)
$$[(4\pi )^{3/2}/5]\beta ^{(2)}(z)=20/\sqrt{70}\left\{v_{222}(z)+v_{224}(z)\right\},$$
(80)
$$[(4\pi )^{3/2}/5]\beta ^{(3)}(z)=35/\sqrt{70}v_{224}(z).$$
(81)
The formulas (77)–(81) relate the parameters of the representation (B) to the coefficient functions in the spherical harmonics expansion (B). The terms $`v_{202}`$ and $`v_{022}`$ describe the coupling between orientation of the molecules and the intermolecular vector, whereas $`v_{220}`$ and $`v_{224}`$ are known as the Maier-Saupe and the quadrupole terms, respectively. For the interaction between NLC molecules, the functions $`\beta ^{(1)}(z)`$ and $`\beta ^{(2)}(z)`$ define the elastic coefficients of NLC that must be positive. This stability condition implies that $`\beta ^{(1)}`$ and $`\beta ^{(2)}`$ are both predominately non-positive Sen and Sullivan (1987); Tjipto-Margo and Sullivan (1988); Teixeira and Sluckin (1992a); Osipov and Hess (1993). |
warning/0507/physics0507060.html | ar5iv | text | # Influence of optical aberrations in an atomic gyroscope
## 1 Introduction
Since the pioneering demonstrations of interferometry with de Broglie atomic waves using resonant light BordePTB ; Bragg and nanofabricated structures Prit91 as atomic beam splitters, a number of new applications have been explored, including measurements of atomic and molecular properties, fundamental tests of quantum mechanics, and studies of various inertial effects Berman . Using atom interferometers as inertial sensors is also of interest for geophysics, tests of general relativity Hyper , and inertial guidance systems.
Atom interferometers based on light-induced beam splitters have already demonstrated considerable sensitivity to inertial forces. Sequences of optical pulses generate the atom optical elements (e.g., mirrors and beam splitters) for the coherent manipulation of the atomic wave packets Borde91 . The sensitivity and accuracy of light-pulse atom interferometer gyroscopes Gustavson00 , gravimeters Peters99 and gravity gradiometers Snadden98 compare favorably with the performances of state-of-the-art instruments. Furthermore, this type of interferometer is likely to lead to a more precise direct determination of the fundamental constant $`\alpha `$ from the measurement of $`\mathrm{}/M`$ wicht2001 . In the case of rotation measurements, the sensitivity reaches that of the best laboratory ring laser gyroscope Stedman . Indeed the Sagnac phase shift, proportional to the total energy of the interfering particle, is much larger for atoms than for photons. This compensates for the smaller interferometer area and the lower flux.
In this paper we focus on the effect of the fluctuations of the atomic trajectory, which might affect the long term stability of atomic gyroscopes when coupled with local phase variations induced by optical aberrations. We will introduce this problem in paragraph 2 and illustrate it quantitatively in the case of our setup in paragraph 3.
Our experiment consists in an almost complete inertial measurement unit Yver03 , using cold Cesium atoms that enable for a drastic reduction of the apparatus dimensions while reaching a sensitivity of $`30`$ nrad.s$`^1.`$Hz<sup>-1/2</sup> to rotation and 4x10<sup>-8</sup> m.s$`^2.`$Hz<sup>-1/2</sup> to acceleration. Its operation is based on recently developed atom interference and laser manipulation techniques. Two interferometers with counter-propagating atomic beams discriminate between rotation and acceleration Gustavson98 . Thanks to the use of a single pair of counter-propagating Raman laser beams, our design is intrinsically immune to uncorrelated vibrations between the three beam splitters, usually limiting such devices. This configuration is made possible by the use of a reduced launch velocity, inducing a reasonable interaction time between the pulses. However, as any atomic gyroscope, our sensor’s scheme remains sensitive to local phase variations, a limitation that has already been encountered in optical atomic clocks Trebst01 .
## 2 Principle
We first briefly review the basic light-pulse method in the case of a symmetric Ramsey-Bordé interferometer scheme Borde02 , where three travelling-wave pulses of light resonantly couple two long-lived electronic states.The two-photon stimulated Raman transitions between ground state hyperfine levels are driven by two lasers with opposite propagation vectors $`𝐤_e`$ and $`𝐤_g`$ ($`𝐤_e𝐤_g`$). First, at $`t=t_1`$ a beam splitting pulse puts the atom into a coherent superposition of its two internal states. Because of conservation of momentum during the atom-light interaction, this pulse introduces a relative momentum $`\mathrm{}𝐤_{\mathrm{eff}}=\mathrm{}𝐤_g\mathrm{}𝐤_e`$ between the atomic wave packets corresponding to each state. These wave packets drift apart for a time $`T`$, after which a mirror pulse is applied at $`t_2=t_1+T`$ to redirect the two wave packets. After another interval of duration $`T`$, the wave packets physically overlap, and a final beam splitting pulse recombines them at $`t_3=t_1+2T`$. The measurement of the probabilities of presence in both internal states at the interferometer output leads to the determination of the difference of accumulated phases along the two paths. In general, atoms are launched with a velocity $`𝐯`$ so that each stimulated Raman transition occurs at a particular position $`\{x_i,y_i,z_i\}_{i=1,2,3}`$ that can be evaluated from the classical trajectories associated with the atomic wave packets Antoine03 , as shown fig. 1. In our setup, Raman laser beams propagate in the *(Ox)* direction and atoms are launched in the $`(y,z)`$ plane. We define $`𝐮_i`$=$`\{y_i,z_i\}`$ the atomic cloud positions in this plane at time $`t_i`$.
In the absence of any external forces, atoms initially prepared in a particular state ($`{}_{}{}^{6}S_{1/2}^{},F=3,m_F=0`$ in the present setup) will return to this state with unit probability. A uniform external acceleration or rotation induces a relative phase shift between the interfering paths. This phase shift modifies the transition probability between the two Cesium internal states $`{}_{}{}^{6}S_{1/2}^{},F=3,m_F=0`$ and $`{}_{}{}^{6}S_{1/2}^{},F=4,m_F=0`$ (noted $`|3`$ and $`|4`$ in the following). Hence the transition probability measurement leads to the determination of the phase shift and finally the evaluation of the perturbing forces.
It can be shown that the only contribution to the phase shift results from the interaction with the laser light fields Antoine03 . In the limit of short, intense pulses, the atomic phase shift associated with a transition $`|3|4`$ (resp. $`|4|3`$) is +$`\varphi _i`$ (resp. -$`\varphi _i`$), where $`\varphi _i`$ is the phase difference between the two Raman laser beams. We then find that the transition probability from $`|3`$ to $`|4`$ at the exit of the interferometer is simply $`\frac{1}{2}\left[1\mathrm{cos}(\mathrm{\Delta }\varphi )\right]`$ where $`\mathrm{\Delta }\varphi =\varphi _12\varphi _2+\varphi _3`$. The three quantities correspond to the phase imparted to the atoms by the initial beam splitting pulse, the mirror pulse, and the recombining pulse where $`\varphi _i=\varphi _g(u_i,t_i)\varphi _e(u_i,t_i)=k_{\mathrm{eff}}x_i+\mathrm{\Phi }(𝐮_i)`$. The sensitivity to rotation and acceleration arises from the first term $`k_{\mathrm{eff}}x_i`$ and simplifies to $`\mathrm{\Delta }\varphi _{\mathrm{acc}}=a_xk_{\mathrm{eff}}T^2`$ and $`\mathrm{\Delta }\varphi _{\mathrm{rot}}=2k_{\mathrm{eff}}v_y\mathrm{\Omega }_zT^2`$ for the present setup. The phase $`\mathrm{\Phi }(𝐮_i)`$ for the pulse at time $`t_i`$ corresponds to the local phase in the $`(y,z)`$ plane due to wavefront distortions of both laser beams<sup>1</sup><sup>1</sup>1The interferometer is also sensitive to time fluctuations of the Raman laser phases Yver03 . These fluctuations are identical for the two interferometers and disappear from the rotation signal. They will be neglected in this paper.. It induces a residual phase error at the exit of the interferometer $`\delta \mathrm{\Phi }=\mathrm{\Phi }(𝐮_1)2\mathrm{\Phi }(𝐮_2)+\mathrm{\Phi }(𝐮_3)`$.
Acceleration cannot be discriminated from rotation in a single atomic beam sensor, as stated above. This limitation can be circumvented by installing a second, counter-propagating, cold atomic beam (fig. 1) Gustavson98 . When both atomic beams perfectly overlap, the area vectors for the resulting interferometer loops have opposite directions. The corresponding rotational phase shifts $`\mathrm{\Delta }\varphi _{\mathrm{rot}}`$ have opposite signs while the acceleration phase shifts $`\mathrm{\Delta }\varphi _{\mathrm{acc}}`$ are identical. Consequently, acceleration is calculated by summing the two interferometer’s phase shifts: $`\mathrm{\Delta }\varphi _+2\mathrm{\Delta }\varphi _{\mathrm{acc}}`$; while taking the difference rejects the contribution of uniform accelerations so that $`\mathrm{\Delta }\varphi _{}2\mathrm{\Delta }\varphi _{\mathrm{rot}}`$. In addition, the residual phase error $`\delta \mathrm{\Phi }`$ vanishes in $`\mathrm{\Delta }\varphi _{}`$, but remains in $`\mathrm{\Delta }\varphi _+`$ as an absolute phase bias $`2\times \delta \mathrm{\Phi }`$.
However, an imperfect overlapping of the two counter-propagating wavepackets trajectories might lead to an imperfect common mode rejection of the residual phase error in $`\mathrm{\Delta }\varphi _{}`$. Thus, a phase bias $`\delta \mathrm{\Phi }_{}=\delta \mathrm{\Phi }^L\delta \mathrm{\Phi }^R`$ will appear, where the notations <sup>L</sup> and <sup>R</sup> concern the left and right atom interferometers. While the phase bias $`\delta \mathrm{\Phi }_+2\times \delta \mathrm{\Phi }`$ depends on the local value of the phase at the average position $`𝐫_i=\frac{𝐮_i^L+𝐮_i^R}{2}`$, the phase bias $`\delta \mathrm{\Phi }_{}`$ depends on the local phase gradient at the average position $`𝐫_i`$ with the position offset $`\delta 𝐫_i=𝐮_i^L𝐮_i^R`$:
$$\begin{array}{ccc}\delta \mathrm{\Phi }_{}& =& \mathrm{\Phi }(𝐫_1)\delta 𝐫_12\mathrm{\Phi }(𝐫_2)\delta 𝐫_2\\ & & +\mathrm{\Phi }(𝐫_3)\delta 𝐫_3.\end{array}$$
(1)
Equation 1 shows that uncorrelated fluctuations of the wavepackets trajectories from shot to shot causes fluctuations of the phase bias, which amplitude depends on the local wavefront slope of the phase. If we consider a perfect control of the launch velocity<sup>2</sup><sup>2</sup>2We can reach a stability of $`10^4`$m.s<sup>-1</sup> or better from shot to shot thanks to the moving molasses technique molasses ., fluctuations of trajectories are only due to fluctuations of the initial positions of the atomic clouds. Consequently, we can consider $`\delta 𝐫_1=\delta 𝐫_2=\delta 𝐫_3`$. The phase fluctuation is then simply proportional to the product of the fluctuations of the cloud initial position $`(y_0,z_0)`$ with the phase gradients $`\mathrm{\Delta }\mathrm{\Phi }_i`$. As the phase gradients are time-independent, the Allan variance of the phase $`\sigma _{\delta \mathrm{\Phi }_{}}^2`$ is simply:
$$\begin{array}{ccc}\sigma _{\delta \mathrm{\Phi }_{}}^2=\sigma _{y_0}^2.\left[_y\left(\mathrm{\Phi }\left(𝐫_1\right)2\mathrm{\Phi }\left(𝐫_2\right)+\mathrm{\Phi }\left(𝐫_3\right)\right)\right]^2& & \\ +\sigma _{z_0}^2.\left[_z\left(\mathrm{\Phi }\left(𝐫_1\right)2\mathrm{\Phi }\left(𝐫_2\right)+\mathrm{\Phi }\left(𝐫_3\right)\right)\right]^2& & \end{array}$$
(2)
where $`\sigma _{y_0}^2`$ and $`\sigma _{z_0}^2`$ are the Allan variances of the initial horizontal and vertical positions. Eq. 2 shows that the fluctuations of the clouds initial positions, as well as the wavefront quality of the Raman beams, have to be systematically investigated in atomic gyroscopes in order to estimate how it affects its performances.
## 3 Experimental results
In our setup, the atomic sources are clouds of Cesium atoms, cooled in magneto-optical traps and launched with a parabolic flight (fig. 2). As the initial angle reaches $`82^o`$, and the launch velocity $`2.4`$m.s<sup>-1</sup>, the horizontal velocity $`v_y`$ is $`0.3`$m.s<sup>-1</sup>.
The single pair of Raman laser beams propagates along the x-axis and is switched on three times at the top of the atomic trajectories. If the three pulses are symmetric with respect to the trajectory apogees, the interferometer oriented enclosed areas are equivalent to their flat horizontal projections: the oriented vertical projection is naught. The time delay between pulses is typically $`45`$ ms. The positions of the atoms during the three Raman pulses are given in fig. 2.
In order to investigate the fluctuations of the atomic initial positions from shot to shot, we image one of the two clouds. The cycling sequence takes about 1.3 s and consists on a trap phase of 500 ms, a molasses phase of 20 ms, a launching phase of 2 ms and a waiting time phase of 800 ms needed to process the image: download of the image, subtraction of a background image and determination of the cloud barycenter position in y- and z-axes. The image is taken just after turning off the trap magnetic field, at the end of the molasses phase. We calculate the Allan standard deviations Allan of the barycenter horizontal and vertical positions (fig. 3) from a one hour acquisition. Two peaks, appearing after 10 s and 150 s of integration time, are characteristic of fluctuations of periods equal to 20 s and 300 s. After about 10 min integration (630 s), the position standard deviations reach 10 $`\mu `$m and 5 $`\mu `$m in the horizontal and vertical directions respectively. This dissymmetry is consistent with the magnetic field gradient configuration, which is twice higher on the Z-direction.
The long-term variations are due to fluctuations of the MOT cooling lasers intensity ratio, which Allan standard deviation is plotted in fig. 3. We see again the oscillation of period 300 s, appearing for 150 s integration time. We analyze this as the period of the air conditioning, creating temperature variations on the fibre splitters delivering the cooling lasers.
This result has to be coupled to the optical aberrations of the Raman lasers. The main contribution to these aberrations comes from the vacuum windows used for the Raman laser beams, which clear diameter is 46 mm. They have been measured with a Zygo wavefront analyzer, which gives the laser phase distortion created by the windows. This distortion is projected on the Zernike polynomial base Zer . As our atomic clouds are about 2 mm wide, the decomposition is pertinent only up to the 36th polynomial. Indeed, the upper numbers correspond to high spatial frequencies, so that their effect will be smoothed by averaging on the atomic cloud dimensions. To reduce the stress on the vacuum windows, essentially due to the mounting, they were glued in place. Thanks to this method, the wavefront quality reaches $`\lambda /50`$ rms over the whole clear diameter of 42 mm.
The wavefront measurement allows for evaluation of the atomic phase shift fluctuations due to the coupling between aberrations and position fluctuations using eq.2 assuming that the two sources are uncorrelated. Their relative position fluctuations are $`\sqrt{2}`$ times greater than these observed for one source. The contribution of this phase fluctuations to the Allan standard deviation of the rotation rate measurement is shown in fig. 4. We compare it with the ultimate stability of our gyroscope, given by the quantum projection noise. It is estimated to $`30/\sqrt{\tau }`$nrad.s<sup>-1</sup> ($`\tau `$ is the integration time) from the ultimate signal-to-noise ratio obtainable with $`10^6`$ atoms.
The rotation noise induced by position fluctuations has a significant contribution for integration times larger than 100 s. At the present stage of the experiment, this limitation is due to the high temperature sensitivity of the fibre splitters. This could be the main limitation of the gyroscope performances.
## 4 Conclusion
In the present paper we studied the stability of a cold atom gyroscope based on two symmetrical Ramsey-Bordé interferometers, with respect to optical phase inhomogeneity. Instability due to aberrations is not a specific problem induced by Raman transitions, but concerns every type of atom interferometer using light beam splitters. We showed that the coupling between wavefront distortions of these lasers and fluctuations of the atomic trajectory becomes predominant at long term, despite a wavefront quality of $`\lambda /50`$ rms obtained thanks to glued windows. In our setup, atomic trajectory fluctuations are mainly due to fluctuations of the intensity ratio of the MOT cooling lasers, induced by the fibre splitters used for their generation.
However several improvements may render their contribution negligible:
\- reduce the atomic trajectory fluctuations, by using discrete optical couplers for the MOT instead of the present fibre splitters,
\- minimize the number of optics which contribute to the interferometer instability. This can be done by including the Raman laser beam imposition optics in the vacuum chamber, in order to remove the aberrations due to the vacuum windows, or by minimizing the number of non-common optics for the two Raman lasers, since only the phase difference between the lasers is imprinted on the atomic phase shift.
Such techniques open large improvement possibilities, which will be confirmed directly on the long-term stability measurement of the atomic signal in our interferometer setup.
## 5 Acknowledgements
The authors would like to thank DGA, SAGEM and CNES for supporting this work, Pierre Petit for the early stage of the experiment and Christian Bordé for helpful discussions. They also thank Thierry Avignon and Lionel Jacubowiez from SupOptique for their help in the wavefront measurement. |
warning/0507/astro-ph0507079.html | ar5iv | text | # Planetary Detection Efficiency of the Magnification 3000 Microlensing Event OGLE-2004-BLG-343
## 1 Introduction
Current microlensing planet searches focus significant effort on high-magnification events, which have great promise for detecting low-mass extrasolar planets. It is therefore crucial to understand the potential for discovering planets and to optimize the early identifications and observational strategy of such events. In previous planetary detection efficiency analyses of high-magnification events, finite-source effects have often been ignored mainly due to computational limitations. However, such effects are intrinsically important in these events because the sources are more likely to be resolved at very low impact parameters. In this study, we improve the method of Yoo et al. (2004b) by incorporating finite-source effects to characterize the planetary detection efficiency of the extremely high-magnification event OGLE-2004-BLG-343, and we develop new efficient algorithms to make the calculations possible. Moreover, we attempt to find useful observational signatures of high-magnification events so as to help alleviate the difficulties in their early recognition.
### 1.1 High-Magnification Events & Earth-Mass Planets
Apart from pulsar timing (Wolszczan & Frail, 1992), microlensing is at present one of the few planet-finding techniques that is sensitive to Earth-mass planets. A planetary companion of an otherwise isolated lens star introduces two kinds of caustics into the magnification pattern: “planetary caustics” associated with the planet itself and a “central caustic” associated with the primary lens. <sup>1</sup><sup>1</sup>1When the planetary companion is close to the Einstein ring, the planetary and central caustics merge into a single “resonant caustic”. When the source passes over or close to one of these caustics, the light curve deviates from its standard Paczyński (1986) form, thus revealing the presence of the planet (Mao & Paczyński, 1991; Liebes, 1964).
Since planetary caustics are generally far larger than central caustics, a “fair sample” of planetary microlensing events would be completely dominated by planetary-caustic events. Nevertheless, central caustics play a crucial role in current microlensing planet searches, particularly for Earth-mass planets (Griest & Safizadeh, 1998), for the simple reason that it is possible to predict in advance that the source of a given event will arrive close to the center of the magnification pattern where it will probe for the presence of these caustics. Hence, one can organize the intensive observations required to characterize the resulting anomalies. By contrast, the perturbations due to planetary caustics occur without any warning. The lower the mass of the planet, the shorter the duration of the anomaly, and so the more crucial is the warning to intensify the observations. This is the primary reason that planet-searching groups give high priority to high-magnification events, i.e., those that probe the central caustics <sup>2</sup><sup>2</sup>2Note that although high-magnification events are guaranteed to have low impact parameters, peak magnification for events with low impact parameters are not necessarily high if they have relatively large source sizes. And large sources will tend to smear out the perturbations induced by the central caustics, thereby decreasing the planetary sensitivity (Griest & Safizadeh, 1998; Chung et al., 2005). So when central caustics are important in producing planetary signals, the maximum magnification of microlensing events serves as a better indicator of planetary detection efficiency than the impact parameter. . As a bonus, high-magnification events are also more sensitive to planetary-caustic perturbations than are typical events (Gould & Loeb, 1992) since their larger images increase the chances that they will be perturbed by planets. However, this enhancement is relatively modest compared to the rich potential of central-caustic crossings.
In principle, it is also possible to search for Earth-mass planets from perturbations due to their larger (and so more common) planetary and resonant caustics, but this would require a very different strategy from those currently being carried out. The problem is that these perturbations occur without warning during otherwise normal microlensing events, and typically last only 1 or 2 hr. Hence, one would have to intensively monitor the entire duration of many events. The only way to do this practically is to intensively monitor an entire field containing many ongoing microlensing events roughly once every 10 minutes in order to detect and properly characterize the planetary deviations. Proposals to make such searches have been advanced for both space-based (Bennett & Rhie, 2002) and ground-based (Sackett, 1997) platforms.
At present, two other microlensing planet-search strategies are being pursued. Both strategies make use of wide-area $`(10\mathrm{deg}^2)`$ searches for microlensing events toward the Galactic bulge. Observations are made once or a few times per night by the OGLE-III<sup>3</sup><sup>3</sup>3 OGLE Early Warning System: http://www.astrouw.edu.pl/~ogle/ogle3/ews/ews.html (Udalski, 2003) and MOA<sup>4</sup><sup>4</sup>4MOA Transient Alert Page: http://www.massey.ac.nz/~iabond/alert/alert.html (Bond et al., 2001) surveys. When events are identified, they are posted as “alerts” on their respective Web sites. In the first approach, these groups check each ongoing event after each observation for signs of anomalous behavior, and if their instantaneous analysis indicates that it is worth doing so, they switch from survey mode to follow-up mode. This approach led to the first reliable detection of a planetary microlensing event, OGLE 2003-BLG-235/MOA-2003-BLG-53 (Bond et al., 2004).
In the second approach, follow-up groups such as the Probing Lensing Anomalies NETwork (PLANET; Albrow et al. 1998) and the Microlensing Follow-Up Network ($`\mu `$FUN; Yoo et al. 2004b) monitor a subset of alerted events many times per day and from locations around the globe. Generally these groups focus to the extent possible on high-magnification events for the reasons stated above. The survey groups can also switch from “survey mode” to “follow-up” mode to probe newly emerging high-magnification events.
Over the past decade several high-magnification events have been analyzed for planets. Gaudi et al. (2002) and Albrow et al. (2001) placed upper limits on the frequency of planets from the analysis of 43 microlensing events, three of which reached magnifications $`A_{\mathrm{max}}100`$, including OGLE-1998-BUL-15 ($`A_{\mathrm{max}}=170\pm 30`$), MACHO-1998-BLG-35 ($`A_{\mathrm{max}}=100\pm 5`$), and OGLE-1999-BUL-35 ($`A_{\mathrm{max}}=125\pm 15`$). However, the first of these was not monitored over its peak. MACHO-1998-BLG-35 was also analyzed by Rhie et al. (2000) and Bond et al. (2002), who incorporated all available data and found modest $`(\mathrm{\Delta }\chi ^2=63)`$ evidence for one, or perhaps two, Earth-mass planets.
Yoo et al. (2004b) analyzed OGLE-2003-BLG-423 ($`A_{\mathrm{max}}=256\pm 43`$), which at the time was the highest magnification event yet recorded. However, because this event was covered only intermittently over the peak, it proved less sensitive to planets than either MACHO-1998-BLG-35 or OGLE-1999-BUL-35.
Abe et al. (2004) analyzed MOA-2003-BLG-32/OGLE-2003-BLG-219, which at $`A_{\mathrm{max}}=525\pm 75`$, is the current record-holder for maximum magnification. Unlike OGLE-2003-BLG-423, this event was monitored intensively over the peak: the Wise Observatory in Israel was able to cover the entire 2.5 hr FWHM during the very brief interval that the bulge is visible from this northern site. The result is that this event has the best sensitivity to low-mass planets to date.
Recently, Udalski et al. (2005) detected a $`3`$-Jupiter mass planet by intensively monitoring the peak of the high-magnification event OGLE-2005-BLG-071. This was the second robust detection of a planet by microlensing and the first from perturbations due to a central caustic.
### 1.2 Planet Detection Efficiencies: Philosophy and Methods
The fundamental aim of microlensing planet searches is to derive meaningful conclusions about the presence of planets (or lack thereof) from these searches. Therefore, it is essential to quantitatively assess what planets could have been detected from the observations of individual non-planetary events if such planets had been present. Actually, this problem is not as easy to properly formulate as it might first appear. For example, the event parameters are measured with only finite precision. Among these, the impact parameter $`u_0`$ (in units of the angular Einstein radius $`\theta _\mathrm{E}`$) is particularly important: if the event really did have a $`u_0`$ equal to its best-fit value, then one could calculate whether a planet at a certain separation and position angle would have given rise to a detectable signal in the observed light curve. But the true value of $`u_0`$ may differ from the best-fit value by, say, $`1\sigma `$, and the same planet may not give rise to a detectable signal for this other, quite plausible geometry. (In principle, an error in the time of maximum, $`t_0`$, would cause a similar problem, displacing the assumed path through the Einstein ring by $`\delta t_0/t_\mathrm{E}`$, where $`t_\mathrm{E}`$ is the Einstein crossing time. However, because $`u_0`$ is strongly correlated with several other parameters while $`t_0`$ is not, the error in $`u_0`$ is substantially larger than the error in $`t_0`$ divided by $`t_\mathrm{E}`$.) Or, as another example, consider finite-source effects. Planetary perturbations have a fairly high probability of exhibiting finite-source effects, which then have a substantial impact on whether the deviation can be detected in a given data stream. If there is such a planetary perturbation, one can measure $`\rho _{}=\theta _{}/\theta _\mathrm{E}`$, the size of the source relative to the Einstein radius. But if there is no planet detected, no finite-source effects are typically detected, and hence there is no direct information on $`\rho _{}`$. Therefore, one cannot reliably determine whether a given planetary perturbation would have been affected by finite-source effects and so whether it would have been detected. Finally, there are technical questions as to what exactly it means that a planet “would have been” detected.
The past decade of microlensing searches has been accompanied by a steady improvement in our understanding of these questions. Gaudi & Sackett (2000) developed the first method to evaluate detection efficiencies, which was later implemented by Albrow et al. (2000) and Gaudi et al. (2002). In this approach, binary models are fitted to the observed data with the three “binary parameters” $`(b,q,\alpha )`$ held fixed and the three “point-lens parameters” $`(t_0,u_0,t_\mathrm{E})`$ allowed to vary. Here $`b`$ is the planet-lens separation in units of $`\theta _\mathrm{E}`$, $`q`$ is the planet-star mass ratio, $`\alpha `$ is the angle of the source trajectory relative to the binary axis, $`t_0`$ is the time of the source’s closest approach to the center of the binary system, $`u_0=u(t_0)`$ is the impact parameter, $`t_E=\theta _\mathrm{E}/\mu `$ is the Einstein timescale, and $`\mu `$ is the source-lens relative proper motion. If a particular $`(b,q,\alpha )`$ yielded a $`\chi ^2`$ improvement $`\mathrm{\Delta }\chi ^2<\chi _{\mathrm{min}}^2=60`$, a planet could be said to be detected. If not, then the ensemble of $`(b,q,\alpha )`$ for which $`\mathrm{\Delta }\chi ^2>\chi _{\mathrm{min}}^2=60`$ was said to be excluded for that event. For each $`(b,q)`$, the fraction of angles $`0\alpha 2\pi `$ that was excluded was designated the “sensitivity” for that geometry.
Gaudi et al. (2002) argued that this method underestimated the sensitivity because it allowed the fit to move $`u_0`$ to values for which the source trajectory would “avoid” the planetary perturbation but still be consistent with the light curve. That is, $`u_0`$ has some definite value, even if it were not known to the modelers exactly what that value should be. Yoo et al. (2004b) followed up on this by holding $`u_0`$ fixed at a series of values and estimated planetary detection efficiency at each. The total efficiency would then be the average of these weighted by the probability of each value of $`u_0`$. In principle, one should also integrate over $`t_0`$ and $`t_\mathrm{E}`$. In practice, Yoo et al. (2004b) found that, at least for the event they analyzed, $`t_0`$ and $`t_{\mathrm{eff}}u_0t_\mathrm{E}`$ were determined very well by the data, so that once $`u_0`$ was fixed, so were $`t_0`$ and $`t_\mathrm{E}`$.
Yoo et al. (2004b) departed from all previous planet-sensitivity estimates by incorporating a Bayesian analysis that accounts for priors derived from a Galactic model of the mass, distance, and velocity properties of source and lens population into the analysis. They simulated an ensemble of events and weighted each by both the prior probability of the various Galactic model parameters (lens mass, lens and source distances, lens and source velocities) and the goodness of fit of the resulting magnification profile to the observed light curve. This approach was essential to enable a proper weighting of different permitted values of $`u_0`$. As a bonus, it allowed one, for the first time, to determine the sensitivities as a function of the physical planetary parameters (such as planet mass $`m_p`$ and planet-star separation $`r_{}`$) as opposed to the microlensing parameters, the planet-star mass ratio $`q`$ and the planet-star projected separation (in the units of $`\theta _\mathrm{E}`$) $`b`$.
Rhie et al. (2000) introduced a procedure for evaluating planet sensitivities that differs qualitatively from that of Gaudi & Sackett (2000). For each trial $`(b,q,\alpha )`$ and observed point-lens parameters $`(t_0,u_0,t_\mathrm{E})`$, they created a simulated light curve with epochs and errors similar to those of the real light curve. They then fitted this light curve to a point-lens model with $`(t_0,u_0,t_\mathrm{E})`$ left as free parameters. If the point-lens model had $`\mathrm{\Delta }\chi ^2>\chi _{\mathrm{min}}^2`$, then this $`(b,q,\alpha )`$ combination was regarded as excluded. That is, they mimicked their planet-detection procedure on simulated planetary events.
Abe et al. (2004) carried out a similar procedure except that they did not fit for $`(t_0,u_0,t_\mathrm{E})`$, but rather just held these three parameters fixed at their point-lens-fit values. Of course, this procedure necessarily yields a higher $`\mathrm{\Delta }\chi ^2`$ than that of Rhie et al. (2000), but Abe et al. (2004) expected that the difference would be small.
While all workers in this field have recognized that finite-source effects are important in principle, they have generally concluded that these did not play a major role in the particular events that they analyzed. This has proved fortunate because the source size is generally unknown, and even a single trial value for the source size typically requires several orders of magnitude more computing time than does a point-source model. Gaudi et al. (2002) estimated angular sizes $`\theta _{}`$ of each of their 43 source stars from their positions on an instrumental color-magnitude diagram (CMD) by adopting $`\mu =12.5\mathrm{km}\mathrm{s}^1\mathrm{kpc}^1`$ for all events and evaluating $`\rho _{}=\theta _{}/(\mu t_\mathrm{E})`$. They made their sensitivity estimates for both this value of $`\rho _{}`$ and for a point source ($`\rho _{}=0`$) and found that generally the differences were small. They concluded that a more detailed finite-source evaluation was unwarranted (and also computationally prohibitive). Using their Monte Carlo technique, Yoo et al. (2004b) were able to evaluate the probability distribution of the parameter combination $`z_0`$, which is equal to impact parameter over source size. This analysis showed that $`z_01`$ (no finite-source effects) with high confidence for their event. This implied that the source did not pass close to the central caustic and hence that finite-source effects were not important. Again, computation for additional values of $`\rho _{}`$ would have been computationally prohibitive.
### 1.3 ”Seeing” the Lens in High-Magnification Events
In the very first paper on microlensing, Einstein (1936) already realized that it might be difficult to observe the magnified source due to “dazzling by the light of the much nearer \[lens\] star.”
Seventy years later, more than 2000 microlensing events have been discovered, but only for two of these has the “dazzling” light of the lens star been definitively observed. The best case is MACHO-LMC-5, for which the lens was directly imaged by the Hubble Space Telescope (HST) (Alcock et al., 2001; Drake et al., 2004), which yielded mass and distance estimates of the lens that agreed to good precision with those derived from the microlensing event itself (Gould et al., 2004).
The next best case is OGLE-2003-BLG-175/MOA-2003-BLG-45, for which Ghosh et al. (2004) showed that the blended light was essentially perfectly aligned with the source. This would be expected if the blend actually was the lens, but it would be most improbable if it were just an ambient field star. In this case, the blend was about 1 mag brighter than the baseline source in $`I`$ and 2 mag brighter in $`V`$, perhaps fitting Einstein’s criterion of a “dazzling” presence.
Intriguingly, the above two events positively identified to harbor luminous lenses are both high-magnification. It is quite plausible that events with luminous lenses are biased towards high magnification since they will more likely be missed if the magnifications are too low. This raises the question of whether OGLE-2003-BLG-423 has a luminous lens. In addition, identifying the lens star would allow us to directly determine the physical properties of the lens, which in turn would help better constrain parameters in the planetary detection analysis.
Here we analyze OGLE-2004-BLG-343, whose maximum magnification $`A_{\mathrm{max}}3000`$ is by far the highest of any observed event and the first to exceed the $`A=1000`$ benchmark initially discussed by Liebes (1964) as roughly the maximum possible magnification for typical Galactic sources and lenses.<sup>5</sup><sup>5</sup>5Liebes (1964) derived that for perfect lens-source alignment, $`A_{u=0}=2\theta _\mathrm{E}/\theta _{}`$ by approximating the source star as having uniform surface brightness, and he evaluated this expression for several typical examples.
As we describe below, the event was alerted as a possibly anomalous, very-high-magnification event in time to trigger intensive observations over the peak, but due to human error, the actual observations caught only the falling side of the peak. We analyze both the actual observations made of this event (in order to evaluate its actual sensitivity to planets) and the sequence of observations that should have been initiated by the trigger. The latter calculation illustrates the potential of state-of-the-art microlensing observations, although unfortunately this potential was not realized in this case.
We analyze the event within the context of the Yoo et al. (2004b) formalism with several major modifications. First, we adopt the Rhie et al. (2000) criterion of planet-sensitivity in place of that of Gaudi & Sackett (2000). That is, we say a planet configuration is ruled out if simulated data generated by this configuration are inconsistent with a point-lens light curve at $`\mathrm{\Delta }\chi ^2>\mathrm{\Delta }\chi _{\mathrm{min}}^2`$ . Second, using a Monte Carlo simulation, we show that for this event, $`z_01`$, and hence finite-source effects are quite important. This requires us to generalize the Yoo et al. (2004b) procedure to include a two-dimensional grid of trial parameters $`(u_0,\rho _{})`$ in place of the one-dimensional $`u_0`$ sequence used by Yoo et al. (2004b). From what was said above, it should be clear that the required computations would be completely prohibitive if they were carried out using previous numerical techniques. Therefore, third, we develop new techniques for finite-source calculations that are substantially more efficient than those used previously.
In § 2, we describe the data. Next, we discuss modeling of the event in § 3. Then in § 4, we present our procedures and results related to planet detection. We explore the possibility that the blended light is due to the lens in § 5. In § 6 we summarize our results and make suggestions on monitoring extremely high-magnification events in the future. Finally, the two new binary-lens finite-source algorithms that we have developed are described in Appendix A.
## 2 Observational Data
The first alert on OGLE-2004-BLG-343 was triggered by the OGLE-III Early Warning System (EWS; Udalski 2003) on 2004 June 16, about 3 days before its peak on HJD$`{}_{}{}^{}`$ HJD$`2450000=3175.7467`$. On June 18, after the first observation of the lens was taken, the OGLE real time lens monitoring system (Early Early Warning System \[EEWS\]; Udalski 2003) triggered an internal alert, indicating a deviation from the single lens light curve based on previous data. Two additional observations were made after that, but the new fits to all of the data were still fully compatible with single-mass lens albeit suggesting high magnification at maximum. Therefore, an alert to the microlensing community was distributed by OGLE on HJD = 3175.1 suggesting OGLE-2004-BLG-343 as a possible high-magnification event. The observation at UT 0:57 (HJD 3175.54508) the next night showed a large deviation from the light-curve prediction based on previous observations, and an internal EEWS alert was triggered again. Usually further observations would have been made soon after such an alert, but unfortunately no observations were performed until UT 6:29 (HJD 3175.77626), about 0.71 hr after the peak. At that time, the event had brightened by almost 3 mag in $`I`$ relative to the previous night’s observation, and therefore it was regarded very likely to be the first caustic crossing of a binary-lens event. As a consequence, no $`V`$-band photometry was undertaken to save observation time and in the hope that observations in $`V`$ could be done when it brightened again. After two post-peak observations confirmed the event’s extremely high magnification, OGLE began maximally intensive observations with a cadence of 4.3 minutes. However, after it was clear that the event was falling in a regular fashion, it was then observed less intensively. A total of 14 observations were performed during 3.39 hr, and a new alert to the microlensing community was immediately released by OGLE as well. However, the next day the event faded drastically, by about 3 mag from the maximum point of the previous night, implying that if the event were a binary, the peak had probably been a cusp crossing rather than a caustic crossing. After being monitored for a few more days, it became clearer that OGLE-2004-BLG-343 was most probably a point-lens event of very high magnification and therefore very sensitive to planets. This recognition prompted OGLE to obtain a $`V`$-band point, but by this time (HJD 3179.51) the source had fallen 6 mag from its peak, so that only a weak detection of the $`V`$ flux was possible. Hence, this yielded only a lower limit on the $`VI`$ color.
By chance, $`\mu `$FUN made one dual-band observation in $`I`$ and $`H`$ 1 day before peak (HJD 3174.74256) solely as a reference point to check on the future progress of the event. After the event peaked, $`\mu `$FUN also concluded that it was uninteresting until OGLE/$`\mu `$FUN email exchanges led to the conclusion that the event was important. Since the source was magnified by A$``$40 at this pre-peak $`\mu `$FUN observation, it enabled a clear $`H`$-band detection and so yielded an $`(IH)`$ color measurement, which can be translated to $`(VI)`$.
The OGLE data are available at the OGLE EWS Web site mentioned above and the $`\mu `$FUN data are available at the $`\mu `$FUN Web site <sup>6</sup><sup>6</sup>6http://www.astronomy.ohio-state.edu/~microfun.
There were 195 images in $`I`$ and eight images in $`V`$, both from OGLE, as well as three images in $`H`$ from $`\mu `$FUN. Since only OGLE $`I`$-band observational data are available near the peak, the following analysis is entirely based on the OGLE $`I`$-band data except that the OGLE $`V`$-band data and $`\mu `$FUN $`H`$-band data are used to constrain the color of the source star. The OGLE errors are renormalized by a factor of 1.42 so that the $`\chi ^2`$ per degree of freedom for the best-fit point-source/point-lens (PSPL) model is close to unity. We also eliminate the two OGLE points that are $`3\sigma `$ outliers. These are both well away from the peak and therefore their elimination has no practical impact on our analysis.
## 3 Event Modeling
Yoo et al. (2004b) introduced a new approach to model microlensing events for which $`u_0`$ is not perfectly measured. As distinguished from previous analyses, this method establishes the prior probability of the event parameters by performing a Monte-Carlo simulation of the event using a Galactic model rather than simply assuming uniform distributions. This approach is not only more realistic but also makes possible the estimation of physical parameters, which are otherwise completely degenerate. Following the procedures of Yoo et al. (2004b), we begin our modeling by fitting the event to a PSPL model, evaluating the finite-source effects and performing a Monte-Carlo simulation. We then improve that method by considering finite-source effects when combining the simulation with the light-curve fits.
### 3.1 Point-Source Point-Lens Model
The PSPL magnification is given by (Paczyński, 1986)
$$A(u)=\frac{u^2+2}{u\sqrt{u^2+4}},u(t)=\sqrt{u_0^2+\frac{(tt_0)^2}{t_\mathrm{E}^2}},$$
(1)
where $`u`$ is the projected lens-source separation in units of the angular Einstein radius $`\theta _\mathrm{E}`$, $`t_0`$ is the time of maximum magnification, $`u_0=u(t_0)`$ is the impact parameter, and $`t_\mathrm{E}`$ is the Einstein timescale.
The predicted flux is then
$$F(t)=F_sA[u(t)]+F_b,$$
(2)
where $`F_\mathrm{s}`$ is the source flux and $`F_\mathrm{b}`$ is the blended-light flux.
The observational data are fitted in the above model with five free parameters ($`t_0`$, $`u_0`$, $`t_\mathrm{E}`$, $`F_\mathrm{s}`$, and $`F_\mathrm{b}`$). The results of the fit are shown in Table 1 (also see Fig. 1). The best-fit $`u_0`$ is remarkably small, $`u_0=0.000333\pm 0.000121`$, which indicates that the maximum magnification is $`A_{\mathrm{max}}=3000\pm 1100`$. As discussed below in § 3.3, the $`3\sigma `$ lower limit is $`A_{\mathrm{max}}1450`$. This is the first microlensing event ever analyzed in the literature with peak magnification higher than 1000. The uncertainties in $`u_0`$, $`t_E`$ and $`F_\mathrm{s}`$ are fairly large, roughly 35%. As pointed out in Yoo et al. (2004b), these errors are correlated, while combinations of these parameters, $`t_{\mathrm{eff}}u_0t_\mathrm{E}`$ and $`F_{\mathrm{max}}F_\mathrm{s}/u_0`$, have much smaller errors:
$$t_{\mathrm{eff}}=0.0141\pm 0.0008\mathrm{days},\mathrm{I}_{\mathrm{min}}=13.805\pm 0.065.$$
(3)
Here $`I_{\mathrm{min}}`$ is the calibrated $`I`$-band magnitude corresponding to $`F_{\mathrm{max}}`$.
### 3.2 Source Properties from Color-Magnitude Diagram
It is by now standard practice to determine the dereddened color and magnitude of a microlensed source by putting the best-fit instrumental color and magnitude of the source on an instrumental $`(I,VI)`$ CMD. The dereddened color and magnitude can then be determined from the offset of the source position from the center of the red clump, which is locally measured to be $`[M_I,(VI)_0]=(0.20,1.00)`$. We adopt a Galactocentric distance $`R_0=8`$kpc. However, at Galactic longitude $`l=+4.21`$, the red clump stars in the OGLE-2004-BLG-343 field are closer to us than the Galactic center by $`0.15`$ mag (Stanek et al., 1997). We derive $`(I,VI)_{0,\mathrm{clump}}=(14.17,1.00)`$. Although the source instrumental color and magnitude are both fit parameters, only the magnitude is generally strongly correlated with other fit parameters. By contrast, the source instrumental color can usually be determined directly by a regression of $`V`$ on $`I`$ flux as the magnification changes. No model of the event is actually required to make this color determination. In the present case, we exploit both $`(VI)`$ and $`(IH)`$ data. Hence, in order to make use of this technique, we must convert the $`(IH)`$ to $`(VI)`$. This will engender some difficulties.
As discussed in § 2, however, $`V`$-band measurements were begun only when the source had fallen nearly to baseline. Hence, the measurement of the $`(VI)`$ color obtained by this standard procedure has very large errors and indeed is consistent with infinitely red ($`F_{\mathrm{s},V}=0`$) at the $`2\sigma `$ level (see Fig. 2). The CMD itself is based on OGLE-II photometry, and we have therefore shifted the OGLE-III-derived fluxes by $`\mathrm{\Delta }I=I_{\mathrm{OGLE}\mathrm{II}}I_{\mathrm{OGLE}\mathrm{III}}=0.26`$ mag. On this now calibrated CMD, the clump is at $`(I,VI)_{\mathrm{clump}}=(15.51,2.04)`$. Hence, the dereddened source color and magnitude are given by $`(I,VI)_0=(I,VI)+(I,VI)_{0,\mathrm{clump}}(I,VI)_{\mathrm{clump}}=(I,VI)(1.34,1.04)`$, the final offset being the reddening vector. This vector corresponds to $`R_{VI}=1+1.34/1.04=2.29`$, which is somewhat high compared to values obtained by Sumi (2004) for typical bulge fields. However, we will present below independent evidence for this or a slightly higher value of $`R_{VI}`$. Figure 2 also shows the position of the blended light, which lies in the so-called reddening sequence of foreground disk main sequence stars. This raises the question as to whether this blended star is actually the lens. We return to this question in § 5.
The source star is substantially fainter than any of the other stars in the OGLE-II CMD. In order to give a sense of the relation between this source CMD position and those of main-sequence bulge stars, we also display the Hipparcos main sequence (ESA, 1997), placed at $`10^{0.15/5}R_0=7.5`$kpc and reddened by the reddening vector derived from the clump. At the best-fit value, $`VI=3.09`$, the source lies well in front of (or to the red of) the bulge main-sequence. However, given the large color error, it is consistent with lying on the bulge main sequence at the $`1\sigma `$ level.
To obtain additional constraints on the color, we consider the $`\mu `$FUN instrumental $`H`$-band data. The single highly-magnified $`(A40)`$ $`H`$-band point (together with a few baseline points) yields $`I_{\mathrm{OGLE}\mathrm{II}}H_{\mu \mathrm{FUN}}=0.59\pm 0.11`$ source color. To be of use, this must be translated to a $`(VI)_{\mathrm{OGLE}\mathrm{II}}`$ color using a $`(VI)/(IH)`$ color-color diagram of the stars in the field.
Unfortunately, there are actually very few field stars in the appropriate color range. This partly results from the small size ($`2`$ arcmin<sup>2</sup>) of the $`H`$-band image and partly from the fact that a large fraction of stars are either too faint to measure in $`V`$-band or saturated in $`H`$-band. We therefore calibrate the $`\mu `$FUN $`H`$-band data by aligning them to Two Micron All Sky Survery (2MASS) data and generate a $`(VI)/(IH_{2\mathrm{M}\mathrm{A}\mathrm{S}\mathrm{S}})`$ color-color diagram by matching stars from the 2MASS $`H`$-band data with OGLE-II $`V,I`$ photometry in a larger field centered on OGLE-2004-BLG-343. We find that $`(H_{2\mathrm{M}\mathrm{A}\mathrm{S}\mathrm{S}}H_{\mu \mathrm{FUN}})=1.99\pm 0.01`$ from 48 stars in common in the field, with a scatter of 0.08 mag. We transform the above $`IH_{\mu \mathrm{FUN}}`$ color to $`IH_{2\mathrm{M}\mathrm{A}\mathrm{S}\mathrm{S}}`$ and plot it as a vertical line on a $`(VI)/(IH_{2\mathrm{M}\mathrm{A}\mathrm{S}\mathrm{S}})`$ color-color diagram (see Fig. 3). From the intersection of the vertical line with the diagonal track of stars in the field, we infer $`VI=2.40\pm 0.15`$.
Since the field stars used to make the alignment are giants, this transformation would be strictly valid only if the source were a giant as well. However, the source star is certainly a dwarf (see Fig. 2). After transforming 2MASS to standard infrared bands (Carpenter, 2001), we use the data from Tables II and III of Bessell & Brett (1988) to construct dwarf and giant tracks on a $`(VI)_0/(IH)_0`$ color-color diagram. These are approximately coincident for blue stars $`(IH)_0<1.6`$ but rapidly separate by 0.28 mag in $`(VI)_0`$ by $`(IH)_0=1.7`$. In principle, we should just adjust our estimate $`(VI)`$ by the difference between these two tracks at the dereddened $`(IH)_0`$ color of the source. Unfortunately, there are two problems with this seemingly straightforward procedure. First, the Bessell & Brett (1988) giant track displays a modest deviation from its generally smooth behavior close to the color of our source, a deviation that is not duplicated by either the giants in our field or the color-color diagram formed by combining Hipparcos and 2MASS data, which both show the same smooth behavior at this location. Second, if the Hipparcos/2MASS diagram or the Bessell & Brett (1988) diagram is reddened using the selective and total extinctions determined above from the position of the clump, then the giant tracks do not align with our field giants. To obtain alignment, one must use $`R_{VI}=2.4`$.<sup>7</sup><sup>7</sup>7Gould et al. (2001) found a similar value using the same method but a different data set. However, the $`R_{VI}`$ we obtained at the beginning of this section is based on the dereddened magnitude of the red clump, which depends on the distance to the Galactic center $`R_0`$. If we were to adopt the new geometric measurement of $`R_0=7.62`$kpc (Eisenhauer et al., 2005), rather than the standard value of $`R_0=8.0`$kpc, we would then have $`I_{0,\mathrm{clump}}=14.07`$, which would give $`R_{VI}=2.39`$. However this value conflicts still more severely with the typical values of $`R_{VI}=1.9`$$`2.1`$ in bulge fields found by Sumi (2004). The conflict among these three determinations of $`R_{VI}`$ ($`1.9`$$`2.1`$ \[Sumi 2004\], 2.29 \[clump\], $`2.4`$ \[Gould et al. 2001\]) is quite a puzzle, but not one that we can explore here.
The bottom line is that there is considerable uncertainty in the dwarf-minus-giant adjustment but only in the upward direction. To take account of this, we add 0.2 mag error in quadrature to the upward error bar and finally adopt $`VI=2.4_{0.15}^{+0.25}`$ for the indirect color determination via the $`(IH)`$ measurement. Finally, we combine this with the direct measurement of $`VI=3.09`$ based on the combined $`V`$ and $`I`$ light curve. Because the errors on the latter measurement are extremely large (and are Gaussian in flux rather than magnitudes), we determine the probability distribution for the combined determination numerically in a flux-based calculation and then convert to magnitudes. We finally find $`VI=(VI)_{\mathrm{best}}\pm \sigma (VI)=2.60\pm 0.20`$, which we show as a magenta point in Figure 2. Hence,
$$(VI)_0=1.56\pm 0.20.$$
(4)
In contrast to most microlensing events that have been analyzed for planets, the color of OGLE-2004-BLG-343 is fairly uncertain. The color enters the analysis in two ways. First, it indicates the surface brightness and so determines the relation between dereddened source flux and angular size. Second, it determines the limb-darkening coefficient.
Given the color error, we consider a range of colors in our analysis and integrate over this range, just as we integrate over a range of impact parameters $`u_0`$ and source sizes (normalized to $`\theta _\mathrm{E}`$) $`\rho _{}`$. We allow colors over the range $`2.2<(VI)<3.0`$ corresponding to $`1.16<(VI)_0<1.96`$. We integrate in steps of 0.1 mag. For each color, we adopt a surface brightness such that the source size $`\theta _{}`$ is given by
$$\theta _{}=\theta _{(VI)}\mathrm{\hspace{0.17em}10}^{0.2(II_{\mathrm{best}})},$$
(5)
where $`I`$ is the (reddened) apparent magnitude in the model, $`I_{\mathrm{best}}=22.24`$, and $`\theta _{2.2}\mathrm{}\theta _{3.0}=(0.350,0.371,0.391,0.421,0.466,0.515,0.546,0.580,0.615)\mu `$as. These values are derived from the color/surface-brightness relations for dwarf stars given in Kervella et al. (2004) using the method as described in Yoo et al. (2004a). In our actual calculations, we use the full distribution of source radii, but for reference we note that the $`1\sigma `$ range of this quantity is
$$\theta _{}=0.47\pm 0.13\mu \mathrm{as}.$$
(6)
We find from the models of Claret (2000) and Hauschildt et al. (1999) that the linear limb-darkening coefficients for dwarfs in our adopted color range vary by only a few hundredths. Therefore, for simplicity, we adopt the mean of these values
$$\mathrm{\Gamma }_I=0.50$$
(7)
for all colors. This corresponds to $`c=3\mathrm{\Gamma }/(2+\mathrm{\Gamma })=0.60`$ in the standard limb-darkening parameterization (Afonso et al., 2000).
Finally, each model specifies not only a color and magnitude for the source, but also a source distance. Evaluation of the likelihood of each specific combination of these requires a color-magnitude relation. We adopt (Reid, 1991)
$$M_I=2.37(VI)_0+2.89$$
(8)
with a scatter of 0.6 mag. The ridge of this relation is shown as a red line segment in Figure 2, with the sources placed at $`10^{0.15/5}R_0=7.5`$kpc and reddened according to the red-clump determination, just as was done for the Hipparcos stars. This track is in reasonable agreement with the Hipparcos stars.
### 3.3 Finite-Source Effects
Yoo et al. (2004b) define $`z_0u_0/\rho _{}`$ (where $`\rho _{}=\theta _{}/\theta _\mathrm{E}`$ is the angular size of the source $`\theta _{}`$ in units of $`\theta _\mathrm{E}`$), which is a useful parameter to characterize the finite-source effects in single-lens microlensing events. We fit the observational data to a set of point-lens models on a grid of ($`u_0`$, $`z_0`$) and then compare the resulting $`\chi ^2`$ with the best-fit PSPL model. As mentioned in § 3.2, we adopt the limb-darkening formalism of Yoo et al. (2004b) and for simplicity choose $`\mathrm{\Gamma }_I=0.50`$.
Figure 4 displays the resulting $`\mathrm{\Delta }\chi ^2`$ contours. It shows that the 1 $`\sigma `$ contour extends from $`z_00.2`$ to arbitrarily large $`z_0`$. This is qualitatively similar to OGLE-2003-BLG-423 as analyzed in Yoo et al. (2004b). However, as we demonstrate in § 3.4, the range of $`z_0`$ that is consistent with the Galactic model is quite different for these two events. This is to be expected since $`z_0=u_0/\rho _{}=(u_0/\theta _{})\theta _\mathrm{E}`$, and $`u_0/\theta _{}`$ is roughly 8 times smaller for this event, while $`\theta _\mathrm{E}`$ is generally of the same order.
Figure 4 shows contours for both $`A_{\mathrm{max}}`$ and $`u_0^1`$. For $`z_0u_0/\rho _{}>2`$, these are very similar, which is expected because in the absence of finite-source effects (and for $`u_01`$), $`A_{\mathrm{max}}=u_0^1`$. Note that the $`A_{\mathrm{max}}`$ contours are roughly rectangular, so that while $`z_0`$ is not well constrained, the $`3\sigma `$ lower limit $`A_{\mathrm{max}}>10^{3.16}1450`$ is quite well defined. This shows that although the blending is very severe, it is also very well constrained, implying that the event’s high magnification is secure.
### 3.4 Monte-Carlo Simulation
We perform a similar Monte-Carlo simulation using a Han & Gould (1996, 2003) model as described in Yoo et al. (2004b) by taking into account all combinations of source and lens distances, $`D_\mathrm{l}<D_\mathrm{s}`$, uniformly sampled along the line of sight toward the source $`(l,b)=(4.21,3.47)`$. The simulation adopts the Gould (2000) mass function taking into account the bulge main sequence stars, white dwarfs (distributed around $`0.6\mathrm{M}_{}`$), neutron stars (narrowly peaked at $`1.35\mathrm{M}_{}`$), and stellar-mass black holes. This mass function is adequate to describe mass distributions of disk lenses except that the disk contains stars with masses greater than $`1\mathrm{M}_{}`$ while the bulge does not. However, a disk main sequence star more massive than the Sun will be too bright to be the lens star for this event (see Fig. 2); so for simplicity, we use this mass function for both disk and bulge in our simulation. In Yoo et al. (2004b), the source flux is determined from the $`t_\mathrm{E}`$ for each Monte-Carlo event since only the PSPL model is considered at this step. However, when finite-source effects are taken into account, each $`t_\mathrm{E}`$ corresponds to a series of $`F_\mathrm{s}`$ depending on the source size $`\rho _{}`$, so there is no longer a 1-1 correspondence between $`F_\mathrm{s}`$ and $`t_\mathrm{E}`$. As discussed in Yoo et al. (2004a), $`\theta _{}`$ can be deduced from the source’s dereddened color and magnitude. Since $`\theta _\mathrm{E}`$ is known for each simulated event, $`\rho _{}`$ is a direct function of $`F_\mathrm{s}`$ and the $`(VI)`$ color of the source,
$$\rho _{}=\frac{\theta _{(VI)}}{\theta _\mathrm{E}}\sqrt{\frac{F_\mathrm{s}}{F_{\mathrm{best}}}},$$
(9)
where $`F_{\mathrm{best}}`$ corresponds to $`I_{\mathrm{best}}`$ in equation (5). Using this constraint, we fit the $`k`$-th Monte-Carlo event to a point-lens model with finite-source effects, holding $`t_{\mathrm{E},k}`$ fixed at the value given by the simulation, for a variety of $`(VI)`$ color values inferred from § 3.2. Hence, for the $`j`$-th $`(VI)`$ color and $`k`$-th Monte-Carlo event, we have best-fit single-lens light-curve parameters $`t_{0,j,k}`$, $`u_{0,j,k}`$, $`\rho _{,j,k}`$, $`F_{\mathrm{s},j,k}`$, $`F_{\mathrm{b},j,k}`$, as well as $`\mathrm{\Delta }\chi _{j,k}^2\chi _{j,k}^2\chi _{\mathrm{PSPL}}^2`$. We construct a three-dimensional table that includes these six quantities as well as the other parameters from the Monte-Carlo simulation ($`t_{\mathrm{E},k}`$, $`\theta _{\mathrm{E},k}`$, $`D_{\mathrm{s},k}`$, $`D_{\mathrm{l},k}`$, $`M_{\mathrm{l},k}`$, $`\mathrm{\Gamma }_k`$), the Einstein timescale and radius, the source and lens distances, the lens mass, and the event rate. From these can also be derived two other important quantities, the source absolute magnitude $`M_{I,j,k}`$ and the physical Einstein radius $`r_{\mathrm{E},k}\theta _{\mathrm{E},k}\times D_{\mathrm{l},k}`$. This three-dimensional table is composed of nine two-dimensional tables, one for each $`(VI)_j`$ color. Each table contains approximately 200,000 rows, one for each simulated event. To make the notation more compact, we refer to the parameters $`a,b,c,\mathrm{}`$ lying in the bin $`a[a_{\mathrm{min}},a_{\mathrm{max}}];b[b_{\mathrm{min}},b_{\mathrm{max}}];c[c_{\mathrm{min}},c_{\mathrm{max}}]\mathrm{}`$ as $`\mathrm{bin}(\{a,b,c,\mathrm{}\})`$.
Similarly to Yoo et al. (2004b), the posterior probability of $`a_i`$ lying in $`\mathrm{bin}(a_i)`$ is given by
$$P[\mathrm{bin}(a_i)]\underset{j,k}{}(P_{VI})_j\times (P_{\mathrm{Reid}})_{j,k}\times \mathrm{exp}[\mathrm{\Delta }\chi _{j,k}^2/2]\times \mathrm{BC}[\mathrm{bin}(\{a_i\}_{j,k})]\times \mathrm{\Gamma }_k,$$
(10)
where $`(P_{\mathrm{Reid}})_{j,k}=\mathrm{exp}(\{(M_I)_{j,k}M_{I,\mathrm{Reid}}[(VI)_{0,j}]\}^2/2[\sigma _{\mathrm{Reid}}^2+(\sigma _{M_I})_{j,k}^2])`$ accounts for the scatter ($`\sigma _{\mathrm{Reid}}=0.6`$) in $`M_I`$ about the Reid relation plus the dispersion $`(\sigma _{M_I})_{j,k}`$ from light-curve fitting, $`(P_{VI})_j=\mathrm{exp}\{[(VI)_j(VI)_{\mathrm{best}}]^2/2\sigma _{(VI)}^2\}`$ reflects the uncertainty in $`VI`$ color, and $`\mathrm{BC}`$ is a boxcar function defined by $`\mathrm{BC}[\mathrm{bin}(a)]\mathrm{\Theta }\left(aa_{\mathrm{min}}\right)\times \mathrm{\Theta }\left(a_{\mathrm{max}}a\right)`$.
Figure 5 shows the posterior probability distributions of various parameters, including $`u_0`$, dereddened apparent $`I`$-band magnitude of the source $`I_0`$, proper motion $`\mu `$, $`z_0`$, source distance modulus, lens distance modulus, absolute $`I`$-band magnitude of the source $`M_I`$, and lens mass. The blue and red histograms represent bulge-disk events and bulge-bulge events, respectively. The relative event rate is normalized in the same way for both bulge-disk and bulge-bulge events. The total rate for bulge-disk events is about $`6`$ times larger than that for bulge-bulge events, which means that the Monte-Carlo simulation tends strongly to favor bulge-disk events.
The Einstein radii are on average smaller for bulge-bulge events than for bulge-disk events, and as a result, the bulge-bulge events tend to have bigger $`\rho _{}`$ and hence smaller $`z_0`$. However, the top right panel of Figure 5 shows that the $`z_0`$ probability distributions have similar shapes for both bulge-bulge and bulge-disk events. This is because the (lack of) finite-source effects constrain $`z_00.7`$ at the $`3\sigma `$ level (see Fig. 4), which cuts off the lower end of the $`z_0`$ distributions for both categories of events. Since bulge-disk events have smaller $`\rho _{}`$ than bulge-bulge events, the $`u_0`$ posterior probability distribution peaks at a lower value for the former. Furthermore, since $`z_00.7`$, the proper motion is constrained to be $`\mu =\theta _{}z_0/t_{\mathrm{eff}}7\mathrm{m}\mathrm{a}\mathrm{s}/\mathrm{yr}`$, which is typical of bulge-disk events but $`2`$ times the proper motion of typical bulge-bulge events.
Figure 5 also shows the distributions of $`u_0`$ and $`I_0`$ from the light-curve data alone by a black solid line. In strong contrast to the corresponding diagrams for OGLE-2003-BLG-423 presented by Yoo et al. (2004b), the light-curve based parameters agree quite well with the Monte-Carlo predictions. In the source distance-modulus panel, the prior distributions for bulge-disk and bulge-bulge events are shown in purple and green histograms, respectively. Again, distinct from OGLE-2003-BLG-423, the most likely source distances of this event agree reasonably well with typical values from the prior distributions. Moreover, from Figure 5, the peak values of source $`M_I`$ distributions are also in good agreement with those derived from the Reid relation ($`M_I=6.59`$ for $`\left[VI\right]_0=1.56`$, see eq. ). Therefore, the source of this event shows very typical characteristics as represented by the Monte-Carlo simulation. Also unlike OGLE-2003-BLG-423, the probability that $`z_01`$ is very high for both bulge-disk and bulge-bulge events. Therefore, finite-source effects must be taken into account in the analysis of this event. In addition to the posterior probability distribution (orange) of the lens mass, the prior distribution (dark green) is displayed in Figure 5 as well. In microlensing analyses, the lens mass is commonly assigned a “typical” value (for example, $`0.3\mathrm{M}_{}`$). However, Figure 5 shows that, lenses with relatively high mass are strongly favored for this event as compared to the prior distribution. Detailed discussions on the lens properties are presented in § 5.
## 4 Detecting Planets
While there are no obvious deviations from point-lens behavior in the light curve of OGLE-2004-BLG-343 at our adopted threshold of $`\mathrm{\Delta }\chi _{\mathrm{min}}^2=60`$, planetary deviations might be difficult to recognize by eye. We must therefore conduct a systematic search for such deviations. Logically, this search should precede the second step of testing to determine what planets we could have detected had they been there. However, as a practical matter it makes more sense to first determine the range of parameter space for which we are sensitive to planets because it is only this range that needs to be searched for planets. We therefore begin with this detection efficiency calculation.
### 4.1 Detection Efficiency
As reviewed in § 1.2, a variety of methods have been proposed to calculate the planetary sensitivities of microlensing events, either in predicting planetary detection efficiencies theoretically or in analyzing real observational data sets. In those methods, $`\mathrm{\Delta }\chi ^2`$ is often calculated by subtracting the $`\chi ^2`$ of single-lens models from that of the binary-lens models to evaluate detection sensitivities. However, the ways in which single-lens and binary-lens models are compared differ from study to study. As noted by Griest & Safizadeh (1998) and Gaudi & Sackett (2000), for real planetary light curves, the lens parameters are not known a priori. Therefore, $`\mathrm{\Delta }\chi ^2`$ will generally be exaggerated if one subtracts from the binary lens model the single-lens model that has the same $`t_0`$, $`u_0`$, and $`t_\mathrm{E}`$ instead of the best-fit single-lens model to the binary light curve. One important factor contributing to this exaggeration is that the center of the magnification pattern (referred to as the center of the caustic in Yoo et al. 2004b) in the binary-lens case is no longer the position of the primary star as it is in the single-lens model (Dominik, 1999b; An & Han, 2002). Therefore light-curve parameters such as $`u_0`$ and $`t_0`$ will shift correspondingly. We find that by not taking into account this effect and directly comparing the simulated binary (i.e., planetary) light curve with the best-fit single-lens model to the data, Abe et al. (2004) exaggerate the planetary sensitivity of MOA 2003-BLG-32/OGLE 2003-BLG-219, although it remains the most sensitive event analyzed to date (see Appendix B).
Following Gaudi & Sackett (2000), planetary systems are characterized by a planet-star mass ratio $`q`$, planet-star separation in Einstein radius $`b`$, and the angle $`\alpha `$ of the source trajectory relative to the planet-star axis. In our binary-lens calculations, $`u_0`$ and $`t_0`$ are defined with respect to the center of magnification discussed above. According to Gaudi & Sackett (2000), the next step is to fit the data to both PSPL models and binary-lens models with a variety of $`(b,q,\alpha )`$ and calculate $`\mathrm{\Delta }\chi ^2(b,q,\alpha )=\chi ^2(b,q,\alpha )\chi _{\mathrm{PSPL}}^2`$. Then $`\mathrm{\Delta }\chi ^2(b,q,\alpha )`$ is compared with a threshold value $`\chi _{\mathrm{thres}}^2`$: if $`\mathrm{\Delta }\chi ^2>\chi _{\mathrm{thres}}^2`$ then a planet with parameters $`b`$, $`q`$, and $`\alpha `$ is claimed to be excluded while it is detected if $`\mathrm{\Delta }\chi ^2<\chi _{\mathrm{thres}}^2`$. The $`(b,q)`$ detection efficiency is then obtained by integrating $`\mathrm{\Theta }(\mathrm{\Delta }\chi ^2\mathrm{\Delta }\chi _{\mathrm{thres}}^2)`$ over $`\alpha `$ in the exclusion region at fixed $`(b,q)`$, where $`\mathrm{\Theta }`$ is a step function. However, Gaudi et al. (2002) point out that for events with poorly constrained light-curve parameters, which is the case for OGLE-2004-BLG-343, this method will significantly underestimate the sensitivities since the binary-lens models will minimize the $`\chi ^2`$ over the relatively large available parameter space. As discussed in Yoo et al. (2004b), the detection efficiency should be evaluated at a series of allowed $`u_0`$ values. To take finite-source effects into account, we generate a grid of permitted $`(u_0,\rho _{})`$, and each $`(u_{0,m},\rho _{,m})`$ bin is associated with the probability $`P[\mathrm{bin}(\{u_{0,m},\rho _{,m}\})]`$ obtained using the following equation:
$$P[\mathrm{bin}(\{u_{0,m},\rho _{,m}\})]\underset{j,k}{}P_{m,j,k}$$
(11)
where
$$P_{m,j,k}=(P_{VI})_j\times (P_{\mathrm{Reid}})_{j,k}\times \mathrm{exp}[\mathrm{\Delta }\chi _{j,k}^2/2]\times \mathrm{BC}[\mathrm{bin}(\{u_{0,m}\}_{j,k})]\times \mathrm{BC}[\mathrm{bin}(\{\rho _{,m}\}_{j,k})]\times \mathrm{\Gamma }_k$$
(12)
If the light-curve parameters were well-constrained, the approaches of Gaudi & Sackett (2000) and Rhie et al. (2000) would be very nearly equivalent, with the former retaining a modest philosophical advantage, since it uses only the observed light curve and does not require construction of light curves for hypothetical events. However, because in our case these parameters are not well constrained, the Gaudi & Sackett (2000) approach would require integration over all binary-lens parameters (except $`F_\mathrm{s}`$ and $`F_\mathrm{b}`$). Regardless of its possible philosophical advantages, this approach is therefore computationally prohibitive in the present case. We therefore do not follow Gaudi & Sackett (2000), but instead construct a binary light curve with the same observational time sequence and photometric errors as the OGLE observations of OGLE-2004-BLG-343, for each $`(b,q,\alpha ;u_0,\rho _{})`$ combination and the associated probability-weighted parameters $`𝒂_{\mathrm{lc}}`$: $`t_0`$, $`t_\mathrm{E}`$, $`F_\mathrm{s}`$ and $`F_\mathrm{b}`$ in the $`m`$-th $`(u_0,\rho _{})`$ bin,
$$𝒂_{\mathrm{lc}(\mathrm{weighted}),\mathrm{m}}=\frac{\underset{j,k}{}P_{m,j,k}𝒂_{\mathrm{lc},j,k}}{\underset{j,k}{}P_{m,j,k}}.$$
(13)
Then each simulated binary light curve $`(b,q,\alpha ;u_{0,m},\rho _{,m})`$ is fitted to a single-lens model with finite-source effects whose best fit yields $`\chi ^2(b,q,\alpha ;u_{0,m},\rho _{,m})`$. Another set of artificial binary light curves is generated under the assumption that OGLE had triggered a dense series of observations following the internal alert at HJD 3175.54508. These cover the peak of the event with the normal OGLE frequency and are used to compare results with those obtained from the real observations.
Magnification calculations for a binary lens with finite-source effects are very time-consuming. Besides $`(b,q,\alpha )`$, our calculations are also performed on $`(u_0,\rho _{})`$ grids, two more dimensions than in any previous search of a grid of models with finite-source effects included. This makes our computations extremely expensive, comparable to those of Gaudi et al. (2002), which equaled several years of processor time. Therefore, we have developed two new binary-lens finite-source algorithms to perform the calculations, as discussed in detail in Appendix A.
In principle, we should consider the full range of $`b`$, i.e., $`0<b<\mathrm{}`$; in practice, it is not necessary to directly simulate $`b<1`$ due to the famous $`bb^1`$ degeneracy (Dominik, 1999a; An, 2005). Instead, we just map the $`b>1`$ results onto $`b<1`$ except for the isolated sensitive zones along the $`x`$axis caused by planetary caustics perturbations.
We define the planetary detection efficiency $`ϵ(b,q)`$ as the probability that an event with the same characteristics as OGLE-2004-BLG-343, except that the lens is a planetary system with configuration of $`(b,q)`$, is inconsistent with the single-lens model (and hence would have been detected),
$`ϵ(b,q,\alpha )`$ $`=`$ $`\{{\displaystyle \underset{m}{}}\mathrm{\Theta }[\chi ^2(b,q,\alpha ;u_{0,m},\rho _{,m})\mathrm{\Delta }\chi _{\mathrm{thres}}^2]`$ (14)
$`\times P[\mathrm{bin}(\{u_{0,m},\rho _{,m}\})]\}`$
$`\times \{{\displaystyle \underset{m}{}}P[\mathrm{bin}(\{u_{0,m},\rho _{,m}\})]\}^1`$
and
$$ϵ(b,q)=\frac{1}{2\pi }_0^{2\pi }ϵ(b,q,\alpha )𝑑\alpha .$$
(15)
### 4.2 Constraints on Planets
Figure 6 shows the planetary detection efficiency of OGLE-2004-BLG-343 for planets with mass ratios $`q=10^3`$, $`10^4`$, and $`10^5`$, as a function of $`b`$, the planet-star separation (normalized to $`\theta _\mathrm{E}`$), and $`\alpha `$, the angle that the moving source makes with the binary axis passing the primary lens star on its left. Different colors indicate 10%, 25%, 50%, 75%, 90% and 100% efficiency. Note that the contours are elongated along an axis that is roughly $`60^{}`$ from the vertical (i.e., the direction of the impact parameter for $`\alpha =0`$). This reflects the fact that the point closest to the peak occurs at $`t=2453175.77626`$ when $`(tt_0)/t_\mathrm{E}=2.16u_0`$, and so when the source-lens separation is at an angle $`\mathrm{tan}^12.16=65^{}`$. For $`q=10^3`$, the region of 100% efficiency extends through $`360^{}`$ within about one octave on either side of the Einstein ring. However, at lower mass ratios there is 100% efficiency only in restricted areas close to the Einstein ring and along the above-mentioned principal axis.
Figure 7 summarizes an ensemble of all figures similar to Figure 6, but with $`q`$ ranging from $`10^{2.5}`$ to $`10^{5.0}`$ in 0.1 increments. To place this summary in a single figure, we integrate over all angles $`\alpha `$ at fixed $`b`$. Comparison of this figure to Figure 8 from Gaudi et al. (2002) shows that the detection efficiency of OGLE-2004-BLG-343 is similar to that of MACHO-1998-BLG-35 and OGLE-1999-BUL-35 despite the fact that their maximum magnifications are $`A_{\mathrm{max}}100`$, roughly 30 times lower than OGLE-2004-BLG-343. Of course, part of the reason is that OGLE-2004-BLG-343 did not actually probe as close as $`u=u_01/3000`$ because no observations were taken near the peak. However, observations were made at $`u1/1200`$, about 12 times closer than in either of the two events analyzed by Gaudi et al. (2002). One problem is that because the peak was not well covered, there are planet locations that do not give rise to observed perturbations at all. But this fact only accounts for the anisotropies seen in Figure 6. More fundamentally, even perturbations that do occur in the regions that are sampled by the data can often be fitted to a point-lens light curve by “adjusting” the portions of the light curve that are not sampled.
Note the central “spike” of reduced detection efficiency plots near $`b=1`$. As first pointed out by Bennett & Rhie (1996), this is due to the extreme weakness of the caustic for nearly resonant ($`b1`$) small mass-ratio ($`q1`$) binary lenses.
### 4.3 No Planet Detected
Based on the detection efficiency levels we obtained in § 4.2, we fit the observational data to binary-lens models to search for a planetary signal in the regions with efficiency greater than zero from $`q=10^5`$ to $`q=10^{2.5}`$. We find no binary-lens models satisfying our detection criteria. In fact, the total $`\chi ^2`$ contributions to the best-fit single-lens model of the observational points over the peak ( HJD $`=2453175.52453176.0`$) are no more than 30, so even if all of these deviations were due to a planetary perturbation, such a binary-lens solution would not easily satisfy our $`\mathrm{\Delta }\chi ^2=60`$ detection criteria. Therefore there are no planet detections in OGLE-2004-BLG-343 data.
### 4.4 Fake Data
Partly to explore further the issue of imperfect coverage of the peak, and partly to understand how well present microlensing experiments can probe for planets, we now ask what would have been the detection efficiency of OGLE-2004-BLG-343 if the internal alert issued on HJD 3175.54508 had been acted upon.
Of course, since the peak was not covered, we do not know exactly what $`u_0`$ and $`\rho `$ for this event are. However, for purposes of this exercise, we assume that they are near the best fit as determined from a combination of the light-curve fitting and the Galactic Monte Carlo, and for simplicity, we choose $`u_0=0.00040,\rho _{}=0.00040`$ which is very close to the best-fit combination. We then form a fake light curve sampled at intervals of 4.3 minutes, starting from the alert and continuing to the end of the actual observations that night. This sampling reflects the intense rate of OGLE follow-up observations actually achieved during this event (see § 2). We assume errors similar to those of the actual OGLE data at similar magnifications. For those points that are brighter than the brightest OGLE point, the minimum actual photometric errors are assigned. We also assume that the color information is known exactly in this case to be $`VI=2.6`$. We then analyze these fake data in exactly the same way that we analyze the real data. In contrast to the real data, however, we do not find a finite range of $`z_0u_0/\rho _{}`$ that are consistent with the fake data. Rather, we find that all consistent parameter combinations have $`z_0=1`$ almost identically. We therefore consider only a one-dimensional set of $`(u_0,\rho _{})`$ combinations subject to this constraint.
Figure 8 is analogous to Figure 6 except that the panels show planet sensitivities for $`q=10^3`$, $`10^4`$, $`10^5`$, and $`10^6`$, that is, an extra decade. In sharp contrast to the real data, these sensitivities are basically symmetric in $`\alpha `$, except for the lowest value of $`q`$. Sensitivities at all mass ratios are dramatically improved. For example, at $`q=10^3`$, there is 100% detection efficiency over 1.7 dex in $`b`$ ($`1/7b7`$). Even at $`q=10^5`$ (corresponding to an Earth-mass planet around an M star), there is 100% efficiency over an octave about the Einstein radius.
Figure 9 is the fake-data analog of Figure 7. It shows that this event would have been sensitive to extremely low mass-ratios, lower than those accessible to any other technique other than pulsar timing.
### 4.5 Detection Efficiency in Physical Parameter Space
One of the advantages of the Monte Carlo approach of Yoo et al. (2004b) is that it permits one to evaluate the planetary detection efficiency in the space of the physical parameters, planet mass and projected physical separation ($`m_p,r_{}`$), rather than just the microlensing parameters $`(b,q)`$. Figures 10 and 11 show this detection efficiency for the real and fake data, respectively. The fraction of Jupiter-mass planets that could have been detected from the actual data stream is greater than $`25\%`$ for $`0.8\mathrm{AU}r_{}10\mathrm{A}\mathrm{U}`$ and is greater than $`90\%`$ for $`2\mathrm{A}\mathrm{U}r_{}6\mathrm{A}\mathrm{U}`$. There is also marginal sensitivity to Neptune-mass planets. However, the detection efficiencies would have been significantly enhanced had the FWHM around the peak been observed, as previously discussed by Rattenbury et al. (2002). For the fake data, more than $`90\%`$ of Jupiter-mass planets in the range $`0.7\mathrm{AU}r_{}20\mathrm{A}\mathrm{U}`$ and more than $`25\%`$ with $`0.3\mathrm{AU}r_{}30\mathrm{A}\mathrm{U}`$ would have been detected. Indeed, some sensitivity would have extended all the way down to Earth-mass planets.
## 5 Luminous Lens?
Understanding the physical properties of their host stars is a major component of the study of extra-solar planets. It is especially important to know the mass and distance of the lens star for planets detected by microlensing because only then can we accurately determine the planet’s mass and physical separation from the star. Obtaining similar information for microlensing events that are unsuccessfully searched for planets enables more precise estimates of the detection efficiency. There are only two known ways to determine the mass and distance of the lens: either measure both the microlensing parallax and the angular Einstein radius (which are today possible for only a small subset of events) or directly image the lens. In most cases the lens is either entirely invisible or is lost in the much brighter light of the source.
A simple argument suggests, however, that in extremely high-magnification events like OGLE-2004-BLG-343, the lens will often be easily visible and, indeed, it is the lens that is unknowingly being monitored, with the source revealing itself only in the course of the event. Events of magnification $`A_{\mathrm{max}}`$ require that the source be much smaller than the Einstein radius, $`\theta _{}2\theta _\mathrm{E}/A_{\mathrm{max}}`$. Since $`\theta _\mathrm{E}=\sqrt{\kappa M\pi _{\mathrm{rel}}}`$, large $`\theta _\mathrm{E}`$ requires a lens that is either massive or nearby, both of which suggest that it is bright. On the other hand, a small $`\theta _{}`$ implies that the source is faint. Generally, if a faint source and a bright potential lens are close on the sky, only the lens will be seen, until it starts to strongly magnify the source. This has important implications for the real time recognition of extreme magnification events, as we discuss in § 6. Here we review the evidence as to whether the blended light in OGLE-2005-BLG-343 is in fact the lens.
As was true for OGLE-2003-BLG-175/MOA-2003-BLG-45 mentioned in § 1.3, the blended light in OGLE-2004-BLG-343 lies in the “reddening sequence” of foreground disk stars. It is certainly “dazzling” by any criterion, being about 50 times brighter than the source in $`I`$ and 150 times brighter in $`V`$ (see Fig. 2). Is the blended light also due to the lens in this case?
There is one argument for this hypothesis and another against. We initiate the first by estimating the mass and distance to the blend as follows. We model the extinction due to dust at a distance $`x`$ along the line of sight by $`dA_I/dx=0.4\mathrm{kpc}^1e^{qx}`$ and set $`q=0.26\mathrm{kpc}^1`$ in order to reproduce the measured extinction to the bulge $`A_I(8\mathrm{kpc})=1.34`$. Using the Reid (1991) color-magnitude relation, we then adjust the distance to the blend until it reproduces the observed color and magnitude of the blend. We find a distance modulus of 12.6 ($`3.3\mathrm{kpc}`$), and with the aid of the Cox (2000) mass-luminosity relation, we estimate a corresponding mass $`M_\mathrm{l}=0.9\mathrm{M}_{}`$. Inspecting Figure 5, we see that this is almost exactly the peak of the lens-distance distribution function predicted by combining light-curve information and the Galactic model. This is quite striking because, in the absence of light-curve information, the lens would be expected to be relatively close to the source. From our Monte-Carlo simulation toward the line of sight of this event, the total prior probability of the bulge-bulge events is about 1.5 times higher than the prior probability of the bulge-disk events, and furthermore, only about $`7\%`$ of all events have lenses less than $`3.3\mathrm{kpc}`$ away (see green and purple histograms in source and lens distance modulus panels of Fig. 5). It is only because the light curve lacks obvious finite-source effects (despite its very high-magnification) that one is forced to consider lenses with large $`\theta _\mathrm{E}`$, which generally drives one toward nearby lenses in the foreground disk. Based on our experience analyzing many blended microlensing events, the blended light is most often from a bulge star rather than a disk star, which simply reflects the higher density of bulge stars. In brief, it is quite unusual for lenses to be constrained to lie in the disk, and it is quite unusual for events to be blended with foreground disk stars. This doubly unusual set of circumstances would be more easily explained if the blend were the lens.
However, if the blend were the lens, then the source and lens would be aligned to better than 1 mas during the event, and one would therefore expect that the apparent position of the source would not change as the source first brightened and then faded. In fact, we find that the apparent position does change by about $`73\pm 9`$mas. However, since the apparent source (i.e., combined source and blended light at baseline) has a near neighbor at 830 mas, which is almost as bright as the source/blend, it is quite possible that the lens actually is the blend, but that this neighbor is corrupting the astrometry.
Thus, the issue cannot be definitively settled at present. However, it could be resolved in principle by, for example, obtaining high-resolution images of the field a decade after the event when the source and lens have separated sufficiently to both be seen. If the blend is the lens, then they will be seen moving directly apart with a proper motion given $`\mu =\theta _\mathrm{E}/t_\mathrm{E}`$, where $`\theta _\mathrm{E}`$ is derived from the estimated mass and distance to the lens and $`t_\mathrm{E}`$ is the event timescale.
Since the blend cannot be positively identified as the lens, we report our main results using a purely probabilistic estimate of the lens parameters. However, for completeness, we also report results here based on the assumption that the lens is the blend. Compared to the previous simulation, in which we considered the full mass function and full range of distances, we sample only the narrow intervals of mass and distance that are consistent with the observed color and magnitude of the lens/blend. To implement these restrictions, we repeat the Monte Carlo, but with the additional constraint that the predicted apparent magnitudes agree with the observed blend magnitude (with an error of 0.5 mag) and that the predicted colors (using the above extinction law and the Reid 1991 color-magnitude relation) also show good agreement with the observed color (with 0.2 mag error). These errors are, of course, much larger than the observational errors. They are included to reflect the fact that the theoretical predictions for color and magnitude at a given mass are not absolutely accurate.
Figure 12 is the resulting version of Figure 10 when the Monte Carlo is constrained to reproduce the blend color and magnitude. The sensitivity contours are narrower and deeper, reflecting the fact that the diagram no longer averages over a broad range of lens masses but rather is restricted effectively to a single mass (and single distance).
## 6 Summary and Discussion
In this paper we present our analysis of microlensing event OGLE-2004-BLG-343, with the highest peak magnification ($`A_{\mathrm{max}}=3000\pm 1100`$) ever analyzed to date. The light curve is consistent with the single-lens microlensing model, and no planet has been detected in this event. We demonstrate that if the peak had been well covered by the observations, the event would have had the best sensitivity to planets to date, and it would even have had some sensitivity to Earth-mass planets (§ 4.4, § 4.5). However, this potential has not been fully realized due to human error (§ 2), and OGLE-2004-BLG-343 turns out to be no more sensitive to planets than a few other high-magnification events analyzed before (§ 4.2, § 4.5). Thus, while ground-based microlensing surveys are technically sufficient to detect very low-mass planets, the relatively short timescale of the sensitive regime of high-magnification microlensing events demands a rapidity of response that is not consistently being achieved. In the final paragraph below, we develop several suggestions to rectify this situation.
In § 3 we show that finite-source effects are important in analyzing this event, so we extend the method of Yoo et al. (2004b) to incorporate such effects in planetary detection efficiency analysis. Moreover, since magnification calculations of binary-lens models with finite-source effects are computationally remarkably expensive, and applying previous finite-source algorithms, it would have taken of order a year of CPU time to do the detection efficiency calculation required by this event. We therefore develop two new binary-lens finite-source algorithms (Appendix A) that are considerably more efficient than previous ones. The “map-making” method (Appendix A.1) is an improvement on the conventional inverse ray-shooting method, which proves to be especially efficient for use in detection efficiency calculations, while the “loop-linking” method (Appendix A.2) is more versatile and could be easily implemented in programs aimed at finding best-fit finite-source binary-lens solutions. Using these algorithms, we were able to complete the computations for this paper in about 4 processor-weeks, roughly an order of magnitude faster than would have been required using previous algorithms.
Finally, we show in § 5 that the blend, which is a Galactic disk star, might very possibly be the lens, and that this case also proves to be highly probable from the Monte-Carlo simulation. However, it seems to contradict the astrometric evidence, and we point out that this issue could in principle be solved by future high-resolution images. Among the high-magnification events discovered by current microlensing survey groups, it is very likely that the lens star, which is also the apparent source, of those events is in the Galactic disk. Thus the blended light is usually far brighter than the source, thereby increasing the difficulty in early identification of such events. This fact motivates the first of several suggestions aimed at improving the recognition of very high-magnification events:
* When events are initially alerted they should be accompanied by instrumental CMD of the surrounding field, with the location of the apparent “source” highlighted. Events whose apparent sources lie on the “reddening sequence” of foreground disk stars (see Fig. 2) have a high probability to actually be lenses of more distant (and fainter) bulge sources. These events deserve special attention even if their initial light curves appear prosaic.
* For each such event it is possible to measure the color (but not immediately the magnitude) of the source by the standard technique of obtaining two-band photometry and measuring the slope of the relative fluxes in the two bands. If the color is different from that of the apparent “source” at baseline, that will prove that this baseline light is not primarily due to the source, and it will increase the probability that this baseline object is the lens. Moreover, if the source color is relatively red, it will show that the source is probably faint and so is (1) most likely already fairly highly magnified (thereby making it possible to detect above the foreground blended light) and (2) capable, potentially at least, of being magnified to very high magnification (see § 5). This would motivate obtaining more data while the event was still faint to help predict its future behavior and would enable a guess as to how to “renormalize” the event’s apparent magnification to its true magnification. This is important because generally one cannot accurately determine this renormalization until the event is within 0.4 mag (when the event is $`t_{\mathrm{eff}}`$ before its peak), at which point it may well be too late to act on this knowledge.
* Both survey groups and follow-up groups should issue alerts on suspected high-magnification events guided by a relatively low threshold of confidence, recognizing that this will lead to more “false alerts” than at present. If such alerts are accompanied by a cautionary note, they will promote intergroup discussions that could lead to more rapid identification of high-magnification events without compromising the credibility of the group.
We would like to thank Jaiyul Yoo and Dale Fields for their generous help. We thank Phil Yock, Ian Bond, Bohdan Paczyński and Juna Kollmeier for their insightful comments on the manuscript. S. D. and A. G. were supported by NSF grant AST 042758. D. D., A. G., and R. P. were supported by NASA grant NNG04GL51G. B. S. G. was supported by a Menzel Fellowship from the Harvard College Observatory. C. H. was supported by the SRC program of Korea Science & Engineering Foundation. B.-G. P. acknowledges support from the Korea Astronomy and Space Science Institute. Support for OGLE was provided by Polish MNII grant 2P03D02124, NSF grant AST-0204908, and NASA grant NAG5-12212. A. U. acknowledges support from the grant “Subsydium Profesorskie” of the Foundation for Polish Science. This publication makes use of data products from the Two Micron All Sky Survey, which is a joint project of the University of Massachusetts and the Infrared Processing and Analysis Center/California Institute of Technology, funded by the National Aeronautics and Space Administration and the National Science Foundation. Any opinions, findings, and conclusions or recommendations expressed in this material are those of the authors and do not necessarily reflect the views of the NSF.
## Appendix A Two New Finite-Source Algorithms
To model planetary light curves, we develop two new binary-lens finite-source algorithms. The first algorithm, called “map-making”, is the main work horse. For a fixed $`(b,q)`$ geometry, it can successfully evaluate the finite-source magnification of almost all data points on the light curve and can robustly identify those points for which it fails. The second algorithm, called “loop-linking”, is much less efficient than map-making but is entirely robust. We use loop-linking whenever the map-making routine decides it cannot robustly evaluate the magnification of a point. In addition, at the present time, map-making does not work for resonant lensing geometries, i.e., geometries for which the caustic has six cusps. For planetary mass ratios, resonant lensing occurs when the planet is very close to the Einstein radius, $`b1`$. We use loop-linking in these cases also.
### A.1 Map-Making
Map-making has two components: a core function that evaluates the magnification and a set of auxiliary functions that test whether the measurement is being made accurately. If a light-curve point fails these tests, it is sent to loop-linking.
Finite-source effects are important when the source passes over or close to a caustic. Otherwise, the magnification can be evaluated using the point-source approximation, which is many orders of magnitude faster than finite-source evaluations. Hence, the main control issue is to ensure that any point that falls close to a caustic is evaluated using a finite-source algorithm or at least is tested to determine whether this is necessary. For very high magnification events, the peak points will always pass close to the central caustic. Hence, the core function of map-making is to “map” an Einstein-ring annulus in the image plane that covers essentially all of the possible images of sources that come close to the central caustic. The method must also take account of the planetary caustic(s), but we address that problem further below.
We begin by inverse ray-shooting an annulus defined by $`A_{\mathrm{PSPL}}>A_{\mathrm{min}}`$, where $`A_{\mathrm{PSPL}}`$ is the Paczyński (1986) magnification due to a point source by a point lens and $`A_{\mathrm{min}}`$ is a suitably chosen threshold. For OGLE-2000-BLG-343, we find that $`A_{\mathrm{min}}=75`$ covers the caustic-approaching points in essentially all cases. The choice of the density of the ray-shooting map is described below. Each such “shot” results in a four-element vector $`(x_i,y_i,x_s,y_s)`$. We divide the portion of the source plane covered by this map into a rectangular grid with $`k=1,\mathrm{},N_g`$ elements. We choose the size of each element to be equal to the smallest source radius being evaluated by the map. Hence, each “shot” is assigned to some definite grid element $`k(x_s,y_s)`$. We then sort the “shots” by $`k`$. For each light-curve point to be evaluated, we first find the grid elements that overlap the source. We then read sequentially through the sorted file<sup>8</sup><sup>8</sup>8Whether this “file” should actually be an external disk file or an array in internal memory depends on both the size of the available internal memory and the total number of points evaluated in each lens geometry. In our case, we used internal arrays. from the beginning of the element’s “shots” to the end. For each “shot”, we ask whether its $`(x_s,y_s)`$ lies within the source. If so, we weight that point by the limb-darkened profile of the source at that radius. Note that a source with arbitrary shape and surface brightness profile would be done just as easily.
For each light-curve point, we first determine whether at least one of the images of the center of the source lies in the annulus. In practice, for our case, the points satisfying this condition are just those on the night of the peak, but for other events this would have to be determined on a point-by-point basis. We divide those points with images outside the annulus into two classes, depending on whether they lie inside or outside two or three rectangles that we “draw”, one around each caustic. Each rectangle is larger than the maximum extent of the caustics by a factor of $`1.5`$ in each direction. If the source center lies outside all of these rectangles, we assume that the point-source approximation applies and evaluate the magnification accordingly. If it lies inside one of the rectangles (and so either near or inside one of the caustics), we perform the following test to see whether the point-source approximation holds. We evaluate the point-source magnifications at five positions, namely, the source center $`A(0,0)`$, two positions along the source $`x`$-axis $`A(\pm \lambda \rho _{},0)`$, and two positions along the source y-axis $`A(0,\pm \lambda \rho _{})`$, where $`\lambda 1`$ is a parameter. We demand that
$$\left|\frac{A(\lambda \rho _{},0)+A(\lambda \rho _{},0)}{2A(0,0)}1\right|+\left|\frac{A(0,\lambda \rho _{})+A(0,\lambda \rho _{})}{2A(0,0)}1\right|<4\sigma ,$$
(A1)
where $`\sigma `$ is the maximum permitted error (defined in § A.3). We use a minimum of five values in order to ensure that the magnification pattern interior to the source is reasonably well sampled; the precise value is chosen empirically and is a compromise between computing speed and accuracy. We require the number of values to be at least $`2[\rho _{}/\sqrt{q}]`$ in order to ensure that small, well-localized perturbations interior to the source caused by low mass ratio companions are not missed.
If a point passes this test, the magnification pattern in the neighborhood of the point is adequately represented by a gradient and so the point-source approximation holds. Points failing this test are sent to loop linking.
The remaining points, those with at least one image center lying in the annulus, are almost all evaluated using the sorted grid as described above. However, we must ensure that the annulus really covers all of the images. We conduct several tests to this end.
First, we demand that no more than one of the three or five images of the source center lie outside the annulus. In a binary lens, there is usually one image that is associated with the companion and that is highly demagnified. Hence, it can generally be ignored, so the fact that it falls outside the annulus does not present a problem. If more than one image center lies outside the annulus, the point is sent to loop-linking. Second, it is possible that an image of the center of the source could lie inside the annulus, but the corresponding image of another point on the source lies outside. In this case, there would be some intermediate point that lay directly on the boundary. To guard against this possibility, we mark the “shots” lying within one grid step of the boundaries of the annulus, and if any of these boundary “shots” fall in the source, we send the point to loop-linking. Finally, it is possible that even though the center of the source lies outside the caustic (and so has only three images), there are other parts of the source that lie inside the caustic and so have two additional images. If these images lay entirely outside the annulus, the previous checks would fail. However, of necessity, some of these source points lie directly on the caustic, and so their images lie directly on the critical curve. Hence, as long as the critical curve is entirely covered by the annulus, at least some of each of these two new images will lie inside the annulus and the “boundary test” just mentioned can robustly determine whether any of these images extend outside the annulus. For each $`(b,q)`$ geometry, we directly check whether the annulus covers the critical curve associated with the central caustic by evaluating the critical curve locus using the algorithm of Witt (1990).
### A.2 Loop-Linking
Loop-linking is a hybrid of two methods: inverse ray-shooting and Stokes’s theorem. In the first method (which was also used above in “map-making”), one finds the source location corresponding to each point in the image plane. Those that fall inside the source are counted (and weighted according to the local surface brightness), while those that land outside the source are not. The main shortcoming of inverse ray-shooting is that one must ensure that the ensemble of “shots” actually covers the entire image of the source without covering so much additional “blank space” that the method becomes computationally unwieldy.
In the Stokes’s theorem approach, one maps the boundary of (a polygon-approximation of) the source into the image plane, which for a binary lens yields either three or five closed polygons. These image polygons form the (interior or exterior) boundaries of one to five images. If one assumes uniform surface brightness, the ratio of the combined areas of these images (which can be evaluated using Stokes’s theorem) to the area of the source polygon is the magnification.
There are two principal problems with the Stokes’s theorem approach. First, sources generally cannot be approximated as having uniform surface brightness. This problem can be resolved simply by breaking the source into a set of annuli, each of which is reasonably approximated as having uniform surface brightness. However, this multiplies the computation time by the number of annuli. Second, there can be numerical problems of several types if the source boundary passes over or close to a cusp. First, the lens solver, which returns the image positions given the source position, can simply fail in these regions. This at least has the advantage that one can recognize that there is a problem and perhaps try some neighboring points. The second problem is that even though the boundary of the source passes directly over a cusp, it is possible that none of the vertices of its polygon approximation lie within the caustic. The polygonal image boundary will then fail to surround the two new images of the source that arise inside the caustic, so these will not be included in the area of the image. Various steps can be taken to mitigate this problem, but the problem is most severe for very low-mass planets (which are of the greatest interest in the present context), so complete elimination of this problem is really an uphill battle.
The basic idea of loop-linking is to map a polygon that is slightly larger than the source onto the image plane, and then to inverse ray shoot the interior regions of the resulting image-plane polygons. This minimizes the image-plane region to be shot compared to other inverse-ray shooting techniques. It is, of course, more time-consuming than the standard Stokes’s theorem technique, but it can accommodate arbitrary surface-brightness profiles and is more robust. As we detail below, loop-linking can fail at any of several steps. However, these failures are always recognizable, and recovery from them is always possible simply by repeating the procedure with a slightly larger source polygon.
Following Gould & Gaucherel (1997), the vertices of the source polygon are each mapped to an array of three or five image positions, each with an associated parity. If the lens solver fails to return three or five image positions, the evaluation is repeated beginning with a larger source polygon. For each successive pair of arrays, we “link” the closest pair of images that has the same parity and repeat this process until all images in these two arrays are linked. An exception occurs when one array has three images and the other has five images, in which case two images are left unmatched. Repeating this procedure for all successive pairs of arrays produces a set of linked “strands”, each with either positive or negative parity. The first element of positive-parity strands and the last element of negative-element strands are labeled “beginnings” and the others are labeled “ends”. Then the closest “beginning” and “end” images are linked and this process is repeated until all “beginnings” and “ends” are exhausted. The result is a set of two to five linked loops. As with the standard Stokes’s theorem approach, it is possible that a source-polygon edge crosses a cusp without either vertex being inside the caustic. Then the corresponding image-polygon edge would pass inside the image of the source, which would cause us to underestimate the magnification. We check for this possibility by inverse ray shooting the image-polygon boundary (sampled with the same linear density as we later sample the images) back into the source plane. If any of these points land in the source, we restart the calculation with a larger source polygon.
We then use these looped links to efficiently locate the region in the image plane to do inverse ray shooting. We first examine all of the links to find the largest difference, $`\mathrm{\Delta }y_{\mathrm{max}}`$, between the y-coordinates of the two vertices of any link. Next, we sort the $`m=1,\mathrm{},n`$ links by the lower y-coordinate of their two vertices, $`y_m^{}`$. One then knows that the upper vertex obeys $`y_m^+y_m^{}+\mathrm{\Delta }y_{\mathrm{max}}`$. Hence, for each $`y`$-value of the inverse ray shooting grid, we know that only links with $`y_m^{}yy_m^{}+\mathrm{\Delta }y_{\mathrm{max}}`$ can intersect this value. These links can quickly be identified by reading through the sorted list from $`y_m^{}=y\mathrm{\Delta }y_{\mathrm{max}}`$ to $`y_m^{}=y`$. The $`x`$ value of each of these crossing links is easily evaluated. Successive pairs of $`x`$’s then bracket the regions (at this value of $`y`$) where inverse rays must be shot. As a check, we demand that the first of each of these bracketing links is an upward-going link and the second is a downward-going link.
### A.3 Algorithm Parameters
Before implementing the two algorithms described above, one must first specify values for certain parameters. Both algorithms involve inverse ray shooting and hence require that a sampling density by specified. Let $`g`$ be the grid size in units of the Einstein radius. For magnification $`A1`$, the image can be crudely approximated as two long strands whose total length is $`\mathrm{}=4A\rho _{}`$ and hence of mean width $`(\pi \rho _{}^2A)/\mathrm{}=(\pi /4)\rho _{}`$. If, for simplicity, we assume that the strand is aligned with the grid, then there will be a total of $`\mathrm{}/g`$ grid tracks running across the strand. Each will have two edges, and on each edge there will be an “error” of $`12^{1/2}`$ in the “proper” number of grid points due to the fact that this number must be an integer, whereas the actual distance across the strand is a real number. Hence, the total number of grid points will be in error by $`[(2\mathrm{}/g)/12]^{1/2}`$, while the total number itself is $`\pi (\rho _{}/g)^2A`$. This implies a fractional error $`\sigma `$,
$$\sigma ^2=\frac{[\pi (\rho _{}/g)^2A]^2}{(2\mathrm{}/g)/12}=\frac{3\pi ^2}{2}A(\rho _{}/g)^3.$$
(A2)
In fact, the error will be slightly smaller than given by equation (A2) in part because the strand is not aligned with the grid, so the total number of tracks across the strands will be lower than $`2\mathrm{}/g`$, and in part because the “discretization errors” at the boundary take place on limb-darkened parts of the star, which have lower surface brightness, so fluctuations here have lower impact. Hence, an upper limit to the grid size required to achieve a fractional error $`\sigma `$ is
$$\frac{g}{\rho _{}}=\left(\frac{3\pi ^2}{2}\right)^{1/3}\sigma ^{2/3}A^{1/3}.$$
(A3)
For each loop-linking point, we know the approximate magnification $`A`$ because we know the weighted parameters of the single-lens model. We generally set $`\sigma `$ at 1/3 of the measurement error, so that the (squared) numerical noise is an order of magnitude smaller than that due to observational error. Using equation (A3) we can then determine the grid size.
For the map-making method, the situation is slightly more complicated. Instead of evaluating one particular point as in the loop-linking method, all of the points on the source plane with the same $`(b,q,u_0,\rho )`$ are evaluated within one map. Therefore, the grid size for this map is the minimum value from equation (A3) to achieve the required accuracy for all of these points. We derive the following equation from equation (A3) to determine the grid size in the map-making method:
$$g=[\frac{3\pi ^2}{2}\frac{Q(u_0)}{u_0}]^{1/3}\rho _{}$$
(A4)
where
$$Q(u_0)=min_{A_i>75}\left\{\frac{F(t_i)F_\mathrm{b}(u_0)}{F_{\mathrm{max}}(t_i)F_b(u_0)}\sigma _i^2\right\}$$
(A5)
We find that $`Q(u_0)=2.65\times 10^7`$ is independent of $`u_0`$ for both the observational and simulated data of OGLE 2004-BLG-343. In principle, one could determine a minimum $`g`$ for all $`(u_0,\rho )`$ combinations and generate only one map for a given $`(b,q)`$ geometry, but this would render the calculation unnecessarily long for most $`(u_0,\rho )`$ combinations. Instead we evaluate $`g`$ for each $`(u_0,\rho )`$ pair and create several maps, one for each ensemble of $`(u_0,\rho )`$ pairs with similar $`g`$’s. The sizes of the ensembles should be set to minimize the total time spent generating, loading and employing maps. Hence, they will vary depending on the application.
## Appendix B MOA-2003-BLG-32/OGLE-2003-BLG-219
MOA-2003-BLG-32/OGLE-2003-BLG-219, with a peak magnification $`A_{\mathrm{max}}=525\pm 75`$ is most sensitive to low-mass planets to date (Abe et al., 2004). However, instead of fitting the simulated binary-lens light curves to single-lens models, Abe et al. (2004) obtain their $`\mathrm{\Delta }\chi ^2`$ by directly subtracting the $`\chi ^2`$ of a simulated binary-lens light curve from that of the light curve that is the best fit to the data. Since the source star of this event could reside in the Sagittarius dwarf galaxy, which makes the Galactic modeling rather complicated, we do not attempt to apply our entire method to this event. We calculate planet exclusion regions with the same $`\mathrm{\Delta }\chi ^2`$ thresholds (60 for $`q=10^3`$ and 40 for the $`q<10^3`$) as Abe et al. (2004) but using our method of obtaining $`\mathrm{\Delta }\chi ^2`$ by fitting the simulated binary-lens light curves to single-lens models. Figure 13 shows our results for the exclusion regions at planet-star mass ratios $`q=10^5,10^4,`$ and $`10^3`$ for MOA-2003-BLG-32/OGLE-2003-BLG-219. The exclusion region we have obtained at $`q=10^3`$ is about $`1/4`$ in vertical direction and $`1/9`$ in horizontal direction relative to the corresponding region in Abe et al. (2004), and the size of our exclusion region at $`10^4`$ is about 60% in each dimension relative to that in Abe et al. (2004). Although according to our analysis, Abe et al. (2004) overestimate the sensitivity of MOA-2003-BLG-32/OGLE-2003-BLG-219 to both Jupiter-mass and Neptune-mass planets, their estimates of sensitivity to Earth-mass planets are basically consistent with our results and MOA-2003-BLG-32/OGLE-2003-BLG-219 nevertheless retains the best sensitivity to planets to date. |
warning/0507/hep-ph0507294.html | ar5iv | text | # Generalized mirror matter models
July 2005
## Abstract
Non-minimal gauge models with exact unbroken improper space-time symmetries are constructed and their cosmological and astrophysical implications explored.
The exact parity modelflv is the minimal extension of the standard model which allows for an exact unbroken parity symmetry \[$`xx,tt`$\]. Each type of ordinary particle (lepton, quark, gauge particles) has a distinct mirror partner. The ordinary and mirror particles form parallel sectors each with gauge symmetry $`G_{SM}SU(3)_cSU(2)U(1)_Y`$, so that the gauge symmetry is $`G_{SM}G_{SM}`$. The interactions of each sector are governed by a Lagrangian of exactly the same form, except with left and right interchanged. That is, in the mirror sector, it is the right-handed chiral fermions which are $`SU(2)`$ doublets, while in the ordinary sector, it is the left-handed chiral fermions. Thus, the Lagrangian has the form:
$`=_{SM}(e_L,e_R,q_L,q_R,W_\mu ,A_\mu ,\mathrm{})+_{SM}(e_R^{},e_L^{},q_R^{},q_L^{},W_\mu ^{},A_\mu ^{},\mathrm{}.)+_{mix}`$ (1)
where $`e,q,W_\mu ,A_\mu `$ denote the leptons, quarks, gauge fields etc and the primed fields are their corresponding mirror partners. The $`_{mix}`$ part describes possible interactions coupling ordinary and mirror particles together.
The exact parity symmetry, $`𝒫`$, has the formflv :
$`xx,tt`$
$`G^\mu G_\mu ^{},W^\mu W_\mu ^{},B^\mu B_\mu ^{}`$
$`\mathrm{}_L\gamma _0\mathrm{}_R^{},e_R\gamma _0e_L^{},q_L\gamma _0q_R^{},`$
$`u_R\gamma _0u_L^{},d_R\gamma _0d_L^{},\varphi \varphi ^{}`$ (2)
where $`G^\mu ,W^\mu ,B^\mu `$ are the standard $`G_{SM}`$ gauge particles, $`\mathrm{}_L,e_R,q_L,u_R,d_R`$ are the standard leptons and quarks (the generation index is implicit) $`\varphi `$ is the standard Higgs doublet and the primes denote the mirror particles. Under this symmetry, $`_{SM}(e,q,W,A,\mathrm{})`$ interchanges with $`_{SM}(e^{},q^{},W^{},A^{},\mathrm{})`$ leaving the full Lagrangian invariant (of course the terms in $`_{mix}`$ must also be invariant under this symmetry). The theory also has an exact unbroken time reversal invariance, $`𝒯`$, with standard CPT identified as the product: $`𝒫𝒯`$flv . In this way, the full Poincare group, containing proper and improper Lorentz transformations, space-time translations etc becomes a fundamental unbroken symmetry – providing strong theoretical motivation for the theory.
Constraints from gauge invariance and renormalizability limit $`_{mix}`$ to just two termsflv :
$`_{mix}=ϵF^{\mu \nu }F_{\mu \nu }^{}+\lambda \varphi ^{}\varphi \varphi ^{}\varphi ^{}`$ (3)
where $`F^{\mu \nu }=^\mu B^\nu ^\nu B^\mu `$ \[$`F^{\mu \nu }=^\mu B^\nu ^\nu B^\mu `$\] is the $`U(1)`$ \[mirror $`U(1)`$\] field strength tensor. The two terms in $`_{mix}`$ provide an important means of experimentally testing this theory (for a recent review of these experimental implications, see Ref.sad ). Note that $`_{mix}`$ can contain other terms if there exists new particles, such as gauge singlet neutrinos. In particular the physics of neutrino mass generation may allow for neutrino - mirror neutrino mass mixing terms in $`_{mix}`$, which would provide another useful way to test the theoryflv2 .
Note that there is a large range of parameters of the Higgs potential for which mirror symmetry is not spontaneously broken by the vacuum (i.e. $`\varphi =\varphi ^{}`$) so that it is an exact, unbroken symmetry of the theory. Explicitly, the most general Higgs potential isflv
$`V=\lambda _1(\varphi ^{}\varphi u^2)^2+\lambda _1(\varphi ^{}\varphi ^{}u^2)^2+\lambda _2(\varphi ^{}\varphi \varphi ^{}\varphi ^{})^2`$ (4)
If $`\lambda _{1,2}>0`$ then $`V0`$. The minimum of the potential, $`V=0`$, occurs when $`\varphi =\varphi ^{}=u`$, demonstrating that the exact parity symmetry is not broken by the vacuum, as advertised.
Importantly, the theory predicts the existence of new particles which are necessarily massive and stable. There is a pressing need for such non-baryonic massive stable particles from astrophysical and cosmological considerations, and a significant amount of work has been done exploring the possibility that mirror matter is the inferred non-baryonic dark matter in the Universea0 ; a1 ; a ; b ; c ; d ; e ; fm ; g ; h ; i . For an up-to-date review, see ref.fr .
The purpose of this article is to examine an obvious generalization of the exact parity model. In the minimal exact parity model there is one ‘copy’ of the standard particles. It is possible that nature may have $`n`$ distinct copies of the standard particles. The copies may be mirror copies or ordinary copies.
Let us introduce the notation, (p,q) to denote $`p`$ ordinary sectors and $`q`$ mirror sectors. This means that there are $`n=p+q1`$ copies of the standard particles, $`q`$ of these of the mirror variety. We assume an exact parity symmetry, $`𝒫`$, which implies equal numbers of ordinary and mirror copies: $`p=q`$ (and hence $`n`$ is odd). If the masses and interactions of the particles in each sector are exactly the same as the standard particles (excepting, of course, that the mirror copies have left and right interchanged), then the Lagrangian would exhibit the discrete symmetry
$`C_pC_p𝒫`$ (5)
where $`C_p`$ is the group of permutations of $`p`$ objects.
These non-minimal mirror models are a straightforward generalization to the exact parity model of ref.flv . The Lagrangian generalizes Eq.(1) in the obvious way:
$`={\displaystyle \underset{i=1}{\overset{p}{}}}_{SM}(e_{iL},e_{iR},q_{iL},q_{iR},W_i^\mu ,A_i^\mu ,\mathrm{})+{\displaystyle \underset{j=1}{\overset{p}{}}}_{SM}(e_{jR}^{},e_{jL}^{},q_{jR}^{},q_{jL}^{},W_j^\mu A_j^\mu ,\mathrm{}.)+_{mix}`$ (6)
where we use the integer subscripts to label the particles from the $`p`$ ordinary sectors and primes plus integer subscripts to label their corresponding mirror partners. In this general case, $`_{mix}`$ has the form:
$`_{mix}=ϵ{\displaystyle \underset{i=1}{\overset{p}{}}}F_i^{\mu \nu }{\displaystyle \underset{j=1}{\overset{p}{}}}F_{j\mu \nu }^{}+ϵ^{}{\displaystyle \underset{k,l=1}{\overset{p}{}}}\left(F_k^{\mu \nu }F_{l\mu \nu }+F_k^{\mu \nu }F_{l\mu \nu }^{}\right)`$
$`+\lambda {\displaystyle \underset{i=1}{\overset{p}{}}}\varphi _i^{}\varphi _i{\displaystyle \underset{j=1}{\overset{p}{}}}\varphi _j^{}\varphi _j^{}+\lambda ^{}{\displaystyle \underset{k,l=1}{\overset{p}{}}}\left(\varphi _k^{}\varphi _k\varphi _l^{}\varphi _l+\varphi _k^{}\varphi _k^{}\varphi _l^{}\varphi _l^{}\right)`$ (7)
where $`kl`$ in the sums and $`F_i^{\mu \nu }^\mu B_i^\nu ^\nu B_i^\mu [F_i^{\mu \nu }^\mu B_i^\nu ^\nu B_i^\mu ]`$.
The most general Higgs potential is given by the straightforward generalization to Eq.(4):
$`V`$ $`=`$ $`\lambda _1{\displaystyle \underset{i=1}{\overset{p}{}}}\left\{[\varphi _i^{}\varphi _iu^2]^2+[\varphi _i^{}\varphi _i^{}u^2]^2\right\}+\lambda _2{\displaystyle \underset{i,j=1}{\overset{p}{}}}[\varphi _i^{}\varphi _i\varphi _j^{}\varphi _j^{}]^2`$ (8)
$`+`$ $`\lambda _3{\displaystyle \underset{i,j=1}{\overset{p}{}}}\left\{[\varphi _i^{}\varphi _i\varphi _j^{}\varphi _j]^2+[\varphi _i^{}\varphi _i^{}\varphi _j^{}\varphi _j^{}]^2\right\}`$
Again, if $`\lambda _{1,2,3}>0`$ then $`V0`$. The minimum of the Higgs potential is then $`V=0`$, which occurs only when each $`\varphi _i=\varphi _i^{}=u`$. This means that the discrete symmetry, Eq.(5), is not spontaneously broken, but is an exact unbroken symmetry of the theory.
One specific motivation for considering such non-minimal models comes from the similarity of the cosmological mass density of non-baryonic dark matter and ordinary matter. Precision measurements of the CMBR from WMAP and other data givecosmo :
$`\mathrm{\Omega }_bh^2=0.0224\pm 0.0009`$
$`\mathrm{\Omega }_mh^2=0.135_{0.009}^{+0.008}`$ (9)
where $`\mathrm{\Omega }_m=\mathrm{\Omega }_b+\mathrm{\Omega }_{dark}`$ is the total matter density (normalized to the critical matter density needed to close the Universe) and $`h`$ is the Hubble parameter measured in units of $`100`$ km/s/Mpc. This means that the cosmological mass density of non-baryonic dark matter is within an order of magnitude of the mass density of baryons:
$`\mathrm{\Omega }_{dark}/\mathrm{\Omega }_b=5.03\pm 0.46.`$ (10)
This interesting result can be explained in principle if the mass and interaction rates of the non-baryonic particles are very similar to ordinary baryons. The exact parity model is one simple and well defined extension of the standard model which has this featuremon ; mon2 . In fact, in the minimal exact parity model, we would expect $`\mathrm{\Omega }_{dark}=\mathrm{\Omega }_b`$ if the evolution of the universe were completely symmetric in the two sectors, i.e. there was no temperature difference between the ordinary and mirror particles during the baryogenesis epoch<sup>1</sup><sup>1</sup>1Actually, $`\mathrm{\Omega }_{dark}=\mathrm{\Omega }_b`$ could also occur – even if there was a temperature asymmetry during baryogenesis – if the baryonic asymmetry was generated by transitions between ordinary and mirror particles as in the scenario of Ref.mon .. During the big bang nucleosynthesis (BBN) epoch, however, the success of standard BBN suggests that $`T^{}`$ is less than $`T`$:
$`T^{}/T\stackrel{<}{}0.6\mathrm{at}T1\mathrm{MeV}`$ (11)
in order for the expansion rate of the universe to have been within an acceptable range. If this temperature asymmetry were induced<sup>2</sup><sup>2</sup>2 The origin of the temperature asymmetry is unknown, however some ideas have been put forward in Ref.inf in the context of inflation. The basic idea is to have an ‘ordinary inflaton’ coupling to ordinary matter, and a ‘mirror inflaton’ coupling to mirror matter. If inflation is triggered by some random fluctuation, then it can occur in the two sectors at different times, leading to $`TT^{}`$ after reheating in the two sectors. Naturally, the bulk of the inflation would be expected to occur prior to baryogenesis (so as not to dilute the baryon number), however this does not exclude the possibility of asymmetric reheating after baryogenesis, since the Universe may have gone through several reheating processes - depending on details such as the the number of weakly coupled scalar fields. before baryogenesis, then we expect $`\mathrm{\Omega }_{dark}\mathrm{\Omega }_b`$, the details would, of course, depend on the precise model of baryogenesis. Obviously if the temperatures of the two sectors were not too different this might explain why $`\mathrm{\Omega }_b\mathrm{\Omega }_{dark}`$ (see also Ref.mon ; mon2 for some related scenarios). Another logical possibility is that the temperature asymmetry required by BBN was induced after baryogenesis. In this case, the abundance of mirror baryons would be exactly the same as the abundance of ordinary baryons, and in the non-minimal mirror models with $`n`$ copies, this would generalize to:
$`\mathrm{\Omega }_{dark}/\mathrm{\Omega }_b=n.`$ (12)
This specific scenario is obviously testable and falsifiable, since it predicts that this ratio is an odd integer. The data are currently consistent with that, suggesting $`n=5`$. Thus we are led to consider the specific particle physics model consisting of the ordinary particles and 5 copies. This is compatible with our hypothesis of exact parity symmetry, implying three ordinary and three mirror sectors.
The successful dark matter features of the mirror matter model would also occur in these non-minimal models, including:
* It would elegantly explaina0 the MACHO population inferred to exist in the galactic halo from numerous microlensing observationsml of nearby galaxies.
* It would be capable of explaining the large scale structure of the universea ; b ; c .
* Provide a straightforward explanationfootdama of the positive annual modulation signal obtained in the DAMA/NaI direct detection experimentdama . Importantly, this explanation is consistent with the null results of the other direct detection experimentsfootdama2 ; sad .
* Explain various solar system anomalies, such as the anomalous accelerationpioneer of the two Pioneer spacecraftfvp , lack of small craters on the asteroid 433 Eros fm etc.
If mirror matter is the non-baryonic dark matter, as the above experiments and observations suggest, then the dark halos inferred to exist in spiral galaxies should be composed predominately of mirror matter. Microlensing studiesml suggest a halo composed of about $`20\%`$ mirror stars with the rest presumably in the form of an ionized mirror gas. Ionized gas, with a typical virial temperature of order $`T100`$ eV radiates energy at a rate per unit volume ofbook :
$`r_{cool}=n_e^{}^2\mathrm{\Lambda }`$ (13)
where $`n_e^{}`$ is the (free) mirror electron number density and $`\mathrm{\Lambda }`$ is a calculable function (which depends on cross section, temperature, composition etc). For a temperature of $`T100eV,\mathrm{\Lambda }10^{23}ergcm^3s^1`$ (see e.g. Ref.book ). In the case of the minimal mirror model, i.e. with $`n=1`$, the halo would have a mirror photon luminosity ofe :
$`L_{halo}={\displaystyle _{R_1}^{100kpc}}n_e^{}^2\mathrm{\Lambda }4\pi r^2𝑑r\left({\displaystyle \frac{3kpc}{R_1}}\right)10^{44}erg/s.`$ (14)
where $`R_1`$ is a phenomenological cutoff. If there are $`n`$ copies, and the particles in each sector are equally abundant in the halo of the galaxy, then $`L_{halo}`$ is reduced by $`n^2`$ for each type of photon. Thus, for $`n=5`$, the luminosity in each type of photon is only about a few times $`10^{42}`$ erg/s. The heating responsible for supporting the halo is not completely clear, however plausible candidates include mirror and/or ordinary supernova which can potentially supply the required energye .
Obviously, since the ordinary particles collapse and form a disk, while the mirror particles are roughly spherically distributed in spiral galaxies <sup>3</sup><sup>3</sup>3 Here, we use the term ‘mirror particles’ as an inclusive term for particles of the n-copies, whether they are mirror copies or not., the evolution is clearly asymmetric. Such an asymmetric evolution is plausible because the ordinary and mirror particles had different temperatures at the epoch of BBN, Eq.(11)<sup>4</sup><sup>4</sup>4 Actually in the case of 5-copies, the BBN bound is slightly more stringent, $`T^{}/T\stackrel{<}{}0.4`$, assuming a common temperature, $`T^{}`$ for each of the 5 copies.. This temperature asymmetry, not only leads to successful large scale structure formationa ; b ; c , but also implies that the chemical composition of the mirror worlds are quite different to the ordinary particle worlda . Specifically, the proportion of mirror helium ($`He^{}`$) to mirror hydrogen ($`H^{}`$) in each of the mirror sectors is expected to be significantly greater than the corresponding ordinary $`He/H`$ ratioa . Because of this major difference, the macroscopic evolution of the ordinary sector will be quite different to that of the mirror sectors. For example, the higher $`He^{}/H^{}`$ ratio implies that mirror stars evolve much faster (an order of magnitude or more) than ordinary starsbrecent , potentially giving a much greater rate of mirror supernova explosions. Evidently, the ordinary-mirror particle asymmetry required to explain a) BBN, b) Large scale structure formation, and c) the disparate distribution of ordinary and mirror particles within spiral galaxies might all be the result of an effective asymmetric boundary condition, with the microscopic interactions, as defined in the quantum field theoretic Lagrangian, remaining completely symmetric. Of course, understanding the complete details of galaxy formation (especially in the non-linear regime) is far from begin fully understood.
In conclusion, we have examined a straightforward generalization of the exact parity model involving $`n`$ copies of the standard particles (with $`n`$ odd, if the fundamental interactions respect an exact parity symmetry). Successful early universe cosmology requires a temperature asymmetry between the ordinary particles and their $`n`$ copies at the BBN epoch. Under the assumption that this temperature asymmetry arose after baryogenesis, the inferred non-baryonic dark matter density $`\mathrm{\Omega }_{dark}/\mathrm{\Omega }_b5`$ suggests that $`n=5`$. Although non-minimal, such models do preserve the successful dark matter features inherent in the minimal exact parity symmetric model and the additional prediction that $`\mathrm{\Omega }_{dark}/\mathrm{\Omega }_b`$ is an odd integer can be further tested by future cosmological observations.
Acknowledgements This work was supported by the Australian Research Council. The author would like to thank H. Georgi for pointing out some deficiencies in a previous version of this paper. |
warning/0507/hep-ex0507045.html | ar5iv | text | # Improved Measurements of Direct 𝐶𝑃 Violation in 𝐵→𝐾⁺𝜋⁻, 𝐾⁺𝜋⁰ and 𝜋⁺𝜋⁰ Decays
The Belle Collaboration
## Abstract
We report an improved measurement of direct $`CP`$ violation in the decay $`B^0K^+\pi ^{}`$, and a search for $`CP`$ violation in the decays $`B^+K^+\pi ^0`$ and $`B^+\pi ^+\pi ^0`$. The measured $`CP`$ violating asymmetries are $`𝒜_{CP}(K^+\pi ^{})=0.113\pm 0.022(\mathrm{stat}.)\pm 0.008(\mathrm{syst}.)`$, $`𝒜_{CP}(K^+\pi ^0)=0.04\pm 0.04(\mathrm{stat}.)\pm 0.02(\mathrm{syst}.)`$ and $`𝒜_{CP}(\pi ^+\pi ^0)=0.02\pm 0.08(\mathrm{stat}.)\pm 0.01(\mathrm{syst}.)`$, where the latter correspond to the intervals $`0.03<𝒜_{CP}(K^+\pi ^0)<0.11`$ and $`0.12<𝒜_{CP}(\pi ^+\pi ^0)<0.15`$ at 90% confidence level. These results are obtained from a data sample that contains 386 million $`B\overline{B}`$ pairs that was collected near the $`\mathrm{{\rm Y}}(4S)`$ resonance, with the Belle detector at the KEKB asymmetric energy $`e^+e^{}`$ collider. All of the results are preliminary.
preprint: BELLE-CONF-0523 LP2005-158 EPS05-495
In the Standard Model (SM) $`CP`$ violation arises via the interference of at least two diagrams with comparable amplitudes but different $`CP`$ conserving and violating phases. Mixing induced $`CP`$ violation in the $`B`$ sector has been established in $`bc\overline{c}s`$ transitions 2phi1 ; 2beta . Direct $`CP`$ violation is expected to be sizeable in the $`B`$ meson system BSS . The first experimental evidence for direct $`CP`$ violation was shown by Belle for the decay mode $`B^0\pi ^+\pi ^{}`$ PIPI . This result suggests large interference between tree and penguin diagrams and the existence of final state interactions FSI . In addition, both Belle belle\_acp\_250 and BABAR babar\_acp\_kpi\_230 have recently reported evidence for direct $`CP`$ violation in the decay $`B^0K^+\pi ^{}`$.
The partial rate $`CP`$ violating asymmetry is defined as:
$`𝒜_{CP}={\displaystyle \frac{N(\overline{B}\overline{f})N(Bf)}{N(\overline{B}\overline{f})+N(Bf)}},`$ (1)
where $`N(\overline{B}\overline{f})`$ is the yield for the $`\overline{B}K\pi /\pi \pi `$ decay and $`N(Bf)`$ denotes that of the charge-conjugate mode. Theoretical predictions with different approaches suggest that $`𝒜_{CP}(K^+\pi ^{})`$ could be either positive or negative acpth . Although there are large uncertainties related to hadronic effects in the theoretical predictions, results for $`𝒜_{CP}(K^+\pi ^{})`$ and $`𝒜_{CP}(K^+\pi ^0)`$ are expected to have the same sign and comparable magnitudes acpth . However, our previous measurements show that $`𝒜_{CP}(K^+\pi ^{})`$ and $`𝒜_{CP}(K^+\pi ^0)`$ are opposite in sign (although $`𝒜_{CP}(K^+\pi ^0)`$ is consistent with no asymmetry), and their central values are found to deviate from each other by $`2.4\sigma `$. These findings are consistent with those reported by BABAR babar\_acp\_kpi\_230 ; babar\_acp\_hpi0\_230 . It is suggested that the disagreement may be due to the contribution of the electroweak penguin process or other mechanisms anom . Therefore, it is important to verify whether the discrepancy persists with improved precision. In this Letter, we report $`𝒜_{CP}`$ measurements using 386 million $`B\overline{B}`$ pairs collected with the Belle detector at the KEKB $`e^+e^{}`$ asymmetric-energy (3.5 on 8 GeV) collider KEKB operating at the $`\mathrm{{\rm Y}}(4S)`$ resonance.
The Belle detector is a large-solid-angle magnetic spectrometer that consists of a silicon vertex detector (SVD), a 50-layer central drift chamber (CDC), an array of aerogel threshold Čerenkov counters (ACC), a barrel-like arrangement of time-of-flight scintillation counters (TOF), and an electro-magnetic calorimeter comprised of CsI(Tl) crystals (ECL) located inside a super-conducting solenoid coil that provides a 1.5 T magnetic field. An iron flux-return located outside of the coil is instrumented to detect $`K_L^0`$ mesons and to identify muons (KLM). The detector is described in detail elsewhere Belle . Two inner detector configurations were used. For the first sample of 152 million $`B\overline{B}`$ pairs (Set I), a 2.0 cm radius beampipe and a 3-layer silicon vertex detector were used; for the latter 234 million $`B\overline{B}`$ pairs (Set II), a 1.5 cm radius beampipe, a 4-layer silicon detector and a small-cell inner drift chamber were usedUshiroda .
The $`B`$ candidate selection is the same as that described in Ref. btohh . Charged tracks are required to originate from the interaction point (IP). Charged kaons and pions are identified using $`dE/dx`$ information and Cherenkov light yields in the ACC. The $`dE/dx`$ and ACC information are combined to form a $`K`$-$`\pi `$ likelihood ratio, $`(K\pi )=_K/(_K+_\pi )`$, where $`_{K/\pi }`$ is the likelihood of kaon/pion. Charged tracks with $`(K\pi )>0.6`$ are regarded as kaons and $`(K\pi )<0.4`$ as pions. Furthermore, charged tracks that are positively identified as electrons are rejected. The $`K/\pi `$ identification efficiencies and misidentification rates are determined from a sample of kinematically identified $`D^+D^0\pi ^+,D^0K^{}\pi ^+`$ decays. Table 1 shows the results. It is clear that the detection efficiency for $`K^{}/\pi ^+`$ is greater than for $`K^+/\pi ^{}`$; this small efficiency bias will be corrected in the $`𝒜_{CP}`$ measurement.
Candidate $`\pi ^0`$ mesons are reconstructed by combining two photons with invariant mass between 115 MeV/$`c^2`$ and 152 MeV/$`c^2`$, which corresponds to $`\pm 2.5`$ standard deviations. Each photon is required to have a minimum energy of 50 MeV in the barrel region ($`32^{}<\theta _\gamma <129^{}`$) or 100 MeV in the end-cap region ($`17^{}<\theta _\gamma <32^{}`$ or $`129^{}<\theta _\gamma <150^{}`$), where $`\theta _\gamma `$ denotes the polar angle of the photon with respect to the beam line. To further reduce the combinatorial background, $`\pi ^0`$ candidates with small decay angles ($`\mathrm{cos}\theta ^{}>0.95`$) are rejected, where $`\theta ^{}`$ is the angle between a $`\pi ^0`$ boost direction from the laboratory frame and its $`\gamma `$ daughters in the $`\pi ^0`$ rest frame.
Two variables are used to identify $`B`$ candidates: the beam-constrained mass, $`M_{\mathrm{bc}}=\sqrt{E_{\text{beam}}^2p_B^2}`$, and the energy difference, $`\mathrm{\Delta }E=E_B^{}E_{\text{beam}}^{}`$, where $`E_{\text{beam}}^{}`$ is the beam energy and $`E_B^{}`$ and $`p_B^{}`$ are the reconstructed energy and momentum of the $`B`$ candidates in the center-of-mass (CM) frame. Events with $`M_{\mathrm{bc}}>5.20`$ GeV/$`c^2`$ and $`0.3\mathrm{GeV}<\mathrm{\Delta }E<0.5\mathrm{GeV}`$ are selected for the final analysis.
The dominant background is from $`e^+e^{}q\overline{q}(q=u,d,s,c)`$ continuum events. To distinguish the signal from the jet-like continuum background, event topology variables and $`B`$ flavor tagging information are employed. We combine a set of modified Fox-Wolfram moments pi0pi0 into a Fisher discriminant. The probability density functions (PDF) for this discriminant, and the cosine of the angle between the $`B`$ flight direction and the $`z`$ axis, are obtained using signal and continuum Monte Carlo (MC) events. These two variables are then combined to form a likelihood ratio $`=_s/(_s+_{q\overline{q}})`$, where $`_{s(q\overline{q})}`$ is the product of signal ($`q\overline{q}`$) probability densities. Additional background discrimination is provided by $`B`$ flavor tagging. The standard Belle flavor tagging algorithm tagging gives two outputs: a discrete variable indicating the flavor of the tagging $`B`$ and the MC-determined dilution factor $`r`$, which ranges from zero for no flavor information to unity for unambiguous flavor assignment. An event that contains a lepton ($`r`$ close to unity) is more likely to be a $`B\overline{B}`$ event so a looser $``$ requirement can be applied. We divide the data into $`r>0.5`$ and $`r<0.5`$ regions. The continuum background is reduced by applying a selection requirement on $``$ for events in each $`r`$ region of Set I and Set II according to the figure-of-merit defined as $`N_s^{exp}/\sqrt{N_s^{exp}+N_{q\overline{q}}^{exp}}`$, where $`N_s^{exp}`$ denotes the expected signal yields based on our previous branching fraction measurements btohh and $`N_{q\overline{q}}^{exp}`$ denotes the expected $`q\overline{q}`$ yields from sideband data ($`M_{\mathrm{bc}}<5.26`$ GeV/$`c^2`$). A typical requirement suppresses 92–99% of the continuum background while retaining 48–67% of the signal.
Backgrounds from $`\mathrm{{\rm Y}}(4S)B\overline{B}`$ events are investigated using a large MC sample. After the $``$ requirement, we find a small charmless three-body background at low $`\mathrm{\Delta }E`$, and reflections from other $`B^0\pi ^+\pi ^{}`$ decays due to $`K`$-$`\pi `$ misidentification.
The signal yields are extracted by applying unbinned two dimensional maximum likelihood (ML) fits to the ($`M_{\mathrm{bc}}`$, $`\mathrm{\Delta }E`$) distributions of the $`B`$ and $`\overline{B}`$ samples. The likelihood for each mode is defined as
$``$ $`=`$ $`\mathrm{exp}({\displaystyle \underset{s,k,j}{}}N_{s,k,j})\times {\displaystyle \underset{i}{}}({\displaystyle \underset{s,k,j}{}}N_{s,k,j}𝒫_{s,k,j,i})`$ (2)
$`𝒫_{s,k,j,i}`$ $`=`$ $`{\displaystyle \frac{1}{2}}[1q_i𝒜_{CP}{}_{j}{}^{}]P_{s,k,j}(M_{\mathrm{bc}i},\mathrm{\Delta }E_i),`$ (3)
where $`s`$ indicates Set I or Set II, $`k`$ distinguishes events in the $`r<0.5`$ or $`r>0.5`$ regions, $`i`$ is the identifier of the $`i`$-th event, $`P(M_{\mathrm{bc}},\mathrm{\Delta }E)`$ is the two-dimensional PDF, $`q`$ indicates the reconstructed $`B`$ meson flavor, $`B(q=+1)`$ or $`\overline{B}(q=1)`$, $`N_j`$ is the number of events for the category $`j`$, which, in turn, corresponds to either signal, $`q\overline{q}`$ continuum, a reflection due to $`K`$-$`\pi `$ misidentification, or background from other charmless three-body $`B`$ decays.
The yields and asymmetries for the signal and backgrounds are allowed to float in all modes. Since the $`K^+\pi ^0`$ and $`\pi ^+\pi ^0`$ reflections are difficult to distinguish with $`\mathrm{\Delta }E`$ and $`M_{\mathrm{bc}}`$, we fit these two modes simultaneously with a fixed reflection-to-signal ratio based on the measured $`K`$-$`\pi `$ identification efficiencies and fake rates. All the signal PDFs ($`P(M_{\mathrm{bc}},\mathrm{\Delta }E)`$) are obtained using MC simulations based on the Set I and Set II detector configurations. The same signal PDFs are used for events in two different $`r`$ regions. No strong correlations between $`M_{\mathrm{bc}}`$ and $`\mathrm{\Delta }E`$ are found for the $`BK^+\pi ^{}`$ signal. Therefore, its PDF is modeled by a product of a single Gaussian for $`M_{\mathrm{bc}}`$ and a double Gaussians for $`\mathrm{\Delta }E`$. For the modes with a $`\pi ^0`$ meson in the final state, there are correlations between $`M_{\mathrm{bc}}`$ and $`\mathrm{\Delta }E`$ in the tails of the signals; hence, their PDFs are described by smoothed two-dimensional histograms. Discrepancies between the signal peak positions in data and MC are calibrated using $`B^+\overline{D}{}_{}{}^{0}\pi _{}^{+}`$ decays, where the $`\overline{D}{}_{}{}^{0}K^+\pi ^{}\pi ^0`$ sub-decay is used for the modes with a $`\pi ^0`$ meson while $`\overline{D}{}_{}{}^{0}K^+\pi ^{}`$ is used for the $`K^+\pi ^{}`$ mode. The MC-predicted $`\mathrm{\Delta }E`$ resolutions are verified using the invariant mass distributions of high momentum $`D`$ mesons. The decay mode $`\overline{D}{}_{}{}^{0}K^+\pi ^{}`$ is used for $`B^0K^+\pi ^{}`$, and $`\overline{D}{}_{}{}^{0}K^+\pi ^{}\pi ^0`$ for the modes with a $`\pi ^0`$ in the final state. The parameters that describe the shape of the signal PDFs are fixed in all of the fits.
The continuum background in $`\mathrm{\Delta }E`$ is described by a first or second order polynomial while the $`M_{\mathrm{bc}}`$ distribution is parameterized by an Argus function $`f(x)=x\sqrt{1x^2}\mathrm{exp}[\xi (1x^2)]`$, where $`x`$ is $`M_{\mathrm{bc}}`$ divided by half of the total center of mass energy. The continuum PDF is the product of an Argus function and a polynomial, where parameters $`\xi `$ and the coefficients of the polynomial are free parameters. These free parameters are $`r`$-dependent. A large MC sample is used to investigate the background from charmless $`B`$ decays and a smoothed two-dimensional histogram is taken as the PDF. The functional forms of the PDFs are the same for the $`B`$ and $`\overline{B}`$ samples.
The efficiency of particle identification is slightly different for positively and negatively charged particles; consequently, the parameter $`𝒜_{CP}`$ in Eq.3 becomes an effective partial rate asymmetry. For the $`K^+\pi ^0`$ and $`\pi ^+\pi ^0`$ modes, this raw asymmetry can be expressed as:
$`𝒜_{CP}^{\mathrm{raw}}={\displaystyle \frac{𝒜_ϵ+𝒜_{CP}}{1+𝒜_ϵ𝒜_{CP}}},`$ (4)
where $`𝒜_{CP}`$ is the true partial rate asymmetry and the efficiency asymmetry $`𝒜_ϵ`$ is the efficiency difference between $`K^{}(\pi ^+`$) and $`K^+(\pi ^{}`$) divided by the sum of their efficiency. The situation is more complicated for the $`K^+\pi ^{}`$ mode because, in addition to the bias due to the efficiency difference between $`K^{}\pi ^+`$ and $`K^+\pi ^{}`$, a $`K^{}\pi ^+`$ signal event can be doubly misidentified as a $`K^+\pi ^{}`$ candidate and dilute $`𝒜_{CP}`$. The efficiency asymmetry results in a $`𝒜_{CP}`$ bias of +0.01, while the small dilution factor due to double misidentification reduces the $`𝒜_{CP}`$ by a factor of 0.98.
Table 2 shows the signal yields and $`𝒜_{CP}`$ values for each mode. The asymmetries for the background components are consistent with zero within errors. Projections of the fits are shown in Figs.1-3. The systematic errors from fitting are estimated from the deviations in $`𝒜_{CP}`$ after varying each parameter of the signal PDFs by 1 standard deviation. The uncertainty in modeling the three-body background is studied by excluding the low $`\mathrm{\Delta }E`$ region ($`<0.12`$ GeV) and repeating the fit. Systematic uncertainty due to particle identification is estimated by repeating the fit after varying the $`K/\pi `$ efficiencies and fake rates by 1 standard deviation. At each step, the deviation in $`𝒜_{CP}`$ is added in quadrature to provide the systematic errors, which are less than 0.01 for all modes. A possible bias from the fitting procedure is checked in MC and a bias due to the $``$ requirement is investigated using the $`B^+\overline{D}{}_{}{}^{0}\pi _{}^{+}`$ samples. No significant bias is observed. The systematic uncertainties due to the detector bias are tested using the fit results for the continuum background listed in Table 2. We find a small background asymmetry dependence on the $``$ requirement for the $`K^+\pi ^{}`$ mode, and assign the uncertainty from the fit result of the $`B^+\overline{D}{}_{}{}^{0}\pi _{}^{+}`$($`\overline{D}{}_{}{}^{0}K^+\pi ^{}`$) sample ($`\pm 0.007`$) as the systematic uncertainty due to detector bias. The final systematic errors are then obtained by quadratically summing the errors due to the detector bias and the fitting systematics.
The partial rate asymmetry $`𝒜_{CP}(K^+\pi ^{})`$ is found to be $`0.113\pm 0.022\pm 0.008`$. The significance including the effect of systematic uncertainty is 4.97$`\sigma `$. This result supersedes our previous measurement belle\_acp\_250 and remains consistent with the value reported by BABAR, $`𝒜_{CP}(K^+\pi ^{})=0.133\pm 0.030\pm 0.009`$ babar\_acp\_kpi\_230 . The observed $`𝒜_{CP}(K^+\pi ^0)`$ value is consistent with zero at the current level of statistical precision. The difference between the results for $`𝒜_{CP}(K^+\pi ^{})`$ and $`𝒜_{CP}(K^+\pi ^0)`$ persists; their central values differ by $`3.1\sigma `$. This suggests a possible contribution from the electroweak penguin process or other mechanisms anom . No evidence of direct $`CP`$ violation is observed in the decay $`B^+\pi ^+\pi ^0`$. We set 90% C.L. intervals: $`0.03<𝒜_{CP}(K^+\pi ^0)<0.11`$ and $`0.12<𝒜_{CP}(\pi ^+\pi ^0)<0.15`$. All of the above results are preliminary.
We thank the KEKB group for the excellent operation of the accelerator, the KEK cryogenics group for the efficient operation of the solenoid, and the KEK computer group and the National Institute of Informatics for valuable computing and Super-SINET network support. We acknowledge support from the Ministry of Education, Culture, Sports, Science, and Technology of Japan and the Japan Society for the Promotion of Science; the Australian Research Council and the Australian Department of Education, Science and Training; the National Science Foundation of China under contract No. 10175071; the Department of Science and Technology of India; the BK21 program of the Ministry of Education of Korea and the CHEP SRC program of the Korea Science and Engineering Foundation; the Polish State Committee for Scientific Research under contract No. 2P03B 01324; the Ministry of Science and Technology of the Russian Federation; the Ministry of Higher Education, Science and Technology of the Republic of Slovenia; the Swiss National Science Foundation; the National Science Council and the Ministry of Education of Taiwan; and the U.S. Department of Energy. |
warning/0507/hep-ex0507029.html | ar5iv | text | # 1 INTRODUCTION
## 1 INTRODUCTION
We report on the preliminary measurement of the branching fraction $`B^0a_1^+(1260)\pi ^{}`$ with $`a_1^+(1260)\pi ^+\pi ^+\pi ^{}`$ . The $`a_1(1260)3\pi `$ decay proceeds mainly through the intermediate states $`(\pi \pi )_\rho \pi `$ and $`(\pi \pi )_\sigma \pi `$ .
The study of this decay mode is complicated by open questions on the parameters of the $`a_1(1260)`$ meson. There are large discrepancies between these parameters when comparing results from analyses involving hadronic interactions and $`\tau `$ decays . Therefore, it is important to verify the theoretical prediction of the branching fraction for this decay mode and have new measurements of the $`a_1(1260)`$ parameters. A theoretical calculation of the branching fraction of this decay mode has been made by Bauer, Stech and Wirbel (BSW) within the framework of the factorisation model. They predict a value of $`38\times 10^6`$, assuming $`\left|\frac{V_{ub}}{V_{cb}}\right|`$ = 0.08. It is also important to note that the $`B^0a_1^+(1260)\pi ^{}`$ channel can be used to measure the Cabibbo-Kobayashi-Maskawa angle $`\alpha `$ of the Unitarity triangle . We presented a preliminary version of this analysis at ICHEP’04 , using an integrated luminosity of $`112fb^1`$ and the measured branching fraction was $`(42.6\pm 4.2\pm 4.1)`$$`\times 10^6`$. For the branching fraction of $`B^0a_1^+(1260)\pi ^{}`$ an upper limit of $`49\times 10^5`$ at the 90% confidence level (C.L.) has been set by CLEO collaboration while the DELPHI collaboration has set the 90% C.L. upper limit of $`28\times 10^5`$ for the branching fraction of $`B^04\pi `$.
Below we present the details of the analysis for the measurement of the branching fraction for $`B^0a_1^+(1260)\pi ^{}2\pi ^+2\pi ^{}`$. Presently, we do not distinguish between the final states $`(\pi \pi )_\rho \pi `$ and $`(\pi \pi )_\sigma \pi `$. Such an analysis would require a study of the angular distributions of the decay products. Possible background contributions from $`B^0`$ decays to $`a_2^+(1320)\pi ^{}`$ and $`\pi ^+(1300)\pi ^{}`$ are studied and taken into account while in the preliminary version presented at ICHEP’04 they were neglected.
## 2 THE BABAR DETECTOR AND DATASET
The results presented in this paper are based on data collected in 1999–2004 with the BABAR detector at the PEP-II asymmetric $`e^+e^{}`$ collider located at the Stanford Linear Accelerator Center. An integrated luminosity of 198 fb<sup>-1</sup>, corresponding to 218 million $`B\overline{B}`$ pairs, was recorded at the $`\mathrm{{\rm Y}}(4S)`$ resonance (“on-resonance”, center-of-mass energy $`\sqrt{s}=10.58\mathrm{GeV}`$). An additional 15 fb<sup>-1</sup> were taken about 40 MeV below this energy (“off-resonance”) for the study of continuum background in which a light or charm quark pair is produced instead of an $`\mathrm{{\rm Y}}(4S)`$.
The asymmetric beam configuration in the laboratory frame provides a boost of $`\beta \gamma =0.56`$ to the $`\mathrm{{\rm Y}}(4S)`$. Charged particles are detected and their momenta measured by the combination of a silicon vertex tracker, consisting of five layers of double-sided silicon microstrip detectors, and a 40-layer central drift chamber, both operating in the 1.5-T magnetic field of a solenoid. The tracking system covers 92% of the solid angle in the center-of-mass frame.
Charged-particle identification is provided by the average energy loss ($`\mathrm{d}E/\mathrm{d}x`$) in the tracking devices and by an internally reflecting ring-imaging Cherenkov detector (DIRC) covering the central region. A $`K/\pi `$ separation of better than four standard deviations ($`\sigma `$) is achieved for momenta below 3 $`\mathrm{GeV}/c`$, decreasing to 2.5 $`\sigma `$ at the highest momenta in the $`B`$ decay final states. Photons and electrons are detected by a CsI(Tl) electromagnetic calorimeter while muons are identified in the magnetic flux return system.
## 3 ANALYSIS METHOD
Monte Carlo (MC) simulations of the signal decay mode, of continuum and $`B\overline{B}`$ backgrounds are used to establish the event selection criteria. We make several particle identification requirements to ensure the identity of all signal pions. For the bachelor charged track we require an associated DIRC Cherenkov angle between $`2\sigma `$ and $`+5\sigma `$ from the expected value for a pion. A $`B`$ meson candidate is characterized kinematically by the energy-substituted mass $`m_{\mathrm{ES}}=\sqrt{(\frac{1}{2}s+𝐩_0𝐩_B)^2/E_0^2𝐩_B^2}`$ and energy difference $`\mathrm{\Delta }E=E_B^{}\frac{1}{2}\sqrt{s}`$, where the subscripts $`0`$ and $`B`$ refer to the initial $`\mathrm{{\rm Y}}(4S)`$ and to the $`B`$ candidate in the lab-frame, respectively, and the asterisk denotes the $`\mathrm{{\rm Y}}(4S)`$ frame. We require $`|\mathrm{\Delta }E|0.2`$ GeV and $`5.25m_{\mathrm{ES}}5.29\mathrm{GeV}/c^2`$. We select $`a_1^+(1260)`$ candidates with the following requirement on the invariant mass: $`0.8<m_{a_1}<1.8`$ $`\mathrm{GeV}/c^2`$. The intermediate dipion state is required to have an invariant mass between 0.51 and 1.1 $`\mathrm{GeV}/c^2`$. The momentum of $`a_1^+(1260)`$ in the center-of-mass frame is required to be between 2.3 and 2.7 $`\mathrm{GeV}/c`$. To reduce fake $`B`$ meson candidates we require p($`\chi ^2`$) $`>`$ 0.01 for the $`B`$ vertex fit. The angular variable $`_{a_1}`$ (cosine of the angle between the direction of the bachelor $`\pi `$ and the flight direction of the $`B`$ in the $`a_1(1260)`$ meson rest frame) is required to be between $`0.85`$ and $`0.85`$ to suppress combinatorics.
To reject continuum background, we make use of the angle $`\theta _T`$ between the thrust axis of the $`B`$ candidate and that of the rest of the tracks and neutral clusters in the event, calculated in the center-of-mass frame. The distribution of $`\mathrm{cos}\theta _T`$ is sharply peaked near $`\pm 1`$ for combinations drawn from jet-like $`q\overline{q}`$ pairs and is nearly uniform for the isotropic $`B`$ meson decays; we require $`|\mathrm{cos}\theta _T|<0.65`$. The remaining continuum background is modelled from “off-resonance” data. We use Monte Carlo simulations of $`B^0\overline{B}^0`$ and $`B^+B^{}`$ decays to look for $`B\overline{B}`$ backgrounds, which can come from both charmless and charm decays. We find that the decay mode $`B^0D^{}\pi ^+`$, with $`D^{}K^+\pi ^{}\pi ^{}`$ and $`D^{}K_S^0\pi ^{}`$, is the only significant background. It is included in the maximum likelihood fit. Final results have been corrected for a small background contribution due to charmless decays.
We use an unbinned multivariate maximum-likelihood fit to extract the signal yields for
$`B^0a_1^+(1260)\pi ^{}`$. The likelihood function incorporates five variables. We describe the $`B`$ decay kinematics using: $`\mathrm{\Delta }E`$, $`m_{\mathrm{ES}}`$, $`m_{a_1}`$, a Fisher discriminant $``$ , and an angular variable A. The Fisher discriminant combines four variables: the angles in the $`\mathrm{{\rm Y}}(4S)`$ frame of the $`B`$ momentum and $`B`$ thrust axis with respect to the beam axis, and the zeroth and second angular moments $`L_{0,2}`$ of the energy flow around the $`B`$ thrust axis. The moments are defined by
$$L_j=\underset{i}{}p_i\left|\mathrm{cos}\theta _i\right|^j,$$
(1)
where $`p_i`$ is the momentum of particle $`i`$, $`\theta _i`$ is the angle between the direction of particle $`i`$ and the trust axis of the B candidate and the sum excludes tracks and clusters used to build the $`B`$ candidate. We have used an angular variable A in order to distinguish $`a_1^+(1260)\pi ^{}`$ from $`a_2^+(1320)\pi ^{}`$ and $`\pi ^+(1300)\pi ^{}`$. If X is our resonance $`a_1(J^P=1^+)`$, $`a_2(J^P=2^+)`$ or $`\pi (1300)(J^P=0^{})`$ that decays into three pions, we evaluate in the X meson rest frame the cosine of the angle between the normal to the plane of the three pions and the flight direction of the bachelor pion. Since we have on average 1.5 B candidates per event, we choose the best one using a $`\chi ^2`$ quantity computed with the $`\rho `$ mass. Since the maximum correlation between the observables in the selected data is 4%, we take the probability density function (PDF) for each event to be a product of the PDFs for the separate observables. The product PDF for event $`i`$ and hypothesis $`j`$, where $`j`$ can be signal (3 types), continuum background or $`B\overline{B}`$ background, is given by
$$𝒫_j^i=𝒫_j(m_{\mathrm{ES}})𝒫_j(\mathrm{\Delta }E)𝒫_j()𝒫_j(m_{a_1})𝒫_j(A).$$
(2)
There is the possibility that a track from a signal candidate is exchanged with a track from the rest of the event. We call these events “self-cross-feed” (SCF) events. The fraction of SCF events with respect to the total number of signal events for each type $`k`$ of signal, $`f_{SCF_k}`$, is fixed to the value found with Monte Carlo signal events (26%). The likelihood function for the event $`i`$ is defined as :
$$^i=\underset{k=1}{\overset{3}{}}\left(n_k(1f_{SCF_k})𝒫_k^i+n_kf_{SCF_k}𝒫_{SCF_k}^i\right)+n_{q\overline{q}}𝒫_{q\overline{q}}^i+n_{B\overline{B}1}𝒫_{B\overline{B}1}^i+n_{B\overline{B}2}𝒫_{B\overline{B}2}^i,$$
(3)
where $`n_k(k=1,3)`$ is the yield for $`a_1^+(1260)\pi ^{}`$, $`a_2^+(1320)\pi ^{}`$, and $`\pi ^+(1300)\pi ^{}`$ respectively, $`n_{q\overline{q}}`$ the number of continuum background events, $`n_{B\overline{B}1}`$ the number of $`B\overline{B}`$ background events $`D^{}\pi ^+`$ with $`D^{}K^+\pi ^{}\pi ^{}`$ and $`n_{B\overline{B}2}`$ the number of $`B\overline{B}`$ background events $`D^{}\pi ^+`$ with $`D^{}K_S^0\pi ^{}`$. The extended likelihood function for all events is :
$$=\frac{\mathrm{exp}(_jn_j)}{N!}\underset{i}{\overset{N}{}}\underset{j}{}n_j𝒫_j^i,$$
(4)
where $`n_j`$ is the yield of events of hypothesis $`j`$ found by the fitter, and $`N`$ is the number of events in the sample. The first factor takes into account the Poisson fluctuations in the total number of events.
We determine the PDFs for signal and $`B\overline{B}`$ backgrounds from MC distributions in each observable. For the continuum background we establish the functional forms and initial parameter values of the PDFs with off-resonance data. We allow the signal $`a_1(1260)`$ PDF parameters and the most important $`q\overline{q}`$ background PDF parameters to float in the final fit. The distributions of invariant mass of $`a_1(1260)`$, $`a_2(1320)`$ and $`\pi (1300)`$ in signal events are parameterized as relativistic Breit-Wigner line-shapes with a mass dependent width which takes into account the effect of the mass dependent $`\rho `$ width. The $`m_{\mathrm{ES}}`$ and $`\mathrm{\Delta }E`$ distributions for signal are parameterized as double gaussian functions. Slowly varying distributions are parameterized by linear functions. The combinatoric background in $`m_{\mathrm{ES}}`$ is described by a phase-space-motivated empirical function . We model the $``$ distribution using a Gaussian function with different widths above and below the mean. The A distributions are modelled using Gaussians in $`a_1^+(1260)\pi ^{}`$and polynomials in $`a_2^+(1320)\pi ^{}`$ and $`\pi ^+(1300)\pi ^{}`$.
## 4 RESULTS
We present the measurement of the branching fraction of the $`B`$ decay to $`a_1^+(1260)\pi ^{}`$, considering $`a_2^+(1320)\pi ^{}`$ and $`\pi ^+(1300)\pi ^{}`$ as sources of background. By generating and fitting simulated samples of signal and background events, we verify that our fitting procedure is working properly. We find that the minimum $`\mathrm{ln}`$ value for the on-resonance data lies well within the $`\mathrm{ln}`$ distribution from these simulated samples. Fits to data show no evidence of $`\pi ^+(1300)\pi ^{}`$, since a negative yield is obtained for this resonance. For this reason the $`\pi ^+(1300)\pi ^{}`$ component has been left out in final fits to the yields.
The reconstruction efficiency is obtained from the fraction of signal MC events passing the selection criteria once corrected for a bias detected in the fit yield. This bias (about 6%) is determined from fits to simulated samples, each equal in size to the data and containing a known number of signal MC events combined with events generated from the background PDFs.
The fitted values of the $`a_1(1260)`$ parameters are: $`m_{a_1}=1.22\pm 0.02`$ $`\mathrm{GeV}/c^2`$ and $`\mathrm{\Gamma }_{a_1}=0.423\pm 0.050`$ $`\mathrm{GeV}/c^2`$ . In Table 1 we show the results of the fits for on-resonance data. The statistical error on the number of events is taken to be the change in the central value when the quantity $`2\mathrm{ln}`$ changes by one unit. The statistical significance is taken as the square root of the difference between the value of $`2\mathrm{ln}`$ for zero signal and the value at its minimum. In Fig. 1 we show the $`m_{\mathrm{ES}}`$, $`\mathrm{\Delta }E`$, $`m_{a_1}`$ and A projections made by selecting events with a signal likelihood (computed without the variable shown in the figure) exceeding a threshold that optimizes the expected sensitivity.
## 5 SYSTEMATIC STUDIES
Most of the systematic errors on the yields that arise from uncertainties in the values of the PDF parameters have already been incorporated into the overall statistical error, since they are floated in the fit. We determine the sensitivity to the other parameters of the signal PDF components by varying these within their uncertainties. The result is shown in the first row of Table 2. This is the only systematic error on the fit yield; the other systematics apply to either the efficiency or the number of $`B\overline{B}`$pairs in the data sample.
The uncertainty in our knowledge of the efficiency is found to be 0.8$`N_t`$%, where $`N_t`$ is the number of signal tracks. We estimate the uncertainty in the number of $`B\overline{B}`$pairs to be 1.1%. The fitting algorithm introduces a systematic bias of 2.8%, which was found from fits to simulated samples with varying background populations. Published world averages provide the $`B`$ daughter branching fraction uncertainties. The systematic error from $`a_1(1260)K`$ cross-feed background is estimated to be 1.4%, while the systematic error due to SCF is found to be 3.5%. We also take into account systematic differences between data and MC for the $`\mathrm{cos}\theta _\mathrm{T}`$ selection (1.8%) and the possibility of interference between the $`a_1`$ and $`a_2`$ amplitudes (4%). The values for each of these contributions are given in Table 2.
## 6 SUMMARY
We have obtained a preliminary measurement of the branching fraction for $`B^0`$ meson decays to $`a_1^+(1260)\pi ^{}`$ with $`a_1^+(1260)`$ $`\pi ^+\pi ^+\pi ^{}`$. The measured branching fraction is:
$$(B^0a_1^+(1260)\pi ^{})=(40.2\pm 3.9\pm 3.9)\times 10^6$$
(5)
The fitted values of the $`a_1(1260)`$ parameters are: $`m_{a_1}=1.22\pm 0.02`$ $`\mathrm{GeV}/c^2`$ and $`\mathrm{\Gamma }_{a_1}=0.423\pm 0.050`$ $`\mathrm{GeV}/c^2`$.
## 7 ACKNOWLEDGMENTS
We are grateful for the extraordinary contributions of our PEP-II colleagues in achieving the excellent luminosity and machine conditions that have made this work possible. The success of this project also relies critically on the expertise and dedication of the computing organizations that support BABAR. The collaborating institutions wish to thank SLAC for its support and the kind hospitality extended to them. This work is supported by the US Department of Energy and National Science Foundation, the Natural Sciences and Engineering Research Council (Canada), Institute of High Energy Physics (China), the Commissariat à l’Energie Atomique and Institut National de Physique Nucléaire et de Physique des Particules (France), the Bundesministerium für Bildung und Forschung and Deutsche Forschungsgemeinschaft (Germany), the Istituto Nazionale di Fisica Nucleare (Italy), the Foundation for Fundamental Research on Matter (The Netherlands), the Research Council of Norway, the Ministry of Science and Technology of the Russian Federation, and the Particle Physics and Astronomy Research Council (United Kingdom). Individuals have received support from CONACyT (Mexico), the A. P. Sloan Foundation, the Research Corporation, and the Alexander von Humboldt Foundation. |
warning/0507/quant-ph0507069.html | ar5iv | text | # Distributing a multiparticle state by entanglement swapping
## Abstract
Using entanglement swapping, we construct a scheme to distribute an arbitrary multiparticle state to remote receivers. Only Bell states and two-qubit collective measurements are required.
The superposition principle is the most novel and least understood feature of the quantum theory. Quantum entanglement arises from the superposition of multiparticle states. For particles located far apart from one another, quantum entanglement gives rise to the mysterious phenomenon of non-local correlation to which there exists no local classical explanation EPR ; Bell ; CHSH ; Aspect . With the advent of quantum information science in recent years, instead of remaining merely a mystery to be solved, quantum entanglement has become a valuable and often indispensable resource in every branch of the field.
A bipartite state $`|\psi _{1,2}`$ is entangled if it cannot be written in a product form:
$$|\psi _{1,2}=|\psi _1^{}|\psi _2^{\prime \prime }.$$
(1)
The two-qubit Bell (or EPR) states,
$`|\varphi _{1,2}^\pm ={\displaystyle \frac{1}{\sqrt{2}}}\left(|0_1|1_2\pm |0_1|1_2\right),`$ (2)
$`|\phi _{1,2}^\pm ={\displaystyle \frac{1}{\sqrt{2}}}\left(|0_1|0_2\pm |1_1|1_2\right),`$ (3)
are the most common entangled states employed in quantum information science. (The $``$ sign will be understood hereafter.) Usually two separated particles are entangled because they were once in contact and interacted with each other. However, by entanglement swapping BVK ; ZZHE ; JWPZ , we can entangle two remotely separated particles which do not have a history of mutual interaction. This is a vivid demonstration of nonlocal correlations in the quantum theory. The idea of entanglement swapping is mathematically very simple. Consider two Bell states, say $`|\phi _{1,2}^+`$ and $`|\varphi _{3,4}^{}`$, where particles 1 and 2 have never been in contact with particles 3 and 4 before. The product of these two states can be rewritten as
$$|\phi _{1,2}^+|\varphi _{3,4}^{}=|\phi _{1,3}^+|\varphi _{2,4}^{}+|\phi _{1,3}^{}|\varphi _{2,4}^+|\varphi _{1,3}^+|\phi _{2,4}^{}|\varphi _{1,3}^{}|\phi _{2,4}^+.$$
(4)
Therefore if we make a Bell measurement on the pair (1,3), then depending on the outcome the (2,4) pair would collapse into one of the corresponding Bell states: $`|\varphi _{2,4}^{}`$, $`|\varphi _{2,4}^+`$, $`|\phi _{2,4}^{}`$, and $`|\phi _{2,4}^+`$. In other words, particles 2 and 4 would become entangled even though they are far apart and have never interacted with each other in the past.
In this paper, we shall use entanglement swapping to distribute an arbitrary $`N`$-particle state to $`MN`$ remote parties. Distribution of quantum information is essential in quantum secret sharing (QSS) and other applications in quantum information science, such as quantum communication network and distributed quantum computation CEHM . QSS is the quantum counterpart of classical secret sharing first discussed by Blakely Blakely and Shamir Shamir in 1979. The idea is as follows. Suppose Alice wants to send a secret message to a remote location, and she can send it to either of her two agents, Bob and Charlie. As a precaution against information leakage and misuse, it is safer for her to split the message into two pieces and send them separately to Bob and Charlie, such that anyone alone has absolutely no knowledge of the message. Bob and Charlie can reconstruct the original secret message only if they collaborate with each other. Clearly the above consideration can be generalized to secret sharing by $`N`$ parties.
QSS refers to the implementation of the secret sharing task outlined above using quantum mechanical resources. Hillary et al. HBB and Karlsson et al. KKI were the first to propose QSS protocols using respectively three-particle Greenberger-Horne-Zeilinger (GHZ) states and two-particle Bell states. Since then, a wide variety of other QSS protocols have been proposed. Although the goal of most QSS protocols is to protect a classical secret message, the notion of QSS has also been generalized to the sharing of a secret quantum state HBB ; KKI ; CGL ; LSBSL , which is also referred to as “quantum state sharing”.
Whether the goal is to share a classical or a quantum secret, in almost all cases, it is necessary to distribute a $`N`$-qubit entangled state among $`M`$ parties, where $`NM2`$. For example, in the QSS protocol of Hillery et al. HBB , Alice uses GHZ states to split a quantum key into two shares such that each of the two agents, Bob and Charlie, gets only one share. To do so, Alice must be able to safely distribute two of the particles in each tripartite GHZ state separately to Bob and Charlie. In the quantum state sharing protocol proposed by Cleve et al. CGL , the secret message to be shared is a single qutrit state. The protocol is similar to an error-correcting code, in which a three-qutrit entangled state is generated from the secret qutrit. The “dealer” then distributes the resulting three qutrits to three different parties. More generally, in a so-called $`(k,n)`$ threshold scheme, the secret (quantum or classical) is divided into $`n`$ shares, such that any $`k`$ of those shares can be used to reconstruct the secret, while any set of less than $`k`$ shares contains absolutely no information about the secret at all. So in general one needs to distribute an arbitrary entangled state to $`n`$ different parties.
Of course Alice could send the particles involved directly over quantum channels to their respective destinations. This practice is however both inefficient and unsafe. First of all, since noise is always present in any available channel, quantum information may get distorted; more seriously decoherence effects may even cause the particle to collapse. The loss of any one particle in an unknown multiparticle state due to decoherence or dissipation will require the whole state to be regenerated again. Hence direct distribution of the particles is not an efficient way to proceed. Furthermore the state to be shared is itself the carrier of information, therefore for security reasons, it is not safe to send them directly over long-distance quantum channels which may be monitored by eavesdroppers. Of course eavesdropping activities can be detected by security testing methods, but they inevitably involve measuring some of the particles going through the channels, which is obviously not acceptable if they are part of the multiparticle state being distributed.
It is therefore more desirable to distribute the particles using quantum entanglement plus local operations and classical communications only. Entanglement resources are usually supplied by two-qubit Bell states or sometimes three-qubit GHZ states shared between the sender and the intended receivers. The advantage of this approach is that the required entanglement can be established and tested independent of the state to be distributed. Once it is securely established, quantum channels are no longer needed, and transmission noise is no longer a problem.
In the following, we show how to faithfully distribute an arbitrary $`N`$-qubit state using entanglement swapping; our scheme generalizes a multipartite QSS protocol discussed in Ref. KKI . Schemes of distributing $`N`$-qubit states by teleportation have also been proposed Rigolin ; IZZ . However, these protocols have the undesirable features that the complexity of the required collective measurements increases with $`N`$. In our scheme, only Bell states are used and only two-qubit collective measurements are required.
To proceed, we first show that entanglement swapping between two Bell states can be generalized to that between an arbitrary $`N`$-qubit state and a Bell state. The setting is that Alice owns an arbitrary $`N`$-qubit state $`|\mathrm{\Psi }_{1,\mathrm{},N}`$, and she shares a Bell state $`|\varphi _{\mu ,\nu }^{}`$ with Bob who is somewhere far away; Alice holds qubit-$`\mu `$ and Bob qubit-$`\nu `$. From
$$|\mathrm{\Psi }_{1,\mathrm{},N}=\left(|0_i0_i|+|1_i1_i|\right)|\mathrm{\Psi }_{1,\mathrm{},N},(1iN),$$
(5)
we see that $`|\mathrm{\Psi }_{1,\mathrm{},N}`$ can always be rewritten as
$$|\mathrm{\Psi }_{1,\mathrm{},N}=a|0_i|\mathrm{\Phi }_{1,\mathrm{},i1,i+1,\mathrm{},N}+b|1_i|\mathrm{\Phi }_{1,\mathrm{},i1,i+1,\mathrm{},N}^{},$$
(6)
where $`|a|^2+|b|^2=1`$, and $`|\mathrm{\Phi }_{1,\mathrm{},i1,i+1,\mathrm{},N}`$ and $`|\mathrm{\Phi }_{1,\mathrm{},i1,i+1,\mathrm{},N}^{}`$ are normalized states of $`(N1)`$ qubits. Note that unless $`a`$ or $`b`$ vanishes, or $`|\mathrm{\Phi }_{1,\mathrm{},i1,i+1,\mathrm{},N}`$ and $`|\mathrm{\Phi }_{1,\mathrm{},i1,i+1,\mathrm{},N}^{}`$ differ only by a phase factor, otherwise qubit-$`i`$ is entangled with the rest of the group.
Similar to Eq. (4), we can rewrite the product of $`|\mathrm{\Psi }_{1,\mathrm{},N}`$ and $`|\varphi _{\mu ,\nu }^{}`$ as
$`|\mathrm{\Psi }_{1,\mathrm{},N}|\varphi _{\mu ,\nu }^{}={\displaystyle \frac{1}{2}}`$ $`[`$ $`|\phi _{i,\mu }^+\left(a|1_\nu |\mathrm{\Phi }_{1,\mathrm{},i1,i+1,\mathrm{},N}b|0_\nu |\mathrm{\Phi }_{1,\mathrm{},i1,i+1,\mathrm{},N}^{}\right)`$ (7)
$`+`$ $`|\phi _{i,\mu }^{}\left(a|1_\nu |\mathrm{\Phi }_{1,\mathrm{},i1,i+1,\mathrm{},N}+b|0_\nu |\mathrm{\Phi }_{1,\mathrm{},i1,i+1,\mathrm{},N}^{}\right)`$
$``$ $`|\varphi _{i,\mu }^+\left(a|0_\nu |\mathrm{\Phi }_{1,\mathrm{},i1,i+1,\mathrm{},N}b|1_\nu |\mathrm{\Phi }_{1,\mathrm{},i1,i+1,\mathrm{},N}^{}\right)`$
$``$ $`|\varphi _{i,\mu }^{}(a|0_\nu |\mathrm{\Phi }_{1,\mathrm{},i1,i+1,\mathrm{},N}+b|1_\nu |\mathrm{\Phi }_{1,\mathrm{},i1,i+1,\mathrm{},N}^{})].`$
Hence a Bell measurement by Alice on the $`(i,\mu )`$ pair will entangle the remote qubit-$`\nu `$ to the local group of ($`N1`$) qubits $`(1,\mathrm{},i1,i+1,\mathrm{},N)`$. This process is depicted diagrammatically in Fig. 1.
The resulting $`N`$-qubit state depends on the outcome of the Bell measurement:
$`(1)|\phi _{i,\mu }^+a|1_\nu |\mathrm{\Phi }_{1,\mathrm{},i1,i+1,\mathrm{},N}b|0_\nu |\mathrm{\Phi }_{1,\mathrm{},i1,i+1,\mathrm{},N}^{},`$ (8)
$`(2)|\phi _{i,\mu }^{}a|1_\nu |\mathrm{\Phi }_{1,\mathrm{},i1,i+1,\mathrm{},N}+b|0_\nu |\mathrm{\Phi }_{1,\mathrm{},i1,i+1,\mathrm{},N}^{},`$ (9)
$`(3)|\varphi _{i,\mu }^+a|0_\nu |\mathrm{\Phi }_{1,\mathrm{},i1,i+1,\mathrm{},N}b|1_\nu |\mathrm{\Phi }_{1,\mathrm{},i1,i+1,\mathrm{},N}^{},`$ (10)
$`(4)|\varphi _{i,\mu }^{}a|0_\nu |\mathrm{\Phi }_{1,\mathrm{},i1,i+1,\mathrm{},N}+b|1_\nu |\mathrm{\Phi }_{1,\mathrm{},i1,i+1,\mathrm{},N}^{}.`$ (11)
Comparing with Eqs. (6), we see that if Alice tells Bob which of the four Bell states she has obtained, then by a local unitary transformation on qubit-$`\nu `$ Bob can rotate the state of the $`N`$-qubit group $`(1,\mathrm{},i1,\nu ,i+1,\mathrm{},N)`$ back to the original state $`|\mathrm{\Psi }_{1,\mathrm{},N}`$ (with qubit-$`i`$ $``$ qubit-$`\nu `$). The required unitary operators are ($`\sigma _z\sigma _x`$, $`\sigma _x`$, $`\sigma _z`$, $`I`$) respectively for the four possible outcomes listed in Eqs. (8$``$11). It is easily seen that the information cost for the whole operation is one e-bit plus two c-bits (classical bits) from Alice to Bob. If qubit-$`i`$ is not entangled with the group $`(1,\mathrm{},i1,i+1,\mathrm{},N)`$, that is, if Eq. (6) can be reduced to
$$|\mathrm{\Psi }_{1,\mathrm{},N}=\left(a|0_i+b|1_i\right)|\mathrm{\Phi }_{1,\mathrm{},i1,i+1,\mathrm{},N},$$
(12)
then what we have done is simply the teleportation BBCJPW of a single qubit $`\left(a|0_i+b|1_i\right)`$ from Alice to Bob. For a general $`|\mathrm{\Psi }_{1,\mathrm{},N}`$ our procedure effectively teleports the entanglement to Bob.
Notice that the above procedure is entirely general, in the sense that it is independent of the state of the inactive qubits $`(1,\mathrm{},i1,i+1,\mathrm{},N)`$. Therefore it can be repeated until all the qubits are distributed to their respective remote locations as desired. We note in passing that, if Alice leaves some of the qubits undistributed, then the distributed qubits form a mixed state.
It is interesting to note that each round of the operation (Bell measurement plus unitary rotation) described above can be summarized by the action of an operator $`U_{i\nu }`$ which interchanges the states of qubit-$`i`$ and qubit-$`\nu `$:
$$U_{i\nu }\left(|\mathrm{\Psi }_{1,\mathrm{},i1,i,i+1,\mathrm{},N}|\varphi _{\mu ,\nu }^{}\right)=|\mathrm{\Psi }_{1,\mathrm{},i1,\nu ,i+1,\mathrm{},N}|\varphi _{\mu ,i}^{},$$
(13)
where without loss of generality we have assumed that the measured Bell state has been rotated back to the singlet Bell state $`|\varphi _{\mu ,i}^{}`$, which Alice can easily do by a local unitary transformation. The effect of $`U_{i\nu }`$ is identical to that of a regular two-qubit swap gate NC except for the fact that here the qubits involved are remotely separated; hence $`U_{i\nu }`$ is a kind of nonlocal swap operator. The whole distribution process can be accomplished by simply repeating this swapping operation $`N`$ times:
$$\underset{i=1}{\overset{N}{}}U_{i\nu _i}\left(|\mathrm{\Psi }_{1,\mathrm{},N}\underset{i=1}{\overset{N}{}}|\varphi _{\mu _i,\nu _i}^{}\right)=|\mathrm{\Psi }_{\nu _1,\mathrm{},\nu _N}\underset{i=1}{\overset{N}{}}|\varphi _{\mu _i,i}^{},$$
(14)
where the $`N`$ qubits $`(\nu _1,\mathrm{},\nu _N)`$ are located at their respective destinations.
From this perspective, it is clear that we could also use the general nonlocal swap operation CLP ; EJPP in place of $`U_{i\nu _i}`$, although it would be rather inefficient to do so. It is known that the implementation of a general nonlocal swap operation requires at least two e-bits plus two c-bits from Alice to Bob plus another two c-bits from Bob to Alice CLP ; EJPP . However in our scheme, only one e-bit and two c-bits from Alice to Bob are needed per swap. This difference in resources consumption is due to the fact that $`U_{i\nu _i}`$ is actually not an universal swap operator, because it cannot be used to swap two arbitrary remote qubits. Our scheme is only good for the special purpose of qubit distribution, where each round of swapping is analogous to the teleportation of a single qubit, hence relatively less resources are required.
In summary, we have constructed a scheme to distribute an arbitrary $`N`$-qubit state (pure or mixed) to $`MN`$ remote parties. The basic operation used is entanglement swapping which involves only Bell states and two-qubit collective measurements. Our scheme is experimentally feasible with currently available technologies LSBSL ; JWPZ . |
warning/0507/math-ph0507008.html | ar5iv | text | # 1 Introduction
## 1 Introduction
In non-relativistic quantum many-body systems, a folk theorem states that a nonvanishing spectral gap above the ground state implies exponentially decaying correlations in the ground state. Perhaps this has been the most popular folk theorem in this field since Haldane predicted a “massive phase” in low dimensional, isotropic quantum systems. Quite recently, this statement was partially proved for quantum lattice systems with a global U(1) symmetry in (fractal) dimensions $`D<2`$. More precisely, a bound which decays to zero at large distance was obtained for correlation functions whose observables behave as a vector under the U(1)-rotation. Unfortunately, the bound is weaker than the expected exponential decay. On the other hand, exponential clustering of the correlations was also proved recently for quantum many-body lattice systems under the gap assumption. This is a non-relativistic version of Fredenhagen’s theorem of relativistic quantum field theory. Clearly the following natural question arises: can this clustering property be combined with the above bound for the decay of the correlations to yield the tighter, exponentially decaying bound for the correlation functions themselves, rather than just for the connected correlation functions? We emphasize that these are different statements; given clustering, the decay of the correlation functions requires also that certain matrix elements vanish in the ground state sector.
In this paper, we address this problem and reexamine the above folk theorem by relying on the exponential clustering of the correlations. Our first step is to provide a rigorous proof of the exponential clustering. We extend the previous results in this case to treat long-range interactions including both power-law and exponentially decaying interactions. In the former case, all the upper bounds for the correlations become power-law bounds.
We then prove that ground state correlation functions of observables which transform as vectors under a U(1) symmetry decay exponentially or with a power law, depending on the form of the interaction, given an additional assumption on a certain self-similarity. In particular, if the system is translationally invariant in one of the spatial directions, this self-similarity condition is automatically satisfied. Therefore the corresponding correlation functions decay exponentially for translationally invariant systems on one-dimensional regular lattices. As a byproduct, we also prove that, if two observables anticommute with each other at large distance, then the corresponding correlation in the ground state decays exponentially under the gap assumption for a wide class of lattice fermion systems with exponentially decaying interactions in any dimensions. In this case, we do not need any other assumption except for those on the interactions and the spectral gap.
This paper is organized as follows: In the next section, we give the precise definitions of the models, and describe our main results. In Section 3, we prove the clustering of generic correlation functions under the gap assumption, and obtain the upper decaying bound for the fermionic correlations. The decay of the bosonic correlations are treated in Section 4. Appendix A is devoted to the proof of the Lieb-Robinson bound for the group velocity of the information propagation in the models with a long-range interaction decaying by power law.
## 2 Models and main results
We consider quantum systems on generic lattices . Let $`\mathrm{\Lambda }_s`$ be a set of the sites, $`x,y,z,w,\mathrm{}`$, and $`\mathrm{\Lambda }_b`$ a set of the bonds, i.e., pairs of sites, $`\{x,y\},\{z,w\},\mathrm{}`$. We call the pair, $`\mathrm{\Lambda }:=(\mathrm{\Lambda }_s,\mathrm{\Lambda }_b)`$, the lattice. If a sequence of sites, $`x_0,x_1,x_2,\mathrm{},x_n`$, satisfies $`\{x_{j1},x_j\}\mathrm{\Lambda }_b`$ for $`j=1,2,\mathrm{},n`$, then we say that the path, $`\{x_0,x_1,x_2,\mathrm{},x_n\}`$, has length $`n`$ and connects $`x_0`$ to $`x_n`$. We denote by $`\mathrm{dist}(x,y)`$ the graph-theoretic distance which is defined to be the shortest path length that one needs to connect $`x`$ to $`y`$. We denote by $`|X|`$ the cardinality of the finite set $`X`$. The Hamiltonian $`H_\mathrm{\Lambda }`$ is defined on the tensor product $`_{x\mathrm{\Lambda }_s}_x`$ of a finite dimensional Hilbert space $`_x`$ at each site $`x`$. We assume $`sup_{\mathrm{\Lambda }_s}sup_x\mathrm{dim}_xN<\mathrm{}`$. For a lattice fermion system, we consider the Fock space.
Consider the Hamiltonian of the form,
$$H_\mathrm{\Lambda }=\underset{X\mathrm{\Lambda }_s}{}h_X,$$
(2.1)
where $`h_X`$ is the local Hamiltonian of the compact support $`X`$. We consider both power-law and exponentially decaying interactions $`h_X`$.
For the power-law decaying interactions $`h_X`$, we require the following conditions:
###### Assumption 2.1
The interaction $`h_X`$ satisfies
$$\underset{Xx,y}{}h_X\frac{\lambda _0}{[1+\mathrm{dist}(x,y)]^\eta }$$
(2.2)
with positive constants, $`\lambda _0`$ and $`\eta `$, and the lattice $`\mathrm{\Lambda }`$ equipped with the metric satisfies
$$\underset{z\mathrm{\Lambda }_s}{}\frac{1}{[1+\mathrm{dist}(x,z)]^\eta }\times \frac{1}{[1+\mathrm{dist}(z,y)]^\eta }\frac{p_0}{[1+\mathrm{dist}(x,y)]^\eta }$$
(2.3)
with a positive constant $`p_0`$.
Remark: If
$$\underset{\mathrm{\Lambda }_s}{sup}\underset{x}{sup}\underset{y\mathrm{\Lambda }_s}{}\frac{1}{[1+\mathrm{dist}(x,y)]^\eta }<\mathrm{},$$
(2.4)
then the inequality (2.3) holds as follows:
$`{\displaystyle \underset{z\mathrm{\Lambda }_s}{}}{\displaystyle \frac{1}{[1+\mathrm{dist}(x,z)]^\eta }}\times {\displaystyle \frac{1}{[1+\mathrm{dist}(z,y)]^\eta }}`$ (2.5)
$`=`$ $`{\displaystyle \frac{1}{[1+\mathrm{dist}(x,y)]^\eta }}{\displaystyle \underset{z\mathrm{\Lambda }_s}{}}{\displaystyle \frac{[1+\mathrm{dist}(x,y)]^\eta }{[1+\mathrm{dist}(x,z)]^\eta [1+\mathrm{dist}(z,y)]^\eta }}`$
$``$ $`{\displaystyle \frac{1}{[1+\mathrm{dist}(x,y)]^\eta }}{\displaystyle \underset{z\mathrm{\Lambda }_s}{}}2^\eta {\displaystyle \frac{[1+\mathrm{dist}(x,z)]^\eta +[1+\mathrm{dist}(z,y)]^\eta }{[1+\mathrm{dist}(x,z)]^\eta [1+\mathrm{dist}(z,y)]^\eta }}`$
$``$ $`{\displaystyle \frac{1}{[1+\mathrm{dist}(x,y)]^\eta }}{\displaystyle \underset{z\mathrm{\Lambda }_s}{}}2^\eta \left\{{\displaystyle \frac{1}{[1+\mathrm{dist}(x,z)]^\eta }}+{\displaystyle \frac{1}{[1+\mathrm{dist}(z,y)]^\eta }}\right\},`$
where we have used the inequality, $`[1+\mathrm{dist}(x,y)]^\eta 2^\eta ([1+\mathrm{dist}(x,z)]^\eta +[1+\mathrm{dist}(z,y)]^\eta )`$. From the assumption (2.2) and the condition (2.4), one has
$$\underset{x}{sup}\underset{Xx}{}h_X|X|s_0<\mathrm{},$$
(2.6)
where $`s_0`$ is a positive constant which is independent of the volume of $`|\mathrm{\Lambda }_s|`$.
Instead of these conditions, we can also require:
###### Assumption 2.2
The interaction $`h_X`$ satisfies
$$\underset{x}{sup}\underset{Xx}{}h_X|X|[1+\mathrm{diam}(X)]^\eta s_1<\mathrm{},$$
(2.7)
where $`\eta `$ is a positive constant, $`\mathrm{diam}(X)`$ is the diameter of the set $`X`$, i.e., $`\mathrm{diam}(X)=\mathrm{max}\{\mathrm{dist}(x,y)|x,yX\}`$, and $`s_1`$ is a positive constant which is independent of the volume of $`|\mathrm{\Lambda }_s|`$.
For exponentially decaying interactions $`h_X`$, we require one of the following two assumptions:
###### Assumption 2.3
There exists a positive $`\eta `$ satisfying the condition (2.4). The interaction $`h_X`$ satisfies
$$\underset{Xx,y}{}h_X\lambda _0\mathrm{exp}[(\mu +\epsilon )\mathrm{dist}(x,y)]$$
(2.8)
with some positive constants, $`\lambda _0,\mu `$ and $`\epsilon `$.
Remark: From the conditions, we have
$$\mathrm{exp}[(\mu +\epsilon )\mathrm{dist}(x,y)]\frac{\lambda _0^{}\mathrm{exp}[\mu \mathrm{dist}(x,y)]}{[1+\mathrm{dist}(x,y)]^\eta }$$
(2.9)
with a positive constant $`\lambda _0^{}`$, and
$$\underset{z\mathrm{\Lambda }_s}{}\frac{\mathrm{exp}[\mu \mathrm{dist}(x,z)]}{[1+\mathrm{dist}(x,z)]^\eta }\times \frac{\mathrm{exp}[\mu \mathrm{dist}(z,y)]}{[1+\mathrm{dist}(z,y)]^\eta }\frac{p_0\mathrm{exp}[\mu \mathrm{dist}(x,y)]}{[1+\mathrm{dist}(x,y)]^\eta }$$
(2.10)
with a positive constant $`p_0`$ in the same way as in the preceding remark.
###### Assumption 2.4
The interaction $`h_X`$ satisfies
$$\underset{x}{sup}\underset{Xx}{}h_X|X|\mathrm{exp}[\mu \mathrm{diam}(X)]s_1<\mathrm{},$$
(2.11)
where $`\mu `$ is a positive constant, and $`s_1`$ is a positive constant which is independent of the volume of $`|\mathrm{\Lambda }_s|`$.
Remark: This assumption is milder than that in by the absence of the factor $`N^{2|X|}`$ in the summand.
Further we assume the existence of a “uniform gap” above the ground state sector of the Hamiltonian $`H_\mathrm{\Lambda }`$. The precise definition of the “uniform gap” is:
###### Definition 2.5
(Uniform gap): We say that there is a uniform gap above the ground state sector if the spectrum $`\sigma (H_\mathrm{\Lambda })`$ of the Hamiltonian $`H_\mathrm{\Lambda }`$ satisfies the following conditions: The ground state of the Hamiltonian $`H_\mathrm{\Lambda }`$ is $`q`$-fold (quasi)degenerate in the sense that there are $`q`$ eigenvalues, $`E_{0,1},\mathrm{},E_{0,q}`$, in the ground state sector at the bottom of the spectrum of $`H_\mathrm{\Lambda }`$ such that
$$\mathrm{\Delta }:=\underset{\mu ,\mu ^{}}{\mathrm{max}}\{|E_{0,\mu }E_{0,\mu ^{}}|\}0\text{as }|\mathrm{\Lambda }_s|\mathrm{}.$$
(2.12)
Further the distance between the spectrum, $`\{E_{0,1},\mathrm{},E_{0,q}\}`$, of the ground state and the rest of the spectrum is larger than a positive constant $`\mathrm{\Delta }E`$ which is independent of the volume $`|\mathrm{\Lambda }_s|`$. Namely there is a spectral gap $`\mathrm{\Delta }E`$ above the ground state sector.
Let $`A_X,B_Y`$ be observables with the support $`X,Y\mathrm{\Lambda }_s`$, respectively. We say that the pair of two observables, $`A_X`$ and $`B_Y`$, is fermionic if they satisfy the anicommutation relation, $`\{A_X,B_Y\}=0`$ for $`XY=\mathrm{}`$. If they satisfy the commutation relation, then we call the pair bosonic.
Define the ground-state expectation as
$$\mathrm{}_{0,\mathrm{\Lambda }}:=\frac{1}{q}\mathrm{Tr}(\mathrm{})P_{0,\mathrm{\Lambda }},$$
(2.13)
where $`P_{0,\mathrm{\Lambda }}`$ is the projection onto the ground state sector. For the infinite volume,
$$\mathrm{}_0:=\mathrm{weak}^{}\text{-}\underset{|\mathrm{\Lambda }_s|\mathrm{}}{lim}\mathrm{}_{0,\mathrm{\Lambda }},$$
(2.14)
where we take a suitable subsequence of finite lattices $`\mathrm{\Lambda }`$ going to the infinite volume so that the expectation converges to a linear functional for a set of quasilocal observables. Although the ground-state expectation thus constructed depends on the subsequence of the lattices and on the observables, our results below hold for any ground-state expectation thus constructed. Further, we denote by
$$\omega (\mathrm{}):=\mathrm{weak}^{}\text{-}\underset{|\mathrm{\Lambda }_s|\mathrm{}}{lim}\mathrm{\Phi }_\mathrm{\Lambda },(\mathrm{})\mathrm{\Phi }_\mathrm{\Lambda }$$
(2.15)
the ground-state expectation in the infinite volume for a normalized vector $`\mathrm{\Phi }_\mathrm{\Lambda }`$ in the sector of the ground state for finite lattice $`\mathrm{\Lambda }`$.
###### Theorem 2.6
(Clustering of fermionic correlations): Let $`A_X,B_Y`$ be fermionic observables with a compact support. Assume that there exists a uniform spectral gap $`\mathrm{\Delta }E>0`$ above the ground state sector in the spectrum of the Hamiltonian $`H_\mathrm{\Lambda }`$ in the sense of Definition 2.5. Let $`\omega `$ be a ground-state expectation (2.15) in the infinite volume limit. Then the following bound is valid:
$`\left|\omega (A_XB_Y){\displaystyle \frac{1}{2}}\left[\omega (A_XP_0B_Y)\omega (B_YP_0A_X)\right]\right|`$ (2.16)
$``$ $`\mathrm{Const}.\times \{\begin{array}{cc}[1+\mathrm{dist}(X,Y)]^{\stackrel{~}{\eta }},\hfill & \text{for power-law decaying }h_X\text{;}\hfill \\ \mathrm{exp}[\stackrel{~}{\mu }\mathrm{dist}(X,Y)],\hfill & \text{for exponentially decaying }h_X\text{,}\hfill \end{array}`$
where $`P_0`$ is the projection onto the sector of the infinite-volume ground state,<sup>1</sup><sup>1</sup>1$`\omega (\mathrm{}P_0\mathrm{})`$ is also defined as a bilinear functional for a set of quasilocal observables in the weak limit. and
$$\stackrel{~}{\eta }=\frac{\eta }{1+2v_\eta /\mathrm{\Delta }E}\text{and}\stackrel{~}{\mu }=\frac{\mu }{1+2v_\mu /\mathrm{\Delta }E}.$$
(2.17)
Here $`v_\eta `$ and $`v_\mu `$ are, respectively, an increasing function of $`\eta `$ and $`\mu `$, and give an upper bound of the group velocity of the information propagation.
Remark: Clearly there exists a maximum $`\mu _{\mathrm{max}}`$ such that the bound (2.8) holds for any $`\mu \mu _{\mathrm{max}}`$. Combining this observation with (2.17), there exists a maximum $`\stackrel{~}{\mu }=\mathrm{max}_{\mu \mu _{\mathrm{max}}}\{\mu /(1+2v_\mu /\mathrm{\Delta }E)\}`$ which gives the optimal decay bound. When the interaction $`h_X`$ is of finite range, one can take any large $`\mu `$. But the upper bound $`v_\mu `$ of the group velocity exponentially increases as $`\mu `$ increases because $`v_\mu `$ depends on $`\lambda _0`$ of (2.8). In consequence, a finite $`\stackrel{~}{\mu }`$ gives the optimal bound.
Formally applying the identity, $`A_XP_0B_Y_0=B_YP_0A_X_0`$, for the bound (2.16), we have the following decay bound for the correlation:<sup>2</sup><sup>2</sup>2See Section 3 for details.
###### Corollary 2.7
Let $`A_X,B_Y`$ be fermionic observables with a compact support. Assume that there exists a uniform spectral gap $`\mathrm{\Delta }E>0`$ above the ground state sector in the spectrum of the Hamiltonian $`H_\mathrm{\Lambda }`$ in the sense of Definition 2.5. Then the following bound is valid:
$$\left|A_XB_Y_0\right|\mathrm{Const}.\times \{\begin{array}{cc}[1+\mathrm{dist}(X,Y)]^{\stackrel{~}{\eta }},\hfill & \text{for power-law decaying }h_X\text{;}\hfill \\ \mathrm{exp}[\stackrel{~}{\mu }\mathrm{dist}(X,Y)],\hfill & \text{for exponentially decaying }h_X\text{,}\hfill \end{array}$$
(2.18)
in the infinite volume limit, where $`\stackrel{~}{\eta },\stackrel{~}{\mu }`$ are as defined above.
###### Theorem 2.8
(Clustering of bosonic correlations): Let $`A_X,B_Y`$ be bosonic observables with a compact support. Assume that there exists a uniform spectral gap $`\mathrm{\Delta }E>0`$ above the ground state sector in the spectrum of the Hamiltonian $`H_\mathrm{\Lambda }`$ in the sense of Definition 2.5. Let $`\omega `$ be a ground-state expectation (2.15) in the infinite volume limit. Then the following bound is valid:
$`\left|\omega (A_XB_Y){\displaystyle \frac{1}{2}}\left[\omega (A_XP_0B_Y)+\omega (B_YP_0A_X)\right]\right|`$ (2.19)
$``$ $`\mathrm{Const}.\times \{\begin{array}{cc}[1+\mathrm{dist}(X,Y)]^{\stackrel{~}{\eta }},\hfill & \text{for power-law decaying }h_X\text{;}\hfill \\ \mathrm{exp}[\stackrel{~}{\mu }\mathrm{dist}(X,Y)],\hfill & \text{for exponentially decaying }h_X\text{,}\hfill \end{array}`$
where $`\stackrel{~}{\eta },\stackrel{~}{\mu }`$ are as defined above.
Remark: Theorem 2.8 is a clustering bound for the connected correlation functions. We now make some additional definitions that will enable us, in certain cases, to prove the decay of $`\left[\omega (A_XP_0B_Y)+\omega (B_YP_0A_X)\right]/2`$ so that Theorem 2.8 can be replaced with a stronger bound below, Theorem 2.10.
###### Definition 2.9
(Self-similarity): Write $`m=q^2`$ with the degeneracy $`q`$ of the ground state sector. We say that the system has self-similarity if the following conditions are satisfied: For any observable $`A`$ of compact support and any given large $`L>0`$, there exist transformations, $`R_1,R_2,\mathrm{},R_m`$, and observables, $`B^{(1)},B^{(2)},\mathrm{},B^{(m)}`$, such that the Hamiltonian $`H_\mathrm{\Lambda }`$ is invariant under the transformations, i.e., $`R_j(H_\mathrm{\Lambda })=H_\mathrm{\Lambda }`$ for any lattice $`\mathrm{\Lambda }`$ with sufficiently large $`|\mathrm{\Lambda }_s|`$, and that the observables satisfy the following conditions:
$$B^{(j)}=R_j(A)\text{and}\left(B^{(j)}\right)^{}=R_j(A^{})\text{for }j=1,2,\mathrm{},m,$$
(2.20)
$$\mathrm{dist}(\mathrm{supp}A,\mathrm{supp}B^{(j)})L\text{for }j=1,2,\mathrm{},m,$$
(2.21)
and
$$\mathrm{dist}(\mathrm{supp}B^{(j)},\mathrm{supp}B^{(k)})L\text{for }jk.$$
(2.22)
In Section 4, we will discuss other conditions similar to this self-similarity condition.
###### Theorem 2.10
Assume that the degeneracy $`q`$ of the ground state sector of the Hamiltonian $`H_\mathrm{\Lambda }`$ is finite in the infinite volume limit, and that there exists a uniform spectral gap $`\mathrm{\Delta }E>0`$ above the ground state sector in the spectrum of the Hamiltonian $`H_\mathrm{\Lambda }`$ in the sense of Definition 2.5. Further assume that the system has self-similarity in the sense of Definition 2.9, and that there exists a subset $`𝒜_b^s`$ of bosonic observables with a compact support such that $`R_j(𝒜_b^s)𝒜_b^s=(𝒜_b^s)^{}`$ for $`j=1,2,\mathrm{},m`$, and that $`A_X^{}B_Y^{}_00`$ as $`\mathrm{dist}(X,Y)\mathrm{}`$ for any pair of bosonic observables, $`A_X^{},B_Y^{}𝒜_b^s`$. Let $`\omega `$ be a ground-state expectation (2.15) in the infinite volume limit, and let $`A_X,B_Y`$ be a pair of bosonic observables satisfying $`A_X𝒜_b^s`$. Then the following bound is valid:
$$\left|\omega (A_XB_Y)\right|\mathrm{Const}.\times \{\begin{array}{cc}[1+\mathrm{dist}(X,Y)]^{\stackrel{~}{\eta }},\hfill & \text{for power-law decaying }h_X\text{;}\hfill \\ \mathrm{exp}[\stackrel{~}{\mu }\mathrm{dist}(X,Y)],\hfill & \text{for exponentially decaying }h_X\text{,}\hfill \end{array}$$
(2.23)
where $`\stackrel{~}{\eta },\stackrel{~}{\mu }`$ are as defined above.
Remark: 1. If the finite system is translationally invariant in one of the spatial directions with a periodic boundary condition, then the self-similarity condition of Definition 2.9 is automatically satisfied by taking the translation as the transformation $`R_j`$. Thus we do not need an additional assumption for such systems.
2. Theorem 2.10 can be extended to a system having infinite degeneracy of the ground state sector in the infinite volume limit if the degeneracy for finite volume is sufficiently small compared to the volume of the system. See Theorem 4.1 in Section 4 for details.
In order to apply this theorem, we need to be able to show that $`A_XB_Y_00`$ as $`\mathrm{dist}(X,Y)\mathrm{}`$ in the infinite volume for any pair of bosonic observables, $`A_X,B_Y𝒜_b^s`$. However, this was proven for quantum spin or fermion systems with a global U(1) symmetry on a class of lattices with (fractal) dimension $`D<2`$ as defined in (2.25) below, so long as the observables behave as a vector under the U(1) rotation.
We first define the dimension for these lattices. The “sphere”, $`S_r(x)`$, centered at $`x\mathrm{\Lambda }_s`$ with the radius $`r`$ is defined as
$$S_r(x):=\{y\mathrm{\Lambda }_s|\mathrm{dist}(y,x)=r\}.$$
(2.24)
Assume that there exists a “(fractal) dimension” $`D1`$ of the lattice $`\mathrm{\Lambda }`$ such that the number $`|S_r(x)|`$ of the sites in the sphere satisfies
$$\underset{x\mathrm{\Lambda }_s}{sup}|S_r(x)|C_0r^{D1}$$
(2.25)
with some positive constant $`C_0`$. This class of the lattices is the same as in .
Consider spin or fermion systems with a global U(1) symmetry on the lattice $`\mathrm{\Lambda }`$ with (fractal) dimension $`1D<2`$, and require the existence of a uniform gap above the ground state sector of the Hamiltonian $`H_\mathrm{\Lambda }`$ in the sense of Definition 2.5. Although the method of can be applied to a wide class of such systems, we consider only two important examples, the Heisenberg and the Hubbard models. We take the set $`𝒜_b^s`$ to be the bosonic observables which behave as a vector under the U(1) rotation. In the rest of this section we use the results of to show as in (2.29,2.12) that the correlation function for this class of observables in these models does decay to zero as $`\mathrm{dist}(X,Y)\mathrm{}`$. The bounds (2.29,2.12) provide only a slow bound on the decay. However, this slow bound on the decay suffices, in conjunction with the self-similarity condition of Definition 2.9 to apply Theorem 2.10. Thus, under the self-similarity assumption as well as the gap assumption, all the upper bounds below (2.29,2.12) are replaced with exponentially decaying bounds by Theorem 2.10. In particular, a system with a translational invariance automatically satisfies the self-similarity condition as mentioned above. Therefore the corresponding correlations show exponential decay for translationally invariant systems on one-dimensional regular lattices.
XXZ Heisenberg model: The Hamiltonian $`H_\mathrm{\Lambda }`$ is given by
$$H_\mathrm{\Lambda }=H_\mathrm{\Lambda }^{XY}+V_\mathrm{\Lambda }(\{S_x^{(3)}\})$$
(2.26)
with
$$H_\mathrm{\Lambda }^{XY}=2\underset{\{x,y\}\mathrm{\Lambda }_b}{}J_{x,y}^{\mathrm{XY}}\left[S_x^{(1)}S_y^{(1)}+S_x^{(2)}S_y^{(2)}\right],$$
(2.27)
where $`𝐒_x=(S_x^{(1)},S_x^{(2)},S_x^{(3)})`$ is the spin operator at the site $`x\mathrm{\Lambda }_s`$ with the spin $`S=1/2,1,3/2,\mathrm{}`$, and $`J_{x,y}^{\mathrm{XY}}`$ are real coupling constants; $`V_\mathrm{\Lambda }(\{S_x^{(3)}\})`$ is a real function of the $`z`$-components, $`\{S_x^{(3)}\}_{x\mathrm{\Lambda }_s}`$, of the spins. For simplicity, we take
$$V_\mathrm{\Lambda }(\{S_x^{(3)}\})=\underset{\{x,y\}\mathrm{\Lambda }_b}{}J_{x,y}^\mathrm{Z}S_x^{(3)}S_y^{(3)}$$
(2.28)
with real coupling constants $`J_{x,y}^\mathrm{Z}`$. Assume that there are positive constants, $`J_{\mathrm{max}}^{\mathrm{XY}}`$ and $`J_{\mathrm{max}}^\mathrm{Z}`$, which satisfy $`|J_{x,y}^{\mathrm{XY}}|J_{\mathrm{max}}^{\mathrm{XY}}`$ and $`|J_{x,y}^\mathrm{Z}|J_{\mathrm{max}}^\mathrm{Z}`$ for any bond $`\{x,y\}\mathrm{\Lambda }_b`$.
Consider the transverse spin-spin correlation, $`S_x^+S_y^{}_0`$, where $`S_x^\pm :=S_x^{(1)}\pm iS_x^{(2)}`$.
###### Theorem 2.11
Assume that the fractal dimension $`D`$ of (2.25) satisfies $`1D<2`$, and that there exists a uniform spectral gap $`\mathrm{\Delta }E>0`$ above the ground state sector in the spectrum of the Hamiltonian $`H_\mathrm{\Lambda }`$ of (2.26) in the sense of Definition 2.5. Then there exists a positive constant $`\gamma `$ such that the transverse spin-spin correlation satisfies the bound,
$$\left|S_x^+S_y^{}_0\right|\mathrm{Const}.\mathrm{exp}\left[\gamma \{\mathrm{dist}(x,y)\}^{1D/2}\right],$$
(2.29)
in the thermodynamic limit $`|\mathrm{\Lambda }_s|\mathrm{}`$.
The proof is given in , and we remark that the result can be extended to more complicated correlations such as the multispin correlation, $`S_{x_1}^+\mathrm{}S_{x_j}^+S_{y_1}^{}\mathrm{}S_{y_j}^{}_0`$. If the system satisfies the self-similarity condition of Definition 2.9, then the upper bound, (2.29), can be replaced with a stronger exponentially decaying one by Theorem 2.10.
Hubbard model : The Hamiltonian on the lattice $`\mathrm{\Lambda }`$ is given by
$$H_\mathrm{\Lambda }=\underset{\{x,y\}\mathrm{\Lambda }_b}{}\underset{\alpha =,}{}\left(t_{x,y}c_{x,\alpha }^{}c_{y,\alpha }+t_{x,y}^{}c_{y,\alpha }^{}c_{x,\alpha }\right)+V(\{n_{x,\alpha }\})+\underset{x\mathrm{\Lambda }_s}{}𝐁_x𝐒_x,$$
(2.30)
where $`c_{x,\alpha }^{},c_{x,\alpha }`$ are, respectively, the electron creation and annihilation operators with the $`z`$ component of the spin $`\mu =,`$, $`n_{x,\alpha }=c_{x,\alpha }^{}c_{x,\alpha }`$ is the corresponding number operator, and $`𝐒_x=(S_x^{(1)},S_x^{(2)},S_x^{(3)})`$ are the spin operator given by $`S_x^{(a)}=_{\alpha ,\beta =,}c_{x,\alpha }^{}\sigma _{\alpha ,\beta }^{(a)}c_{x,\beta }`$ with the Pauli spin matrix $`(\sigma _{\alpha ,\beta }^{(a)})`$ for $`a=1,2,3`$; $`t_{i,j}𝐂`$ are the hopping amplitude, $`V(\{n_{x,\alpha }\})`$ is a real function of the number operators, and $`𝐁_x=(B_x^{(1)},B_x^{(2)},B_x^{(3)})𝐑^3`$ are local magnetic fields. Assume that the interaction $`V(\{n_{x,\alpha }\})`$ is of finite range in the sense of the graph theoretic distance.
###### Theorem 2.12
Assume that the fractal dimension $`D`$ of (2.25) satisfies $`1D<2`$, and that there exists a uniform spectral gap $`\mathrm{\Delta }E>0`$ above the ground state sector in the spectrum of the Hamiltonian $`H_\mathrm{\Lambda }`$ of (2.30) in the sense of Definition 2.5. Then the following bound is valid:
$$\left|c_{x,}^{}c_{x,}^{}c_{y,}c_{y,}_0\right|\mathrm{Const}.\mathrm{exp}\left[\gamma \{\mathrm{dist}(x,y)\}^{1D/2}\right]$$
with some constant $`\gamma `$ in the thermodynamic limit $`|\mathrm{\Lambda }_s|\mathrm{}`$. If the local magnetic field has the form $`𝐁_x=(0,0,B_x)`$, then we further have
$$\left|S_x^+S_y^{}_0\right|\mathrm{Const}.\mathrm{exp}\left[\gamma ^{}\{\mathrm{dist}(x,y)\}^{1D/2}\right]$$
(2.31)
with some constant $`\gamma ^{}`$.
The proof is given in . Clearly the Hamiltonian $`H_\mathrm{\Lambda }`$ of (2.30) commutes with the total number operator $`𝒩_\mathrm{\Lambda }=_{x\mathrm{\Lambda }_s}_{\mu =,}n_{x,\mu }`$ for a finite volume $`|\mathrm{\Lambda }_s|<\mathrm{}`$. We denote by $`H_{\mathrm{\Lambda },N}`$ the restriction of $`H_\mathrm{\Lambda }`$ onto the eigenspace of $`𝒩_\mathrm{\Lambda }`$ with the eigenvalue $`N`$. Let $`P_{0,\mathrm{\Lambda },N}`$ be the projection onto the ground state sector of $`H_{\mathrm{\Lambda },N}`$, and we denote the ground-state expectation by
$$\mathrm{}_{0,\nu }=\mathrm{weak}^{}\text{-}\underset{|\mathrm{\Lambda }_s|\mathrm{}}{lim}\frac{1}{q_N}\mathrm{Tr}(\mathrm{})P_{0,\mathrm{\Lambda },N},$$
(2.32)
where $`q_N`$ is the degeneracy of the ground state, and $`\nu `$ is the limit of the filling factor $`N/|\mathrm{\Lambda }_s|`$ of the electrons. Since the operators $`S_x^\pm `$ do not connect the sectors with the different eigenvalues $`N`$, we have
###### Theorem 2.13
Assume that the fractal dimension $`D`$ of (2.25) satisfies $`1D<2`$, and that there exists a uniform spectral gap $`\mathrm{\Delta }E>0`$ above the ground state sector in the spectrum of the Hamiltonian $`H_{\mathrm{\Lambda },N}`$ in the sense of Definition 2.5. Then the following bound is valid for the filling factor $`\nu `$ of the electrons:
$$\left|S_x^+S_y^{}_{0,\nu }\right|\mathrm{Const}.\mathrm{exp}\left[\gamma ^{}\{\mathrm{dist}(x,y)\}^{1D/2}\right]$$
(2.33)
with some constant $`\gamma ^{}`$ in the infinite volume limit.
The proof is given in . If the system satisfies the self-similarity condition of Definition 2.9, then these three upper bounds, (2.12), (2.31) and (2.33), can be replaced with a stronger exponentially decaying one by Theorem 2.10.
## 3 Clustering of correlations
In order to prove the power-law and the exponential clustering, Theorems 2.6 and 2.8, we follow the method . The key tools of the proof are Lemma 3.1 below and the Lieb-Robinson bound for the group velocity of the information propagation. The sketch of the proof is that the static correlation function can be derived from the time-dependent correlation function by the lemma, and the large-distance behavior of the time-dependent correlation function is estimated by the Lieb-Robinson bound. As a byproduct, we obtain the decay bound (2.18) for fermionic observables.
Consider first the case of the bosonic observables. Let $`A_X,B_Y`$ be bosonic observables with compact supports $`X,Y\mathrm{\Lambda }_s`$, respectively, and let $`A_X(t)=e^{itH_\mathrm{\Lambda }}A_Xe^{itH_\mathrm{\Lambda }}`$, where $`t𝐑`$ and $`H_\mathrm{\Lambda }`$ is the Hamiltonian for finite volume. Let $`\mathrm{\Phi }`$ be a normalized vector in the ground state sector. The ground state expectation of the commutator is written as
$`\mathrm{\Phi },[A_X(t),B_Y]\mathrm{\Phi }`$ $`=`$ $`\mathrm{\Phi },A_X(t)(1P_{0,\mathrm{\Lambda }})B_Y\mathrm{\Phi }\mathrm{\Phi },B_Y(1P_{0,\mathrm{\Lambda }})A_X(t)\mathrm{\Phi }`$ (3.1)
$`+`$ $`\mathrm{\Phi },A_X(t)P_{0,\mathrm{\Lambda }}B_Y\mathrm{\Phi }\mathrm{\Phi },B_YP_{0,\mathrm{\Lambda }}A_X(t)\mathrm{\Phi }.`$
In terms of the ground state vectors $`\mathrm{\Phi }_{0,\nu },\nu =1,2,\mathrm{},q`$, with the energy eigenvalues, $`E_{0,\nu }`$, and the excited state vectors $`\mathrm{\Phi }_n`$ with $`E_n,n=1,2,\mathrm{}`$, one has
$$\mathrm{\Phi },A_X(t)(1P_{0,\mathrm{\Lambda }})B_Y\mathrm{\Phi }=\underset{\nu ,\nu ^{}}{}\underset{n0}{}a_\nu ^{}a_\nu ^{}\mathrm{\Phi }_{0,\nu },A_X\mathrm{\Phi }_n\mathrm{\Phi }_n,B_Y\mathrm{\Phi }_{0,\nu ^{}}e^{it(E_nE_{0,\nu })},$$
(3.2)
$$\mathrm{\Phi },B_Y(1P_{0,\mathrm{\Lambda }})A_X(t)\mathrm{\Phi }=\underset{\nu ,\nu ^{}}{}\underset{n0}{}a_\nu ^{}a_\nu ^{}\mathrm{\Phi }_{0,\nu },B_Y\mathrm{\Phi }_n\mathrm{\Phi }_n,A_X\mathrm{\Phi }_{0,\nu ^{}}e^{it(E_nE_{0,\nu ^{}})},$$
(3.3)
$$\mathrm{\Phi },A_X(t)P_{0,\mathrm{\Lambda }}B_Y\mathrm{\Phi }=\underset{\nu ,\nu ^{}}{}\underset{\mu }{}a_\nu ^{}a_\nu ^{}\mathrm{\Phi }_{0,\nu },A_X\mathrm{\Phi }_{0,\mu }\mathrm{\Phi }_{0,\mu },B_Y\mathrm{\Phi }_{0,\nu ^{}}e^{it(E_{0,\mu }E_{0,\nu })}$$
(3.4)
and
$$\mathrm{\Phi },B_YP_{0,\mathrm{\Lambda }}A_X(t)\mathrm{\Phi }=\underset{\nu ,\nu ^{}}{}\underset{\mu }{}a_\nu ^{}a_\nu ^{}\mathrm{\Phi }_{0,\nu },B_Y\mathrm{\Phi }_{0,\mu }\mathrm{\Phi }_{0,\mu },A_X\mathrm{\Phi }_{0,\nu ^{}}e^{it(E_{0,\mu }E_{0,\nu ^{}})},$$
(3.5)
where we have written
$$\mathrm{\Phi }=\underset{\nu =1}{\overset{q}{}}a_\nu \mathrm{\Phi }_{0,\nu }.$$
(3.6)
In order to get the bound for $`\mathrm{\Phi },A_X(t=0)B_Y\mathrm{\Phi }`$, we want to extract only the “negative frequency part” (3.2) from the time-dependent correlation functions (3.1). For this purpose, we use the following lemma :
###### Lemma 3.1
Let $`E𝐑`$, and $`\alpha >0`$. Then
$`\underset{T\mathrm{}}{lim}\underset{ϵ0}{lim}{\displaystyle \frac{i}{2\pi }}{\displaystyle _T^T}{\displaystyle \frac{e^{iEt}e^{\alpha t^2}}{t+iϵ}}𝑑t`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\sqrt{{\displaystyle \frac{\pi }{\alpha }}}{\displaystyle _{\mathrm{}}^0}𝑑\omega \mathrm{exp}[(\omega +E)^2/(4\alpha )]`$ (3.7)
$`=`$ $`\{\begin{array}{cc}1+𝒪(\mathrm{exp}[\mathrm{\Delta }E^2/(4\alpha )])\hfill & \text{for }E\mathrm{\Delta }E\text{;}\hfill \\ 𝒪(\mathrm{exp}[\mathrm{\Delta }E^2/(4\alpha )])\hfill & \text{for }E\mathrm{\Delta }E\text{.}\hfill \end{array}`$
Proof: Write
$$I(E)=\frac{i}{2\pi }_T^T\frac{e^{iEt}e^{\alpha t^2}}{t+iϵ}𝑑t.$$
(3.8)
Using the Fourier transformation,
$$e^{iEt}e^{\alpha t^2}=\frac{1}{2\pi }\sqrt{\frac{\pi }{\alpha }}_{\mathrm{}}^{\mathrm{}}\mathrm{exp}[(\omega +E)^2/(4\alpha )]e^{i\omega t}𝑑\omega ,$$
(3.9)
we decompose the integral $`I(E)`$ into three parts as
$$I(E)=I_{}(E)+I_0(E)+I_+(E),$$
(3.10)
where
$$I_{}(E)=\frac{i}{2\pi }\frac{1}{2\pi }\sqrt{\frac{\pi }{\alpha }}_T^T𝑑t\frac{1}{t+iϵ}_{\mathrm{}}^{\mathrm{\Delta }\omega }𝑑\omega \mathrm{exp}[(\omega +E)^2/(4\alpha )]e^{i\omega t},$$
(3.11)
$$I_0(E)=\frac{i}{2\pi }\frac{1}{2\pi }\sqrt{\frac{\pi }{\alpha }}_T^T𝑑t\frac{1}{t+iϵ}_{\mathrm{\Delta }\omega }^{\mathrm{\Delta }\omega }𝑑\omega \mathrm{exp}[(\omega +E)^2/(4\alpha )]e^{i\omega t},$$
(3.12)
and
$$I_+(E)=\frac{i}{2\pi }\frac{1}{2\pi }\sqrt{\frac{\pi }{\alpha }}_T^T𝑑t\frac{1}{t+iϵ}_{\mathrm{\Delta }\omega }^{\mathrm{}}𝑑\omega \mathrm{exp}[(\omega +E)^2/(4\alpha )]e^{i\omega t},$$
(3.13)
where we choose $`\mathrm{\Delta }\omega =bT^{1/2}`$ with some positive constant $`b`$.
First let us estimate $`I_0(E)`$. Note that
$$\frac{1}{t+iϵ}=\frac{t}{t^2+ϵ^2}\frac{iϵ}{t^2+ϵ^2}.$$
(3.14)
Using this identity, one has
$$I_0(E)=\frac{i}{2\pi }\frac{1}{2\pi }\sqrt{\frac{\pi }{\alpha }}_{\mathrm{\Delta }\omega }^{\mathrm{\Delta }\omega }𝑑\omega \mathrm{exp}[(\omega +E)^2/(4\alpha )]_T^T𝑑t\left[\frac{t\mathrm{sin}\omega t}{t^2+ϵ^2}\frac{iϵ\mathrm{cos}\omega t}{t^2+ϵ^2}\right],$$
(3.15)
where we have interchanged the order of the double integral by relying on $`|t|T<\mathrm{}`$. Since the integral about $`t`$ can be bounded by some constant, one obtains
$$|I_0(E)|\mathrm{Const}.\times \alpha ^{1/2}\mathrm{\Delta }\omega \mathrm{Const}.\times \alpha ^{1/2}T^{1/2}.$$
(3.16)
Therefore the corresponding contribution is vanishing in the limit $`T\mathrm{}`$.
Note that
$$\frac{i}{2\pi }_T^T𝑑t\frac{e^{i\omega t}}{t+iϵ}=\{\begin{array}{cc}𝒪(\omega ^1T^1)\hfill & \text{for }\omega >0\text{;}\hfill \\ e^{ϵ\omega }+𝒪(\omega ^1T^1)\hfill & \text{for }\omega <0\text{.}\hfill \end{array}$$
(3.17)
Using this, the function $`I_+(E)`$ of (3.13) can be evaluated as
$$|I_+(E)|\mathrm{Const}.\times T^{1/2}.$$
(3.18)
This is also vanishing in the limit.
Thus it is enough to consider only the integral $`I_{}(E)`$. In the same way as the above, one has
$$I_{}(E)=\frac{1}{2\pi }\sqrt{\frac{\pi }{\alpha }}_{\mathrm{}}^{\mathrm{\Delta }\omega }𝑑\omega \mathrm{exp}[(\omega +E)^2/(4\alpha )]e^{ϵ\omega }+𝒪(T^{1/2}).$$
(3.19)
Since $`e^{ϵ\omega }1`$ for $`\omega <0`$, one has
$$\underset{T\mathrm{}}{lim}\underset{ϵ0}{lim}I_{}(E)=\frac{1}{2\pi }\sqrt{\frac{\pi }{\alpha }}_{\mathrm{}}^0𝑑\omega \mathrm{exp}[(\omega +E)^2/(4\alpha )].$$
(3.20)
Note that, for $`E\mathrm{\Delta }E`$,
$$\frac{1}{2\pi }\sqrt{\frac{\pi }{\alpha }}_{\mathrm{}}^0𝑑\omega \mathrm{exp}[(\omega +E)^2/(4\alpha )]\frac{1}{2}\mathrm{exp}[\mathrm{\Delta }E^2/(4\alpha )],$$
(3.21)
and, for $`E\mathrm{\Delta }E`$,
$`{\displaystyle \frac{1}{2\pi }}\sqrt{{\displaystyle \frac{\pi }{\alpha }}}{\displaystyle _{\mathrm{}}^0}𝑑\omega \mathrm{exp}[(\omega +E)^2/(4\alpha )]`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\sqrt{{\displaystyle \frac{\pi }{\alpha }}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega \mathrm{exp}[(\omega +E)^2/(4\alpha )]`$ (3.22)
$``$ $`{\displaystyle \frac{1}{2\pi }}\sqrt{{\displaystyle \frac{\pi }{\alpha }}}{\displaystyle _0^{\mathrm{}}}𝑑\omega \mathrm{exp}[(\omega +E)^2/(4\alpha )]`$
$`=`$ $`1+𝒪(\mathrm{exp}[\mathrm{\Delta }E^2/(4\alpha )]).`$
Clearly these imply (3.7).
From Lemma 3.1 and the expression (3.1) of the correlation function with (3.2) and (3.3), one has
$`\underset{T\mathrm{}}{lim}\underset{ϵ0}{lim}{\displaystyle \frac{i}{2\pi }}{\displaystyle _T^T}𝑑t{\displaystyle \frac{1}{t+iϵ}}\mathrm{\Phi },[A_X(t),B_Y]\mathrm{\Phi }e^{\alpha t^2}`$ (3.23)
$`=`$ $`\mathrm{\Phi },A_X(1P_{0,\mathrm{\Lambda }})B_Y\mathrm{\Phi }+𝒪(\mathrm{exp}[\mathrm{\Delta }E^2/(4\alpha )])`$
$`+`$ $`\underset{T\mathrm{}}{lim}\underset{ϵ0}{lim}{\displaystyle \frac{i}{2\pi }}{\displaystyle _T^T}𝑑t{\displaystyle \frac{1}{t+iϵ}}\left[\mathrm{\Phi },A_X(t)P_{0,\mathrm{\Lambda }}B_Y\mathrm{\Phi }\mathrm{\Phi },B_YP_{0,\mathrm{\Lambda }}A_X(t)\mathrm{\Phi }\right]e^{\alpha t^2}`$
for finite volume.
In the following, we treat only the power-law decaying interaction $`h_X`$ because one can treat the exponentially decaying interactions in the same way. See also refs. in which the exponential clustering of the correlations is proved for finite-range interactions under the gap assumption along the same line as below.
In order to estimate the left-hand side, we recall the Lieb-Robinson estimate (A.1) in Appendix A,
$$\frac{1}{t}[A_X(t),B_Y]\mathrm{Const}.\times \frac{1}{(1+r)^\eta }\frac{e^{v|t|}1}{|t|},$$
(3.24)
for $`r>0`$, where we have written $`r=\mathrm{dist}(X,Y)`$. Using this estimate, the integral can be evaluated as
$`\left|{\displaystyle _T^T}𝑑t{\displaystyle \frac{\mathrm{\Phi },[A_X(t),B_Y]\mathrm{\Phi }}{t+iϵ}}e^{\alpha t^2}\right|`$ (3.25)
$``$ $`\left|{\displaystyle _{|t|c\mathrm{}}}𝑑t{\displaystyle \frac{\mathrm{\Phi },[A_X(t),B_Y]\mathrm{\Phi }}{t+iϵ}}e^{\alpha t^2}\right|+\left|{\displaystyle _{|t|>c\mathrm{}}}𝑑t{\displaystyle \frac{\mathrm{\Phi },[A_X(t),B_Y]\mathrm{\Phi }}{t+iϵ}}e^{\alpha t^2}\right|`$
$``$ $`\mathrm{Const}.\times {\displaystyle \frac{1}{(1+r)^{\eta cv}}}+{\displaystyle \frac{\mathrm{Const}.}{\sqrt{\alpha }\mathrm{}}}\mathrm{exp}[\alpha c^2\mathrm{}^2],`$
where $`c`$ is a positive, small parameter, and $`\mathrm{}=\mathrm{log}(1+r)`$, and we have used
$$_{|t|c\mathrm{}}\frac{e^{v|t|}1}{|t|}𝑑t2e^{cv\mathrm{}}.$$
(3.26)
In order to estimate the integral in the right-hand side of (3.23), we consider the matrix element $`\mathrm{\Phi }_{0,\nu },A_X(t)P_{0,\mathrm{\Lambda }}B_Y\mathrm{\Phi }_{0,\nu ^{}}`$ because the other matrix elements in the ground state can be treated in the same way. Using Lemma 3.1, one has
$`\underset{T\mathrm{}}{lim}\underset{ϵ0}{lim}{\displaystyle \frac{i}{2\pi }}{\displaystyle _T^T}𝑑t{\displaystyle \frac{1}{t+iϵ}}\mathrm{\Phi }_{0,\nu },A_X(t)P_{0,\mathrm{\Lambda }}B_Y\mathrm{\Phi }_{0,\nu ^{}}e^{\alpha t^2}`$
$`=`$ $`{\displaystyle \underset{\mu =1}{\overset{q}{}}}\mathrm{\Phi }_{0,\nu },A_X\mathrm{\Phi }_{0,\mu }\mathrm{\Phi }_{0,\mu },B_Y\mathrm{\Phi }_{0,\nu ^{}}{\displaystyle \frac{1}{2\pi }}\sqrt{{\displaystyle \frac{\pi }{\alpha }}}{\displaystyle _{\mathrm{}}^0}𝑑\omega \mathrm{exp}[(\omega +\mathrm{\Delta }_{\mu ,\nu })^2/(4\alpha )],`$
where $`\mathrm{\Delta }_{\mu ,\nu }=E_{0,\mu }E_{0,\nu }`$. Using the assumption (2.12) and the dominated convergence theorem, we have that, for any given $`\epsilon >0`$, there exists a sufficiently large volume of the lattice $`\mathrm{\Lambda }_s`$ such that
$$\left|\underset{T\mathrm{}}{lim}\underset{ϵ0}{lim}\frac{i}{2\pi }_T^T𝑑t\frac{e^{\alpha t^2}}{t+iϵ}\mathrm{\Phi }_{0,\nu },A_X(t)P_{0,\mathrm{\Lambda }}B_Y\mathrm{\Phi }_{0,\nu ^{}}\frac{1}{2}\mathrm{\Phi }_{0,\nu },A_XP_{0,\mathrm{\Lambda }}B_Y\mathrm{\Phi }_{0,\nu ^{}}\right|<\epsilon .$$
(3.28)
Combining this observation, (3.23) and (3.25), and choosing $`\alpha =\mathrm{\Delta }E/(2c\mathrm{})`$, one obtains
$`\left|\omega (A_XB_Y){\displaystyle \frac{1}{2}}\left[\omega (A_XP_0B_Y)+\omega (B_YP_0A_X)\right]\right|`$ (3.29)
$``$ $`\mathrm{Const}.\times {\displaystyle \frac{1}{(1+r)^{\eta cv}}}+\mathrm{Const}.\times \mathrm{exp}[{\displaystyle \frac{c\mathrm{\Delta }E}{2}}\mathrm{}]`$
in the infinite volume limit, where the ground-state expectation $`\omega `$ is given by (2.15). Choosing $`c=\eta /(v+\mathrm{\Delta }E/2)`$, we have
$$\left|\omega (A_XB_Y)\frac{1}{2}\left[\omega (A_XP_0B_Y)+\omega (B_YP_0A_X)\right]\right|\frac{\mathrm{Const}.}{[1+\mathrm{dist}(X,Y)]^{\stackrel{~}{\eta }}},$$
(3.30)
with $`\stackrel{~}{\eta }=\eta /(1+2v/\mathrm{\Delta }E)`$. In the same way, we have
$$|\omega (A_XB_Y)\frac{1}{2}[\omega (A_XP_0B_Y)+\omega (B_YP_0A_X)]|\mathrm{Const}.\times \mathrm{exp}[\stackrel{~}{\mu }\mathrm{dist}(X,Y)]$$
(3.31)
for the exponentially decaying interaction $`h_X`$, where $`\stackrel{~}{\mu }=\mu /(1+2v/\mathrm{\Delta }E)`$. This proves Theorem 2.8. The corresponding bound for finite-range interactions was already obtained in . Using the definition (2.13) of the expectation $`\mathrm{}_{0,\mathrm{\Lambda }}`$ and the identity,
$$A_XP_0B_Y_{0,\mathrm{\Lambda }}=B_YP_0A_X_{0,\mathrm{\Lambda }},$$
(3.32)
for the integral in the right-hand side of (3.23), we obtain
$`\left|A_XB_Y_{0,\mathrm{\Lambda }}A_XP_{0,\mathrm{\Lambda }}B_Y_{0,\mathrm{\Lambda }}\right|`$ (3.33)
$``$ $`\mathrm{Const}.\times \{\begin{array}{cc}[1+\mathrm{dist}(X,Y)]^{\stackrel{~}{\eta }},\hfill & \text{for power-law decaying }h_X\text{;}\hfill \\ \mathrm{exp}[\stackrel{~}{\mu }\mathrm{dist}(X,Y)],\hfill & \text{for exponentially decaying }h_X\hfill \end{array}`$
for any finite lattice $`\mathrm{\Lambda }_sX,Y`$ in the same way as in the above.
Next consider the case that the pair, $`A_X,B_Y`$, is fermionic. Note that
$`\mathrm{\Phi }_{0,\nu },\{A_X(t),B_Y\}\mathrm{\Phi }_{0,\nu ^{}}`$ (3.34)
$`=`$ $`\mathrm{\Phi }_{0,\nu },A_X(t)(1P_{0,\mathrm{\Lambda }})B_Y\mathrm{\Phi }_{0,\nu ^{}}+\mathrm{\Phi }_{0,\nu },B_Y(1P_{0,\mathrm{\Lambda }})A_X(t)\mathrm{\Phi }_{0,\nu ^{}}`$
$`+`$ $`\mathrm{\Phi }_{0,\nu },A_X(t)P_{0,\mathrm{\Lambda }}B_Y\mathrm{\Phi }_{0,\nu ^{}}+\mathrm{\Phi }_{0,\nu },B_YP_{0,\mathrm{\Lambda }}A_X(t)\mathrm{\Phi }_{0,\nu ^{}}.`$
Since the difference between bosonic and fermionic observables is in the signs of some terms, one has
$`\left|\omega (A_XB_Y){\displaystyle \frac{1}{2}}\left[\omega (A_XP_0B_Y)\omega (B_YP_0A_X)\right]\right|`$ (3.35)
$``$ $`\mathrm{Const}.\times \{\begin{array}{cc}[1+\mathrm{dist}(X,Y)]^{\stackrel{~}{\eta }},\hfill & \text{for power-law decaying }h_X\text{;}\hfill \\ \mathrm{exp}[\stackrel{~}{\mu }\mathrm{dist}(X,Y)],\hfill & \text{for exponentially decaying }h_X\text{.}\hfill \end{array}`$
In particular, thanks to the identity (3.32), we obtain
$`\left|A_XB_Y_0\right|`$ $``$ $`\mathrm{Const}.\times \{\begin{array}{cc}[1+\mathrm{dist}(X,Y)]^{\stackrel{~}{\eta }},\hfill & \text{for power-law decaying }h_X\text{;}\hfill \\ \mathrm{exp}[\stackrel{~}{\mu }\mathrm{dist}(X,Y)],\hfill & \text{for exponentially decaying }h_X\text{.}\hfill \end{array}`$
This is nothing but the desired bound. We stress that, for infinite degeneracy of the infinite-volume ground state, this upper bound is also justified in the same argument with the dominated convergence theorem.
## 4 Vanishing of the matrix elements in the ground state
The aim of this section is to prove the bound (2.23) for the correlation and discuss an extension of Theorem 2.10 to a system having infinite degeneracy of infinite-volume ground state. The latter result is summarized as Theorem 4.1 below. We will give only the proof of Theorem 4.1 because Theorem 2.10 is proved in the same way. By the clustering bounds (3.30) and (3.31), it is sufficient to show that all the matrix elements, $`\mathrm{\Phi }_{0,\nu ^{}}A_X\mathrm{\Phi }_{0,\nu }`$, in the sector of the ground state are vanishing. The key idea of the proof is to estimate the absolute values of the matrix elements by using the self-similarity condition and the decay bound (4.2) below of the correlations at a sufficiently large distance.
We denote by $`q_\mathrm{\Lambda }`$ the degeneracy of the sector of the ground state for the finite lattice $`\mathrm{\Lambda }`$, and we allow $`q_\mathrm{\Lambda }\mathrm{}`$ as $`|\mathrm{\Lambda }_s|\mathrm{}`$. We write $`m=q_\mathrm{\Lambda }^2`$ for short. To begin with, we write the bound (3.33) as
$$\left|A_XB_Y_{0,\mathrm{\Lambda }}A_XP_{0,\mathrm{\Lambda }}B_Y_{0,\mathrm{\Lambda }}\right|G_0(\mathrm{dist}(X,Y)),$$
(4.1)
where we have written the upper bound of the right-hand side by the function $`G_0`$ of the distance. We assume that the following bound holds:
$$\left|A_XB_Y_{0,\mathrm{\Lambda }}\right|G_1(\mathrm{dist}(X,Y))$$
(4.2)
with an upper bound $`G_1`$ which is vanishing at the infinite distance. Further we define $`\stackrel{~}{G}_\mathrm{\Lambda }`$ as
$$\stackrel{~}{G}_\mathrm{\Lambda }(A_X,B_Y):=\mathrm{max}\{G_0(\mathrm{dist}(X,Y)),G_1(\mathrm{dist}(X,Y))\}.$$
(4.3)
###### Theorem 4.1
Let $`\omega `$ be a ground-state expectation (2.15) in the infinite volume limit, and let $`A_X,B_Y`$ be a pair of bosonic observable with compact supports $`X,Y`$. Assume that there exists a uniform spectral gap $`\mathrm{\Delta }E>0`$ above the ground state sector in the spectrum of the Hamiltonian $`H_\mathrm{\Lambda }`$ in the sense of Definition 2.5. Suppose that, for any given $`ϵ>0`$, there exists $`M_0>0`$ such that, for any large lattice $`\mathrm{\Lambda }`$ satisfying $`|\mathrm{\Lambda }_s|M_0`$, there exists a set of observables, $`B^{(j)}`$, $`j=1,2,\mathrm{},m`$, and a set of transformations, $`R_j`$, $`j=1,2,\mathrm{},m`$, satisfying the following conditions: Any pair of the observables, $`A_X,B^{(1)},\mathrm{},B^{(m)}`$, is bosonic,
$$B^{(j)}=R_j(A),(B^{(j)})^{}=R_j(A_X^{})\text{and}R_j(H_\mathrm{\Lambda })=H_\mathrm{\Lambda },$$
(4.4)
and
$$q_\mathrm{\Lambda }^3\underset{i,j\{0,1,\mathrm{},m\}:}{\mathrm{max}}_{ij}\left\{\stackrel{~}{G}_\mathrm{\Lambda }(\left(B^{(i)}\right)^{},B^{(j)})\right\}<ϵ,$$
(4.5)
where we have written $`B^{(0)}=A_X^{}`$. Then we have the bound,
$$\left|\omega (A_XB_Y)\right|\mathrm{Const}.\times \{\begin{array}{cc}[1+\mathrm{dist}(X,Y)]^{\stackrel{~}{\eta }},\hfill & \text{for power-law decaying }h_X\text{;}\hfill \\ \mathrm{exp}[\stackrel{~}{\mu }\mathrm{dist}(X,Y)],\hfill & \text{for exponentially decaying }h_X\text{,}\hfill \end{array}$$
(4.6)
in the infinite volume limit.
Proof: From the bound (3.30) or (3.31) and the Schwarz inequality,
$$\left|\omega (A_XP_0B_Y)\right|^2\omega (A_XP_0A_X^{})\omega (B_Y^{}P_0B_Y),$$
(4.7)
it is sufficient to show $`\omega (A_XP_0A_X^{})=\omega (A_X^{}P_0A_X)=0`$. Further, we have
$$\mathrm{\Phi },AP_{0,\mathrm{\Lambda }}A^{}\mathrm{\Phi }q_\mathrm{\Lambda }AP_{0,\mathrm{\Lambda }}A^{}_{0,\mathrm{\Lambda }}=q_\mathrm{\Lambda }A^{}P_{0,\mathrm{\Lambda }}A_{0,\mathrm{\Lambda }}$$
(4.8)
for any ground state vector $`\mathrm{\Phi }`$ with norm one and any observable $`A`$ on the finite lattice $`\mathrm{\Lambda }`$. Therefore we estimate $`q_\mathrm{\Lambda }A_XP_{0,\mathrm{\Lambda }}A_X^{}_{0,\mathrm{\Lambda }}`$.
Note that, from the clustering bound (4.1), (4.2) and (4.3), we have
$`\left|A_XP_{0,\mathrm{\Lambda }}B_Y_{0,\mathrm{\Lambda }}\right|`$ $``$ $`\left|A_XB_Y_{0,\mathrm{\Lambda }}\right|+\left|A_XB_Y_{0,\mathrm{\Lambda }}A_XP_{0,\mathrm{\Lambda }}B_Y_{0,\mathrm{\Lambda }}\right|`$ (4.9)
$``$ $`2\stackrel{~}{G}_\mathrm{\Lambda }(A_X,B_Y).`$
We define
$$B_i^{(j)}:=\mathrm{\Phi }_{0,\nu ^{}},B^{(j)}\mathrm{\Phi }_{0,\nu }$$
(4.10)
for $`j=0,1,\mathrm{},m`$ and for the finite lattice $`\mathrm{\Lambda }`$, where we have written $`i=(\nu ^{},\nu )`$ with $`i=1,2,\mathrm{},m`$ for short. Since $`(B_1^{(j)},B_2^{(j)},\mathrm{},B_m^{(j)})`$ is an $`m`$-dimensional vector, there exist complex numbers, $`C_j,j=0,1,\mathrm{},m`$, such that, at least, one of $`C_j`$ is nonvanishing and that
$$\underset{j=0}{\overset{m}{}}C_jB_i^{(j)}=0.$$
(4.11)
Let $`\mathrm{}`$ be the index which satisfies $`|C_{\mathrm{}}|=\mathrm{max}\{|C_0|,|C_1|,\mathrm{},|C_m|\}`$. Clearly, we have
$$B_i^{(\mathrm{})}=\underset{j\mathrm{}}{}\frac{C_j}{C_{\mathrm{}}}B_i^{(j)}.$$
(4.12)
Therefore
$`(B^{(\mathrm{})})^{}P_{0,\mathrm{\Lambda }}B^{(\mathrm{})}_{0,\mathrm{\Lambda }}={\displaystyle \frac{1}{q_\mathrm{\Lambda }}}{\displaystyle \underset{i=1}{\overset{m}{}}}\left|B_i^{(\mathrm{})}\right|^2`$ $`=`$ $`{\displaystyle \underset{j\mathrm{}}{}}{\displaystyle \frac{C_j}{C_{\mathrm{}}}}{\displaystyle \frac{1}{q_\mathrm{\Lambda }}}{\displaystyle \underset{i=1}{\overset{m}{}}}(B_i^{(\mathrm{})})^{}B_i^{(j)}`$ (4.13)
$``$ $`m\underset{j\mathrm{}}{\mathrm{max}}\left\{\left|(B^{(\mathrm{})})^{}P_{0,\mathrm{\Lambda }}B^{(j)}_{0,\mathrm{\Lambda }}\right|\right\}`$
$``$ $`2q_\mathrm{\Lambda }^2\underset{j\mathrm{}}{\mathrm{max}}\left\{\stackrel{~}{G}_\mathrm{\Lambda }((B^{(\mathrm{})})^{},B^{(j)})\right\},`$
where we have used the inequality (4.9) for getting the last bound. When $`\mathrm{}=0`$, we obtain
$$q_\mathrm{\Lambda }A_XP_{0,\mathrm{\Lambda }}A_X^{}_{0,\mathrm{\Lambda }}2ϵ$$
(4.14)
from $`B^{(0)}=A_X^{}`$ and the assumption (4.5). When $`\mathrm{}0`$, we reach the same conclusion by using the relation,
$$A_X^{}P_{0,\mathrm{\Lambda }}A_X_{0,\mathrm{\Lambda }}=R_{\mathrm{}}(A_X^{})P_{0,\mathrm{\Lambda }}R_{\mathrm{}}(A)_{0,\mathrm{\Lambda }}=\left(B^{(\mathrm{})}\right)^{}P_{0,\mathrm{\Lambda }}B^{(\mathrm{})}_{0,\mathrm{\Lambda }},$$
(4.15)
which is derived from the assumption (4.4).
Remark: 1. The advantage of Theorem 4.1 is that it is easier to find $`B^{(j)}`$ and $`R_j`$ because of the finiteness of the lattice. Actually one can construct $`B^{(j)}`$, $`R_j`$ and $`\mathrm{\Lambda }`$ satisfying the requirement by connecting $`m`$ copies of a small, finite lattice to each other at their boundaries. But, if the degeneracy $`q_\mathrm{\Lambda }`$ exceeds $`\sqrt{|\mathrm{\Lambda }_s|}`$, we cannot find the observables, $`B^{(j)}`$, and the transformations, $`R_j`$. Therefore our argument does not work in such cases.
2. Under the weaker assumption,
$$q_\mathrm{\Lambda }^2\underset{i,j\{0,1,\mathrm{},m\}:}{\mathrm{max}}_{ij}\left\{\stackrel{~}{G}_\mathrm{\Lambda }(\left(B^{(i)}\right)^{},B^{(j)})\right\}<ϵ,$$
(4.16)
than (4.5), we can obtain the bound,
$$\left|A_XB_Y_0\right|\mathrm{Const}.\times \{\begin{array}{cc}[1+\mathrm{dist}(X,Y)]^{\stackrel{~}{\eta }},\hfill & \text{for power-law decaying }h_X\text{;}\hfill \\ \mathrm{exp}[\stackrel{~}{\mu }\mathrm{dist}(X,Y)],\hfill & \text{for exponentially decaying }h_X\text{,}\hfill \end{array}$$
(4.17)
in the infinite volume limit.
3. Consider the situation of the above Remark 2 or the case with a finite degeneracy of the infinite-volume ground state. Then, instead of introducing the transformations $`R_j`$, we can directly require
$$A_X^{}P_0A_X_0=\left(B^{(j)}\right)^{}P_0B^{(j)}_0\text{for }j=1,2,\mathrm{},m,$$
(4.18)
in the infinite volume limit, and at infinite distance between the observables $`A_X`$ and $`B^{(j)}`$.
## Appendix A Lieb-Robinson bound for group velocity
Quite recently, Nachtergaele and Sims have extended the Lieb-Robinson bound to a wide class of models with long-range, exponentially decaying interactions. In this appendix, we further extend the bound to the power-law decaying interactions. We also tighten the bound on the exponentially decaying case. (See Assumption 2.4 compared to that in ). However, in our proof, the time $`t`$ must be real.
In the following, we treat only the case with bosonic observables and with the power-law decaying interaction $`h_X`$ because the other cases including the previous results can be treated in the same way.
###### Theorem A.1
Let $`A_X,B_Y`$ be a pair of bosonic observables with the compact support, $`X,Y`$, respectively. Assume that the system satisfies the conditions in Assumption 2.1 or 2.2. Then
$$[A_X(t),B_Y]CA_XB_Y|X||Y|\frac{e^{v|t|}1}{[1+\mathrm{dist}(X,Y)]^\eta }\text{for }\mathrm{dist}(X,Y)>0,$$
(A.1)
where the positive constants, $`C`$ and $`v`$, depend only on the interaction of the Hamiltonian and the metric of the lattice.
Remark: The same bound for fermionic observables is obtained by replacing the commutator with the anticommutator in the left-hand side.
For exponentially decaying interaction $`h_X`$, the following bound is valid:
###### Theorem A.2
Let $`A_X,B_Y`$ be a pair of bosonic observables with the compact support, $`X,Y`$, respectively. Assume that the system satisfies the conditions in Assumption 2.3 or 2.4. Then
$$[A_X(t),B_Y]CA_XB_Y|X||Y|\mathrm{exp}[\mu \mathrm{dist}(X,Y)]\left[e^{v|t|}1\right]\text{for }\mathrm{dist}(X,Y)>0,$$
(A.2)
where the positive constants, $`C`$ and $`v`$, depend only on the interaction of the Hamiltonian and the metric of the lattice.
Remark: For the proof under Assumption 2.3, we rely on the inequalities, (2.9) and (2.10). Assumption 2.4 is milder than that in ref. as remarked in Section 2.
We assume that the volume $`|\mathrm{\Lambda }_s|`$ of the lattice $`\mathrm{\Lambda }`$ is finite. If it is necessary to consider the infinite volume limit, we take the limit after deriving the desired Lieb-Robinson bounds which hold uniformly in the size of the lattice. Let $`A,B`$ be observables supported by compact sets, $`X,Y\mathrm{\Lambda }_s`$, respectively. The time evolution of $`A`$ is given by $`A(t)=e^{itH_\mathrm{\Lambda }}Ae^{itH_\mathrm{\Lambda }}`$. First, let us derive the inequality (A.12) below for the commutator $`[A(t),B]`$. We assume $`t>0`$ because the negative $`t`$ can be treated in the same way. Let $`ϵ=t/N`$ with a large positive integer $`N`$, and let
$$t_n=\frac{t}{N}n\text{for}n=0,1,\mathrm{},N.$$
(A.3)
Then we have
$$[A(t),B][A(0),B]=\underset{i=0}{\overset{N1}{}}ϵ\times \frac{[A(t_{n+1}),B][A(t_n),B]}{ϵ}.$$
(A.4)
In order to obtain the bound (A.12) below, we want to estimate the summand in the right-hand side. To begin with, we note that the identity, $`U^{}OU=O`$, holds for any observable $`O`$ and for any unitary operator $`U`$. Using this fact, we have
$`[A(t_{n+1}),B][A(t_n),B]`$ $`=`$ $`[A(ϵ),B(t_n)][A,B(t_n)]`$
$``$ $`[A+iϵ[H_\mathrm{\Lambda },A],B(t_n)][A,B(t_n)]+𝒪(ϵ^2)`$
$`=`$ $`[A+iϵ[I_X,A],B(t_n)][A,B(t_n)]+𝒪(ϵ^2)`$
with
$$I_X=\underset{Z:ZX\mathrm{}}{}h_Z,$$
(A.6)
where we have used
$$A(ϵ)=A+iϵ[H_\mathrm{\Lambda },A]+𝒪(ϵ^2)$$
(A.7)
and the triangle inequality. Further, by using
$$A+iϵ[I_X,A]=e^{iϵI_X}Ae^{iϵI_X}+𝒪(ϵ^2),$$
(A.8)
we have
$`[A+iϵ[I_X,A],B(t_n)]`$ $``$ $`[e^{iϵI_X}Ae^{iϵI_X},B(t_n)]+𝒪(ϵ^2)`$ (A.9)
$`=`$ $`[A,e^{iϵI_X}B(t_n)e^{iϵI_X}]+𝒪(ϵ^2)`$
$``$ $`[A,B(t_i)iϵ[I_X,B(t_n)]]+𝒪(ϵ^2)`$
$``$ $`[A,B(t_n)]+ϵ[A,[I_X,B(t_n)]]+𝒪(ϵ^2).`$
Substituting this into the right-hand side in the last line of (LABEL:difnorm), we obtain
$`[A(t_{n+1}),B][A(t_n),B]`$ $``$ $`ϵ[A,[I_X,B(t_n)]]+𝒪(ϵ^2)`$ (A.10)
$``$ $`2ϵA[I_X(t_n),B]+𝒪(ϵ^2).`$
Further, substituting this into the right-hand side of (A.4) and using (A.6), we have
$`[A(t),B][A(0),B]`$ $``$ $`2A{\displaystyle \underset{n=0}{\overset{N1}{}}}ϵ\times [I_X(t_n),B]+𝒪(ϵ)`$ (A.11)
$``$ $`2A{\displaystyle \underset{Z:ZX\mathrm{}}{}}{\displaystyle \underset{n=0}{\overset{N1}{}}}ϵ\times [h_Z(t_n),B]+𝒪(ϵ).`$
Since $`h_Z(t)`$ is the continuous function of the time $`t`$ for a finite volume, the sum in the right-hand side converges to the integral in the limit $`ϵ0`$ ($`N\mathrm{}`$) for any fixed finite lattice $`\mathrm{\Lambda }`$. In consequence, we obtain
$$[A(t),B][A(0),B]2A\underset{Z:ZX\mathrm{}}{}_0^{|t|}𝑑s[h_Z(s),B].$$
(A.12)
We define
$$C_B(X,t):=\underset{A𝒜_X}{sup}\frac{[A(t),B]}{A},$$
(A.13)
where $`𝒜_X`$ is the set of observables supported by the compact set $`X`$. Then we have<sup>3</sup><sup>3</sup>3Since the local interaction $`h_Z`$ with $`ZX`$ does not change the support $`X`$ of $`A`$ in the time evolution, we can expect that the sum in the right-hand side of (A.14) can be restricted to the set $`Z`$ satisfying $`ZX\mathrm{}`$ and $`ZX\mathrm{}`$. However, this restriction does not affect the resulting Lieb-Robinson bound. Therefore we omit the discussion.
$$C_B(X,t)C_B(X,0)+2\underset{Z:ZX\mathrm{}}{}h_Z_0^{|t|}𝑑sC_B(Z,s)$$
(A.14)
from the above bound (A.12).
We recall that the observables, $`A`$ and $`B`$, are, respectively, supported by the compact sets, $`X,Y\mathrm{\Lambda }_s`$. Assume $`\mathrm{dist}(X,Y)>0`$. Then we have $`C_B(X,0)=0`$ from the definition of $`C_B(X,t)`$, and note that
$$C_B(Z,0)\{\begin{array}{cc}2B,\hfill & \text{for }ZY\mathrm{}\text{;}\hfill \\ 0,\hfill & \text{otherwise.}\hfill \end{array}$$
(A.15)
Using these facts and the above bound (A.14) iteratively, we obtain
$`C_B(X,t)`$ $``$ $`2{\displaystyle \underset{Z_1:Z_1X\mathrm{}}{}}h_{Z_1}{\displaystyle _0^{|t|}}𝑑s_1C_B(Z_1,s_1)`$
$``$ $`2{\displaystyle \underset{Z_1:Z_1X\mathrm{}}{}}h_{Z_1}{\displaystyle _0^{|t|}}𝑑s_1C_B(Z_1,0)`$
$`+`$ $`2^2{\displaystyle \underset{Z_1:Z_1X\mathrm{}}{}}h_{Z_1}{\displaystyle \underset{Z_2:Z_2Z_1\mathrm{}}{}}h_{Z_2}{\displaystyle _0^{|t|}}𝑑s_1{\displaystyle _0^{|s_1|}}𝑑s_2C_B(Z_2,s_2)`$
$``$ $`2B(2|t|){\displaystyle \underset{Z_1:Z_1X\mathrm{},Z_1Y\mathrm{}}{}}h_{Z_1}`$
$`+`$ $`2B{\displaystyle \frac{(2|t|)^2}{2!}}{\displaystyle \underset{Z_1:Z_1X\mathrm{}}{}}h_{Z_1}{\displaystyle \underset{Z_2:Z_2Z_1\mathrm{},Z_2Y\mathrm{}}{}}h_{Z_2}`$
$`+`$ $`2B{\displaystyle \frac{(2|t|)^3}{3!}}{\displaystyle \underset{Z_1:Z_1X\mathrm{}}{}}h_{Z_1}{\displaystyle \underset{Z_2:Z_2Z_1\mathrm{}}{}}h_{Z_2}{\displaystyle \underset{Z_3:Z_3Z_2\mathrm{},Z_3Y\mathrm{}}{}}h_{Z_3}+\mathrm{}`$
Proof of Theorem A.1 under Assumption 2.1 The first sum in the power series (LABEL:CBexpabound) is estimated as
$`{\displaystyle \underset{Z_1:Z_1X\mathrm{},Z_1Y\mathrm{}}{}}h_{Z_1}`$ $``$ $`{\displaystyle \underset{xX}{}}{\displaystyle \underset{yY}{}}{\displaystyle \underset{Z_1x,y}{}}h_{Z_1}`$ (A.17)
$``$ $`{\displaystyle \frac{\lambda _0|X||Y|}{[1+\mathrm{dist}(X,Y)]^\eta }}`$
from the assumption (2.2). The second, double sum is estimated as
$`{\displaystyle \underset{Z_1:Z_1X\mathrm{}}{}}h_{Z_1}{\displaystyle \underset{Z_2:Z_2Z_1\mathrm{},Z_2Y\mathrm{}}{}}h_{Z_2}`$ (A.18)
$``$ $`{\displaystyle \underset{xX}{}}{\displaystyle \underset{yY}{}}{\displaystyle \underset{z_{12}\mathrm{\Lambda }_s}{}}{\displaystyle \underset{Z_1x,z_{12}}{}}h_{Z_1}{\displaystyle \underset{Z_2z_{12},y}{}}h_{Z_2}`$
$``$ $`{\displaystyle \underset{xX}{}}{\displaystyle \underset{yY}{}}{\displaystyle \underset{z_{12}\mathrm{\Lambda }_s}{}}{\displaystyle \frac{\lambda _0}{[1+\mathrm{dist}(x,z_{12})]^\eta }}{\displaystyle \frac{\lambda _0}{[1+\mathrm{dist}(z_{12},y)]^\eta }}`$
$``$ $`{\displaystyle \frac{\lambda _0^2p_0|X||Y|}{[1+\mathrm{dist}(X,Y)]^\eta }},`$
where we have used the assumptions (2.2) and (2.3). Similarly, the third, triple sum can be estimated as
$`{\displaystyle \underset{Z_1:Z_1X\mathrm{}}{}}h_{Z_1}{\displaystyle \underset{Z_2:Z_2Z_1\mathrm{}}{}}h_{Z_2}{\displaystyle \underset{Z_3:Z_3Z_2\mathrm{},Z_3Y\mathrm{}}{}}h_{Z_3}`$ (A.19)
$``$ $`{\displaystyle \underset{xX}{}}{\displaystyle \underset{yY}{}}{\displaystyle \underset{z_{12}\mathrm{\Lambda }_s}{}}{\displaystyle \underset{z_{23}\mathrm{\Lambda }_s}{}}{\displaystyle \underset{Z_1x,z_{12}}{}}h_{Z_1}{\displaystyle \underset{Z_2z_{12},z_{23}}{}}h_{Z_2}{\displaystyle \underset{Z_3z_{23},y}{}}h_{Z_3}`$
$``$ $`{\displaystyle \underset{xX}{}}{\displaystyle \underset{yY}{}}{\displaystyle \underset{z_{12}\mathrm{\Lambda }_s}{}}{\displaystyle \underset{z_{23}\mathrm{\Lambda }_s}{}}{\displaystyle \frac{\lambda _0}{[1+\mathrm{dist}(x,z_{12})]^\eta }}{\displaystyle \frac{\lambda _0}{[1+\mathrm{dist}(z_{12},z_{23})]^\eta }}{\displaystyle \frac{\lambda _0}{[1+\mathrm{dist}(z_{23},y)]^\eta }}`$
$``$ $`{\displaystyle \underset{xX}{}}{\displaystyle \underset{yY}{}}{\displaystyle \underset{z_{12}\mathrm{\Lambda }_s}{}}{\displaystyle \frac{\lambda _0}{[1+\mathrm{dist}(x,z_{12})]^\eta }}{\displaystyle \frac{\lambda _0^2p_0}{[1+\mathrm{dist}(z_{12},y)]^\eta }}`$
$``$ $`{\displaystyle \frac{\lambda _0^3p_0^2|X||Y|}{[1+\mathrm{dist}(X,Y)]^\eta }}.`$
From these observations, we have
$`C_B(X,t)`$ $``$ $`{\displaystyle \frac{2B|X||Y|}{[1+\mathrm{dist}(X,Y)]^\eta }}\left\{2\right|t|\lambda _0+{\displaystyle \frac{(2|t|)^2}{2!}}\lambda _0^2p_0+{\displaystyle \frac{(2|t|)^3}{3!}}\lambda _0^3p_0^2+\mathrm{}\}`$ (A.20)
$`=`$ $`{\displaystyle \frac{2p_0^1B|X||Y|}{[1+\mathrm{dist}(X,Y)]^\eta }}\left\{\mathrm{exp}[2\lambda _0p_0|t|]1\right\}.`$
Consequently, we obtain
$$[A(t),B]\frac{2p_0^1AB|X||Y|}{[1+\mathrm{dist}(X,Y)]^\eta }\left\{\mathrm{exp}[2\lambda _0p_0|t|]1\right\}$$
(A.21)
from (A.13).
Proof of Theorem A.1 under Assumption 2.2 The first sum in the power series (LABEL:CBexpabound) is estimated as
$`{\displaystyle \underset{Z_1:Z_1X\mathrm{},Z_1Y\mathrm{}}{}}h_{Z_1}`$ $``$ $`{\displaystyle \underset{xX}{}}{\displaystyle \underset{yY}{}}{\displaystyle \underset{Z_1x,y}{}}h_{Z_1}`$ (A.22)
$``$ $`{\displaystyle \underset{xX}{}}{\displaystyle \underset{yY}{}}{\displaystyle \underset{Z_1x,y}{}}h_{Z_1}[1+\mathrm{dist}(x,y)]^\eta [1+\mathrm{diam}(Z_1)]^\eta `$
$``$ $`[1+\mathrm{dist}(X,Y)]^\eta |X||Y|s_0,`$
where
$$s_0=\underset{x}{sup}\underset{Zx}{}h_Z[1+\mathrm{diam}(Z)]^\eta .$$
(A.23)
Clearly, this constant $`s_0`$ is finite from the assumption (2.7). The second, double sum is estimated as
$`{\displaystyle \underset{Z_1:Z_1X\mathrm{}}{}}h_{Z_1}{\displaystyle \underset{Z_2:Z_2Z_1\mathrm{},Z_2Y\mathrm{}}{}}h_{Z_2}`$ (A.24)
$``$ $`{\displaystyle \underset{xX}{}}{\displaystyle \underset{yY}{}}{\displaystyle \underset{z_{12}\mathrm{\Lambda }_s}{}}{\displaystyle \underset{Z_1x,z_{12}}{}}h_{Z_1}{\displaystyle \underset{Z_2z_{12},y}{}}h_{Z_2}`$
$``$ $`{\displaystyle \underset{xX}{}}{\displaystyle \underset{yY}{}}{\displaystyle \underset{z_{12}\mathrm{\Lambda }_s}{}}[1+\mathrm{dist}(x,z_{12})]^\eta [1+\mathrm{dist}(z_{12},y)]^\eta `$
$`\times `$ $`{\displaystyle \underset{Z_1x,z_{12}}{}}h_{Z_1}[1+\mathrm{diam}(Z_1)]^\eta {\displaystyle \underset{Z_2z_{12},y}{}}h_{Z_2}[1+\mathrm{diam}(Z_2)]^\eta `$
$``$ $`[1+\mathrm{dist}(X,Y)]^\eta {\displaystyle \underset{xX}{}}{\displaystyle \underset{yY}{}}{\displaystyle \underset{z_{12}\mathrm{\Lambda }_s}{}}{\displaystyle \underset{Z_1x,z_{12}}{}}h_{Z_1}[1+\mathrm{diam}(Z_1)]^\eta `$
$`\times `$ $`{\displaystyle \underset{Z_2z_{12},y}{}}h_{Z_2}[1+\mathrm{diam}(Z_2)]^\eta `$
$``$ $`[1+\mathrm{dist}(X,Y)]^\eta {\displaystyle \underset{xX}{}}{\displaystyle \underset{yY}{}}{\displaystyle \underset{Z_1x}{}}h_{Z_1}[1+\mathrm{diam}(Z_1)]^\eta `$
$`\times `$ $`{\displaystyle \underset{Z_2y}{}}h_{Z_2}|Z_2|[1+\mathrm{diam}(Z_2)]^\eta `$
$``$ $`[1+\mathrm{dist}(X,Y)]^\eta |X||Y|s_0s_1,`$
where we have used the assumption (2.7) and the inequality,
$`[1+\mathrm{dist}(x,z)]^\eta [1+\mathrm{dist}(z,y)]^\eta `$ (A.25)
$`=`$ $`[1+\mathrm{dist}(x,z)+\mathrm{dist}(z,y)+\mathrm{dist}(x,z)\mathrm{dist}(z,y)]^\eta `$
$``$ $`[1+\mathrm{dist}(x,z)+\mathrm{dist}(z,y)]^\eta `$
$``$ $`[1+\mathrm{dist}(x,y)]^\eta ,`$
for any $`z\mathrm{\Lambda }_s`$. Similarly, the third, triple sum can be estimated as
$`{\displaystyle \underset{Z_1:Z_1X\mathrm{}}{}}h_{Z_1}{\displaystyle \underset{Z_2:Z_2Z_1\mathrm{}}{}}h_{Z_2}{\displaystyle \underset{Z_3:Z_3Z_2\mathrm{},Z_3Y\mathrm{}}{}}h_{Z_3}`$ (A.26)
$``$ $`{\displaystyle \underset{xX}{}}{\displaystyle \underset{yY}{}}{\displaystyle \underset{z_{12}\mathrm{\Lambda }_s}{}}{\displaystyle \underset{z_{23}\mathrm{\Lambda }_s}{}}{\displaystyle \underset{Z_1x,z_{12}}{}}h_{Z_1}{\displaystyle \underset{Z_2z_{12},z_{23}}{}}h_{Z_2}{\displaystyle \underset{Z_3z_{23},y}{}}h_{Z_3}`$
$``$ $`[1+\mathrm{dist}(X,Y)]^\eta {\displaystyle \underset{xX}{}}{\displaystyle \underset{yY}{}}{\displaystyle \underset{z_{12}\mathrm{\Lambda }_s}{}}{\displaystyle \underset{z_{23}\mathrm{\Lambda }_s}{}}{\displaystyle \underset{Z_1x,z_{12}}{}}h_{Z_1}[1+\mathrm{diam}(Z_1)]^\eta `$
$`\times `$ $`{\displaystyle \underset{Z_2z_{12},z_{23}}{}}h_{Z_2}[1+\mathrm{diam}(Z_2)]^\eta {\displaystyle \underset{Z_3z_{23},y}{}}h_{Z_3}[1+\mathrm{diam}(Z_3)]^\eta `$
$``$ $`[1+\mathrm{dist}(X,Y)]^\eta |X||Y|s_0s_1^2.`$
From these observations, we have
$`C_B(X,t)`$ $``$ $`{\displaystyle \frac{2s_0s_1^1B|X||Y|}{[1+\mathrm{dist}(X,Y)]^\eta }}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(2s_1|t|)^n}{n!}}`$ (A.27)
$`=`$ $`{\displaystyle \frac{2s_0s_1^1B|X||Y|}{[1+\mathrm{dist}(X,Y)]^\eta }}\left\{\mathrm{exp}[2s_1|t|]1\right\}.`$
As a result, we obtain
$$[A(t),B]\frac{2s_0s_1^1AB|X||Y|}{[1+\mathrm{dist}(X,Y)]^\eta }\left\{\mathrm{exp}[2s_1|t|]1\right\}$$
(A.28)
from (A.13).
Acknowledgements: We would like to thank Jens Eisert and Tobias Osborne for useful comments. TK thanks Bruno Nachtergaele and Hal Tasaki for helpful discussions. MBH was supported by US DOE W-7405-ENG-36. |
warning/0507/nucl-th0507004.html | ar5iv | text | # Polishing the Lens: I Pionic Final State Interactions and HBT Correlations– Distorted Wave Emission Function (DWEF) Formalism and Examples
## I Introduction
Measurements of the two-particle momentum correlations between pairs of identical particles have been used to study the space-time structure of the “fireball” produced in the collision between two heavy ions moving relativistically. The quantum statistical effects of symmetrization cause an enhancement of the two-boson coincidence rate at small momentum differences that can be related to the space-time extent of the particle source. This method, called HBT interferometry, has been applied extensively in recent experiments at the Relativistic Heavy Ion Collider (RHIC) by the STAR and PHENIX collaborations. See the reviewsPratt:wm ; Wiedemann:1999qn ; Kolb:2003dz ; Lisa:2005dd .
The invariant ratio of the cross section for the production of two pions of momenta $`𝐩_1,𝐩_2`$ to the product of single particle production cross sections is analyzed as the correlation function $`C(𝐩_1,𝐩_2)`$. We define $`𝐪`$=$`𝐩_1`$$`𝐩_2`$ and $`𝐊`$=$`(𝐩_1`$+$`𝐩_2)/2`$, with $`𝐊_T`$ as the component perpendicular to the beam direction. (We focus on mid-rapidity data, where $`𝐊=𝐊_T`$.) The correlation function can be parameterized for small $`𝐪`$ as $`C(𝐪,𝐊)1\lambda \mathrm{exp}(R_O^2q_O^2R_S^2q_S^2R_L^2q_L^2)\lambda (1R_O^2q_O^2R_S^2q_S^2R_L^2q_L^2)(q_iR_i1),`$ where $`O,S,L`$ represent directions parallel to $`𝐊_T`$, perpendicular to $`𝐊_T`$ and the beam direction, and parallel to the beam directionBP-HBT . Early Rischke:1996em and recent Kolb:2003dz hydrodynamic calculations predicted that a fireball evolving through a quark-gluon-hadronic phase transitions would emit pions over a long time period, causing a large ratio $`R_O/R_S`$. The puzzling experimental result that $`R_O/R_S1`$ Adler:2001zd is part of what has been called “the RHIC HBT puzzle” Heinz:2002un .
Another part of the puzzle is that the measured radii depend strongly on the average momentum $`K`$, typically decreasing in size by about 50% over the measured range. This dependence of a geometrical parameter on the probe momentum shows immediately that the radii are not simply a property of a static source. The influence of the interactions between the pionic probe and the medium, as well as the effects of transverse flow, must be taken into account when extracting the radii. The medium at RHIC seems to be a very high density, strongly interacting plasmaGyulassy:2004zy , so that any pions made in its interior would be expected to interact strongly.
We studied the effects of including the pionic interactions in previous workCramer:2004ih , finding that the only way to simultaneously describe the measured HBT radii and pionic spectra is to include the effects of pion-medium final state interactions by solving the relevant relativistic wave equation. These interactions are so strongly attractive that the pions can be taken as propagating through a system with a restored chiral symmetry.
The principal aim of the present work is to present a detailed treatment of the formalism that will allow wide application of our technique, and we present some new applications here. We also provide specific simple examples to demonstrate that the effects of pionic interactions cause the measured sizes of the medium to be different than the true sizes. Furthermore, we shall explicitly demonstrate that classical treatments of the pion-medium interactions, based on using the eikonal approximation for solutions of the wave equation are not valid.
An outline of the remainder of this paper follows. Previous standard formalisms that use plane wave pions are briefly reviewed in Sect. II. The technical method of incorporating the influence of final state interactions between the pion and the medium is described in Sect. III. Pionic emission in the absence of final state interaction is described with an emission function $`S_0`$ that is of a form motivated by hydrodynamics. The connection between the symmetries of this function (for the case of head-on collisions) and the form of the pion optical potential is described in Sect. IV. The role of the complex optical potential and the use of chiral symmetry to constrain its form at low energies is explained in Sect. V. The the specific numerical algorithm necessary to incorporate the optical potential is presented in Sect. VI. Once our approach is defined, there is only one approximation we need to make. The validity of this requires only that the source size be much larger than the inverse of the temperature. This large source approximation, LSA, is explained in Sect. VII. The resulting distorted wave emission function, DWEF, is evaluated using two different equivalent methods in Sec. VIII. In Sect. IX we apply these methods to STAR central Au+Au data at $`\sqrt{s}=200`$ GeV. Instead of treating the temperature of the system as a free parameter, as was done in Cramer:2004ih , we fix its value at the transition temperature $`T_c=`$ 193 MeV, obtained in the most recent lattice QCD calculationKatz:2005br .
For conventional calculations of the spectra, chemical equilibrium analyzes yield lower temperatures $`T_{ch}=174`$ MeV Braun-Munzinger:2001ip . A large difference between $`T_c`$ and $`T_{ch}`$ implies that the hadrons interact after the deconfinement transition occurs. This notion is entirely consistent with our treatment of pionic distortions which has as its fundamental assumption that pions interact in a hot dense medium before escaping to freedom. We also fix the pion chemical potential at the pion mass ($`\mu _\pi =`$ 139.57 MeV). The resulting procedure reduces the number of free parameters by two to a total of nine. Then we are able to reproduce a total of 32 data points with high accuracy. Section IX also extends earlier numerical resultsCramer:2004ih by computing the dependence of the HBT radii on the centrality in Au+Au collisions and by making predictions for Cu+Cu collisions vs. centrality. The eikonal approximation to solving the wave equation is discussed in Sect. X in which the important influence of the opacity and the vanishing of the effects of the real potential are also described. The importance of the real part of the optical potential in obtaining oscillating radii at low energies is illustrated through two examples in Sect. XI. A brief summary is presented in Sect. XII. Finally, a short appendix verifies an approximation to an integral.
## II Previous Formalism – Plane wave pions
The aim is to include the effects of final state interactions of outgoing pions. We’ll begin with a brief review of the formalism previously used to describe HBT correlations for situations in which the pions do not interact with the medium.
The relevant observables are the covariant single- and two-particle emission functions defined as appropriately normalized ratios of cross sectionsGKW79
$`𝒫_1(𝐩)=E{\displaystyle \frac{1}{\sigma _\pi }}{\displaystyle \frac{d\sigma _\pi }{d^3𝐩}}`$ (1)
$`𝒫_2(𝐩_1,𝐩_2)=E_1E_2{\displaystyle \frac{1}{\sigma _\pi }}{\displaystyle \frac{d\sigma _{\pi \pi }}{d^3𝐩_1d^3𝐩_2}}`$ (2)
$`{\displaystyle \frac{d^3p}{E}𝒫_1(𝐩)}=N`$ (3)
$`{\displaystyle \frac{d^3p_1}{E_1}\frac{d^3p_2}{E_2}𝒫_2(𝐩_1,𝐩_2)}=N(N1)`$ (4)
$`C(𝐩_1,𝐩_2)={\displaystyle \frac{N^2}{N(N1)}}{\displaystyle \frac{𝒫_2(𝐩_1,𝐩_2)}{𝒫_1(𝐩_1)𝒫_1(𝐩_2)}}{\displaystyle \frac{𝒫_2(𝐩_1,𝐩_2)}{𝒫_1(𝐩_1)𝒫_1(𝐩_2)}}.`$ (5)
The last step of Eq. (5) is obtained because $`N^2N`$ for the very high energy collisions of interest here.
These observables can be expressed in terms of an emission function $`S(x,K)`$ that is the Wigner transform of the density matrix associated with the currents that emit the pions. The literature Wiedemann:1999qn ; GKW79 presents the single-particle emission function as
$`S_0(x,p)={\displaystyle \frac{d^4y}{2(2\pi )^3}\mathrm{exp}(ipy)J^{}(x+y/2)J(xy/2)}`$
$`={\displaystyle \frac{d^4y}{2(2\pi )^3}e^{ip(x+y/2)}J^{}(x+y/2)e^{ip(xy/2)}J(xy/2)},`$ (6)
where $`p`$ is the four-momentum of an on-shell emitted pion and $`J(x)`$ is the current that acts as a source of pion fields. We shall assume that $`J(x)`$ is not altered by the presence of the produced pions. The brackets indicate that one takes an ensemble average over chaotic sources. Thus we assume that the emission process is initially uncorrelated: the pions are emitted from chaotic sources that have random phasesed73 ; GKW79 ; Kapusta:2005pt . This assumption seems consistent with analysis of data produced at SPS and RHICLisa:2005dd The second term of Eq. (6) is written to show that the emission function is determined by the product of a plane wave factor and the current that emits the pions. This is the standard result of scattering theory, if final state interactions are ignored. Here and throughout the paper we use natural units in which $`\mathrm{}`$ and $`c`$ are unity. Our use of the subscript 0 denotes that the effects of final state interactions are ignored.
Next we discuss the relationship between the function $`S_0(x,p)`$ and the single-pion production cross section. If the source current couples linearly with the pion field (with a chosen interaction Lagrangian, $`L_I`$), and no final state interactions occur, the matrix element of the interaction $`L_I`$ between the vacuum and single-pion final state is
$`\stackrel{~}{J}(p){\displaystyle d^4xe^{ipx}J(x)}={\displaystyle d^4x\varphi _𝐩(x)J(x)},`$ (7)
where $`\varphi _𝐩=e^{ipx}`$ is the solution of
$`\left({\displaystyle \frac{^2}{t^2}}^2+m_\pi ^2\right)\varphi _𝐩=0.`$ (8)
In this plane wave (pw) approximation, the covariant single-particle emission function can be expressed as
$`𝒫_1^{\mathrm{pw}}(p)=\left|\stackrel{~}{J}(p)\right|^2={\displaystyle d^4xS_0(x,p)},`$ (9)
The literature Wiedemann:1999qn presents the two-particle emission function as
$`S_0(x,K)={\displaystyle \frac{d^4y}{2(2\pi )^3}\mathrm{exp}(iKy)J^{}(x+y/2)J(xy/2)}.`$ (10)
The quantity that enters in the HBT correlation is $`S_0(x,K,q)`$ with
$`S_0(x,K,q)S_0(x,K)e^{iqx}={\displaystyle \frac{d^4y}{2(2\pi )^3}e^{ip_1(x+y/2)}J^{}(x+y/2)e^{ip_2(xy/2)}J(xy/2)}S_0(x,p_1,p_2),`$ (11)
where $`p_1,p_2`$ are the four-momenta of the two on-shell emitted pions, and
$`q=p_1p_2,K={\displaystyle \frac{1}{2}}(p_1+p_2).`$ (12)
The brackets again indicate that one takes an ensemble average over chaotic sources. The second form, Eq. (11) shows again that products of plane-wave functions and currents determine the observables. The notation $`S_0(x,p_1,p_2)`$ is introduced to emphasize the two-particle nature of the emission function.
In plane wave approximation the correlation function is given by the expression
$`𝒫_1^{\mathrm{pw}}(p_1)𝒫_1^{\mathrm{pw}}(p_2)C_0(p_1,p_2)=\left|{\displaystyle \frac{d^4x_1}{2(2\pi )^3}\frac{d^4x_2}{2(2\pi )^3}\frac{1}{\sqrt{2}}\left(e^{ip_1x_1}e^{ip_2x_2}+e^{ip_2x_1}e^{ip_1x_2}\right)J(x_1)J(x_2)}\right|^2,`$ (13)
in which the effects of the Bose statistics is explicit. One proceeds by taking the absolute square to find that the expression contains four terms. The direct terms involve $`\left(e^{ip_1x_1}e^{ip_1x_1^{}}J(x_1)J^{}(x_1^{})e^{ip_2x_2}e^{ip_2x_2^{}}J(x_2)J^{}(x_2)\mathrm{plus}\mathrm{\hspace{0.33em}1}2\right)`$ and the exchange terms involve $`(e^{ip_1x_1}e^{ip_2x_1^{}}J(x_1)J^{}(x_1^{}))e^{ip_2x_2}e^{ip_1x_2^{}}J(x_2)J^{}(x_2))\mathrm{plus}\mathrm{\hspace{0.33em}1}2).`$ The placement of the brackets arises from the use of chaotic sources: the effects of pions produced by two different chaotic sources average to 0ed73 ; GKW79 . Using the absolute square and the previous definitions in Eq. (13) yields the result
$`C_0(p_1,p_2)=1+{\displaystyle \frac{\left|d^4xS_0(x,p_2,p_1)\right|^2}{\left|\stackrel{~}{J}(p_1)\right|^2\left|\stackrel{~}{J}(p_2)\right|^2}}.`$ (14)
## III Introducing Final state interactions
Our emission function, $`S_0`$, is not meant to be the same as that obtained in conventional hydrodynamic calculations of hadronic spectral functions. In such calculationsWiedemann:1999qn the quantity $`S_0(x,K)`$ is assumed to be a local equilibrium Bose-Einstein distribution localized on a 3-dimensional freeze-out hypersurface that separates the thermalized interior of the hot dense medium from the free-streaming particles on its exteriorCooper:1974mv . Here we assume that the emitted pions (and other hadrons) undergo significant interactions while escaping the hot dense medium. The emitted pions interact both with the hot dense matter and other hadrons as they escape the system. This is emission from a coexistence phase. The purpose of the present section is to provide a formalism that allows such effects to take place, while also including the usual effects of emission from a freeze out surface.
### III.1 Distorted Waves
We wish to include the effect that an escaping pion has to ”fight” its way through the medium. Our formalism relies heavily on earlier derivations in Ref. GKW79 . The effects of the fighting distort the wave away from its plane wave form dictated by Eq. (8). These interactions of the pions with the medium, represented by the optical potential $`U`$, lead to a modified equation (approximated as a one-body equation)
$`\left({\displaystyle \frac{^2}{t^2}}^2+U+m_\pi ^2\right)\psi =0.`$ (15)
This equation provides an approximate treatment of the complicated final state interactions. In principle, the optical potential should be related to the underlying pion source currents. We adopt a phenomenological approach. The resulting S-matrix depends on the exact solution, $`\psi _𝐩^{()}(x)`$, to this equation with the out boundary condition of approaching at $`t\mathrm{}`$, the free wave $`\varphi _𝐩(x)`$BJD .
Ref. GKW79 uses standard S-matrix reduction techniques to show that matrix element of the pion field between the vacuum and a one-pion final state $`p`$ is simply the complex conjugate of $`\psi _𝐩^{()}(x)`$, if the optical potential is introduced as in Eq. (15). This means that in calculating the matrix elements of $`L_I`$ between the vacuum and multi-pion final states one simply replaces the plane wave solutions of Eq. (8) by the out-going wave solutions of Eq. (15). Thus the simple replacement:
$`e^{ipx}\psi _𝐩^{()}(x)`$ (16)
is sufficient to include final state interaction effects within the optical model approximation. Therefore the single-particle emission function is given by
$`S(x,p_1)={\displaystyle \frac{d^4y}{2(2\pi )^3}J^{}(x+y/2)J(xy/2)\psi _{𝐩_1}^{()}(x+y/2)\psi _{𝐩_1}^{()}(xy/2)}.`$ (17)
This expression reduces to the plane wave form $`S_0(x,p_1)`$ Eq. (10) if the distorted wave $`\psi _{𝐩_1}`$ are replaced by plane waves. The two-particle emission function that includes the effects of final state interactions is obtained by replacing Eq. (10) by the result
$`S(x,p_2,p_1)=S(x,K,q)={\displaystyle \frac{d^4y}{2(2\pi )^3}J^{}(x+y/2)J(xy/2)\psi _{𝐩_1}^{()}(x+y/2)\psi _{𝐩_2}^{()}(xy/2)},`$ (18)
which applies for calculating the two-particle emission function. In the plane wave limit, $`S(x,K,q)S_0(x,K)e^{iqx}.`$
Physical observables for the emission of two pions of momenta $`p_1,p_2`$ are determined by the correlation function $`C(q,K)(K=\frac{1}{2}(p_1+p_2),q=p_1p_2)`$, which is given by
$`C(q,K)=1+{\displaystyle \frac{|d^4xS(x,K,q)|^2}{d^4xS(x,p_1)d^4xS(x,p_2)}}.`$ (19)
One obtains the usual expression (14) if the wave functions $`\psi _𝐩^{()}`$ are replaced by plane wave functions ($`\psi _𝐩^{()}(x)e^{ipx}`$).
The expression (18) contains the ensemble average of the currents. This may be expressed in terms of $`S_0(x,K)`$ by taking the Fourier transform of Eq. (6). We then obtain the convolution formula:
$`S(x,K,q)={\displaystyle d^4K^{}S_0(x,K^{})\frac{d^4y}{(2\pi )^4}e^{iK^{}y}\psi _{𝐩_1}^{()}(x+y/2)\psi _{𝐩_2}^{()}(xy/2)},`$ (20)
where the subscripts indicate the momenta $`p_1,p_2`$ of the detected pions and $`K=\frac{1}{2}(p_1+p_2)`$. The quantity $`S(x,K,q)`$ is used to compute experimental observables in the same way that $`S_0`$ was previously used. This is the distorted wave emission function DWEF formalism. The expression (20) is the coordinate space version of Eq. (5.25) of Ref. GKW79 .
We emphasize that $`S_0`$ (which reflects the true properties of the source) is no longer directly related to observables–the appearance of distorted waves obscures the relationship between the data and the true properties of the source.
### III.2 Extracting HBT radii
The correlation function of Eq. (19) are related to HBT radii in two different ways. In the first method, one treats the momentum differences $`q_{O,S,L}`$ as small quantities and then expands keeping terms to second order so that:
$`C(q,K)11q_O^2R_O^2q_S^2R_S^2q_L^2R_L^2,`$ (21)
where $`q_O`$ is the transverse component that is parallel to the direction of $`𝐊`$, $`q_S`$ is the transverse component that is perpendicular to the direction of $`𝐊`$, and $`q_L`$ is the longitudinal component. Here and below, because we are focusing on central collisions we ignore the $`q_Oq_LR_{OL}^2`$ cross term. Another parameterization is
$`C(q,K)1\mathrm{exp}(q_O^2R_O^2q_S^2R_S^2q_L^2R_L^2),`$ (22)
In practice, describing data and extracting radii require using:
$`C(q,K)1\lambda \mathrm{exp}(q_O^2R_O^2q_S^2R_S^2q_L^2R_L^2),`$ (23)
$`C(q,K)1\lambda (1q_O^2R_O^2q_S^2R_S^2q_L^2R_L^2).`$ (24)
Here the reduction factor $`\lambda `$ (typically about 1/2) is the fraction of pairs that originate in the space-time region relevant for correlations, see the review Lisa:2005dd . Current understanding Lisa:2005dd is that the HBT data are consistent with incoherent emission, and accounting for the many pions that are produced by the decays of resonances far outside the collision region can reproduce the factor $`\lambda .`$ These “halo” pions do not have a BE-enhanced correlation in the $`q`$ region measured with the pions emitted from the core of hot dense matter, but they cannot be experimentally separated from the latter.
Our calculations employ the core-halo model core-halo in which pions are assumed to arise either from the hot, dense core or from the halo. This model is a simplified version of more detailed treatments of resonance decays that are discussed in the review Wiedemann:1999qn .
The influence of resonances affects the extraction of radii and the measurements of the pion spectrum in different ways. We account for this in our phenomenological analysis, see the erratum of Ref. Cramer:2004ih and Sect. IX-A. Here we note only that effects of pions produced by short-lived resonances (such as $`\mathrm{\Delta },\rho ,\mathrm{}`$) that are not explicitly included in the pionic wave equation are included in the function $`S_0`$. The pions resulting from long lived resonances are not included in $`S_0`$. Pions resulting from the decay of slowly moving $`\mathrm{\Omega }`$ mesons that occur inside the dense matter system are included in $`S_0`$, but the pions from rapidly moving $`\omega `$ mesons that decay outside of the system are not included in $`S_0`$.
The approximate forms (23,24) provide two ways to extract radii. The former suggests that
$`R_i^2={\displaystyle \frac{1}{q_i^2}}\mathrm{ln}{\displaystyle \frac{\lambda }{C(q_i,K)1}},i=O,S,L`$ (25)
while the latter implies
$`R_i^2={\displaystyle \frac{1}{q_i^2}}{\displaystyle \frac{C(q_i,K)1}{\lambda }},i=O,S,L`$ (26)
Eq. (25) can be used for values of $`q_I`$ such that $`(C1)/\lambda `$ can be approximated by a Gaussian function. The use of Eq. (24) requires that $`q_iR_i1`$. We show below in Sect. VIII that, while the correlation functions are not Gaussians (so that the squared radii are not moments of the correlation function), the Gaussian parameterization is quite accurate at the modest (but not very small) values of $`q_i30`$ MeV/c that dominate the experimental extraction of radii.
## IV Symmetries of $`S_0(x,K)`$ and the form of the pion distorted waves
We use the hydrodynamic parameterization of the source of Ref. CL96a ; H96 ; Tomasik:1997eq . Generally based on the Bjorken tube model, it is given by
$`S_0(x,K)d^4x`$ $`=`$ $`{\displaystyle \frac{M_{}\mathrm{cosh}(\eta Y)}{(2\pi )^3}}{\displaystyle \frac{1}{\mathrm{exp}\left[\frac{(Ku(x)\mu _\pi )}{T(x)}\right]1}}\rho (b)\mathrm{exp}\left[{\displaystyle \frac{(\eta \eta _0)^2}{2(\mathrm{\Delta }\eta )^2}}\right]`$ (27)
$`\times \tau d\tau \left[{\displaystyle \frac{1}{\sqrt{2\pi (\mathrm{\Delta }\tau )^2}}}\mathrm{exp}\left({\displaystyle \frac{(\tau \tau _0)^2}{2(\mathrm{\Delta }\tau )^2}}\right)\right]d\eta bdbd\varphi ,`$
$`d^4x`$ $`=`$ $`\tau d\tau d\eta bdbd\varphi .`$ (28)
Here, $`\mu _\pi `$ is the pion chemical potential, the variables $`(b,\varphi )=𝐛`$ are equivalent to $`𝐱_{}`$,$`\eta =\frac{1}{2}\mathrm{ln}\frac{t+z}{tz},\tau =\sqrt{t^2z^2},M_T=\sqrt{K_{}^2+m_\pi ^2},Y=\frac{1}{2}\mathrm{ln}\frac{E_K+K_z}{E_KK_z}.`$ The factor in the brackets involving $`\tau `$ has the same normalization as the delta function: $`\delta (\tau \tau _0)`$. We use a Bose-Einstein distribution instead of the Boltzmann distribution of Ref. Tomasik:1997eq and also allow the transverse density $`\rho (b)`$ to have a general form instead of a Gaussian. Note that the model contains the parameters: $`R,\eta _0,\mathrm{\Delta }\eta ,\mathrm{\Delta }\tau ,\tau _0,\eta _f`$. As stated originally, temperature gradients were included, but we treat the temperature as a constant for our numerical calculations. We shall see below that our formalism can easily be generalized to allow the temperature to vary as a function of $`b`$.
The longitudinal $`\eta ,\tau `$ and transverse $`(𝐛,𝐊)`$ variables can be separated in the following form:
$`S_0(x,K)`$ $`=`$ $`𝒮_0(\eta ,\tau ,Y)B_\eta (𝐛,𝐊)`$ (29)
$`B_\eta (𝐛,𝐊)`$ $``$ $`{\displaystyle \frac{M_{}}{\mathrm{exp}[(Ku\mu _\pi )/T]1}}\rho (b)`$ (30)
$`𝒮_0(\eta ,\tau ,Y)`$ $``$ $`{\displaystyle \frac{\mathrm{cosh}(\eta Y)}{(2\pi )^3}}\mathrm{exp}\left[{\displaystyle \frac{\eta ^2}{2(\mathrm{\Delta }\eta )^2}}\right]{\displaystyle \frac{1}{\sqrt{2\pi (\mathrm{\Delta }\tau )^2}}}\mathrm{exp}\left[{\displaystyle \frac{(\tau \tau _0)^2}{2(\mathrm{\Delta }\tau )^2}}\right].`$ (31)
Here $`\rho (b)`$ is a function, normalized as $`\rho (0)=1`$ that represents the transverse density. To be specific we use
$`\rho (b)=[1/(\mathrm{exp}((bR_{WS})/a_{WS})+1)]^2.`$ (32)
This distribution has a correct exponential fall-off at large values of $`b`$, and different choices of the parameters $`R_{WS},a_{WS}`$ allow a variety of shapes to be assumed.
We concentrate on the kinematics of the STAR experiment which detects pions within half a unit of rapidity of moving perpendicular to the beam, so we take $`Y=0`$. This means that the average momentum is transverse: $`𝐊=𝐊_{}=𝐊_T`$. With colliding beams of equal mass and energy there is a fore-aft symmetry along the longitudinal axis, so that we use $`\eta _0=0.`$ Therefore it is useful to define
$`𝒮_0(\eta ,\tau )𝒮_0(\eta ,\tau ,Y=0).`$ (33)
The velocity field $`u(x)`$ describing the dynamics of the expanding source is parameterized by WSH96
$$u^\mu (x)=(\mathrm{cosh}\eta \mathrm{cosh}\eta _t(b),\mathrm{cos}\varphi \mathrm{sinh}\eta _t(b),\mathrm{sin}\varphi \mathrm{sinh}\eta _t(b),\mathrm{sinh}\eta \mathrm{cosh}\eta _t(b)),$$
(34)
with $`\varphi `$ the angle between $`𝐊_{}`$ and $`𝐮`$ (or when appearing in the single-particle emission function, it is the angle between $`𝐩_i`$ and $`𝐮`$). Eq. (34) implements a boost-invariant longitudinal flow profile $`v_L=z/t`$, with a linear radial profile of strength $`\eta _f`$ for the transverse flow rapidity:
$$\eta _t(b)=\eta _f\frac{b}{R_{WS}}.$$
(35)
The exponent that enters in the Bose-Einstein distribution is given by the four-vector dot product:
$$Ku(x)=M_{}\mathrm{cosh}(\eta Y)\mathrm{cosh}\eta _t(b)K_{}\mathrm{sinh}\eta _t(b)\mathrm{cos}\varphi .$$
(36)
The presence of a non-vanishing value of $`\eta _f`$ causes the emission function to depend on both $`K`$ and $`M_{}`$. The Bose-Einstein distribution of Eq. (30) is evaluated as a sum of Boltzmann distributions:
$`{\displaystyle \frac{1}{\mathrm{exp}\left[\frac{Ku(x)\mu _\pi }{T}\right]1}}={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\mathrm{exp}({\displaystyle \frac{Ku+\mu _\pi }{T_n}}),T_n{\displaystyle \frac{T}{n}}.`$ (37)
The parameterization (27) is motivated by hydrodynamical models with approximately boost-invariant longitudinal dynamics. It uses thermodynamic and hydrodynamic parameters and appropriate coordinates. The “emission function” given above was originally intended to parameterize the distribution of points of last interaction in the source. In conventional treatments, the Cooper-FryeCooper:1974mv matching procedure is used to obtain distributions of detected particles. Our approach is more general. We assume that pions can be formed at any space-time point (including but not limited to the freeze-out surface) during the collision, and propagate through the dense medium while interacting before being detected. Thus we take $`S_0(x,K)`$ to be the emission function in the absence of final state interactions.
We note that the emission function of Eq. (27) has been previously used in the “blast-wave model” Retiere:2003kf , and we would like to comment on the differences between our formalism and that model. These include:
* The blast-wave model does not attempt to reproduce the normalization of the spectrum, and sets the chemical potential $`\mu _\pi `$ to 0. We find that taking $`\mu _\pi `$ to be the pion mass allows us to reproduce both the normalization and the shape of the pion momentum spectrum.
* The blast-wave model uses the smoothness approximation and computes HBT radii as moments of the emission function.
* The blast-wave model uses plane waves and therefore omits the effects of the optical potential.
* The blast-wave model uses a parameterized emission function to describe the six-dimensional distribution of pions on some freeze-out hypersurface, while the emission function of the DWEF model describes the initial emission of pions within the hyper-volume of the hot, dense medium produced by the collision. Therefore, it is not appropriate to compare emission-function parameters like $`T`$ and $`R`$ that are derived from blast-wave fits with those of the present model.
## V Final state interactions and the optical potential
The salient feature of the 200 GeV data is the high density of the produced matter, so we treat the effects of pion final state interactions with a dense medium. We adopt a single-channel approach that uses the interaction–distorted incoming wave $`\mathrm{\Psi }_{𝐩_1}^{()}(x_1)`$.
We assume that the matter produced in the central region of the collision is cylindrically symmetric with a very long axis, so that an expression of the form (29) is valid. In that case, the optical potential $`U`$ representing the interaction between a pion and the medium is a complex, azimuthally-symmetric function depending on pion momentum and local density. Within our formalism the influence of some time-dependent effects in $`U`$ introduced by the time-dependent source $`S_0`$ is incorporated in the energy dependence of the optical potential, and the pion-medium interaction time is restricted by $`S_0`$.
The optical potential accounts for the interaction between each pion and the surrounding medium, but does not include the interaction between the two pions. In particular, the Coulomb interaction is known to be important. The experimental analysis removes this effect before extracting radii from the data. The Coulomb interaction between pions is of long range and its important effects occur when the pions are outside the medium. Indeed, specific analyzes show that Coulomb effect occur only at very low relative momenta STARHBT . In contrast, the optical potential acts only when pions are inside. We therefore expect that presence of the optical potential would not influence the removal of the Coulomb interaction. A quantitative treatment of the effect of the optical potential on the removal of Coulomb effects would involve solving a three-body problem. This is not warranted at the present stage of development, but might eventually become worthwhile.
The optical potential accounts for situations in which the pion changes energy or disappears entirely due to its interactions with the dense medium. We do not assume to know the content of the dense medium, and therefore will use a phenomenological optical potential. But it is worthwhile to consider a simple example to get an idea about how large the optical potential can be. Suppose, e.g., that the medium is a gas of pions. Then $`\pi \pi `$ scattering would be the origin of $`U`$. In the impulse approximation, the central optical potential would be $`U_0=4\pi f\rho _0`$, where $`f`$ is the complex forward scattering amplitude and $`\rho _0`$ the central density. For low energy pion-pion interactions, $`4\pi `$Im$`[f(p)]=p\sigma ,`$ with $`\sigma 1`$mb. At a momentum $`p=1\mathrm{fm}^1=197.3`$ MeV/c, using a pion density about ten times the baryon density of ordinary nuclear matter, Im$`[U(0)]0.15\mathrm{fm}^2`$, representing significant opacity.
The optical potential must be an analytic function of energy, and therefore the existence of an imaginary part mandates the existence of a real part. Thus, any analysis needs to treat $`U`$ as a complex function. Under certain circumstances the real part can be very large. For example, if two interacting pions each have less energy than half of the rho meson mass, the final state interactions caused by virtual transitions to a rho meson would be strongly attractive. Additionally, the influence of chiral symmetry restoration can lead to a strong real part. This is discussed next.
### V.1 Chiral Symmetry Restoration
Suppose the dense medium is one in which chiral symmetry is restored. This means that the value of the quark condensate vanishes, an effect that could be caused by an increase in temperature or density. The pion mass is proportional to the quark condensate via the GMOR relation Gell-Mann:1968rz
$`m_\pi ^2f_\pi ^2={\displaystyle \frac{m_u+m_d}{2}}0|\overline{u}u+\overline{d}d|0,`$ (38)
where $`f_\pi `$ is the weak pion decay constant $`93`$ MeV. It is believed that $`m_\pi ,f_\pi `$ and the condensate $`0|\overline{u}u+\overline{d}d|0`$ all depend on temperature and density. If one takes the Brown-RhoBrown:1991kk scaling relation for $`f_\pi `$ and the perturbatively calculated temperature dependence of the condensate, the pion mass is proportional to the cube root of the condensate, and therefore vanishes for sufficiently large temperatures. See the reviews Koch:1996vt . Suppose the optical potential arises only from the temperature dependence of the pion mass. Then the Klein-Gordon equation would take the form:
$`(^2+m_\pi ^2(T))\psi =(p^2+m_\pi ^2)\psi ,`$ (39)
for regions inside the medium. The effects of the medium are incorporated through the difference between $`m\pi ^2(T)`$ and $`m_\pi ^2`$. If one re-writes Eq. (39) as a Klein-Gordon equation $`(^2+U)\psi =p^2\psi ,`$ then the optical potential takes the form:
$`U(b)=(m_\pi ^2(T)m_\pi ^2)\rho (b),`$ (40)
in which the finite extent of the medium is accounted for by the factor $`\rho (b)`$. If $`m_\pi (T)`$ approaches zero, the optical potential is attractive with magnitude $`m_\pi ^2.`$
A more recent study by Son & StephenovSon:2002ci provides a more detailed treatment of the effects of chiral restoration in which the general $`p`$-wave nature of the low-energy interaction between pions and any target is included. We are guided by this work. Son & Stephenov use the dispersion relation for low momentum pions in infinite nuclear matterboy ; Son:2002ci :
$`\omega ^2=u_\pi ^2(\widehat{p}^2+m_\pi (T)^2),`$ (41)
where $`\widehat{p}^2`$ is the infinite-sized matter version of $`^2`$. The quantity $`u_\pi `$ is termed the pion velocity, even though it is only that when $`m_\pi (T)`$ vanishes. The term $`m_\pi (T)`$ is denoted the pion screening mass. This quantity appears in the expression for the static Euclidean pion correlatorSon:2002ci . The energy of a pion at $`𝐩=0`$ is termed the pion pole mass. The free pion mass is $`m_\pi `$. In Ref. Son:2002ci Eq. (41) applies only for $`T<T_c`$. For larger temperatures, chiral symmetry is restored and pions are massless.
Defining $`t(T_cT)/T_c`$, SS find $`m_\pi ^2(T)t^{\beta \nu },u^2t^\beta `$, with $`\beta <\nu `$, e.g $`\nu =.73,\beta =.38`$Baker:1977hp . These equations are valid for temperature close to (but not too close to) the critical point. Another view about the dispersion relation can be found in Ref. sas . This general discussion about the influence of chiral restoration provides some guidance, but does not tell us exactly to use.
We wish to obtain an equivalent optical potential and see if it is attractive or repulsive. Use the Klein Gordon equation in the form
$`\widehat{p}^2+U+m_\pi ^2=p^2+m_\pi ^2.`$ (42)
In regions outside the dense medium where $`U=0`$, the operator $`\widehat{p}^2`$ is simply the square of the momentum, $`p^2`$. Subtract (42) from (41) to obtain
$`U=u_\pi ^2m_\pi ^2(T)m_\pi ^2+(u_\pi ^21)\widehat{p}^2,`$ (43)
an expression that is the sum of two negative definite terms. This form can be simplified by using the wave equation (42) to remove term $`\widehat{p}^2`$. Then one finds a momentum-dependent optical potential:
$`U={\displaystyle \frac{u_\pi ^2m_\pi ^2(T)m_\pi ^2+(u_\pi ^21)p^2}{u_\pi ^2}}.`$ (44)
Note that if $`u_\pi `$ becomes really small the optical potential becomes very strongly attractive.
For matter of finite size, the term $`\widehat{p}^2`$ can also be interpreted as $`c`$ which is the Kisslinger term lsk . For infinite nuclear matter only forward scattering occurs and the two terms are identical, but differences may arise for scattering from media of finite size. One can not tell the difference between the two terms at the start, so that we find a general form
$`U(b)=(w_0+w_2(1ϵ)p^2)\rho (b)ϵw_2\rho (b),`$ (45)
with both the real and imaginary parts of $`w_0,w_2,ϵ`$ positive for attractive interactions, and $`\rho (b)`$ taken from Eq. (32). This simple form is strictly valid only for low-energy pions. The limit of an infinite tube is used to write $`U`$ as independent of $`z`$. We find for the 200 GeV data, that for temperatures below our assumed value of 193 MeV the fitting prefers a very small value of the Kisslinger term, so we simply set $`ϵ=0`$. We also take $`w_0`$ real (so that there is no opacity at p=0). We shall see below, that if the temperature is set to a value greater than 193 MeV, the fit is improved by including the gradient terms at about the 20% level ($`ϵ0.2`$). The simple form Eq. (45) is sufficient to account for the data we study, $`K_T600`$ MeV/c, but we remind the reader that the interaction strength does not grow as $`p^2`$ for $`p`$ much greater than about 400 MeV/c.
We shall discuss the precise parameters of the optical potential in subsequent sections. For now we may simply proceed using the assumption that $`U`$ is a complex, azimuthally-symmetric function depending on pion momentum and local density of the form of Eq. (45).
## VI Finding $`\psi ^{(\pm )}(𝐱)`$
The evaluation of the emission function (20) requires performing an eight dimensional integral using a distorted wave. We shall use symmetries to reduce the number of numerical evaluations. Doing this depends on obtaining a compact expression for the distorted wave. In the present section, we show how the distorted waves are evaluated.
The first step is to realize that the function $`\psi _𝐩^{()}(x)`$ represents an energy-eigenfunctionGKW79 provided the optical potential does not change with time. We shall show below that the value of $`\delta \tau `$ is not large. Thus the the time-independent optical potential that we use can be thought of as a time-averaged optical potential. So we have
$`\psi _𝐩^{()}(x)=e^{i\omega _px^0}\mathrm{\Psi }_𝐩^{()}(𝐱).`$ (46)
To proceed we need to examine the properties of the wave function $`\psi _𝐩^{()}(𝐱)`$. It is conventional to compute $`\psi _𝐩^{(+)}(𝐱)`$, and we will follow this convention. Then we use time reversal invariance in the form
$`\mathrm{\Psi }_𝐩^{()}(𝐱)=\mathrm{\Psi }_{𝐩}^{(+)}{}_{}{}^{}(𝐱)`$ (47)
to obtain the desired wave function.
The next step is to realize that for central collisions, $`Y=0`$, the emission function (31) has a cylindrical symmetry. This means that the expected optical potential is azimuthally symmetric. If we take the matter to have the form of a very long tube, the optical potential will be independent of $`z`$. Then one obtains a solution that takes a product form
$`\mathrm{\Psi }_{𝐩_{1,2}}^{()}(𝐱)`$ $`=`$ $`e^{iq_Lz/2}\psi _{𝐩_{1,2}}^{()}(𝐱_{}=𝐛),`$ (48)
$`𝐩_{1,2}`$ $`=`$ $`𝐊\pm 𝐪/2\pm \widehat{𝐳}q_L/2,`$ (49)
where the vector $`𝐪`$ is defined as a transverse vector $`𝐪\widehat{𝐳}=0`$.
We may obtain the wave function $`\psi _𝐩^{(+)}(𝐛)`$ by solving the wave equation
$`\left(_{}^2+U(b)\right)\psi _𝐩^{()}(𝐛)=p^2\psi _𝐩^{()}(𝐛).`$ (50)
If $`U=0,\psi _𝐩^{(+)}(𝐛)=e^{i𝐩𝐛}`$. Many previous treatments of opacity can be understood as using the eikonal approximation to obtain solutions to Eq. (50).
We take the optical potential to have the azimuthally-symmetric form of Eq. (45) so that the solution for $`\psi _𝐩^{(\pm )}(𝐛)`$ can be expanded in partial wave form in plane polar coordinates ($`b\sqrt{𝐱_{}^2},\varphi ,\mathrm{cos}\varphi \widehat{𝐩}\widehat{𝐛}`$):
$`\psi _𝐩^{(+)}(𝐛)={\displaystyle \underset{m=\mathrm{},\mathrm{}}{}}f_m(p,b)i^me^{im\varphi },`$ (51)
$`\psi _𝐩^{(+)}(𝐛)=f_0(p,b)+2{\displaystyle \underset{m=1,\mathrm{}}{}}f_m(p,b)i^m\mathrm{cos}m\varphi ,`$ (52)
with (52) taking into account the invariance of the differential equation for $`f_m`$ under the interchange $`bb`$. Note that we may use Eq. (47) to find
$`\psi _𝐩^{()}(𝐱_{})=f_0(p,b)+2{\displaystyle \underset{m=1,\mathrm{}}{}}f_m(p,b)(i)^m\mathrm{cos}m\varphi .`$ (53)
In practice a finite number of terms is needed, with $`mm_{max}2pR/\mathrm{}`$.
Use Eq. (52) in (50) to find
$`\left({\displaystyle \frac{d^2}{db^2}}+{\displaystyle \frac{1}{b}}{\displaystyle \frac{d}{db}}+(p^2{\displaystyle \frac{m^2}{b^2}})\right)f_m(p,b)U({\displaystyle \frac{b^2}{R^2}})f_m(p,b)=0.`$ (54)
Note that for large enough $`b`$, $`U`$ vanishes and the $`f_m`$ are linear combinations of Bessel $`J_m`$ and Neumann $`N_m`$ functions or Hankel functions $`H_m^{(1,2)}`$.
This differential equation can be solved numerically using the Runge-Kutta technique. One determines the wave function by matching the numerical solutions to the analytic solution
$`f_m(p,b)=A_m\left(J_m(pb)+T_mH_m^{(1)}(pb)\right).`$ (55)
One matches the numerical function and its derivative (at large enough $`b`$ so that $`U=0`$) so as to determine the constants $`A_m,T_m`$. The normalization of $`\psi _𝐩^{(\pm )}`$ is such that $`A_m=1`$–asymptotically the wave is a sum of an ordinary plane wave and an outgoing wave, and the functions $`f_m(p,b)`$ can be regarded as phase shifted Bessel functions. Using the partial-wave form of two-dimensional wave function (52) will simplify the evaluation of $`S(x,K)`$ of Eq. (20).
## VII The Distorted Wave Emission Function (DWEF) and the Large Source Approximation (LSA)
Using (46) in Eq.(20) allows the integrals over $`y^0`$ and $`K^0^{}`$ to be evaluated as:
$`S(x,K,q)={\displaystyle \frac{d^3K^{}}{(2\pi )^3}S_0(x;K^0,𝐊^{})e^{i(\omega _2\omega _1)x^0}d^3ye^{i𝐊^{}𝐲}\mathrm{\Psi }_{𝐩_1}^{()}(𝐱+𝐲/2)\mathrm{\Psi }_{𝐩_2}^{()}(𝐱𝐲/2)},`$ (56)
and using (49) allows the integrals over $`y^3`$ and $`K_{}^{3}{}_{}{}^{}`$ to be evaluated yielding a four-dimensional integral:
$`S(x,K,q)`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^2}}𝒮_0(\tau ,\eta )e^{iq^0tiq_lz}{\displaystyle d^2b^{}\stackrel{~}{B}_\eta (𝐛,𝐛^{})\psi _{𝐩_1}^{()}(𝐛+𝐛^{}/2)\psi _{𝐩_2}^{()}(𝐛𝐛^{}/2)},`$ (57)
$$\stackrel{~}{B}_\eta (𝐛,𝐛^{})d^2K_T^{}B_\eta (𝐛,𝐊_T^{})\mathrm{exp}\left[i𝐊_T^{}𝐛^{}\right].$$
The result (57) still requires the evaluation of an 6-dimensional integral (over $`\tau ,\eta ,𝐛^{},𝐊_T^{}`$) to obtain the correlation function. We search for simplifications. The integral (57) simplifies if we ignore the effects of transverse flow rapidity. So to gain insight, let’s set
$`\eta _f=0,`$ (58)
consider a fixed value of $`\eta `$, and take one of the terms in the series (37). Then
$`\stackrel{~}{B}_\eta (𝐛,𝐛^{})=\rho (b)g(𝐛^2)`$ (59)
$`g(𝐛^2)=2{\displaystyle d^2K_{}M_{}\mathrm{exp}\left[\frac{M_{}\mathrm{cosh}\eta }{T}\right]\mathrm{exp}\left[i𝐊_{}𝐛^{}\right]}`$ (60)
Fig. 1 shows $`g(b)`$ for the case $`T/\mathrm{cosh}\eta =m_\pi `$. It is clear from (60) that the quantity $`T`$ controls the range of allowed values of $`K_{}`$, so that the extent of $`b^{}`$ is of order $`1/T1\mathrm{fm}`$ for $`T`$ = 200 MeV. This is much smaller than the size of the presumed fireball that controls the extent of $`b`$ in (57). Thus it is natural to think of neglecting the terms involving $`\pm 𝐛^{}/2`$ of (57). This however, is too extreme an approximation because it would not lead to the formalism of Sect. 1 in the plane wave limit. Instead, we use the approximation
$`\psi ^{()}(𝐛\pm 𝐛^{}/2)e^{i𝐩_{}𝐛^{}/2}\psi ^{()}(𝐛)`$ (61)
This is exact in the plane wave limit, but its wider validity relies on the replacement
$`\psi _{𝐩_i}^{()}(𝐛+𝐛^{}/2)\psi _{𝐩_j}^{()}(𝐛𝐛^{}/2)g(𝐛^2)\psi _{𝐩_i}^{()}(𝐛)\psi _{𝐩_j}^{()}(𝐛)g(𝐛^2)\mathrm{exp}(i𝐊_{}𝐛^{}),`$ (62)
that requires that the size of the source be much larger than $`n/T`$, where $`n`$ is the expansion order in the Bose-Einstein distribution (37). Our sources typically have a diameter of about 25 fm, and $`T1\mathrm{fm}^1`$, so that $`n`$ must be no bigger than about 25. We achieve numerical convergence with $`n<10`$, so that the source is truly large enough for our approximation. We denote Eq. (62) to be the “Large Source Approximation” (LSA) and use it to immediately integrate over $`𝐛^{}`$ and to obtain a simpler version of (57):
$`S(x,K)=𝒮_0(\tau ,\eta ,Y)e^{i(\omega _2\omega _1)\tau \mathrm{cosh}\eta }e^{iq_L\tau \mathrm{sinh}\eta }B_\eta (𝐛,𝐊)\psi _{𝐩_1}^{()}(𝐛)\psi _{𝐩_2}^{()}(𝐛),`$ (63)
that is obtained for any value of $`\eta _f`$. The expansion (37) gives
$`B_\eta (𝐛,𝐊)={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\mathrm{exp}({\displaystyle \frac{Ku+\mu _\pi }{T_n}})M_T\rho (b)`$
$`={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\mathrm{exp}({\displaystyle \frac{M_T\mathrm{cosh}\eta \mathrm{cosh}\eta _t(b)+\mu _\pi }{T_n}})\mathrm{exp}({\displaystyle \frac{K_T\mathrm{sinh}\eta _t(b)\mathrm{cos}\varphi }{T_n}})M_T\rho (b),`$ (64)
so that Eq. (63) becomes
$`S(x,K,q)=𝒮_0(\tau ,\eta ,Y)e^{i(\omega _2\omega _1)\tau \mathrm{cosh}\eta iq_L\tau \mathrm{sinh}\eta }{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\mathrm{exp}(\gamma _n(b)\mathrm{cosh}\eta )B_n(𝐛,𝐊)\psi _{𝐩_1}^{()}(𝐛)\psi _{𝐩_2}^{()}(𝐛)`$ (65)
$`B_n(𝐛,𝐊)\mathrm{exp}(\mu _\pi /T_n)\mathrm{exp}(K\mathrm{sinh}\eta _t(b)\mathrm{cos}\varphi /T_n)M_T\rho (b)`$ (66)
$`S(x,p_i)=𝒮_0(\tau ,\eta ,Y){\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\mathrm{exp}(\gamma _n(b)\mathrm{cosh}\eta )B_n(𝐛,𝐩_i)\left|\psi _{𝐩_i}^{()}(𝐛)\right|^2,`$ (67)
$`\gamma _n(b){\displaystyle \frac{M_T\mathrm{cosh}\eta _t(b)}{T_n}}.`$ (68)
The number of terms in the summation over $`n`$ required to achieve an accurate result depends on the values of $`\mu _\pi ,𝐊`$. If desired the quantities $`T,\mu _\pi `$ can be treated as functions of $`b`$ CL96a in Eq. (66), rather than as constants, as we have done here.
## VIII Correlation function
The specific form of the wave function enables us to express the correlation function as $`C(q_L\widehat{𝐳}+𝐪,𝐊)`$:
$`C(q_L\widehat{𝐳}+𝐪,𝐊)=1+{\displaystyle \frac{|d^4xS(x,K,q)|^2}{d^4xS(x,p_1)d^4xS(x,p_2)}}`$ (69)
$`𝐩_{1,2}=𝐊\pm q_L/2\widehat{𝐳}\pm 𝐪/\mathrm{𝟐}.`$ (70)
The dependence on $`q_L`$ occurs only in the numerator.
The remaining task is to perform the integral over $`d^4x`$. We use two separate approaches. The first Cramer:2004ih involves expanding $`S(x,K)`$, Eq. (66), as double power series in $`q_L`$ and $`\tau `$, keeping all terms up to second order. The second involves exact numerical integration. The two methods yield nearly identical results, with the second being more accurate and taking only slightly more computer time. We present both methods here. First, we take $`𝐪`$ to be small and make expansions Sec. VIII.1. This formalism is applied to obtain numerical results for radii and spectra in Sects. IX.1,IX.5 and Sect. IX.6. The formalism to compute the correlation functions without the use of expansion is contained in Sec. VIII.2, and this formalism is applied to compute correlation functions in Sec. IX.7.
### VIII.1 Evaluation of correlation function by expansion
Now we make the above mentioned expansion keeping terms to order $`q_L^2`$, (and anticipate the integration over $`\eta ,\tau `$). The term linear in $`q_L`$ is an odd function of $`\eta `$ so that it vanishes when the integration over $`\eta `$ is carried out. Use $`\omega _2\omega _1=q_o\beta ,\beta K_T/M_T`$. Thus the expansion of Eq. (66) is
$`S(x,K)=𝒮_0(\tau ,\eta ,Y)(1iq_o\beta \tau \mathrm{cosh}\eta {\displaystyle \frac{1}{2}}q_o^2\beta ^2\tau ^2\mathrm{cosh}^2\eta {\displaystyle \frac{1}{2}}q_L^2\tau ^2\mathrm{sinh}^2\eta )`$
$`\times {\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\mathrm{exp}(\gamma _n(b)\mathrm{cosh}\eta )B_n(𝐛,𝐊_T)\psi _{𝐩_1}^{()}(𝐛)\psi _{𝐩_2}^{()}(𝐛),`$ (71)
The effects of the $`Y`$ dependence of $`𝒮_0(\tau ,\eta ,Y)`$ that appears for $`q_L0`$ is the same for the numerators and denominators of the correlation functions and vanish. Hence we do not keep this explicit dependence in the intermediate steps in the following calculations.
The next step in evaluating (69) is to integrate over all $`\tau ,\eta `$ using the measure $`d\eta \tau d\tau `$. The first integral to appear arises from the factor of unity appearing inside the parenthesis of (71). It is
$`I_0{\displaystyle \frac{1}{\sqrt{2\pi }\mathrm{\Delta }\tau }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\eta \mathrm{cosh}\eta {\displaystyle _{\mathrm{}}^{\mathrm{}}}\tau 𝑑\tau \mathrm{exp}({\displaystyle \frac{\eta ^2}{2(\mathrm{\Delta }\eta )^2}})\mathrm{exp}({\displaystyle \frac{(\tau \tau _0)^2}{2(\mathrm{\Delta }\tau )^2}})\mathrm{exp}(\gamma _n(b)\mathrm{cosh}\eta )`$
$`2\tau _0\mathrm{exp}({\displaystyle \frac{1}{\mathrm{\Delta }\eta ^2}})K_1(\gamma _n+{\displaystyle \frac{1}{\mathrm{\Delta }\eta ^2}})\tau _0f_0(\xi _n)\mathrm{exp}({\displaystyle \frac{1}{\mathrm{\Delta }\eta ^2}}),f_0(\xi _n)=2K_1(\xi _n),\xi _n\gamma _n+{\displaystyle \frac{1}{\mathrm{\Delta }\eta ^2}}.`$ (72)
The approximation involves the replacement
$`\mathrm{exp}({\displaystyle \frac{\eta ^2}{2(\mathrm{\Delta }\eta )^2}})\mathrm{exp}({\displaystyle \frac{1}{\mathrm{\Delta }\eta ^2}})\mathrm{exp}({\displaystyle \frac{1}{\mathrm{\Delta }\eta ^2}}\mathrm{cosh}\eta ).`$ (73)
Similar replacements are made below. The dominant contributions to the integral involve small values of $`\eta `$, so the approximation is expected to be very good. The appendix shows that the error involved is less than about 2%.
We proceed to use the same replacement to evaluate the remaining integrals and find
$`I_2{\displaystyle \frac{1}{\sqrt{2\pi }\mathrm{\Delta }\tau }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{cosh}\eta d\eta {\displaystyle _{\mathrm{}}^{\mathrm{}}}\tau 𝑑\tau \mathrm{exp}({\displaystyle \frac{\eta ^2}{2(\mathrm{\Delta }\eta )^2}})\mathrm{exp}({\displaystyle \frac{(\tau \tau _0)^2}{2(\mathrm{\Delta }\tau )^2}})\mathrm{exp}(\gamma _n(b)\mathrm{cosh}\eta )(iq_0\beta )\tau \mathrm{cosh}\eta `$
$`(iq_0\beta )(\tau _0^2+(\mathrm{\Delta }\tau )^2)\mathrm{exp}({\displaystyle \frac{1}{\mathrm{\Delta }\eta ^2}})f_2(\xi _n),f_2(\xi _n)=K_0(\xi _n)+K_2(\xi _n)=2(K_0(\xi _n)+{\displaystyle \frac{K_1(\xi _n)}{\xi _n}}),`$ (74)
$`I_3{\displaystyle \frac{1}{\sqrt{2\pi }\mathrm{\Delta }\tau }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{cosh}\eta d\eta {\displaystyle _{\mathrm{}}^{\mathrm{}}}\tau 𝑑\tau \mathrm{exp}({\displaystyle \frac{\eta ^2}{2(\mathrm{\Delta }\eta )^2}})\mathrm{exp}({\displaystyle \frac{(\tau \tau _0)^2}{2(\mathrm{\Delta }\tau )^2}})\mathrm{exp}(\gamma _n(b)\mathrm{cosh}\eta )(q_0\beta )^2\tau ^2\mathrm{cosh}^2\eta `$
$`(3\tau _0\mathrm{\Delta }\tau ^2+\tau _0^3)(q_0\beta )^2\mathrm{exp}({\displaystyle \frac{1}{\mathrm{\Delta }\eta ^2}})f_3(\xi _n),f_3(\xi _n)=2{\displaystyle \frac{K_2(\xi _n)}{\xi _n}}=2\left(K_0(\xi _n)/\xi _n+K_1(\xi _n)(1+{\displaystyle \frac{2}{\xi _n^2}})\right),`$ (75)
$`I_1{\displaystyle \frac{1}{\sqrt{2\pi }\mathrm{\Delta }\tau }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\eta \mathrm{cosh}\eta {\displaystyle _{\mathrm{}}^{\mathrm{}}}\tau 𝑑\tau \mathrm{exp}({\displaystyle \frac{\eta ^2}{2(\mathrm{\Delta }\eta )^2}})\mathrm{exp}({\displaystyle \frac{(\tau \tau _0)^2}{2(\mathrm{\Delta }\tau )^2}})\mathrm{exp}(\gamma _n(b)\mathrm{cosh}\eta )\tau ^2\mathrm{sinh}^2\eta `$
$`\tau _0(3\mathrm{\Delta }\tau ^2+\tau _0^2)\mathrm{exp}({\displaystyle \frac{1}{\mathrm{\Delta }\eta ^2}})f_1(\xi _n),f_1(\xi _n)=2(K_2(\xi _n)/\xi _n)=2K_0(\xi _n)/\xi _n+4K_1(\xi _n)/\xi _n^2.`$ (76)
Then
$`{\displaystyle }d^4xS(x,K,q)=\tau _0\mathrm{exp}({\displaystyle \frac{1}{\mathrm{\Delta }\eta ^2}})\times `$
$`\left[\mathrm{\Phi }_{12}iq_o\beta (\tau _0+\mathrm{\Delta }\tau ^2/\tau _0)F_2(K_T){\displaystyle \frac{1}{2}}q_o^2\beta ^2(3\mathrm{\Delta }\tau ^2+\tau _0^2)F_3(K_T)q_L^2/2F_1(K_T)\right],`$ (77)
where
$`\mathrm{\Phi }_{12}={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle d^2bf_0(\xi _n)B_n(𝐛,𝐊_T)\psi _{𝐩_1}^{()}(𝐛)\psi _{𝐩_2}^{()}(𝐛)}`$ (78)
$`F_2(K_T)={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle d^2bf_2(\xi _n)B_n(𝐛,𝐊_T)|\psi _{K_T}^{()}(𝐛)|^2}`$ (79)
$`F_3(K_T)={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle d^2bf_3(\xi _n)B_n(𝐛,𝐊_T)|\psi _{K_T}^{()}(𝐛)|^2}`$ (80)
$`F_1(K_T)={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle d^2bf_1(\xi _n)B_n(𝐛,𝐊_T)|\psi _{K_T}^{()}(𝐛)|^2}`$ (81)
The denominator is obtained by evaluating the emission function (68) and the following functions enter:
$`F_0(K_T)={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle d^2bf_0(\xi _n)B_n(𝐛,𝐊_T)|\psi _{K_T}^{()}(𝐛)|^2}`$ (82)
$`\mathrm{\Phi }_{ii}={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle d^2bf_0(\xi _n)B_n(𝐛,𝐊_T)|\psi _{𝐩_i}^{()}(𝐛)|^2}.`$ (83)
With these definitions, one may show:
$`C(𝐪+q_L\widehat{𝐳},𝐊_T)=1q_l^2R_l^2q_o^2\beta ^2\stackrel{~}{\mathrm{\Delta }\tau }^2+{\displaystyle \frac{|\mathrm{\Phi }_{12}|^2}{\mathrm{\Phi }_{11}\mathrm{\Phi }_{22}}}`$ (84)
$`R_l^2=(3\mathrm{\Delta }\tau ^2+\tau _0^2)F_1(K_T)/F_0(K_T)`$ (85)
$`\stackrel{~}{\mathrm{\Delta }\tau }^2=(3\mathrm{\Delta }\tau ^2+\tau _0^2)F_3(K_T)/F_0(K_T)(\tau _0+\mathrm{\Delta }\tau ^2/\tau _0)^2\left|{\displaystyle \frac{F_2(K_T)}{F_0(K_T)}}\right|^2.`$ (86)
The spectra are given by Eq. (3) Wiedemann:1999qn
$`E_p{\displaystyle \frac{dN}{d^3p}}={\displaystyle \frac{dN}{dYM_{}dM_{}d\varphi _p}}={\displaystyle d^4xS(x,p)},`$ (87)
so that in general
$`{\displaystyle \frac{dN}{dM_{}^2}}={\displaystyle \frac{1}{2}}(2\pi )){\displaystyle }dY{\displaystyle }d^4xS(x,p),`$ (88)
in which the azimuthal symmetry of the angular distribution is used.
The STAR detector receives pions for values of $`Y`$ between $`\pm 0.5`$ and presents its results in terms of $`\frac{dN}{2\pi M_{}dM_{}dY}_{|Y|<0.5}`$ which is given by
$`{\displaystyle \frac{dN}{2\pi M_{}dM_{}dY}}={\displaystyle _{0.5}^{0.5}}𝑑Y{\displaystyle d^4xS(x,K)}.`$ (89)
Numerical studies showed us that the average over $`Y`$ is extremely well approximated by simply replacing $`Y`$ by 0. This requires $`\delta \eta >1/2,`$ and we use $`\delta \eta 1`$. Thus we find
$`{\displaystyle \frac{dN}{2\pi M_{}dM_{}dY}}_{|Y|<0.5}{\displaystyle \frac{\tau _0}{8\pi ^3}}\mathrm{exp}({\displaystyle \frac{1}{\mathrm{\Delta }\eta ^2}})F_0(𝐊_T).`$ (90)
The evaluation proceeds by reducing the two dimensional integrals of Eqs. (78-83) to those of one dimension by an analytic evaluation of the angular integrals. We need the partial wave expansions
$`\psi _{𝐩_1}^{()}{}_{}{}^{}(𝐛)=f_0(p_1,b)+2{\displaystyle \underset{m=1,\mathrm{}}{}}f_m(p_1,b)(i)^m\mathrm{cos}m\varphi _1`$ (91)
$`\psi _{𝐩_2}^{()}(𝐛)=f_0^{}(p_2,b)+2{\displaystyle \underset{m=1,\mathrm{}}{}}f_{m}^{}{}_{}{}^{}(p_2,b)(i)^m\mathrm{cos}m\varphi _2,`$ (92)
where $`\mathrm{cos}\varphi _i=\widehat{𝐩}_i\widehat{𝐛}.`$
We encounter integrals of the form
$`A_{mn}(z)={\displaystyle _0^{2\pi }}𝑑\varphi e^{z\mathrm{cos}\varphi }\mathrm{cos}m\varphi _1\mathrm{cos}n\varphi _2.`$ (93)
For $`𝐪𝐊,𝐪=𝐪_o,\varphi _1=\varphi _2=\varphi `$, so that we define
$`A_{mn}(,z){\displaystyle _0^{2\pi }}d\varphi e^{z\mathrm{cos}\varphi }\mathrm{cos}m\varphi \mathrm{cos}n\varphi `$ (94)
$`=\pi \left(I_{m+n}(z)+I_{mn}(z)\right),`$ (95)
where $`I_{m\pm n}(z)`$ are modified Bessel functions:
$`I_n(z)=(i)^nJ_n(iz),`$ (96)
for real $`z`$, and $`I_n(z)=I_n(z)`$. An integral representation is
$`I_n(z)={\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{2\pi }}𝑑\varphi e^{z\mathrm{cos}\varphi }\mathrm{cos}n\varphi .`$ (97)
For $`𝐪𝐊=0`$, $`𝐪=𝐪_s`$, we define
$`\mathrm{cos}\alpha ={\displaystyle \frac{K}{\sqrt{K^2+q^2/4}}},\mathrm{sin}\alpha ={\displaystyle \frac{q/2}{\sqrt{K^2+q^2/4}}}.`$ (98)
Then
$`A_{mn}(,z)=\pi \left(I_{m+n}(z)\mathrm{cos}(mn)\alpha +I_{mn}(z)\mathrm{cos}(m+n)\alpha \right).`$ (99)
Note that for the denominators, one gets integrals in which the two angles of (93) are the same. Then the expression (95) is to be used.
The use of (95) and (99) in (65) yields
$`\mathrm{\Phi }_{12}(,)`$ $`=`$ $`M_{}(K){\displaystyle _0^{\mathrm{}}}b𝑑b\rho (b)e^{M_{}(K)\mathrm{cosh}\eta _t(b)/T}`$ (100)
$`\times \left[{\displaystyle \underset{m,n=0}{\overset{\mathrm{}}{}}}ϵ_mϵ_nf_m(p_1,b)f_n^{}(p_2,b)(i)^{nm}A_{mn}((,),{\displaystyle \frac{K}{T}}\mathrm{sinh}\eta _t(b))\right]`$
$`ϵ_0`$ $`=`$ $`1,ϵ_{n>0}=2.`$ (101)
The notation $`(,)`$ denotes either of the two possibilities $`𝐪𝐊,𝐪𝐊`$. Similarly,
$`\mathrm{\Phi }_{ii}(,)=M_{}(p_i){\displaystyle _0^{\mathrm{}}}bdb\rho (b)e^{M_{}(p_i)\mathrm{cosh}\eta _t(b)/T}`$
$`\times \left[{\displaystyle \underset{m,n=0}{\overset{\mathrm{}}{}}}ϵ_mϵ_nf_m(p_i,b)f_n^{}(p_i,b)(i)^{nm}A_{mn}((,),K\mathrm{sinh}\eta _t(b)/T)\right]`$ (102)
### VIII.2 Exact numerical evaluation of correlation function
We present an alternate calculational technique that gives the correlation function for any value of $`𝐪`$. To do this, start by recalling Eqs. (69) and (63) and note the appearance of the integral
$`I_{\tau \eta }`$
$`{\displaystyle \frac{1}{\sqrt{2\pi }\mathrm{\Delta }\tau }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\eta \mathrm{cosh}\eta \mathrm{exp}(\gamma _n(b)\mathrm{cosh}\eta )\mathrm{exp}({\displaystyle \frac{\eta ^2}{2(\mathrm{\Delta }\eta )^2}})`$
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}\tau 𝑑\tau \mathrm{exp}({\displaystyle \frac{(\tau \tau _0)^2}{2(\mathrm{\Delta }\tau )^2}})e^{i\tau \left((\omega _2\omega _1)\mathrm{cosh}\eta iq_L\mathrm{sinh}\eta \right)}`$
$`={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\eta \mathrm{cosh}\eta \mathrm{exp}(\gamma _n(b)\mathrm{cosh}\eta +i\alpha \tau _0)\mathrm{exp}({\displaystyle \frac{\eta ^2}{2(\mathrm{\Delta }\eta )^2}})(\tau _0+i\alpha \mathrm{\Delta }\tau ^2)\mathrm{exp}(\alpha ^2\mathrm{\Delta }\tau ^2/2),`$
(103)
$`\alpha (\omega _2\omega _1)\mathrm{cosh}\eta q_L\mathrm{sinh}\eta .`$ (104)
The procedure of this section is to evaluate the integral over $`\eta `$ numerically.
Let’s set up the full calculation. Integrating over $`\tau `$ and using the result (103) yields
$`{\displaystyle }d^4xS(x,K,q)={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle }d^2bB_n(𝐛,𝐊_T))\psi _{𝐩_1}^{()}(𝐛)\psi _{𝐩_2}^{()}(𝐛)\times `$
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\eta \mathrm{cosh}\eta \mathrm{exp}(\gamma _n(b)\mathrm{cosh}\eta +i\alpha \tau _0)\mathrm{exp}({\displaystyle \frac{\eta ^2}{2(\mathrm{\Delta }\eta )^2}})(\tau _0+i\alpha \mathrm{\Delta }\tau ^2)\mathrm{exp}(\alpha ^2\mathrm{\Delta }\tau ^2/2),`$
There are some remarks to be made here: there is no need to make the approximation of Eq. (144), and the present expression is actually much more compact than (77). Possible cross terms Chapman:1994yv involving $`q_Oq_L`$ are small for the experiments with $`Y0`$ that we analyze, so we neglect the cross terms and take either $`q_L=0`$ or $`q_O=0`$. Then the integrands have terms either even or odd in $`\eta `$. The odd terms cancel. So define (with $`\mathrm{\Delta }\omega \omega _2\omega _1)`$)
$`I(\gamma _n(b),\mathrm{\Delta }\omega ,q_L,\mathrm{\Delta }\eta ,\mathrm{\Delta }\tau ,\tau _0)`$
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\eta \mathrm{cosh}\eta \mathrm{exp}(\gamma _n(b)\mathrm{cosh}\eta +i\alpha \tau _0)\mathrm{exp}({\displaystyle \frac{\eta ^2}{2(\mathrm{\Delta }\eta )^2}})(\tau _0+i\alpha \mathrm{\Delta }\tau ^2)\mathrm{exp}(\alpha ^2\mathrm{\Delta }\tau ^2/2).`$ (106)
We use various specific values of the arguments of the function $`I`$ to compute the different observables. According to Eqs. (25,26) a radius $`R_i`$ can be computed using $`q_i0,q_{ji}=0`$. Thus to compute $`R_O`$ we take $`q_{L,S}=0`$, so that
$`I(\gamma ,\mathrm{\Delta }\omega ,0,\mathrm{\Delta }\eta ,\mathrm{\Delta }\tau ,\tau _0)I_O(\gamma ,\mathrm{\Delta }\omega ,\mathrm{\Delta }\eta ,\mathrm{\Delta }\tau ,\tau _0)`$
$`2{\displaystyle _0^{\mathrm{}}}𝑑\eta \mathrm{cosh}\eta \mathrm{exp}(\gamma \mathrm{cosh}\eta )\mathrm{exp}({\displaystyle \frac{\eta ^2}{2(\mathrm{\Delta }\eta )^2}})e^{iz\tau _0}(\tau _0+iz\mathrm{\Delta }\tau ^2)\mathrm{exp}(z^2\mathrm{\Delta }\tau ^2/2)`$
$`z\mathrm{\Delta }\omega \mathrm{cosh}\eta .`$ (107)
To compute $`R_L`$ we take $`\mathrm{\Delta }\omega =0`$, so that
$`I(\gamma ,0,q_L,\mathrm{\Delta }\eta ,\mathrm{\Delta }\tau ,\tau _0)I_L(\gamma ,q_L,\mathrm{\Delta }\eta ,\mathrm{\Delta }\tau ,\tau _0)`$
$`2{\displaystyle _0^{\mathrm{}}}𝑑\eta \mathrm{cosh}\eta \mathrm{exp}(\gamma \mathrm{cosh}\eta )\mathrm{exp}({\displaystyle \frac{\eta ^2}{2(\mathrm{\Delta }\eta )^2}})(\tau _0\mathrm{cos}(y\tau _0)y\mathrm{\Delta }\tau ^2\mathrm{sin}(y\tau _0))\mathrm{exp}(y^2\mathrm{\Delta }\tau ^2/2)`$
$`yq_L\mathrm{sinh}\eta `$ (108)
In computing $`R_S`$, the spectra, or the denominator of the correlation function we have $`\mathrm{\Delta }\omega =0,q_L=0`$. Then we use:
$`I(\gamma ,0,0,\mathrm{\Delta }\eta ,\mathrm{\Delta }\tau ,\tau _0)2\tau _0{\displaystyle _0^{\mathrm{}}}𝑑\eta \mathrm{cosh}\eta \mathrm{exp}(\gamma \mathrm{cosh}\eta )\mathrm{exp}({\displaystyle \frac{\eta ^2}{2(\mathrm{\Delta }\eta )^2}}).`$ (109)
Then the correlation function is given by
$`C(𝐪+q_L\widehat{𝐳},𝐊)=1+{\displaystyle \frac{\left|\chi _{12}\right|^2}{\chi _{11}\chi _{22}}}`$ (110)
$`\chi _{12}={\displaystyle d^4xS(x,K,q)}`$
$`={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle d^2bB_n(𝐛,𝐊_T)I(\gamma _n(b),\mathrm{\Delta }\omega ,q_L,\mathrm{\Delta }\eta ,\mathrm{\Delta }\tau ,\tau _0)\psi _{𝐩_1}^{()}(𝐛)\psi _{𝐩_2}^{()}(𝐛)}`$ (111)
$`\chi _{ii}={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle d^2bB_n(𝐛,𝐩_i)I(\gamma _n(b),0,0,\mathrm{\Delta }\eta ,\mathrm{\Delta }\tau ,\tau _0)\left|\psi _{𝐩_i}^{()}(𝐛)\right|^2}`$ (112)
For very small values of $`𝐪`$ this correlation function reduces to the one obtained in second order, (84). There are two changes. If $`𝐪`$ is not small then one obtains the correct correlation function. A second change is that the approximation (144) is not used. The differences between using the procedure of this subsection and that of the previous subsection are negligible ($`<2\%`$ and indistinguishable) for the cases we have studied. This provides a further verification of the approximations used in Sect. (VIII.1), but the present formalism avoids those approximations.
## IX Applications
The description of the formalism is essentially complete. The plane wave emission function $`S_0`$ is defined in Sect. II and its symmetry discussed in Sect. IV. The distorted wave emission function function $`S`$ is defined in Sect. III, and its evaluation elaborated in Sect. VII. The optical potential is defined in Sec. V, and its use in a wave equation to obtain the distorted wave is discussed in Sec. VI. The correlation function is evaluated in Sect. VIII.
Our technique may be compared with the Buda-Lund model, an efficient representation of the dataBuda , in which the temperature and fugacity are taken as position-dependent functions appearing in a Boltzmann distribution. The effects of our optical potential could provide an explanation of some of those deduced dependencies.
We are now ready to confront the data.
### IX.1 DWEF Fits to Central Au+Au Collisions
Table 1 gives the parameters of our best fit to the STAR data for Au+Au central collisions at 200 GeV, which we will refer to as “F193”. For this fit to the data, as shown in Figs. 3 and 4, the $`\chi ^2`$ is 56.45, and the $`\chi ^2`$ per degree of freedom is 2.45. Table I also gives the estimated variances of those fit parameters that were varied, as calculated by determining the parameter variation required to increase the $`\chi ^2`$ value by one unit. Correlations between different parameters are not considered, but, as will be discussed below, more than one set of parameters can produce a quality fit to the data.
To illustrate the importance of the various effects of the DWEF model for computing the radii $`R_O,R_S`$ and the spectrum, we have done calculations using the F193 fit parameters with various effects switched on separately. This is shown in Fig. 2. The curves labeled DWEF show the full calculation. Those labeled PWEF are computed using plane waves, i.e., the optical potential ($`w_{0,2}`$) and flow ($`\eta _F`$) are set to zero. The curves labeled $`Re(w_0)`$ and $`Re(w_2)`$ use only the real constant or momentum-dependent parts of the optical potential, respectively, set to the F193 values, with no flow. The curves labeled $`Im(w_2)`$ use only the imaginary momentum-dependent part of the optical potential set to the F193 value, with no flow. Finally, the curves labeled $`flow`$ uses the F193 value of $`\eta _F`$, with no optical potential. All other parameters in these calculations are set to the F193 values of Table 1. This study indicates that both flow and the optical potential modify the HBT radii, but only the momentum dependent parts of the optical potential ($`w_2`$) affect the spectrum.
Figs. 3 and 4 show DWEF calculations using F193, as compared with the STAR data STARHBT ; STARspec used in the fit. We see that both the HBT radii and the spectrum (including the spectrum normalization) are reproduce very well indeed. However, some clarification is needed here on the issue of resonance pions. Our model describes only those “direct” pions that are emitted directly within the fireball and that participate fully in the HBT correlation. It does not predict the “halo” fraction of pions originating from resonance decays later in the process that will not participate in the HBT correlation (or that would make an unmeasurable “spike” in the correlation function near $`q=0`$). These latter pions are present in the measured spectrum and must be removed before fitting the spectrum with a DWEF calculation.
We have done this by using the HBT $`\lambda `$ parameter previously extracted as part of the HBT analysis, which we take to represent the probability that both correlated pions are direct rather than halo pions. We fit the extracted $`\lambda `$ parameter with a straight line in transverse momentum to obtain the function $`\lambda (p_T)`$. Then we remove the 12% weak decay correction of the published pion spectrumSTARspec and multiply it by $`\sqrt{\lambda (p_T)}`$, thereby producing an estimate of the momentum spectrum of direct pions only. This corrected spectrum, as indicated by the black small-dashed line in Fig. 4, is then used in the data fitting. Following the DWEF calculation, the predicted theoretical spectrum is “uncorrected” by dividing it by $`\sqrt{\lambda (p_T)}`$ and then reducing it by 12%, so that it can be compared directly with the published data. These predictions, now including halo pions, are the curves shown in the spectrum plots. We note that the effects of the two corrections to the spectrum tend cancel each other, and that their net effect on the fits is primarily an increase in temperature and flow parameters. We note also that the Phobos Collaboration has published spectrum data in the low $`p_T`$ regionBack:2004uh , but because the value of $`\lambda `$ in this region are not known, we cannot make similar corrections and therefore have not included these data in our current analysis. The DWEF fit parameters, particularly the temperature, flow, chemical potential, and Woods-Saxon diffuseness, are sensitive to the details of our treatment of resonances, thus introducing a procedure-dependent systematic error in DWEF parameter extraction.
Our DWEF calculations predict the correlation function of two pions with a particular vector momentum difference $`𝐪`$. There are at least two ways of extracting HBT radii from such correlation functions. One is to use the quadratic form (24) of the correlation function at very small magnitudes of $`𝐪`$. The other is to use the Gaussian form (23) at larger values of $`𝐪`$ that correspond to the falloff region of the correlation function. In our previous publicationCramer:2004ih we used the former method with $`q=K_T/40`$, a technique which essentially extracts HBT radii from the calculated curvature of the correlation function near $`q=0`$. In the present work we use the Gaussian form with $`q=0.15`$ fm<sup>-1</sup>, a momentum difference that evaluates the correlation function near its half-maximum point, a procedure that bears more resemblance to the extraction of radii from experimental data with Gaussian fitting. We have found that in the medium and high momentum regions the two procedures gives very similar descriptions of the data and extract similar radii. However, we have found that for values of $`K_T<175`$ MeV/c, the two methods diverge, and that the Gaussian form gives more reliable radii.
### IX.2 Wave Function Plots
We can understand more about how HBT measurements in various momentum regions probe the system under study by examining the computed wave functions from the DWEF calculation. Such wave functions calculated with the F193 fit are shown in Fig. 5. The complete DWEF wave functions (left column) differ significantly from the semi-classical eikonal approximation calculated with the same absorption (right column) at each of the energies we consider. The importance of the real part of the optical potential is illustrated by comparison with the center column, showing wave functions calculated with the real part of the optical potential set to zero. The maxima and minima of these plots result from interference effects caused by the optical potential as incorporated by solving the quantum mechanical wave equation. Their spacing gives a qualitative indication of the pion wavelength in the medium. The differences between the quantum calculations and the semi-classical eikonal model grow smaller as the value of $`K_T`$ increases. We note that while the high momentum wave functions sample only a “bright ring” at the edge of the fireball, the low momentum wave functions is also non-vanishing in most parts of the fireball.
### IX.3 Temperature Ambiguities
In our previous publicationCramer:2004ih the temperature parameter used in the pion emission function was $`T=222`$ MeV. This is an uncomfortably large value for the temperature of the medium, high enough that pions would be expected to “melt” in such an environment and lose their identity. Therefore, we explored (not shown) the extent to which the DWEF model requires such large values of temperature and showed that there is a parameter ambiguity involving the parameters $`T`$ and $`a_{WS}`$, such that a fit of reasonable quality could be obtained for temperature values over a fairly broad range.
Therefore, it is desirable to choose a temperature appropriate to the medium. If one takes quite literally the expectation that the DWEF model describes the initial emission of pions and that the first pions are produced directly in the strongly interacting quark-gluon plasma as it make a phase transition to a hadronic phase, then the emission function should have the temperature of the QGP transition. Similarly, if the pions are produced as massless objects due to chiral symmetry restoration in the medium, then the chemical potential should equal the free pion mass. Recent lattice gauge calculations reported at Quark Matter 2005Katz:2005br give the critical QGP transition temperature to be 193 MeV. Therefore, we adopt (Table I) T=193 MeV and $`\mu _\pi `$=139.57 MeV and search for a new fit to the STAR $`\sqrt{S_{NN}}`$=200 GeV Au+Au data. To our surprise, the fit we obtain (see Figs. 3 and 4) with parameters fixed at these values is the best we have found, with an overall $`\chi ^2`$ of 56.45 and a $`\chi ^2`$ per degree of freedom of 2.45. Further, subsequent searches in which the fixed temperature was set to other values between 173 MeV and 220 MeV show that there is a definite minimum in $`\chi ^2`$ at just the value of temperature given by the lattice gauge calculation. Table II shows nine different set of fit parameters, all of which give good fits to the central STAR $`\sqrt{S_{NN}}`$=200 GeV Au+Au data. In Table II, the fixed parameters are indicated in bold face. We see that F193, the parameter set of Table I, gives the best fit, and that for temperatures greater than 193 MeV it is necessary to use a non-zero value of $`ϵ`$, invoking a Kisslinger-type wave equation Eq. (45), to obtain a quality fit. However, fits F193a and F193b indicate that searching on $`ϵ`$ as part of a search at $`T=`$ 193 MeV gives values near zero and did not improve the fit.
Fig. 6 shows $`\chi ^2`$ per degree of freedom as a function of temperature for the nine fits listed in Table II. The solid line is a parabolic fit to the points. We do not assert that there is physics driving the preference of the STAR data for the temperature predicted by lattice gauge calculations, but we believe this is more than a coincidence.
Fig. 7 shows the ratio of the parameters of Table II to the F193 parameters, and indicates the range of variations and the correlations of parameters. Note in particular the correlation between $`T`$, $`a_{WS}`$, $`w_0`$, $`\mathrm{\Delta }\tau `$, and $`\eta _F`$.
Fig. 8 shows the predictions of all nine fits as compared with the STAR data. As can be seen, there are no striking differences. In order to make the differences more visible, Fig. 9 shows the ratio of the predictions of the eight other fits to those of F193. as c an be seen the differences in the radius predictions are less than 1%, while those of the spectrum predictions are as large as 10%.
However, we note that in the momentum region below 50 MeV/c, where no HBT or spectrum data is available, the different fits make different predictions. Fig. 10 shows the predictions of the fits of Table II in this region. We note that the “wiggles” arise mainly for the extreme fits, and in particular that these variations are relatively small for the F193 fit that provided the best fit to the higher momentum data.
### IX.4 Meaning of optical potential parameters
Table I shows that $`w_0=0.14`$ and that $`w_2`$ is 0.85 +0.12 i. Let us try to understand these values using the ideas of Sect. V. First consider the real parts. Comparing Eq. (45) and Eq. (44) shows that
$`w_0={\displaystyle \frac{m_\pi ^2}{u_\pi ^2}}m_\pi ^2(T),\mathrm{Re}(w_2)={\displaystyle \frac{1u_\pi ^2}{u_\pi ^2}}`$ (113)
which allows us to extract the values:
$`u_\pi ^2=0.54,u_\pi =0.74,u_\pi m_\pi (T)=0.65\mathrm{fm}^1.`$ (114)
These values are comparable to the estimates of Son:2002ci .
We also compare with the results of ShuryakShuryak:1991hb for pions propagating in hot dilute matter. His results for $`T=200`$ MeV (the stated limit of validity of the calculation) are most relevant for our case. His Fig. 1 shows results for $`V=U/(2m_\pi )`$ as a function of pion momentum. The imaginary potential is approximately proportional to $`p^2`$ for $`p400MeV/c`$. For $`p=1`$ fm<sup>-1</sup> Shuryak obtains $`Im(U)0.1`$ fm<sup>-2</sup>, which is close to our value of 0.12 fm<sup>-2</sup>. Our real potential is about seven times larger in magnitude than our imaginary potential, while Shuryak finds that the strengths are comparable. Thus our real potential is much stronger than expected from a dilute gas approximation and consistent (via the formalism of Son:2002ci ) with the occurance of a chiral phase transition.
### IX.5 Non-central Au+Au
Our analysis so far has focused on the central (0-5%) STAR $`\sqrt{S_{NN}}`$=200 GeV Au+Au data. However, the STAR collaboration performed measurements of pion correlations and spectra at $`\sqrt{S_{NN}}`$=200 GeV Au+Au as a function of centrality.
For non-central events, our optical potential would depend on the direction of $`𝐛`$ as well as its magnitude. The simple dependence on $`b`$ was exploited heavily in previous sections to simplify the calculations, so, in principle, non-cental collisions are not a part of our model. However, we can make simplifying assumptions that can allow us to predict the observables for non-central collisions. In particular, we assume that a non-central collision resembles a central collision with the same number of participants. This assumption allows us to extrapolate our results to systems that do not have perfect centrality by using participant scaling. In particular, we take the space-time parameters $`R_{WS},a_{WS}`$, and $`\tau _0`$ to scale as the centrality-dependent number of participant particles to the one-third power: $`N_{\mathrm{part}}^{1/3}`$. The values of $`N_{\mathrm{part}}`$ are taken from Glauber-model calculationsMillerthesis . For the Au+Au system, the value of $`\mathrm{\Delta }\tau `$ is kept at the value of F193, because this is a dynamic quantity describing the proper-time duration during which pions are emitted in the collision. Table III gives the parameters $`R_{WS},\tau _0,a_{WS}`$ and $`\mathrm{\Delta }\tau `$ vs. centrality, as scaled from the F193 fit of Table I. Parameters not listed in the table are the same as those of Table I.
In Fig. 11 we see that the participant-scaled predictions work fairly well for the first few centrality bins, but agreement diminishes for the more peripheral collisions. We do not expect this simple scaling procedure to be accurate very far from the purely central collision assumption of the model, so this loss of predictive power is to be expected.
In particular, looking for differences between peripheral and central collisions is a standard way to test for the possible existence of QGP effects. This means that the strength of the optical potential parameters of our model should change as the collision becomes peripheral. We have kept the potential depth parameters constant with centrality because we have no way to predict their centrality dependence. Fig. 11 shows an excellent reproduction of $`R_L`$ with centrality, a fairly good reproduction of $`R_O`$, but only a fair description of $`R_S`$. We take this to be consistent with the disappearance of a chiral phase transition in the more peripheral collisions.
### IX.6 Predictions for Cu+Cu collisions
The STAR collaboration has also measured (but not yet published) radii and spectra vs. centrality for Cu+Cu collisions at $`\sqrt{S_{NN}}`$=200 GeV, so it is of interest to make predictions for this system. Here again use participant scaling, but have also assumed that the emission duration parameter $`\mathrm{\Delta }\tau `$ scales as $`A^{1/3}`$, where $`A`$ is that atomic mass number of the colliding nuclei. Table IV gives the scaled space-time parameters used to predict the HBT radii for the Cu+Cu system. Parameters not present in the table are the same as those of Table I.
### IX.7 Correlation functions and the Gaussian approximation
We apply the present formalism to obtain correlation function $`C`$ (with $`\lambda =1`$), using the parameters of Table I, for $`K_T=158,316`$ MeV/c. The results are shown in Figs. 13,14. We see that the correlation functions $`C(K_T,q)`$ are fairly well represented by Gaussians, for the relevant range of small values of $`q`$, with widths that are approximately independent of $`K_T`$. For a typical radius of about 7 fm, and $`q=0.15`$ fm<sup>-1</sup>, (where data are measured) $`q^2R^21`$. Therefore, using the approximation (24) is not very accurate. On the other hand, the correlation functions are close to Gaussian in shape in this region, suggesting that approximation (23) is more appropriate. A detailed look at the ratio of computed correlation functions ($`C1`$) to its Gaussian fit in Fig. 14 shows that the Gaussian curves represent the correlation functions fairly well in the region $`0<q<0.22`$ fm<sup>-1</sup> where $`C1`$ is large and radii are extracted, but that the correlation functions are systematically larger than the Gaussian fit for $`R_{O,L}`$ and smaller for $`R_S`$ in the “tail” region ($`q>0.22`$ fm).
### IX.8 Possible Extensions of the DWEF Model
The DWEF model presented here uses the empirical “hydrodynamics-inspired” emission function of Eq. (27). However, we note that the application of distorted waves to an emission function is more generally applicable, and that the formalism we present here can be applied to any emission function that has the same symmetry properties as Eq. (27). One extension of the model would be to calculate the emission function as a multi-dimensional numerical table directly from hydrodynamics and use this with the DWEF model to calculate spectra and radii. We also note that, while it has not been not implemented here, the DWEF model can allow the temperature and chemical potential to depend on $`b`$, as is done in the Buda-Lund model Buda . Preliminary investigations of this extension of the model, however, did not lead to an improved description of the data.
## X The Eikonal Approximation
Several previous calculations Heiselberg:1997vh Wong:2003dh of the effects of opacity have used the eikonal approximation to solve the wave equation (50). Here we explain the nature of this approximation and discuss its weaknesses and (fewer) strengths when applied to the current situation.
The basic idea is that if the momentum is large one may say approximately that the wave propagate in a given direction (here the out direction, which is taken as along the $`x`$ axis). Then one assumes a solution of the form $`\psi _{}^{()}{}_{}{}^{}(𝐛)=e^{ipx}\mathrm{\Phi }(𝐛)`$ that is inserted into the wave equation. Taking the Laplacian of the approximate wave function gives a term proportional to $`p^2`$ that is canceled, a term proportional to $`p`$ that is kept, and another term that is ignored. Then one finds
$`\psi _{}^{()}{}_{}{}^{}(𝐛=x,y)=e^{ipx}\mathrm{exp}\left[{\displaystyle \frac{i}{2p}}{\displaystyle _x^{\mathrm{}}}U(x^{},y)𝑑x^{}\right].`$ (115)
for propagation in the $`x`$ (out) direction. The corrections to this solution are of order $`\frac{1}{2iK}\frac{1}{U}\frac{U}{x}\frac{U}{4K^2}`$ times the terms that are kept. The first correction can be large in the surface region in which $`U`$ varies greatly and the second term can be large in the interior region in which $`U`$ reaches its full value. We are concerned with pions of momentum ranging from 30 to 600 MeV/c, so that the eikonal approximation be expected to be poor. However, the ease of application, and its wide use makes it worthwhile for us to assess the use of Eq. (115). We consider a purely imaginary potential first and then use a general complex potential.
### X.1 Strong Absorption at High K – Purely Imaginary potential
In the impulse approximation $`U=4\pi f\rho `$ where $`f`$ is the projectile-target scattering amplitude and $`\rho `$ is the density of scatterers. The optical theorem relates the imaginary part of $`f`$ to the total cross section $`\sigma `$ so that $`Im[U]=p\sigma \rho `$. If we keep only the imaginary potential $`i/(2p)U=1/2\sigma \rho `$ and for a constant density the intensity of the wave falls as $`e^{x\lambda _{\mathrm{mfp}}}`$ with the mean free path, $`\lambda _{\mathrm{mfp}}=\frac{1}{\sigma \rho }`$. More generally the wave function for a purely imaginary optical potential is given by
$`\psi ^{()}(𝐩_i,𝐛)=e^{i𝐩_i𝐛}e^{l_i/2\lambda _{\mathrm{mfp}}}`$ (116)
where $`i=1,2`$ for the two wave functions and $`l_i(𝐛,𝐊)`$ is the direct line path length (parallel to the direction of $`𝐩_i`$ from the emission point $`𝐛`$ to the edge of the medium. For a purely imaginary optical potential $`U=iK\sigma \rho `$ and $`\lambda _{\mathrm{mfp}}`$ is the resulting mean free path.
We have, see Fig. 15,
$`𝐑=𝐥_𝐢+𝐛,`$ (117)
$`R^2=l_i^2+b^2+2bl_i\mathrm{cos}(\theta \theta _s).`$ (118)
$`\theta `$ is angle between $`𝐛`$ and $`x`$-axis (direction of $`𝐊`$) $`\theta _s`$ is angle between $`𝐩_{\mathrm{𝟏},\mathrm{𝟐}}`$ and $`𝐊`$. This was called $`\alpha `$ in previous sections. The angle between $`𝐛`$ and $`𝐥_𝐢`$ is $`\theta \theta _s`$. Solving we find
$`l=b\mathrm{cos}(\theta \theta _s)+\sqrt{(b\mathrm{cos}(\theta \theta _s))^2+R^2b^2}.`$ (119)
For the R-out case $`\theta _s=0`$, and
$`l_i=l=x+\sqrt{R^2y^2}`$ (120)
In this approximation the correlation function minus one is given by the absolute square of the ratio of integrals: $`\frac{I_,(q)}{I_,(q=0)}.`$ In particular,
$`I_{}(q)`$ $`=`$ $`{\displaystyle d^2be^{i𝐪𝐛}e^{l/\lambda _{\mathrm{mfp}}}},𝐪\widehat{𝐱}`$ (121)
$`=`$ $`{\displaystyle \frac{\lambda }{1+iq\lambda }}{\displaystyle _R^R}𝑑y\left(e^{iq\sqrt{R^2y^2}}e^{(iq+2/\lambda _{\mathrm{mfp}})\sqrt{R^2y^2}}\right).`$ (122)
If $`\lambda _{\mathrm{mfp}}<<R`$ we neglect the second term. Since we are interested in radii, we expand the remaining term in powers of the exponential (keeping $`\lambda _{\mathrm{mfp}}<<R`$):
$`I_{}(q)={\displaystyle \frac{\lambda _{\mathrm{mfp}}}{1+iq\lambda _{\mathrm{mfp}}}}\left(2Rq^2R^2/2(2R2R/3)+iq\pi R/2\right)`$ (123)
$`C_{}(q)2q^2R^2(2/3\pi ^2/16)`$ (124)
$`R_O^2=R^2(2/3\pi ^2/16)=0.0498R^2;R_O=R/4.48`$ (125)
In the plane wave approximation (Eq. (122) with $`\lambda _{\mathrm{mfp}}\mathrm{}`$) one would find $`R_O^{PW}=R/\sqrt{8}`$. The striking result of Heiselberg & VischerHeiselberg:1997vh indicated that the measured radius should be 40% smaller than the radius obtained in plane wave approximation. This result is confirmed for the highest value of $`K_T=3\mathrm{fm}^1600`$ MeV/c by the calculations shown above in Fig. 2. The value of $`R_o`$ is reduced by approximately 40% by the influence of the imaginary optical potential. Fig. 16 shows the effects of increasing the imaginary potential (by varying $`Im[w_2]`$ from 0.1 to 0.5 in steps of 0.1). The computed value of $`R_O`$ does not change when $`Im[w_2]`$ is large enough and $`K_T`$ is high enough. Small deviations between our results for highest $`K_T`$ and the resultHeiselberg:1997vh can be attributed to the non-zero value of the diffuseness $`a_{WS}`$.
In the model of the present section, the correlation function and pion intensity would each be proportional to $`\lambda _{\mathrm{mfp}}`$, and a very small value would yield a very small pionic spectrum. More generally, the extraction of the chemical potential from the pionic spectrum depends on the mean free path parameter $`\lambda _{\mathrm{mfp}}`$.
Now let us do the R-side case. Then
$`\mathrm{cos}\theta _s={\displaystyle \frac{K}{\sqrt{K^2+q^2/4}}};\mathrm{sin}\theta _s=\pm {\displaystyle \frac{q/2}{\sqrt{K^2+q^2/4}}}`$ (126)
$`\mathrm{cos}(\theta \theta _s)=\mathrm{cos}\theta (1{\displaystyle \frac{q^2}{8K^2}}){\displaystyle \frac{q}{2K}}\mathrm{sin}\theta .`$ (127)
We find that
$`l_ix+\sqrt{R^2y^2}\pm 𝒪({\displaystyle \frac{q}{K}})+𝒪({\displaystyle \frac{q^2}{K^2}})x+\sqrt{R^2y^2}.`$ (128)
The terms $`\pm 𝒪(\frac{q}{K})`$ cancel in computing the term $`l_1+l_2`$ that enters in computing the present correlation function. Thus the correction is of order $`1/K^2`$. The eikonal approximation works only if $`KR1,`$ so the corrections must be presumed to be small. Then
$`I_{}(q)={\displaystyle _R^R}𝑑ye^{iqy}{\displaystyle _{\sqrt{R^2y^2}}^{\sqrt{R^2y^2}}}𝑑xe^{x/\lambda _{\mathrm{mfp}}}e^{\frac{1}{\lambda _{\mathrm{mfp}}}\sqrt{R^2y^2}}`$ (129)
$`\lambda {\displaystyle _R^R}𝑑y(1q^2y^2/2)=\lambda _{\mathrm{mfp}}(2Rq^2R^3/3)`$ (130)
$`C_{}(q)=2q^2R^2/3`$ (131)
$`R_s^2=R^2/3,`$ (132)
Without distortion would be $`C_{}^{PW}(q)=2q^2R^2/4`$, so in this case the strong absorption increases the radius by a factor of $`\sqrt{4/3}.`$ This result is qualitatively obeyed in our realistic solutions of the wave equation, see Fig. 16.
### X.2 Complex potential with non-vanishing real part
If there is an attractive real potential the wave function of Eq.(116) becomes:
$`\psi ^{()}(𝐩_1,𝐛)=e^{i𝐩_1𝐛}e^{l_1(1+i\alpha )/2\lambda _{\mathrm{mfp}}}`$ (133)
with $`\alpha `$ dimensionless, real and positive. If $`\alpha =0`$ one obtains the purely absorptive model of the previous sub-section. Conversely, the limit of a purely real potential occurs when $`1/\lambda _{\mathrm{mfp}}0,\alpha /\lambda _{\mathrm{mfp}}1/\lambda _0`$. We also have
$`\psi _{}^{}{}_{}{}^{()}(𝐩_2,𝐛)=e^{i𝐩_2𝐛}e^{l_2(1i\alpha )/2\lambda _{\mathrm{mfp}}}`$ (134)
In the product $`\psi _{𝐩_1}\psi _{𝐩_2}^{}`$ enters in computing the correlation function, so unlike the previous case of pure absorption, a term of the form $`l_1l_2`$ enters. This difference is of order $`q/K`$ compared to other terms, but its influence in computing radii must lead eventually to a term of order $`(\frac{q}{K})^2`$. The validity of the eikonal approximation depends on the ability to disregard such terms. Thus a valid eikonal approximation means that $`l_1=l_2`$ so the factors of $`\alpha `$ cancel out in the product $`\psi ^{()}(𝐩_1,𝐛)\psi _{}^{}{}_{}{}^{()}(𝐩_2,𝐛)`$. Thus, if one assumes the eikonal approximation is valid at all values of $`K_T`$, one would find erroneously that the real potential never has an influence on the calculation of HBT radii.
Unlike the effects of the imaginary potential, which are qualitatively captured by the eikonal approximation (even if applied wrongly at low values of $`K_T`$), the effects of the real potential are completely lost. Thus the eikonal approximation can not be used for values of $`K_T`$ such that the real potential contributes. As shown in Fig. 3 the real potential is important for all values of $`K_T`$ less than 600 MeV/c. Conversely, for much larger values of $`K_T`$, for which the eikonal approximation does accurately reproduce the solution of the wave equation, the real potential will not play a role in determining radii.
## XI Oscillations – a simple square well example
Our numerical results indicate that the radii may have significant oscillations for small values of $`K_T=K`$. The purpose of this sub-section is to provide a simple example that also yields oscillating radii. Consider a cylindrical source of radius $`R`$ of infinite extent in the longitudinal direction. Suppose this leads to a real, attractive square well potential of radius $`R`$ that is proportional to the square of the momentum, as motivated by Eq. (45). Then (50) becomes
$`(^2U_0p^2)\psi =p^2\psi (bR)`$ (135)
$`^2\psi =p^2\psi (b>R),`$ (136)
with $`U_0>0`$. This equation is easily solved using the partial wave expansion (53). To provide a simple analytic example we further specify to the case of the lowest partial wave, $`m=0`$. In this case we find
$`\psi (𝐛)=J_0(p\sqrt{1+U_0}b),(bR)`$ (137)
This shows immediately that the effect of the interaction is to scale each momenta $`p_1,p_2`$ that appear in the wave functions by a factor $`\sqrt{1+U_0}.`$ We define
$`\stackrel{~}{p}p\sqrt{1+U_0}.`$ (138)
Eq. (137) is a valid approximation to the full wave function (for $`bR`$) only if $`p\sqrt{1+U_0}R1`$. However, it is interesting to also consider larger values of $`p`$.
To compute radii, recall the correlation function (69). For simplicity we take $`\eta _f=0`$, and consider fixed values of $`\eta `$ and $`\tau `$, with $`\mathrm{\Delta }\tau =0`$. This corresponds to evaluating $`S_0`$ at its peak and neglecting the influence of time duration. In this case many factors in the numerator and denominator of Eq. (70) cancel. Then the correlation function is given by a simple expression that provides some insight:
$`C(K,q)1={\displaystyle \frac{\varphi _R^2(\stackrel{~}{p_1},\stackrel{~}{p_2})}{\varphi _R(\stackrel{~}{p_1})\varphi _R(\stackrel{~}{p}_2)}}`$ (139)
$`\varphi _R(p_i,p_j){\displaystyle \frac{1}{R^2}}{\displaystyle _0^R}b𝑑bJ_0(p_ib)J_0(p_jb)={\displaystyle \frac{p_iJ_0(p_jR)J_1(p_iR)p_jJ_0(p_iR)J_1(p_jR)}{R(p_i^2p_j^2)}}`$ (140)
$`\varphi _R(p_i){\displaystyle \frac{1}{R^2}}{\displaystyle _0^R}b𝑑bJ_0^2(p_ib)={\displaystyle \frac{1}{2}}\left(J_0^2(p_iR)+J_1^2(p_iR)\right)`$ (141)
If we are concerned with the side radius, inside the well $`p_1=p_2=\sqrt{K^2+q^2/4}`$ for the arguments of $`J_0`$ that enter. In that case, $`\varphi _R(p_1,p_2)\varphi _R(p_1,p_1)=\varphi _R(p_1)`$ and the correlation function takes on the value of 2 and the side-radius vanishes. This is a specific consequence of the approximation of taking only $`m=0`$–all directions of $`𝐊`$ are equivalent, so there is no influence of vectors $`𝐪`$ that are perpendicular to $`𝐊`$.
The out case is more interesting because the energies and magnitudes of momenta $`p_{1,2}=(K\pm q/2)`$ of the two pions are different. A non-zero radius is obtained by evaluating $`C(K,q)`$ for very small values of $`q`$ and using Eq. (26). The quantity $`R_O(K)/R`$ is displayed in Fig. 17.
We see that an oscillatory pattern similar to our realistic results, such as Fig. 3 for small values of $`K_T`$, emerges from the present calculation. This occurs simply because Bessel functions oscilate. Indeed the out-radius obtained without the influence of final state distortions ($`U_0=0`$) also has oscillations, but with a different frequency. Therefore another quantity of interest is the ratio of radii computed with $`R_O(K)`$ and without $`R_O^0(K)`$ the influence of distortions:
$`{\displaystyle \frac{R_O(K)}{R_O^0(K)}}\mathrm{Ratio}`$ (142)
displayed in Fig. 17. The oscillations seen in the right part of Fig. 17 demonstrate the significance of final state interactions that cause enhancements by factors of greater than six and suppressions by factors of 2! Furthermore, comparing the left and right hand sides of the figure shows the quantity $`R_O^0(K)/R`$ also oscillates. This is caused by the sharp edge of the square well and because we limit ourselves to $`m=0`$.
## XII Summary
A complete formal treatment of the distorted waves treatment of HBT correlations is presented here. The need for incorporating the influence of an optical potential $`U`$ and the resulting distorted wave emission function is explained. The partial wave formalism necessary to compute the pionic distorted waves and the resulting emission function is detailed. Two different methods with equivalent results to evaluate the necessary eight-dimensional integral are described. Chiral symmetry restricts the form of $`U`$ boy ; Son:2002ci at low energy and the necessary constraints are implemented here and in Cramer:2004ih . An excellent description of the STAR Au+Au HBT and spectrum data is achieved for central collisions and the use of an average area formulation leads to a very good description of these observables for non-central collisions. We also use four different versions of geometrical scaling to predict the results of central Cu+Cu collisions. The ability to calculate the absolute magnitude of the spectrum as well as the radii Eq. (25),Eq. (26) (which involve ratios of functions of emmission functions) required for the computation of radii is a principal advantage of our method. The Blast Wave Model is discussed in Sec. IV.
The above results are achieved using a temperature $`T_c`$ fixed at the recently determined critical value of 193 MeV Katz:2005br . Such a value could present difficulties for conventional calculations of the spectra because chemical equilibrium analyses yield lower temperatures $`T_{ch}=174`$ MeV Braun-Munzinger:2001ip . A large difference between $`T_c`$ and $`T_{ch}`$ implies that the hadrons interact after the deconfinement transition occurs. This notion is entirely consistent with our treatment of pionic distortions which has as its fundamental assumption that pions interact in a hot dense medium before escaping to freedom.
We find that the necessary real optical potential is so strongly attractive that the pion can be said to lose its mass inside the medium. That chiral symmetry seems to be restored is the conclusion of our earlier workCramer:2004ih . The RHIC-HBT puzzle is therefore replaced by the need to investigate this restoration.
Explicit evaluation of wave functions obtained by exact numerical solutions of the wave equation show some interesting features of the strong interaction and also display differences with the solutions obtained using the eikonal approximation. A critical discussion of the eikonal approximation as applied to computing HBT radii shows that its use causes the crucial influence of the real part of the optical potential to be entirely lost. The huge importance of the real part of the optical potential is explicitly illustrated through a simple
There are many immediate applications of this formalism. In particular, a treatment of HBT data obtained at all energies is in progress and will be presented elsewhere.
## Acknowledgments
This work is partially supported by the USDOE grants Nos. DE-FG-02-97ER41014 and DE-FG-02-97ER41020. GAM thanks LBL, TJNAF and BNL for their hospitality during the course of this work. We thank W. Busza, T. Csörgő, J. Draper, M. Lisa, M. Luzum, S. Pratt, J. Rafelski, S. Reddy, E. Shuryak and D. Son for useful discussions.
## Appendix
Consider the integral appearing in Sect. VIII.1:
$`f(\beta ,\mathrm{\Delta }\eta )={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\eta \mathrm{exp}(\eta ^2/(2\mathrm{\Delta }\eta ^2))\mathrm{cosh}(\eta )\mathrm{exp}\left(\beta \mathrm{cosh}\eta \right),`$ (143)
that can be approximated Retiere:2003kf using Eq. (73)as
$`f(\beta ,\mathrm{\Delta }\eta )=2\mathrm{exp}({\displaystyle \frac{1}{\mathrm{\Delta }\eta ^2}})K_1(\beta +{\displaystyle \frac{1}{\mathrm{\Delta }\eta ^2}}).`$ (144)
Figure 18 shows that the approximation is excellent. |
warning/0507/quant-ph0507259.html | ar5iv | text | # Stronger Two-Observer All-Versus-Nothing Violation of Local Realism
## Abstract
We introduce a two-observer all-versus-nothing proof of Bell’s theorem which reduces the number of required quantum predictions from 9 \[A. Cabello, Phys. Rev. Lett. 87, 010403 (2001); Z.-B. Chen et al., Phys. Rev. Lett. 90, 160408 (2003)\] to 4, provides a greater amount of evidence against local realism, reduces the detection efficiency requirements for a conclusive experimental test of Bell’s theorem, and leads to a Bell’s inequality which resembles Mermin’s inequality for three observers \[N. D. Mermin, Phys. Rev. Lett. 65, 1838 (1990)\] but requires only two observers.
The Greenberger-Horne-Zeilinger (GHZ) proof GHZ89 ; Mermin90a of Bell’s theorem Bell64 provided a direct contradiction between the Einstein-Podolsky-Rosen (EPR) EPR35 local elements of reality (LERs) by considering 4 perfect correlations predicted by quantum mechanics. However, while Bell’s theorem required only two spacelike separated observers, the GHZ proof required three. A two-observer all-versus-nothing (AVN) proof Mermin90b based on 9 perfect correlations was introduced in Refs. Cabello01a ; Cabello01b , then adapted for two-photon systems CPZBZ03 , and recently tested in the laboratory YZZYZZCP05 . Despite requiring only two observers, this AVN proof has two disadvantages when compared to GHZ: it is more complex, in the sense that it requires a higher number of quantum predictions, and provides less evidence against local realism than GHZ VGG03 ; thus a conclusive experiment based on the two-observer proof would require higher detection efficiencies than one based on GHZ. Both disadvantages are indeed related: in the two-observer proof, 8 out of these 9 perfect correlations (i.e., 89%) can be explained by LERs; in GHZ, only 3 out of these 4 perfect correlations (i.e., only 75%) can be explained by LERs.
Here we introduce an AVN proof requiring two observers and only 4 perfect correlations. We then show that conclusive experiments based on this new proof would require lower detection efficiencies than those needed for the previous two-observer AVN proof, and derive the corresponding Bell inequality for real experiments.
Let us consider a two-photon system entangled both in polarization and path degrees of freedom CPZBZ03 ; YZZYZZCP05 ; Kwiat97 prepared in the state
$`|\psi `$ $`=`$ $`{\displaystyle \frac{1}{2}}(|Hu_1|Hu_2+|Hd_1|Hd_2`$ (1)
$`+|Vu_1|Vu_2|Vd_1|Vd_2),`$
where $`|H_j`$ and $`|V_j`$ represent horizontal and vertical polarization, and $`|u_j`$ and $`|d_j`$ denote two orthonormal path states for photon $`j`$. The state (1) can be viewed as a two-photon version of the four-qubit cluster state BR01 . Lu Lu05a and Lu Lu05b have described procedures for preparing state (1).
Let us also consider 6 local observables on photon $`j`$: 3 for polarization degrees of freedom, defined by the operators
$`X_j`$ $`=`$ $`|H_jV|+|V_jH|,`$ (2)
$`Y_j`$ $`=`$ $`i\left(|V_jH||H_jV|\right),`$ (3)
$`Z_j`$ $`=`$ $`|H_jH||V_jV|,`$ (4)
and 3 for path degrees of freedom, defined by the operators
$`x_j`$ $`=`$ $`|u_jd|+|d_ju|,`$ (5)
$`y_j`$ $`=`$ $`i\left(|d_ju||u_jd|\right),`$ (6)
$`z_j`$ $`=`$ $`|u_ju||d_jd|.`$ (7)
Each of these observables can have only one of two possible values: $`1`$ or $`1`$.
The proof has two steps. First, we will show that the 7 local observables $`X_1`$, $`Y_1`$, $`x_1`$, $`X_2`$, $`Y_2`$, $`y_2`$, and $`z_2`$ satisfy the EPR condition for LER; namely, “if, without in any way disturbing a system, we can predict with certainty (i.e., with probability equal to unity) the value of a physical quantity, then there exists an element of physical reality corresponding to this physical quantity” EPR35 . The observables $`X_1`$, $`Y_1`$, and $`x_1`$ of photon 1 are LERs because their values can be predicted with certainty from spacelike separated measurements on photon 2. Spacelike separation guarantees that photon 1 has not been disturbed in any way Einstein48 . These predictions with certainty follow from the fact that state (1) satisfies the following equations:
$`X_1X_2z_2|\psi `$ $`=`$ $`|\psi ,`$ (8)
$`Y_1Y_2z_2|\psi `$ $`=`$ $`|\psi ,`$ (9)
$`x_1Z_2x_2|\psi `$ $`=`$ $`|\psi .`$ (10)
Equation (8) tells us that, from the results of measuring $`X_2`$ and $`z_2`$ on photon 2, we can predict with certainty the result $`v(X_1)`$ of measuring $`X_1`$ on photon 1. Equation (9) tells us that, from the results of measuring $`Y_2`$ and $`z_2`$ on photon 2, we can predict with certainty the result $`v(Y_1)`$ of measuring $`Y_1`$ on photon 1. Equation (10) tells us that, from the results of measuring $`Z_2`$ and $`x_2`$ on photon 2, we can predict with certainty the result $`v(x_1)`$ of measuring $`x_1`$ on photon 1.
Analogously, the observables $`X_2`$, $`Y_2`$, $`y_2`$, and $`z_2`$ of photon 2 are also LERs because state (1) also satisfies the following equations:
$`X_1z_1X_2|\psi `$ $`=`$ $`|\psi ,`$ (11)
$`Y_1z_1Y_2|\psi `$ $`=`$ $`|\psi ,`$ (12)
$`Z_1y_1y_2|\psi `$ $`=`$ $`|\psi ,`$ (13)
$`z_1z_2|\psi `$ $`=`$ $`|\psi .`$ (14)
Therefore, the results $`v(X_2)`$, $`v(Y_2)`$, $`v(y_2)`$, and $`v(z_2)`$ of measuring $`X_2`$, $`Y_2`$, $`y_2`$, and $`z_2`$ on photon 2 can be predicted with certainty from spacelike separated measurements on photon 1.
Moreover, we can prove that two observables on the same photon, but corresponding to different degrees of freedom, like $`X_2`$ and $`z_2`$, are independent LERs in the sense that measuring one of them does not change the value of the other (thus there is no need for any additional assumptions beyond the EPR condition; see Marinatto03 for a similar discussion). For instance, a suitable measurement of $`X_2`$ does not change $`v(z_2)`$ because $`v(z_2)`$ can be predicted with certainty from a spacelike separated measurement on photon 1 \[see Eq. (14)\], and this prediction is not affected by whether $`X_2`$ is measured before measuring $`z_2`$, or $`X_2`$ and $`z_2`$ are jointly measured. Therefore, the EPR criterion is enough to guarantee that $`z_2`$ has a value $`v(z_2)`$, which does not change by measuring $`X_2`$. A similar reasoning applies to all 7 local observables involved in the proof.
The second step of the proof consists of showing the contradiction between the predictions of quantum mechanics and those of local realistic theories. For this purpose, note that state (1) also satisfies the following equations:
$`X_1x_1Y_2y_2|\psi `$ $`=`$ $`|\psi ,`$ (15)
$`Y_1x_1X_2y_2|\psi `$ $`=`$ $`|\psi .`$ (16)
To be consistent with Eqs. (8), (9), (15), and (16), local realistic theories predict the following relations between the values of the LERs:
$`v(X_1)`$ $`=`$ $`v(X_2)v(z_2),`$ (17)
$`v(Y_1)`$ $`=`$ $`v(Y_2)v(z_2),`$ (18)
$`v(X_1)v(x_1)`$ $`=`$ $`v(Y_2)v(y_2),`$ (19)
$`v(Y_1)v(x_1)`$ $`=`$ $`v(X_2)v(y_2),`$ (20)
respectively. However, it is impossible to assign the values $`1`$ or $`1`$ in a way consistent with all Eqs. (17)–(20). This can be proved as follows: In Eqs. (17)–(20) every value appears twice; therefore, the product of Eqs. (17)–(20) gives $`1=1`$. We therefore conclude that the 4 predictions of quantum mechanics given by Eqs. (8), (9), (15), and (16) cannot be reproduced by EPR LERs.
The remarkable property of this AVN proof is that the contradiction appears after considering only 4 quantum predictions, while the previous two-observer AVN proof Cabello01a ; Cabello01b ; CPZBZ03 ; YZZYZZCP05 required 9 quantum predictions. Moreover, while a local realistic model can reproduce 8 out of the 9 predictions of the previous two-observer AVN proof, it can be easily seen that it can reproduce only 3 out of the 4 predictions in the proof presented here. In more practical terms, this means that an experimental realization of the previous AVN proof YZZYZZCP05 requires photodetectors of a higher efficiency to avoid the detection loophole Pearle70 than an experimental realization of the proposed proof. To show this, we will estimate the detector efficiency required for both proofs. For this purpose, it is useful to see both proofs as games Vaidman99 .
Let us start by translating the new proof into a game in which a quantum-based strategy beats any classical strategy. Consider a team of two players, Alice and Bob, each of them isolated in a booth. Alice is asked one out of two possible questions: (I) “What are $`v(X_1)`$ and $`v(x_1)`$?” or (II) “What are $`v(Y_1)`$ and $`v(x_1)`$?” Bob is asked one out of 4 possible questions: (i) “What are $`v(X_2)`$ and $`v(y_2)`$?,” (ii) “What are $`v(X_2)`$ and $`v(z_2)`$?,” (iii) “What are $`v(Y_2)`$ and $`v(y_2)`$?,” or (iv) “What are $`v(Y_2)`$ and $`v(z_2)`$?” Each of them must give one of the following answers: “$`1`$ and $`1`$,” “$`1`$ and $`1`$,” “$`1`$ and $`1`$,” or “$`1`$ and $`1`$.” When Alice is asked (I), Bob is asked (ii) or (iii); when Alice is asked (II), Bob is asked (i) or (iv). Since $`v(X_2)`$ represents a LER, Bob’s answer to $`v(X_2)`$ must be independent on whether $`v(X_2)`$ is asked together with $`v(y_2)`$ or with $`v(z_2)`$. The same applies for Bob’s answers to $`v(Y_2)`$, $`v(y_2)`$, and $`v(z_2)`$. The team wins if their answers satisfy the corresponding equation of the set (17)–(20).
Assuming that the 4 possible combinations of questions are asked with the same frequency, no classical strategy allows the players to win in more than $`3/4`$ of the rounds. For instance, a simple optimal classical strategy is that each player always answers $`1`$ to any question. The hidden set of local instructions
$`\left\{\begin{array}{ccc}v(X_1)& v(X_2)& v(z_2)\\ v(Y_1)& v(Y_2)& v(y_2)\end{array}\right\},`$ (23)
where the vertical bar separates Alice’s instructions from Bob’s, is then
$`G:=\left\{\begin{array}{ccc}\hfill 1& \hfill 1& \hfill 1\\ \hfill 1& \hfill 1& \hfill 1\end{array}\right\}.`$ (26)
This strategy only fails whenever Bob is asked (iv) (i.e., in $`1/4`$ of the rounds). However, there is a quantum strategy that never fails: the players can win all the rounds if they share pairs of photons in the state (1), and give as answers the results of the corresponding measurements on their photons.
In a real experiment for testing Bell’s theorem, the low efficiency of detectors opens the possibility that non-detections could correspond to local instructions such as “if $`X`$ is measured, then the photon will not activate the detector.” This allows local realistic theories to simulate the observed results. To estimate the detection efficiency required to experimentally discard these theories, let us introduce a modification to the rules of the game. Let us allow each player to give no answer whatsoever in a fraction $`1\eta `$ of the rounds. If any of the players does not answer, then that round is not valid. This modification opens the possibility of the players also using a fraction of sets of local instructions like
$`B_1:=\left\{\begin{array}{ccc}\hfill 1& \hfill 1& \hfill 1\\ \hfill 0& \hfill 1& \hfill 1\end{array}\right\}`$ (29)
or
$`B_2:=\left\{\begin{array}{ccc}\hfill 1& \hfill 1& \hfill 1\\ \hfill 1& \hfill 0& \hfill 1\end{array}\right\},`$ (32)
where $`0`$ means that the corresponding player will not answer the corresponding question. For instance, if the players share $`B_1`$, Alice will not answer question (II); and if they share $`B_2`$, Bob will not answer questions (iii) and (iv).
Let us suppose that the players are using sets of predefined answers (or, equivalently, that the observed data can be described by a local realistic theory). For instance, sets like $`G`$ with a frequency $`1p`$, sets like $`B_1`$ with a frequency $`p/2`$, and sets like $`B_2`$ with a frequency $`p/2`$, where $`p`$ depends on the efficiency of the photodetector corresponding to photon $`j`$,
$$\eta _j=1p+\frac{p}{2}f_j+\frac{p}{2},$$
(33)
where $`f_1`$ ($`f_2`$) is the probability that Alice (Bob) answers \[i.e., she (he) does not get the instruction $`0`$ in her (his) set\] when they are using a $`B_1`$ ($`B_2`$) set. In our case $`f_j=1/2`$.
Let us calculate the minimum detection efficiency required to discard the possibility that the players are using this particular set of predefined answers (or, equivalently, that the observed data can be described by a local realistic theory). Then, to simulate the quantum probability of winning the game, the minimum value of $`p`$ is obtained by solving the equation
$$P_Q=(1p)P_G+\frac{p}{2}P_{B_1}+\frac{p}{2}P_{B_2},$$
(34)
where $`P_Q`$ is the quantum probability of winning the game, $`P_G`$ is the probability of winning the game when the players use a $`G`$ set, and $`P_{B_j}`$ is the probability of winning when the players use a $`B_j`$ set and both answer the questions. In our case, $`P_Q=1`$, $`P_G=3/4`$ and $`P_{B_j}=1`$. Introducing these values in Eqs. (33) and (34), we arrive at the conclusion that our local model can simulate the quantum predictions if $`\eta _j3/4=0.75`$. An exhaustive examination of all possible sets of local instructions shows that the previously presented model is indeed optimal and therefore we conclude that local realistic theories cannot simulate the quantum predictions if
$$\eta _j>3/4,$$
(35)
which is the same efficiency needed for a loophole-free experiment based on the three-observer version of GHZ’s proof Larsson98 . Indeed, the efficiency required for a two-observer AVN proof based on the state (1) can be lowered to $`\eta _j>11/160.69`$ if we consider 12 quantum predictions Cabello05 .
Let us compare these efficiencies with that required for a loophole-free experiment based on the two-photon version CPZBZ03 ; YZZYZZCP05 of the the two-observer AVN proof Cabello01a ; Cabello01b . To be consistent with 9 specific predictions of quantum mechanics (for details, see Cabello01b ), local realistic theories predict the following 9 relations between the values of the LERs:
$`v(Z_1)`$ $`=`$ $`v(Z_2),`$ (36)
$`v(z_1)`$ $`=`$ $`v(z_2),`$ (37)
$`v(X_1)`$ $`=`$ $`v(X_2),`$ (38)
$`v(x_1)`$ $`=`$ $`v(x_2),`$ (39)
$`v(Z_1z_1)`$ $`=`$ $`v(Z_2)v(z_2),`$ (40)
$`v(X_1x_1)`$ $`=`$ $`v(X_2)v(x_2),`$ (41)
$`v(Z_1)v(x_1)`$ $`=`$ $`v(Z_2x_2),`$ (42)
$`v(X_1)v(z_1)`$ $`=`$ $`v(X_2z_2),`$ (43)
$`v(Z_1z_1)v(X_1x_1)`$ $`=`$ $`v(Z_2x_2)v(X_2z_2).`$ (44)
It is impossible to assign the values $`1`$ or $`1`$ in a way consistent with all Eqs. (36)–(44Cabello01b . To calculate the detection efficiency required for a conclusive test based on this proof we will convert this impossibility into a game in which a quantum-based strategy beats any classical strategy. Consider again a team of two players, each of them isolated in a booth. Alice is asked one out of three possible questions: (I) “What are $`v(Z_1)`$ and $`v(x_1)`$?,” (II) “What are $`v(X_1)`$ and $`v(z_1)`$?,” or (III) “What are $`v(Z_1z_1)`$ and $`v(X_1x_1)`$?” Analogously, Bob is asked one out of three possible questions: (i) “What are $`v(Z_2)`$ and $`v(z_2)`$?,” (ii) “What are $`v(X_2)`$ and $`v(x_2)`$?,” or (iii) “What are $`v(Z_2x_2)`$ and $`v(X_2z_2)`$?” Each of them must give one of the following answers: “$`1`$ and $`1`$,” “$`1`$ and $`1`$,” “$`1`$ and $`1`$,” or “$`1`$ and $`1`$.” An interesting feature of this game is that it does not require a promise: all nine possible combinations of questions are legitimate. Assuming that all questions are asked with the same frequency, no classical strategy allows the players to win in more than $`8/9`$ of the rounds. However, the players can win all the rounds using a quantum strategy Cabello01b . An optimal classical strategy consists on using sets of instructions
$`\left\{\begin{array}{cccc}v(Z_1)& v(z_1)& v(Z_2)& v(z_2)\\ v(X_1)& v(x_1)& v(X_2)& v(x_2)\\ v(Z_1z_1)& v(X_1x_1)& v(Z_2x_2)& v(X_2z_2)\end{array}\right\}`$ (48)
like
$`G:=\left\{\begin{array}{cccc}\hfill 1& \hfill 1& \hfill 1& \hfill 1\\ \hfill 1& \hfill 1& \hfill 1& \hfill 1\\ \hfill 1& \hfill 1& \hfill 1& \hfill 1\end{array}\right\},`$ (52)
which gives correct answers except when Alice is asked (III) and Bob is asked (iii).
Now let us suppose that each player is also allowed to give no answer in a fraction $`1\eta `$ of the rounds. This is equivalent to assuming that they can use sets of local instructions like
$`B_1:=\left\{\begin{array}{cccc}\hfill 1& \hfill 1& \hfill 1& \hfill 1\\ \hfill 1& \hfill 1& \hfill 1& \hfill 1\\ \hfill 0& \hfill 0& \hfill 1& \hfill 1\end{array}\right\},`$ (56)
or
$`B_2:=\left\{\begin{array}{cccc}\hfill 1& \hfill 1& \hfill 1& \hfill 1\\ \hfill 1& \hfill 1& \hfill 1& \hfill 1\\ \hfill 1& \hfill 1& \hfill 0& \hfill 0\end{array}\right\}.`$ (60)
In this case, an optimal ensemble of sets of local instructions turns out to be $`G`$ sets with frequency $`1p`$, $`B_1`$ sets with frequency $`p/2`$, and $`B_2`$ sets with frequency $`p/2`$. Therefore, $`P_Q=1`$, $`P_G=8/9`$, $`P_{B_j}=1`$, and $`f_j=2/3`$, which leads us to conclude that, to avoid the detection loophole in the two-observer AVN proof in Cabello01a ; Cabello01b ; CPZBZ03 , we need
$$\eta _j>5/60.83,$$
(61)
which is higher than the value, given by inequality (35), required for the two-observer AVN proof introduced in this Letter.
AVN proofs follow directly from perfect correlations predicted by quantum mechanics. However, in a laboratory realization of the experiment, the observed correlations will not be as perfect as the proof requires. It is therefore convenient to derive Bell’s inequalities whose validity is necessary for the observed correlations to be consistent with a very general probabilistic local realistic theory, which are violated by the quantum predictions by an amount allowing significant room for the blurring effect of the imperfections in real experiments. The relevant features of the AVN proof derive from the fact that, for the state (1),
$$\psi |X_1X_2z_2Y_1Y_2z_2+X_1x_1Y_2y_2+Y_1x_1X_2y_2|\psi =4,$$
(62)
while, as can be easily checked, in any local realistic theory, this expected value must satisfy
$$|X_1X_2z_2Y_1Y_2z_2+X_1x_1Y_2y_2+Y_1x_1X_2y_2|2.$$
(63)
Inequality (63) resembles Mermin’s inequality for three observers Mermin90b , in the sense that quantum mechanics predicts a violation of 4, which is indeed the maximum possible violation of inequality (63), while the local realistic limit is 2. However, inequality (63) requires only two observers, as in the original Bell inequality Bell64 .
Summing up, we have introduced an AVN proof which combines the most interesting features appearing, separately, in previous AVN proofs: it is simple (the contradiction to EPR LERs follows from only 4 quantum predictions), provides a greater amount of evidence against local realism (only 3 out of these 4 predictions can be reproduced by LERs), and requires only two observers. In addition, we have shown that a conclusive experimental test of this proof would require lower efficiency detectors than those needed for the previous two-observer AVN proof, and we have derived the corresponding Bell inequality for real experiments, which resembles Mermin’s inequality for three observers but requires only two observers.
The author thanks A. Broadbent, E. Galvão, A. Lamas-Linares, J.-Å. Larsson, C.-Y. Lu, S. Lu, E. Santos, and H. Weinfurter for useful comments, and acknowledges support by Projects No. BFM2002-02815 and No. FQM-239. |
warning/0507/astro-ph0507603.html | ar5iv | text | # Contribution of HD molecules in cooling of the primordial gas
## 1 Introduction
Formation of stars is intimately connected with the ability of gas to cool. In the metal-free primordial medium the radiative cooling is mainly provided by molecular hydrogen and its isotope analogue HD. In the expanding universe H<sub>2</sub> and HD form after the epoch of recombination \[1-5\]. Due to a non-zero dipole moment and lower excitation energy HD molecules can efficiently cool gas at temeperature $`T200`$ K, where the rate of H<sub>2</sub> cooling decreases sharply. If HD abundance is small, cooling of primordial gas stops practically at $`T200`$ K.
Thus, thermal evolution of the primordial gas at low temperature, and as a consequence characteristics of the first stars in the universe critically depend on the abundance of HD molecules. At present, an analysis of conditions when HD can be thermodynamically important in the primordial gas is absent. Conclusions about the role of HD molecules in a prestellar universe are contradictory. In particular, it is shown in that HD cooling never dominates in a collapsing spherical cloud. At the same time the authors \[6-8\] point out that HD cooling can be important in primordial clouds, although their calculations are restricted only by initial stages of the collapse. It is obvious though, that the abundance of HD crucially depends on thermal evolution of gas in the temperature range $`T>500`$ K. The apparent contradiction about the role of HD is connected with differences in the initial conditions used in from one side and in \[6-8\] from the other. In addition, the H<sub>2</sub> cooling function adopted in is overestimated. Therefore, in order to make firm conclusions about the role of HD molecules in a prestellar universe it is necessary to investigate the efficiency of HD formation in a wider range of initial conditions.
It is known that formation of H<sub>2</sub> and HD molecules is largely boosted behind shock waves \[10-13\]. It is connected mostly with the fact that temperature and fractional ionization behind the shock waves increase, and as a consequence the rates of molecular reactions are enhanced. In a postshock gas cooled down to temperature $`10^4`$ the fraction of electrons remains sufficiently high, $`x\stackrel{>}{}0.001`$, which favours rapid formation of H<sub>2</sub> molecules in the catalystic reactions with H<sup>-</sup> ions. At such conditions H<sub>2</sub> fraction can reach $`10^2`$. Further cooling is mainly provided by H<sub>2</sub> molecules efficient in the temperature range $`2007000`$ K. When lower temperatures are reached, $`T200`$ K, deuterium begins to convert into HD molecules due to chemical fractionation . The contribution of HD cooling in energy losses increases when temperature decreases, and if it becomes dominant the gas temperature can fall down to several tens of degrees. One can expect that at least in a restricted range of shock parameters formation of HD molecules is as efficient as to provide such a predominance. We aim to study this possibility.
In Section 2 we describe a thermochemical model of a shocked gas; the results are presented in Section 3; in Section 4 discussion and in Section 5 summary are given. Throughout the paper we assume a $`\mathrm{\Lambda }`$CDM cosmology with the parameters $`(\mathrm{\Omega }_0,\mathrm{\Omega }_\mathrm{\Lambda },\mathrm{\Omega }_m,\mathrm{\Omega }_b,h)=(1.0,0.71,0.29,0.047,0.72)`$ and deuterium abundance $`2.62\times 10^5`$ .
## 2 Molecular kinetics behind shock waves
In the absence of thermal conductivity and diffusion the thermochemical evolution of gas behind a shock wave can be described by a system of ordinary differential equations for a Lagrangian fluid element, which include equations of chemical kinetics
$$\dot{x}_i=F(x_i,T,n)D(x_i,T,n),$$
(1)
and energy equation
$$\dot{T}=\frac{2}{3}\underset{i}{}[\mathrm{\Gamma }_i(x_i,T,n)\mathrm{\Lambda }_i(x_i,T,n)]+\frac{2}{3}\frac{T}{n}\dot{n},$$
(2)
where $`x_i`$ is the fraction of $`i`$-th species, $`F_i(𝐱,T,n)`$, $`D_i(𝐱,T,n)`$ are the corresponding formation and destruction rates, $`\mathrm{\Gamma }_i(𝐱,T,n)`$, $`\mathrm{\Lambda }_i(𝐱,T,n)`$ are the heating and cooling rates. Chemical kinetics of primordial gas include the following main species: H, H<sup>+</sup>, H<sup>-</sup>, He, He<sup>+</sup>, He<sup>++</sup>, H<sub>2</sub>, H$`{}_{}{}^{+}{}_{2}{}^{}`$, D, D<sup>+</sup>, D<sup>-</sup>, HD, HD<sup>+</sup>, $`e`$. The rates for collisional and radiative processes are taken from , and for D<sup>-</sup> ion from . The energy equation accounts cooling processes connected with H, He, He<sup>+</sup>, He<sup>++</sup>, such as collisional excitation and ionization, bremsstrahlung radiation, recombination, dielectronic recombination, molecular cooling by H<sub>2</sub> and HD, and Compton cooling. In the absence of external ionizing radiation the abundances of chemical species and radiative cooling are determined by collisions. Thus, the right-hand side of (1) and first term of the right-hand side of (2) are proportional to the gas density, and for convenience we can introduce the fluence $`\eta `$ through
$$d\eta =ndt,$$
(3)
where $`t`$ is time, $`n`$ is number density; below the results are presented as functions of the fluence.
For the H<sub>2</sub> cooling functions we adopted the expressions given in , the HD cooling is taken from , the other rates are from . In addition, the effects of interaction between molecules and the CMB photons are accounted which imply that when gas temperature is close to the CMB temperature, H<sub>2</sub> and HD molecules are populated by the CMB photons, and then collisionally transfer the energy to kinetic temperature, thus heating the gas. Therefore, gas cannot cool lower than the CMB temperature. We assume that the gas temperature behind the front jumps and reach the value
$$T_0=\alpha ^2\frac{m_pv_c^2}{k}1.2\times 10^2\alpha ^2\left(\frac{v_c}{1\mathrm{k}\mathrm{m}\mathrm{c}^1}\right)^2,$$
(4)
where $`\alpha ^2=3/16`$ is for a shock propagating in a static gas, and $`\alpha ^2=1/3`$ is for a shock wave formed in a head-on collision of two flows (clouds) with equal velocities $`v_c`$ . Neglecting thermal conduction the evolution of each Lagrangian volume of gas behind the shock is isobaric , so that the density is described by
$$n=\frac{p}{\mu kT},$$
(5)
where $`\mu `$ is the average molecular weight.
## 3 Formation of HD behind shock waves
In the contemporary scenarious of structure formation in the universe the first protogalaxies form at the epoch $`z=1030`$. For specifity, we consider thermochemical evolution of baryons behind a shock at the redshift $`z=20`$. Formation of dark haloes (future protogalaxies) and their subsequent virialization are accompanied by shock waves in the baryon component. The duration of this process is close to the Hubble time $`t_H(z)`$ , so we restrict computations by $`t_H(z)`$, which means the ending redshift $`z_e12`$ for the initial $`z_i=20`$. One can expect that collisions of baryonic flows during the virialization of dark haloes result in considerable density variations. In order to study possible influence of such density variations on thermochemical evolution we conduct the calculations for a wide range of density in colliding flows: from the background to the virial value. As a characteristic value we adopt the virial density, $`18\pi ^2n_b(1+z)^3`$ (see e.g. ), where $`n_b`$ is the background baryon density today, and consider the dependence of the thermochemical evolution on the initial gas density. Gas before the shock is assumed to be cold compared to the gas just after the front, so we consider strong shock waves.
HD molecules form efficiently at low teperature in the presence of sufficient fraction of molecular hydrogen through the reaction
$$\mathrm{D}^++\mathrm{H}_2\mathrm{HD}+\mathrm{H}^+.$$
(6)
For this reason we briefly discuss formation and destruction of H<sub>2</sub> molecules behind shock fronts .
In the primordial gas H<sub>2</sub> forms in interactions of neutral hydrogen with H<sup>-</sup> and H$`{}_{}{}^{+}{}_{2}{}^{}`$ ions efficiently born at high temperature. As shock waves significantly increase temperature it enhances formation of H<sub>2</sub> and increases its abundunce . It is known that HI cooling becomes inefficient at temperature $`\stackrel{<}{}10^4`$ K, and in the primordial gas only the molecular hydrogen can provide further cooling to lower temperature. One can estimate a minimum fraction of H<sub>2</sub> needed for effective cooling: the cooling time must be shorter then the Hubble time which condition fulfills when H<sub>2</sub> fraction becomes greater than the critical value $`x_{\mathrm{H}_2}=5\times 10^4`$ ; an increase of H<sub>2</sub> abundance shortens the cooling time. In a shocked gas at temperature $`\stackrel{>}{}8\times 10^3`$ K collisions increase fractional ionization, and this further enhances formation of H<sup>-</sup>, H$`{}_{}{}^{+}{}_{2}{}^{}`$ and H<sub>2</sub>. Figures 1-3 show the evolutionary paths of the thermochemical state of a gas element behind the shock: the evolution begins at high temperature and follows a monotonous cooling accompanied by grownig H<sub>2</sub> and HD fractions. It is seen that at high collisional velocities the electron fraction significantly increases at initial stages, which stimulates formation of H<sup>-</sup> ions and H<sub>2</sub> molecules and a strong decrease of temperature. The H<sub>2</sub> abundance grows rapidly in the temperature range $`T10^310^4`$ K, while at lower temperatures formation of H<sub>2</sub> practically exhausts . Already formed H<sub>2</sub> molecules provide further cooling down to $`T200400`$ K depending on the exact value of their abundance. If gas temperature falls down to $`T\stackrel{<}{}150`$ K deuterium rapidly converts to HD molecules due to chemical fractionation (Fig.4.).
Fig. 5 presents the dependence of electron, H<sub>2</sub> and HD fractions on the fluence $`\eta =n𝑑t`$ for several values of the shock velocity. Note that for the shock waves with $`v_s3.5\alpha ^1`$ km s<sup>-1</sup> the final abundance of H<sub>2</sub> is greater than the limit $`5\times 10^4`$, and as a cosequence the gas behind the shock can lose significant fraction of its thermal energy in one Hubble time. Since larger shock velocity corresponds to higher temperature, the maximum density behind the shock increases with the velocity $`\rho v_s^2`$ (eqs. 4 and 5). In Fig.5 this corresponds to larger fluence. It is obvious, that the characteristic time of the thermal evolution decreases as $`v_s^2`$. For $`v_s4.6\alpha ^1`$ km s<sup>-1</sup> H<sub>2</sub> fraction equals $`7\times 10^4`$ and becomes sufficient for cooling down to $`T130`$ K; HD fraction at these conditions is $`4\times 10^7`$. Further increase of the shock velocity results in an increase of molecular fractions and a decrease of temperature on a shorter time, with the dominant cooling provided by HD molecules (Fig.6). Equating the cooling rates from H<sub>2</sub> and HD one can estimate the critical temperature below which the contribution of HD into thermodynamics becomes dominant. For the adopted cooling functions it occurs at $`T_{\mathrm{cr}}130`$ K. For the shock wave velocity $`4.6\alpha ^1`$ km s<sup>-1</sup> gas cools down to this limit, and for larger velocities temperature falls below $`T_{\mathrm{cr}}`$ where the HD cooling dominates.
In the velocity range $`4.6\alpha ^1\stackrel{<}{}v_s\stackrel{<}{}8.7\alpha ^1`$ km s<sup>-1</sup> the initial temperature behind the shock is insufficient for collisional ionization, and in the subsequent evolution the electron fraction can only decrease (Figs.1 and 5). At such conditions H<sub>2</sub> and HD abundances increase with velocity only because the reaction rates grow with temperature. As seen in Fig.2 for shock velocities in this range the maximum fraction of H<sub>2</sub> is approximately equal to $`x_{\mathrm{H}_2}8\times 10^4`$. This is enough for an efficient formation of HD molecules and successful cooling down to $`T130`$ K. At $`v\stackrel{>}{}7\alpha ^1`$ km s<sup>-1</sup> the HD fraction becomes sufficient for cooling down to the CMB temperature $`T_{\mathrm{CMB}}=2.7(1+z)`$: due to the strong emission in rotational lines of HD gas temperature falls to several tens, $`30`$ K, approaching $`T_{\mathrm{CMB}}(z)`$ at a given redshift (Fig.5). This is because HD molecules provide an efficient exchange of energy between the CMB and baryons through absorption of CMB photons and subsequent collisional de-exitation . It is obvios that similar picture is valid at all redshift, and the final temperature of cold baryons is $`T_{\mathrm{min}}T_{\mathrm{CMB}}`$.
For velocities $`v_s\stackrel{>}{}9.2\alpha ^1`$ km s<sup>-1</sup> gas temperature behind the shock becomes greater than $`10^4`$ K, which results in an increase of fractional ionization immediately after the gas element crosses the shock front: for $`v_s=10.4\alpha ^1`$ km s<sup>-1</sup> $`x_e`$ increases by factor of 2, and for $`v_s\stackrel{>}{}11.6\alpha ^1`$ km s<sup>-1</sup> more than an order (Fig.1 and 5). At such conditions chemical kinetics changes qualitatively – the enhancement of H<sub>2</sub> formation is caused in this case by the two factors: increasing reaction rates and a higher ionization fraction, resulting in more frequent catalystic processes H$`+e`$H<sup>-1</sup>, H$`{}_{}{}^{1}+`$H$``$H$`{}_{2}{}^{}+e`$. The corresponding evolution of $`x_e(t)`$, $`x_{\mathrm{H}_2}(t)`$ and $`x_{\mathrm{HD}}(t)`$ looks qualitively different as seen in Fig.5: while for $`v_s<9.2\alpha ^1`$ km s<sup>-1</sup> an increase of velocity by 1 km s<sup>-1</sup> produces insignificant changes in $`x_{\mathrm{H}_2}(t)`$ and $`x_{\mathrm{HD}}(t)`$, for greater velocities such an increase results in a considerable (half order of magnutude) increase of $`x_{\mathrm{H}_2}(t)`$ and $`x_{\mathrm{HD}}(t)`$. Thus, $`x_{\mathrm{H}_2}(t)`$ and $`x_{\mathrm{HD}}(t)`$ fractions at equal temepratures are higher for larger velocities (Figs. 2 and 3). For lower velocities the evolutionary paths $`x_{\mathrm{H}_2}(T)`$ and $`x_{\mathrm{HD}}(T)`$ for different $`v_s`$ practically coincide at temperature $`T10^3`$ K. One should stress, that in the considered range of velocities H<sub>2</sub> molecules form primarely through H<sup>-</sup> ions, the contribution from H$`{}_{}{}^{+}{}_{2}{}^{}`$ ions is as a rule negligible.
Thus, for velocities $`v4.6\alpha ^1`$ km s<sup>-1</sup> HD molecules behind shock waves provide lower temperatures than H<sub>2</sub> can do. It is worth noting that at $`z45`$ the CMB temperature is higher than the critical value $`T_{\mathrm{cr}}`$ at which HD cooling dominates. Therefore, under these conditions HD molecules can only heat gas, and at larger redshifts become unimportant.
Everywhere above we adopted gas density in colliding flows equal to the virial value at the corresponding redshift. However, one can assume that in the process of merging of haloes a fraction of baryonic mass can be lost. In the intergalactic medium such ”separated” baryonic flows can greatly expand, and their final density depends on collision velocity $`v_c`$, masses of the merging haloes, details of separation and so on. Subsequently such baryonic clumps can collide with gaseous components of other haloes or with each other. In these conditions the thermal evolution differs from that of denser baryonic flows. Let us consider how the thermochemical evolution depends on the density. Fig.7 shows the temperature and the HD fraction versus the density for several shock velocities at two redshifts. It is clearly seen that at $`z=20`$ and in the low velocity range $`v5.8\alpha ^1`$ km s<sup>-1</sup> only for densities close to the virial value gas temperature drops below the critical value where contribution from HD dominates. However, for higher velocities HD cooling remains efficient even for densities of one order of magnitude lower than the virial value. At $`v5.8\alpha ^1`$ km s<sup>-1</sup> and the density close to the virial value only $`0.25`$ of deuterium converts in HD, however it becomes sufficient to cool the gas down to the CMB temperature. For higher velocities this can occur for several times lower densities than the virial value. Collisions with higher velocities $`v11.6\alpha ^1`$ km s<sup>-1</sup> change H<sub>2</sub> kinetics: formation of H<sub>2</sub> becomes more efficient due to significant increase of fractional ionization behind the front, and behaves similar to the collisions with virial density at $`v_s\stackrel{>}{}10.4\alpha ^1`$ km s<sup>-1</sup> shown in Figs. 1 and 5. These features are seen in Fig.7. In general, one can conclude that HD molecules can also play a significant role in cooling of baryonic flows of low density.
## 4 Discussion: formation of protostellar fragments
Birth of stars is always accompanied with shock waves. This is unconditionally true for the first stars in the universe. The origin of shock waves can be connected both with the formation of the first protogalaxies in merging flows, and with supernovae explosions in already formed galaxies. Formation of the first protogalaxies implies separation of highly overcritical density perturbations from the Hubble expansion and a predominantly one-dimensional compression . This is a source of shock waves in baryonic component. Similar processes take place in the hierarchical scenario of structure formation, where massive objects form in collisions and following mergings of less massive haloes. These processes can be treated as collisions of the gas and dark matter flows. A collisionless dark matter reaches the virial state apparently through the violent relaxation, while in gas component shock waves form. The shock wave velocity depends on mass of a forming protogalaxy: the velocity amplitude in a perturbation of mass $`M`$ is close to the value
$$v_c=\sqrt{3}\sigma ,$$
(7)
where $`\sigma `$ is the one-dimensional velocity dispersion
$$\sigma ^2=\frac{GM}{2R}.$$
(8)
Thus, the parameters of the shock waves are determined by the total mass of matter involved in motion, by the redshift at which the object forms, and so on. The efficiency of HD formation is sensitive to these parameters. Therefore, one can expect that the characteristcs of stellar population vary in galaxies of different mass. As the maximum abundance of HD depends on the shock velocity, and HD molecules cool the gas to much lower temperature than H<sub>2</sub> molecules do, a typical mass of protostellar molecular clouds is expected to decrease with increasing mass of a forming galaxy .
Shock waves in the epoch of galaxy formation can be connected with explosions of the first supernovae. In these events much more powerful shock waves form: typical velocities can be greater $`\stackrel{>}{}100`$ km s<sup>-1</sup>, the corresponding temperature behind the front is $`\stackrel{>}{}2.8\times 10^5`$ K. At radiative stages when the gas temperature reaches $`\stackrel{<}{}10^4`$ K, the H<sub>2</sub> fraction becomes sufficiently high $`\stackrel{>}{}5\times 10^3`$ due to high ionization fraction at the preceding stages. As a consequence, the gas temperature definitely falls to the lower values at which the cooling is essentially determined by HD molecules. In these conditions fragmentation of an expanding shell can occur .
Due to isobaric compression the gas density behind the shock increases considerably compared to the initial value. For instance, for the velocity $`v_c7\alpha ^1`$ km s<sup>-1</sup> the initial temperature is $`T5.8\times 10^3`$ K (4), and when the gas cools down to $`T_{\mathrm{CMB}}2.7(1+z)`$, its density increases more than 200 times. At such conditions fragmentation and formation of stars become possible behind the shock . Gravitationally unstable fragments can give rise to protostars or protostellar clusters. It follows that when cooling is determined by HD molecules, formation of low mass protostellar clouds becomes possible . Indee, under these conditions the gas temperature falls down to $`2.7(1+z)`$ K which at $`z=20`$ is 4 times smaller than can be provided only by H<sub>2</sub> molecules. The Jeans mass $`M_J30T^{3/2}n^{1/2}M_{}`$ behind the front is $`M_J15T^2n_0^{1/2}T_0^{1/2}M_{}`$, where $`n_0`$, the gas density in a flow (a cloud) before collision, $`T_0`$ is the temperature at the shock; here density at the front is taken 4$`n_0`$ as for strong shock waves. If we assume that gas density in flows (clouds) before collision is equal to the virial value (see e.g. ), $`n_0=18\pi ^2n_b(1+z)^3`$, then for velocities $`v_c7\alpha ^1`$ km s<sup>-1</sup> (or initial temperature $`T_05.8\times 10^3`$ K) the Jeans mass is
$`M_J\stackrel{<}{}2.4\times 10^5M_{}\left({\displaystyle \frac{1+z}{T_0\delta _c}}\right)^{0.5}=7.2\times 10^3M_{}\left({\displaystyle \frac{\alpha v}{1\mathrm{km}\mathrm{c}^1}}\right)^1\left({\displaystyle \frac{\delta _c}{18\pi ^2}}\right)^{0.5}\left({\displaystyle \frac{1+z}{20}}\right)^{0.5}`$ (9)
where $`\delta _c`$ is the ratio of gas density before collision to the background baryonic density. At the same time, when cooling is determined only by H<sub>2</sub> molecules the typical gas temperature is of $`200`$ K, and the corresponding Jeans mass $`13.5[(1+z)/20]^2`$ times exceeds the value given by (9)
$`M_J\stackrel{<}{}1.3\times 10^9M_{}\left({\displaystyle \frac{1}{T_0\delta _c(1+z)^3}}\right)^{0.5}=10^5M_{}\left({\displaystyle \frac{\alpha v}{1\mathrm{km}\mathrm{c}^1}}\right)^1\left({\displaystyle \frac{\delta _c}{18\pi ^2}}\right)^{0.5}\left({\displaystyle \frac{1+z}{20}}\right)^{1.5}`$ (10)
In other words, the question of how massive are the fragments formed behind shock fronts depends on whether the cooling is determined by H<sub>2</sub> or HD molecules. The density of the fragments is $`10300`$ cm<sup>-3</sup> depending on the redshift and the initial temperature. Subsequent collapse is isothermal until the fragment becomes opaque in H<sub>2</sub> and HD lines, which takes place at the density $`10^910^{10}`$ cm <sup>-3</sup>. At this stage, if the cooling is dominated by HD molecules, the Jeans mass is $`M_J30T_{\mathrm{CMB}}^{3/2}n^{1/2}M_{}10^3(1+z)^{3/2}M_{}`$ , however when HD is underabundant the Jeans mass can be 2-3 orders greater . Further evolution is determined by the accretion of gas onto the central core . If the accretion rate is below the Eddigton limit, the mass of a forming star is comparable to the initial mass of a protostellar cloud, in the opposite case it can be much lower . Thus, one can expect that a typical mass of the first stars born in protogalaxies of higher masses (corresponding to higher collisional velocities) is shifted towards the lower end due to cooling by HD molecules.
It is readily seen that since the overall thermal evolution of low density flows differs from that of higher densities, the final value of the Jeans mass and its dependence on redshift will differ from the above value. Fig.7 shows the Jeans mass versus the gas density in the flow. In the range of collisional velocities $`5.8\alpha ^1v_c<8.6\alpha ^1`$ km s<sup>-1</sup> only flows with the initial density very close to the virial can have Jeans mass $`M_J`$ smaller than $`10^4M_{}`$, which may correspond to the mass of a protostellar cloud. However for $`v8.6\alpha ^1`$ km s<sup>-1</sup> the Jeans mass becomes $`10^4M_{}`$ for the initial density 4 times lower than the virial. As mentioned above, for higher velocities gas cools down to the CMB temperature, and the Jeans mass is $`10^3M_{}`$, what is seen also from (9) – on Fig.7 flat parts of the lines for the velocity $`8.6\alpha ^1`$ km s<sup>-1</sup> reflects this circumstance. Thus, the Jeans mass is considerably higher than (9) only for flows with a low collisional velocity and a low density. High-velocity collisions $`v_c8.6\alpha ^1`$ km s<sup>-1</sup> provide cooling down to the temperature $`TT_{\mathrm{CMB}}`$ even for low density flows.
Let us consider collision of flows whose density is equal to the background value, i.e. $`\delta _c1`$, $`\rho /\rho _{\mathrm{vir}}6\times 10^3`$. It is seen from Fig.7 that even for high-velocity collisions $`v_c=11.6\alpha ^1`$ km s<sup>-1</sup> gas cannot cool sufficiently in one Hubble time, and the Jeans mass is quite high: $`10^610^7M_{}`$. Moreover, the free-fall time for such low densities is greater than the comoving Hubble time. Under these conditions baryonic objects cannot be formed. However, further increase of the collisional velocity, $`v_c>11.6\alpha ^1`$ km s<sup>-1</sup>, makes the gas behind the shock able to cool down to the temperature $`1000`$ K during one Hubble time. For instance, for low density flows collided with the velocity $`v_c19.2\alpha ^1`$ km s<sup>-1</sup> the final temperature behind the shock is $`200`$ K, and can reach lower values for higher velocities. Under these conditions HD molecules will dominate in gas colling, and the Jeans mass becomes as small as $`7\times 10^4M_{}`$.
## 5 Conclusions
The influence of HD molecules on the thermochemical evolution of the primordial gas behind shock waves possibly formed during the epoch of galaxy formation has been studied.
1. We showed that deuterium converts efficiently to HD molecules and the contribution of HD to coolig becomes dominant for the shock waves with velocities $`\stackrel{>}{}4.6\alpha ^1`$ km s<sup>-1</sup> ($`\alpha 0.5`$). Behind such shock waves the conditions are favourable for fragmentation and, as a consequence, for formation of protostellar clusters.
2. For shock velocities $`\stackrel{>}{}7\alpha ^1`$ km s<sup>-1</sup> gas is able to cool down to the CMB temperature. Under these conditions Jeans mass depends only on the redshift and the initial density: $`M_J\stackrel{<}{}2.4\times 10^5M_{}(1+z)^{0.5}(T_0\delta _c)^{0.5}`$, for virial haloes ($`\delta _c=18\pi ^2`$) at $`z=20`$ this corresponds to $`M_J\stackrel{<}{}10^3M_{}`$.
3. At $`z\stackrel{>}{}45`$ the CMB temperature is close to the critical value $`T_{\mathrm{cr}}`$, at which the contribution from HD molecules to the total cooling is comparable to that from H<sub>2</sub>. At these conditions HD molecules begin to heat gas, and at higher redshifts become unimportant in thermal history of baryons.
4. For densities of colliding flows smaller than the virial value the efficiency of HD molecule formation decreases. In particular, at $`z=20`$ gas temperature behind the shock with $`v5.8\alpha ^1`$ km s<sup>-1</sup> drops substantially only for a density close to the virial value. However, for the shock velocities $`8.6\alpha ^1`$ km s<sup>-1</sup> HD molecules are important in cooling for densities of 2-3 times lower than the virial value. For $`v11.6\alpha ^1`$ km s<sup>-1</sup> HD cooling is effective for densities close to the background value, and almost all deuterium converts to HD for collisions of less dense than if the flows were virial. For the gas density equal to the background value and for the velocity $`v\stackrel{>}{}19.2\alpha ^1`$ km s<sup>-1</sup> temperature drops to $`200`$ K and HD molecules begin to dominate radiative cooling.
Note added in manuscript 2005 July 26.- After acceptance of this paper we have been informed about the paper by A. Lipovka, R. Núñez-López, and V. Avila-Reese (MNRAS, 2005, in press, astro-ph/0503682), where new calculations of the HD cooling function are reported. In the temperature range of interest ($`T<10^3`$ K) this function coincides with that given in , while at higher temperatures, where Lipovka et al predict an order of magnitude enhanced HD cooling rate, the abundance of HD is too low to contribute. |
warning/0507/math0507164.html | ar5iv | text | # Extensions of the linear bound in the Füredi–Hajnal conjecture
## 1 Introduction
Füredi and Hajnal asked in whether for every fixed 0-1 permutation matrix $`P`$ the maximum number of 1’s in an $`n\times n`$ 0-1 matrix $`M`$ avoiding $`P`$ is $`O(n)`$; the avoidance here means that $`P`$ cannot be obtained from $`M`$ by a series of row deletions, column deletions, and replacements of $`1`$’s with $`0`$’s (in particular, permuting rows or columns of $`M`$ is not allowed). The Füredi–Hajnal conjecture was settled by Marcus and Tardos in where they proved that if $`M`$ avoids a $`k\times k`$ permutation matrix, then the number of 1’s in $`M`$ is at most $`2k^4\left(\genfrac{}{}{0pt}{}{k^2}{k}\right)n`$. In this paper we present extensions of their linear bound to more general structures.
The Marcus–Tardos bound can be reformulated in the language of graph theory, since matrices with entries 0 and 1 can be viewed as the incidence matrices of bipartite graphs. Thus if $`P=([2k],E(P))`$ is a graph on the vertex set $`[2k]=\{1,2,\mathrm{},2k\}`$ with $`k`$ mutually disjoint edges, each of which connects the sets $`[k]`$ and $`[k+1,2k]=\{k+1,k+2,\mathrm{},2k\}`$, and $`M=([2n],E(M))`$ is a graph on $`[2n]`$ which only has edges connecting $`[n]`$ and $`[n+1,2n]`$ and does not contain $`P`$ as an ordered subgraph, then $`M`$ has only linearly many edges, i.e. $`|E(M)|=O(n)`$. It is easy to modify the proof in so that it gives a linear bound for all $`P`$-avoiding graphs $`G`$ (not necessarily bipartite) on the vertex set $`[2n]`$ (and therefore $`[n]`$). In Section 2 we extend this bound further to hypergraphs with edges of arbitrary size. We also discuss exponential enumerative bounds which follow from the linear extremal bounds as corollaries.
In yet another light, 0-1 matrices can be viewed as the (characteristic matrices of) binary relations. In Section 3 we generalize the original proof of Marcus and Tardos to $`d`$ dimensions and show that every $`d`$-ary relation on $`[n]`$ which avoids a fixed $`d`$-dimensional permutation has at most $`O(n^{d1})`$ elements.
## 2 Extensions to hypergraphs
For a graph $`G^{}=([k],E^{})`$, we define $`\mathrm{gex}_<(n,G^{})`$ to be the maximum number $`|E|`$ of edges in a graph $`G=([n],E)`$ that does not contain $`G^{}`$ as an ordered subgraph. We represent a permutation $`\pi =a_1a_2\mathrm{}a_k`$ of $`[k]`$ by the graph
$$P(\pi )=([2k],\{\{i,k+a_i\}:i[k]\}).$$
As we mentioned in Section 1, it is easy to modify the proof in to obtain the bound
$$\mathrm{gex}_<(n,P(\pi ))=O(n)$$
(1)
where the constant in $`O`$ depends only on $`\pi `$.
For the hypergraph extension we need a few more definitions. A hypergraph is a finite collection $`H=(E_i:iI)`$ of finite nonempty edges $`E_i`$ which are subsets of $`𝐍=\{1,2,\mathrm{}\}`$. The vertex set is $`V(H)=_{iI}E_i`$. For simplicity we do not allow (unlike in the graph case) isolated vertices; for our extremal problems this restriction is immaterial (isolated vertices in graphs can be represented by singleton edges in our extension). In general we will allow multiple edges, and will denote a hypergraph as simple if it has no multiple edges. We say that $`H^{}=(E_i^{}:iI^{})`$ is contained in $`H=(E_i:iI)`$, written $`H^{}H`$, if there exists an increasing injection $`f:V(H^{})V(H)`$ and an injection $`g:I^{}I`$ such that $`f(E_i^{})E_{g(i)}`$ for every $`iI^{}`$; otherwise we say that $`H`$ avoids $`H^{}`$. To put it differently, $`H^{}H`$ means that $`H^{}`$ can be obtained from $`H`$ by deleting some edges, deleting vertices from the remaining edges, and relabeling the vertices so that their ordering is preserved. This containment generalizes the ordered subgraph relation. Note that a simple hypergraph can contain a non-simple hypergraph.
The order $`v(H)`$ of $`H`$ is the number of vertices $`v(H)=|V(H)|`$, the size $`e(H)`$ is the number of edges $`e(H)=|I|`$, and the weight $`i(H)`$ is the number of incidences $`i(H)=_{iI}|E_i|`$. We define two hypergraph extremal functions.
###### Definition.
Let $`F`$ be any hypergraph. We associate with $`F`$ the functions $`\mathrm{ex}_e(,F),\mathrm{ex}_i(,F):𝐍𝐍`$,
$`\mathrm{ex}_e(n,F)`$ $`=`$ $`\mathrm{max}\{e(H):HF\&H\text{ is simple}\&v(H)n\}`$
$`\mathrm{ex}_i(n,F)`$ $`=`$ $`\mathrm{max}\{i(H):HF\&H\text{ is simple}\&v(H)n\}.`$
Obviously, $`\mathrm{ex}_e(n,F)\mathrm{ex}_i(n,F)`$ for every $`F`$ and $`n`$. If $`F`$ has at least two edges and has no two separated edges (edges $`E_1`$ and $`E_2`$ satisfying $`E_1<E_2`$), Klazar’s Theorem 2.3 in gives an inequality in the opposite direction:
$$\mathrm{ex}_i(n,F)(2v(F)1)(e(F)1)\mathrm{ex}_e(n,F).$$
So in particular, for every permutation $`\pi `$ of $`[k]`$,
$$\mathrm{ex}_i(n,P(\pi ))(4k1)(k1)\mathrm{ex}_e(n,P(\pi )).$$
(2)
Thus a linear bound on $`\mathrm{ex}_i(n,P(\pi ))`$ follows directly from one on $`\mathrm{ex}_e(n,P(\pi ))`$.
The latter bound can be derived using the techniques in along with the graph bound in (1). To explain the reduction we need the notion of the blow-up of a graph. A graph $`G^{}`$ is an $`m`$-blow-up of a graph $`G`$ if for every edge coloring of $`G^{}`$ by colors from $`𝐍`$ such that every color is used at most $`m`$ times, there exists a subgraph of $`G^{}`$ that is order-isomorphic to $`G`$ and no two of its edges have the same color. Let $`G`$ be a graph with $`k`$ vertices, $`H`$ be a $`\left(\genfrac{}{}{0pt}{}{k}{2}\right)`$-blow-up of $`G`$, and $`f:𝐍𝐍`$ be a function such that $`\mathrm{gex}_<(n,H)<nf(n)`$ for every $`n𝐍`$. Then Theorem 3.1 in states that, for every $`n𝐍`$,
$$\mathrm{ex}_e(n,G)<e(G)\mathrm{gex}_<(n,G)\mathrm{ex}_e(2f(n)+1,G).$$
(3)
Combining the bounds in (1), (2), and (3) we obtain the following result:
###### Theorem 2.1.
For every permutation $`\pi `$,
$$\mathrm{ex}_i(n,P(\pi ))=O(n).$$
###### Proof.
For $`m𝐍`$ and a $`k`$-permutation $`\pi `$, consider a permutation graph $`P(\pi ^{})`$ that arises from $`P(\pi )`$ by replacing every edge in $`P(\pi )`$ with a bundle of $`k(m1)+1`$ edges so that the initial vertices of the edges in each bundle form an interval in $`P(\pi ^{})`$ and the same holds for the final vertices. The positions of the intervals are as in $`P(\pi )`$, that is, for every selection of one edge from each bundle the resulting graph is order-isomorphic to $`P(\pi )`$. There are many such graphs $`P(\pi ^{})`$ ($`\pi ^{}`$ is a $`(k^2(m1)+k)`$-permutation) and each of them is, by the pigeonhole principle, an $`m`$-blow-up of $`P(\pi )`$.
We set $`m=\left(\genfrac{}{}{0pt}{}{2k}{2}\right)`$. By the graph bound in (1), there are constants $`c_\pi `$ and $`c_\pi ^{}`$ such that
$$\mathrm{gex}_<(n,P(\pi ))<c_\pi n\text{ and }\mathrm{gex}(n,P(\pi ^{}))<c_\pi ^{}n$$
for every $`n`$. We set $`H=P(\pi ^{})`$ and $`f(n)=c_\pi ^{}`$ and apply the bound in (3) to get the linear bound
$$\mathrm{ex}_e(n,P(\pi ))<kc_\pi \mathrm{ex}_e(2c_\pi ^{}+1,P(\pi ))n.$$
By the bound in (2),
$$\mathrm{ex}_i(n,P(\pi ))<k(k1)(4k1)c_\pi \mathrm{ex}_e(2c_\pi ^{}+1,P(\pi ))n,$$
proving the claim. ∎
Klazar posed the following six extremal and enumerative conjectures in :
1. The number of simple $`H`$ such that $`v(H)=n`$ and $`HP(\pi )`$ is $`<c_1^n`$.
2. The number of maximal simple $`H`$ with $`v(H)=n`$ and $`HP(\pi )`$ is $`<c_2^n`$.
3. For every simple $`H`$ with $`v(H)=n`$ and $`HP(\pi )`$, we have $`e(H)<c_3n`$.
4. For every simple $`H`$ with $`v(H)=n`$ and $`HP(\pi )`$, we have $`i(H)<c_4n`$.
5. The number of simple $`H`$ with $`i(H)=n`$ and $`HP(\pi )`$ is $`<c_5^n`$.
6. The number of $`H`$ with $`i(H)=n`$ and $`HP(\pi )`$ is $`<c_6^n`$.
He showed that all six conjectures hold for a large class of permutations $`\pi `$ and that they hold for every $`\pi `$ in weaker forms: with almost linear and almost exponential bounds (respectively). Conjecture C4, however, is precisely the statement of Theorem 2.1, and it is easy to extend the proof given in this paper to affirm that all six conjectures hold for every permutation $`\pi `$.
We shall show how to amend the proofs in to prove C1, and then note that C1 implies C2, C3, C5 and C6 via Lemma 2.1 of .
###### Corollary 2.2.
For every permutation $`\pi `$ there exists a constant $`c_1>1`$ (depending on $`\pi `$) so that the number of simple hypergraphs on the vertex set $`[n]`$ avoiding $`P(\pi )`$ is $`<c_1^n`$.
###### Proof.
Theorems 2.4 and 2.5 in show that the number of hypergraphs with a given weight $`i(H)`$ that avoid $`P(\pi )`$ is at most $`9^{(3^{2k}+2k)i(H)}`$. Thus by Theorem 2.1, we are done. ∎
The Stanley–Wilf conjecture (see Bóna ), proved by Marcus and Tardos in as a corollary of their linear extremal bound, claimed that for every permutation $`\pi `$ there is a constant $`c=c(\pi )`$ such that the number of permutations $`\sigma `$ of $`[n]`$ avoiding $`\pi `$ is $`<c^n`$; the avoidance of permutations here means that $`P(\sigma )`$ is not an ordered subgraph of $`P(\pi )`$. In view of the reformulation from permutations to bipartite graphs mentioned in Section 1, Corollary 2.2 is an extension of the Stanley–Wilf conjecture. A related extension was proposed by Brändén and Mansour in Section 5 of : they conjectured that the number of words over the ordered alphabet $`[n]`$ which have length $`n`$ and avoid $`\pi `$ is at most exponential in $`n`$. These words can be represented by simple graphs $`G`$ on $`[2n]`$ in which every edge connects $`[n]`$ and $`[n+1,2n]`$ and every $`x[n]`$ has degree exactly $`1`$; the containment of ordered words is then just the ordered subgraph relation. Hence this extension is subsumed in Corollary 2.2.
Corollary 2.2 subsumes yet another extension of the Stanley–Wilf conjecture to set partitions proposed by Klazar . This extension is related to $`k`$-noncrossing and $`k`$-nonnesting set partitions whose exact enumeration was recently investigated by Chen et al. and Bousquet-Mélou and Xin . Consider, for a set partition $`H`$ of $`[n]`$, the graph $`G(H)=([n],E)`$ in which an edge connects two neighboring elements of a block (not separated by another element of the same block). Thus $`H`$ is represented by increasing paths which are spanned by the blocks. $`H`$ is a $`k`$-noncrossing (resp. $`k`$-nonnesting) partition iff $`P(12\mathrm{}k)`$ (resp. $`P(k(k1)\mathrm{}1)`$) is not an ordered subgraph of $`G(H)`$. Thus Corollary 2.2 provides an exponential bound: for fixed $`k`$, the numbers of $`k`$-noncrossing and $`k`$-nonnesting partitions of $`[n]`$ grow at most exponentially.
## 3 An extension to $`d`$-dimensional matrices
We now generalize the original Füredi–Hajnal conjecture from ordinary 0-1 matrices to $`d`$-dimensional 0-1 matrices. As was mentioned in Section 1, these are just $`d`$-ary relations (or, as we will discuss later, $`d`$-uniform, $`d`$-partite hypergraphs). We keep the matrix terminology, however, both for the sake of consistency and to highlight the similarities with the original Marcus–Tardos proof in .
###### Definition.
We will call a $`(d+1)`$-tuple $`M=(M;n_1,\mathrm{},n_d)`$ where $`M[n_1]\times \mathrm{}\times [n_d]`$ a $`d`$-dimensional (0-1) matrix, and will refer to the elements of $`M`$ as edges.
If $`F=(F;k_1,\mathrm{},k_d)`$ and $`M=(M;n_1,\mathrm{},n_d)`$ are two $`d`$-dimensional matrices, we say that $`F`$ is contained in $`M`$, $`FM`$, if there exist $`d`$ increasing injections $`f_i:[k_i][n_i]`$, $`i=1,2,\mathrm{},d`$, such that for every $`(x_1,\mathrm{},x_d)F`$ we have $`(f_1(x_1),\mathrm{},f_d(x_d))M`$; otherwise we say that $`M`$ avoids $`F`$.
###### Definition.
We set $`f(n,F,d)`$ to be the maximum size $`|M|`$ of a $`d`$-dimensional matrix $`(M;n,\mathrm{},n)`$ that avoids a $`d`$-dimensional matrix $`F`$.
For $`i[d]`$, we will denote the projection mapping from $`[n_1]\times \mathrm{}\times [n_d]`$ to $`[n_i]`$ as $`\pi _i`$. For $`t[d]`$, we define the $`t`$-remainder of $`M=(M;n_1,\mathrm{},n_d)`$ to be the $`(d1)`$-dimensional matrix $`N=(N;n_1^{},\mathrm{},n_{d1}^{})`$ where $`n_1^{}=n_1,\mathrm{},n_{t1}^{}=n_{t1}`$, $`n_t^{}=n_{t+1},\mathrm{},n_{d1}^{}=n_d`$ and the edge $`(e_1,\mathrm{},e_{d1})N`$ if and only if $`(e_1,\mathrm{},e_{t1},x,e_t,e_{t+1},\mathrm{},e_{d1})M`$ for some $`x[n_t]`$.
Let $`I_1<I_2<\mathrm{}<I_r`$ be a partition of $`[n]`$ into $`r`$ intervals and $`M=(M;n,\mathrm{},n)`$ a $`d`$-dimensional matrix. We define the contraction of $`M`$ (with respect to the intervals) to be the $`d`$-dimensional matrix $`N=(N;r,\mathrm{},r)`$ given by $`(e_1,\mathrm{},e_d)N`$ iff $`M(I_{e_1}\times \mathrm{}\times I_{e_d})\mathrm{}`$ (we could define the contraction operation for a general $`d`$-dimensional matrix and with distinct and general partitions in each coordinate but we will not need such generality).
We say that $`P=(P;k,\mathrm{},k)`$ is a $`d`$-dimensional permutation of $`[k]`$ if for every $`i[d]`$ and $`x[k]`$ there is exactly one edge $`eP`$ with $`\pi _i(e)=x`$. Note that $`|P|=k`$ and that there are exactly $`(k!)^{d1}`$ $`d`$-dimensional permutations of $`[k]`$. For $`d=1`$, the only 1-dimensional permutation $`(P;k)`$ is $`[k]`$, and for $`d=2`$ the 2-dimensional permutations $`P=(P;k,k)`$ are precisely the $`k\times k`$ 0-1 permutation matrices. A $`d`$-dimensional permutation of $`[k]`$ can be thought of as a $`d\times k`$ matrix with the first row normalized to $`1,2,\mathrm{},k`$ and with each row being a permutation of $`1,2,\mathrm{},k`$. The columns would then give the coordinates of the $`k`$ edges in $`P`$.
It is also possible to view the structure $`M=(M;n_1,\mathrm{},n_d)`$ as an ordered, $`d`$-uniform, $`d`$-partite hypergraph with partitions $`[n_i]`$. In this interpretation, the image of $`M`$ by the projection $`\pi _i`$ is obtained by intersecting the edges with the $`i^{th}`$ partition, while the intersections with the union of all partitions except the $`t^{th}`$ one give the $`t`$-remainder of $`M`$ (in both cases we disregard multiplicity of edges). Furthermore, the set of $`d`$-dimensional permutations of $`[k]`$ would be the set of perfect matchings of the complete $`d`$-uniform, $`d`$-partite hypergraph on $`kd`$ vertices.
We will make use of two observations, analogous to those made in :
1. For any $`t[d]`$, the $`t`$-remainder of a $`d`$-dimensional permutation of $`[k]`$ is a $`(d1)`$-dimensional permutation of $`[k]`$. Furthermore, each edge of the resulting $`t`$-remainder can be completed (by adding the $`t`$-th coordinate) in a unique way to an edge of the original permutation.
2. If $`M=(M;n,\mathrm{},n)`$ avoids a $`d`$-dimensional permutation, then so does any contraction of $`M`$.
###### Theorem 3.1.
For every fixed $`d`$-dimensional permutation $`P`$,
$$f(n,P,d)=O(n^{d1}).$$
On the other hand it is clear that for a $`d`$-dimensional permutation $`P`$ with $`|P|>1`$ we have $`f(n,P,d)n^{d1}`$ ($`f(n,P,d)=0`$ if $`|P|=1`$). Thus, for a $`d`$-dimensional permutation $`P`$ with $`|P|>1`$,
$$f(n,P,d)=\mathrm{\Theta }(n^{d1}).$$
This bound can be given an equivalent formulation. We say that a matrix $`M=(M;n_1,\mathrm{},n_d)`$ is a $`d`$-dimensional $`k`$-grid if each $`[n_i]`$ can be partitioned in $`k`$ intervals $`I_{i,1}<I_{i,2}<\mathrm{}<I_{i,k}`$ so that $`|M(I_{1,j_1}\times I_{2,j_2}\times \mathrm{}\times I_{k,j_k})|=1`$ for every $`d`$-tuple $`(j_1,j_2,\mathrm{},j_d)[k]^d`$ (thus, in particular, $`|M|=k^d`$). Let $`g(n,k,d)`$ be the maximum size of a $`d`$-dimensional $`n\times n\times \mathrm{}\times n`$ matrix that contains no $`d`$-dimensional $`k`$-grid. Then
$$g(n,k,d)=\mathrm{\Theta }(n^{d1}).$$
It is clear that $`g(n,k,d)n^{d1}`$. The bound $`g(n,k,d)=O(n^{d1})`$ implies $`f(n,P,d)=O(n^{d1})`$ for every $`P`$ because every $`d`$-dimensional $`k`$-grid contains every $`d`$-dimensional permutation of $`[k]`$. In the other way, it is easy to see that there exist $`d`$-dimensional $`k`$-grids that are $`d`$-dimensional permutations of $`[k^d]`$. Thus $`f(n,P,d)=O(n^{d1})`$ implies $`g(n,k,d)=O(n^{d1})`$.
To prove Theorem 3.1, we will show that a $`d`$-dimensional matrix of big enough size must contain every $`d`$-dimensional permutation of $`k`$. We set
$$f(n,k,d)=\underset{P}{\mathrm{max}}f(n,P,d)$$
where $`P`$ runs through all $`d`$-dimensional permutations of $`[k]`$.
###### Lemma 3.2.
Let $`d2`$, $`m,n_0𝐍`$. Then
$$f(mn_0,k,d)(k1)^df(n_0,k,d)+dn_0m^d\left(\genfrac{}{}{0pt}{}{m}{k}\right)f(n_0,k,d1).$$
###### Proof.
Let $`M=(M,mn_0,\mathrm{},mn_0)`$ be a $`d`$-dimensional matrix that avoids $`P`$, a $`d`$-dimensional permutation of $`[k]`$. We aim to bound the size of $`M`$.
We split $`[mn_0]`$ into $`n_0`$ intervals $`I_1<I_2<\mathrm{}<I_{n_0}`$, each of length $`m`$, and define, for $`i_1,\mathrm{},i_d[n_0]`$,
$$S(i_1,\mathrm{},i_d)=\{eM:\pi _j(e)I_{i_j}\text{ for }j=1,\mathrm{},d\}.$$
Note that this partitions the set of edges of $`M`$ into $`n_0^d`$ pieces. We will call these sets of edges blocks and we define a cover of these blocks by a total of $`dn_0+1`$ sets $`\{U_0\}\{U(t,j):t[d],j[n_0]\}`$ as follows:
* $`U(t,j)=\{S(i_1,\mathrm{},i_d):i_t=j\text{ and }|\pi _t(S(i_1,\mathrm{},i_d))|k\}`$
* $`U_0`$ consists of the blocks which are not in any $`U(t,j)`$
Note that the total number of non-empty blocks is exactly the number of edges in the contraction of $`M`$ with respect to the partition $`\{I_i\}`$. Since $`M`$ does not contain $`P`$, the contraction of $`M`$ can not contain $`P`$, so the number of non-empty blocks is at most $`f(n_0,k,d)`$. Also note that any block $`B`$ in $`U_0`$ has at most $`(k1)^d`$ edges in it (because $`BX_1\times \mathrm{}\times X_d`$ for some $`X_i[mn_0]`$ with $`|X_i|<k`$). Hence
$$|U_0|(k1)^df(n_0,k,d).$$
Now we fix $`t[d]`$ and $`j[n_0]`$. Clearly,
$$|U(t,j)|m^d|U(t,j)|.$$
We assume, for a contradiction, that $`|U(t,j)|>\left(\genfrac{}{}{0pt}{}{m}{k}\right)f(n_0,k,d1)`$. By the definition of $`U(t,j)`$ and the pigeonhole principle, there are $`k`$ numbers $`c_1<c_2<\mathrm{}<c_k`$ in $`I_j`$ and $`r`$ blocks $`S_1,S_2,\mathrm{},S_r`$ in $`U(t,j)`$ where $`r>f(n_0,k,d1)`$ such that for every $`S_a`$ and every $`c_b`$ there is an $`eS_a`$ with $`\pi _t(e)=c_b`$. Let $`P^{}`$ be the $`t`$-remainder of $`P`$ and $`M^{}=(M^{};n_0,\mathrm{},n_0)`$ be the $`(d1)`$-dimensional matrix arising from contracting $`(_{i=1}^rS_i,n,\mathrm{},n)`$ with respect to the intervals $`I_i`$ and then taking the $`t`$-remainder. Since $`|M^{}|=r>f(n_0,k,d1)`$, $`M^{}`$ contains $`P^{}`$. Thus among the blocks $`S_1,S_2,\mathrm{},S_r`$ there exist $`k`$ of them—call them $`S_1,S_2,\mathrm{},S_k`$—so that for any selection of $`k`$ edges $`e_1S_1,\mathrm{},e_kS_k`$ their $`t`$-remainders form a copy of $`P^{}`$. Furthermore, due to the property of the blocks $`S_i`$, it is possible to select $`e_1,\mathrm{},e_k`$ so that their $`t`$-th coordinates agree with $`P`$. Then $`e_1,\mathrm{},e_k`$ form a copy of $`P`$, a contradiction. Therefore
$$|U(t,j)|m^d|U(t,j)|m^d\left(\genfrac{}{}{0pt}{}{m}{k}\right)f(n_0,k,d1)$$
and
$$|_{t,j}U(t,j)|dn_0m^d\left(\genfrac{}{}{0pt}{}{m}{k}\right)f(n_0,k,d1).$$
Combining this with the bound for $`U_0`$ gives the stated bound. ∎
Theorem 3.1 will be a direct consequence of the following lemma:
###### Lemma 3.3.
For $`m=k^{d/(d1)}`$, $`f(n,k,d)k^d\left(dm\left(\genfrac{}{}{0pt}{}{m+1}{k}\right)\right)^{d1}n^{d1}`$.
###### Proof.
We will proceed by induction on $`d+n`$. For $`d=1`$ this holds since $`f(n,k,1)=k1`$ and for $`n<k^2`$, this holds trivially. Now given $`n`$ and $`d2`$, assume that the hypothesis is true for all $`d^{},n^{}`$ such that $`d^{}+n^{}<d+n`$.
Let $`n_0=n/m`$ and
$$c_d=k^d\left(dm\left(\genfrac{}{}{0pt}{}{m+1}{k}\right)\right)^{d1}.$$
Using the inequality $`f(n,k,d)<f(mn_0,k,d)+dmn^{d1}`$, Lemma 3.2, the inductive hypotheses, and $`n_0n/m`$, we get
$$f(n,k,d)<\left(\frac{(k1)^d}{m^{d1}}c_d+dm\left(\left(\genfrac{}{}{0pt}{}{m}{k}\right)c_{d1}+1\right)\right)n^{d1}.$$
Since $`\frac{(k1)^d}{m^{d1}}(1\frac{1}{k})^d1\frac{1}{k}`$ and $`\left(\genfrac{}{}{0pt}{}{m}{k}\right)c_{d1}+1\left(\genfrac{}{}{0pt}{}{m+1}{k}\right)c_{d1}`$, it follows that $`f(n,k,d)<c_dn^{d1}`$ with the above defined $`c_d`$. ∎
## 4 Concluding remarks
We were informed recently that Balogh, Bollobás and Morris derived Theorem 2.1 (their Theorem 2) and Corollary 2.2 (their Theorem 1) independently. The proofs in are self-contained (not appealing to the results in ) and their approach is different from ours. In fact, they are able to prove stronger statements, which in turn imply Theorem 2.1 and Corollary 2.2 from this paper.
It should be noted that we make no effort to optimize any of the constants in Section 3, however it would be interesting to see if any of the constants could be drastically reduced. The constant achieved in this paper is double exponential in $`k`$, whereas we conjecture the true constant is in fact much smaller.
## 5 Acknowledgments
The authors would like to thank Gábor Tardos for, among other things, his enlightening conversations and helpful remarks. |
warning/0507/astro-ph0507286.html | ar5iv | text | # MACHOs in M31?Based on observations made with the Isaac Newton Telescope operated on the island of La Palma by the Isaac Newton Group in the Spanish Observatorio del Roque de los Muchachos of the Instituto de Astrofisica de Canarias Absence of evidence but not evidence of absence
## 1 Introduction
Compact objects that emit little or no radiation form a class of plausible candidates for the composition of dark matter halos. Examples include black holes, brown dwarfs, and stellar remnants such as white dwarfs and neutron stars. These objects, collectively known as Massive Astrophysical Compact Halo Objects or MACHOs, can be detected indirectly through gravitational microlensing wherein light from a background star is amplified by the space-time curvature associated with the object (Paczyński, 1986).
The first microlensing surveys were performed by the MACHO (Alcock et al., 2000) and EROS (Lasserre et al., 2000; Afonso et al., 2003) collaborations and probed the Milky Way halo by monitoring stars in the Large and Small Magellanic Clouds. While both collaborations detected microlensing events they reached different conclusions. The MACHO collaboration reported results that favour a MACHO halo fraction of 20%. On the other hand, the results from EROS are consistent with no MACHOs and imply an upper bound of 20% on the MACHO halo fraction. The two surveys are not inconsistent with each other since they probe different ranges in MACHO masses. They do leave open the question of whether MACHOs make up a substantial fraction of halo dark matter and illustrate an inherent difficulty with microlensing searches for MACHOs, namely that they must contend with a background of self-lensing events (i.e., both lens and source stars in the Milky Way or Magellanic clouds), variable stars, and supernovae. The Magellanic Cloud surveys are also hampered by having only two lines of sight through the Milky Way halo.
Microlensing surveys towards M31 have important advantages over the Magellanic Cloud surveys (Crotts, 1992). The microlensing event rate for M31 is greatly enhanced by the high density of background stars and the availability of lines-of-sight through dense parts of the M31 halo. Furthermore, since lines of sight toward the far side of the disk pass through more of the halo than those toward the near side, the event distribution due to a MACHO population should exhibit a near-far asymmetry (Gyuk & Crotts, 2000; Kerins et al., 2001; Baltz et al., 2003).
Unlike stars in the Magellanic Clouds, those in M31 are largely unresolved, a situation that presents a challenge for the surveys but one that can be overcome by a variety of techniques. To date microlensing events toward M31 have been reported by four different collaborations, VATT-Columbia (Uglesich et al., 2004), MEGA (de Jong et al., 2004), POINT-AGAPE (Paulin-Henriksson et al., 2003; Calchi Novati et al., 2003, 2005) and WeCAPP (Riffeser et al., 2003).
Recently, the POINT-AGAPE collaboration presented an analysis of data from three seasons of INT observations in which they concluded that “at least 20% of the halo mass in the direction of M31 must be in the form of MACHOs” (Calchi Novati et al., 2005). Their analysis is significant because it is the first for M31 to include a model for the detection efficiency.
The MEGA collaboration is conducting a microlensing survey in order to quantify the amount of MACHO dark matter in the M31 halo. Observations are carried out at a number of telescopes including the 2.5m Isaac Newton Telescope (INT) on La Palma, and, on Kitt Peak, the 1.3m McGraw-Hill, 2.4m Hiltner, and 4m Mayall telescopes. The observations span more than 4 seasons. The first three seasons of INT data were acquired jointly with the POINT-AGAPE collaboration though the data reduction and analysis have been performed independently.
In de Jong et al. (2004) (hereafter Paper I) we presented 14 candidate microlensing events from the first two seasons of INT data. The angular distribution of these events hinted at a near-far asymmetry albeit with low statistical significance. Recently An et al. (2004a) pointed out that the distribution of variable stars also shows a near-far asymmetry raising questions about the feasibility of the M31 microlensing program. However, the asymmetry in the variable stars is likely caused by extinction which can be modelled.
In this paper, we present our analysis of the 4-year INT data set. This extension of the data by two observing seasons compared to Paper I is a significant advance, but this data set is still only a subset of the MEGA survey. The forthcoming analysis of the complete data set will feature a further increase in time-sampling and baseline coverage and length. But there are more significant advances from Paper I. We improve upon the photometry and data reduction in order to reduce the number of spurious variable-source detections. We fully automate the selection of microlensing events and model the detection efficiency through extensive Monte Carlo simulations. Armed with these efficiencies, we compare the sample of candidate microlensing events with theoretical predictions for the rate of events and their angular and timescale distributions. These predictions are based on new self-consistent disk-bulge-halo models (Widrow & Dubinski, 2005) and a model for differential extinction across the M31 disk. The models are motivated by photometric and kinematic data for M31 as well as a theoretical understanding of galactic dynamics.
Our analysis shows that the observed number of events can be explained by self-lensing due to stars in the disk and bulge of M31, contrary to the findings of Calchi Novati et al. (2005). Our results are consistent with a no MACHO hypothesis, although we cannot rule out a MACHO fraction of 30%.
Data acquisition and reduction methods are discussed in Sect. 2. The construction of a catalogue of artificial microlensing events is described in Sect. 3. This catalogue provides the basis for a Monte Carlo simulation of the survey and is used, in Sect. 4, to set the selection criteria for microlensing events. Our candidate microlensing events are presented in Sect. 5. The artificial event catalogue is then used in Sect. 6 to calculate the detection efficiency. Our extinction model is presented in Sect. 7. In Sect. 8 the theoretical models are described and the predictions for event rate and distribution are presented. A discussion of the results and our conclusions are presented in Sects. 9 and 10.
## 2 Data acquisition and reduction
Observations of M31 were carried out using the INT Wide Field Camera (WFC) and spread equally over the two fields of view shown in Fig. 1. The WFC field of view is approximately 0.25$`\mathrm{}^o`$ and consists of four 2048x4100 CCDs with a pixel scale of 0.333″. The chosen fields cover a large part of the far side (SE) of the M31 disk and part of the near side. Observations span four observing seasons each lasting from August to January. Since the WFC is not always mounted on the INT, observations tend to cluster in blocks of two to three weeks with comparable-sized gaps during which there are no observations.
Exposures during the first (1999/2000) observing season were taken in three filters, r, g and i, which correspond closely to Sloan filters. For the remaining seasons (2000/01, 2001/02, 2002/03), only the r and i filters were used. Nightly exposure times for the first season were typically 10 minutes in duration but ranged from 5 to 30 minutes. For the remaining seasons the default exposure time was 10 minutes per field and filter. Standard data reduction procedures, including bias subtraction, trimming and flatfielding were performed in IRAF.
### 2.1 Astrometric registration and image subtraction
We use Difference Image Photometry (DIP) (Tomaney & Crotts, 1996) to detect variable objects in the highly crowded fields of M31. Individual images are subtracted from a high quality reference image to yield difference images in which variable objects show up as residuals. Most operations are carried out with the IRAF package DIFIMPHOT.
Images are transformed to a common astrometric reference frame. A high signal-to-noise (S/N) reference image is made by stacking high-quality images from the first season. Exposures from a given night are combined to produce a single “epoch” with Julian date taken to be the weighted average of the Julian dates of the individual exposures.
Average point spread functions (PSFs) for each epoch and for the reference image are determined from bright unsaturated stars. A convolution kernel is calculated by dividing the Fourier transform of the PSF from an epoch by the PSF transform from the reference image. This kernel is used to degrade the image with better seeing (usually the reference image) before image subtraction is performed (Tomaney & Crotts, 1996).
Image subtraction does not work well in regions with very high surface brightness because of a lack of suitable, unsaturated stars. For this reason we exclude a small part of the south field located in a high-surface brightness region of the bulge (see Fig. 1).
### 2.2 Variable source detection
Variable sources show up in the difference images as residuals which can be positive or negative depending on the flux of the source in a given epoch relative to the average flux of the source as measured in the reference image. However, difference images tend to be dominated by shot noise. The task at hand is to differentiate true variable sources from residuals that are due to noise.
The program SExtractor (Bertin & Arnouts, 1996) is used to detect “significant residuals” in r epochs, defined as groups of 4 or more connected pixels that are all at least 3$`\sigma `$ above or below the background. Residuals from different epochs are cross-correlated and those that appear in two or more consecutive epochs are catalogued as variable sources. (Because of fringing, the i difference images are of poorer quality than the r ones and we therefore use r data to make the initial identification.)
### 2.3 Lightcurves and Epoch quality
The difference images for a number of epochs are discarded for a variety of reasons. Epochs with poor seeing do not give clean difference images. We require better than 2″ seeing and discard 7 epochs and parts of 12 epochs where this condition is not met. PSF-determination fails if an image is over-exposed. We discard 7 epochs and parts of another 7 epochs for this reason. Finally 2 epochs from the second and third seasons are discarded because of guiding errors.
Lightcurves for the variable sources are obtained by performing PSF-fitting photometry on the residuals in the difference images. For every pixel the Poisson-noise is evaluated as well as the fractional flux error due to photometric inaccuracies in the matching and subtraction steps for the difference image in question. Fluxes and their error bars are derived by optimal weighting of the individual pixel values. Lightcurves are also produced at positions where no variability is identified and fit to a flat line. These lightcurves serve as a check on the contribution to the flux error bars derived from the photometric accuracy of each difference image. For each epoch, we examine by eye the distribution of the deviations from the flat-line fits normalised by the photometric error bar. Epochs where this distribution shows broad non-Gaussian wings are discarded since wings in the distribution are likely caused by guiding errors or highly variable seeing between individual exposures. For epochs where this normalised error distribution is approximately Gaussian but with a dispersion greater than one, the error bars are renormalised.
Approximately 19% of the 209 r epochs and 22% of the 183 i epochs are discarded. The number of epochs that remain for each season, filter, and field are tabulated in Table 1. Though variable objects are detected in r, lightcurves are constructed in both r and i. In total, 105,447 variable source lightcurves are generated.
## 3 Artificial microlensing events
This section describes the construction of a catalogue of artificial microlensing lightcurves which forms the basis of our Monte Carlo simulations. We add artificial events to the difference images and generate lightcurves in the same manner as is done with the actual data. The details of this procedure follow a review of microlensing basics and terminology.
### 3.1 Microlensing lightcurves
The lightcurve for a single-lens microlensing event is described by the time-dependent flux (Paczyński, 1986):
$$F(t)=F_0\frac{u^2+2}{u\sqrt{u^2+4}}F_0A(t)$$
(1)
where $`F_0`$ is the unlensed source flux and $`A`$ is the amplification. $`u=u(t)`$ is the projected separation of the lens and the source in units of the Einstein radius,
$$R_\text{E}=\sqrt{\frac{4Gm}{c^2}\frac{D_{\mathrm{OL}}D_{\mathrm{LS}}}{D_{\mathrm{OS}}}},$$
(2)
where $`m`$ is the lens mass and the $`D`$’s are the distances between observer, lens and source. If the motions of lens, source, and observer are uniform for the duration of the lensing event we can write
$$u(t)=\sqrt{\beta ^2+\left(\frac{tt_{\mathrm{max}}}{t_\mathrm{E}}\right)^2}$$
(3)
where $`\beta `$ is the impact parameter in units of $`R_\mathrm{E}`$, that is, the mimimum value attained by $`u`$. $`t_{\mathrm{max}}`$ is the time of maximum amplification and $`t_\mathrm{E}`$ is the Einstein time, defined as the time it takes the source to cross the Einstein radius.
In classical microlensing the measured lightcurves contain contributions from unlensed sources. Blending, as this effect is known, changes the shape of the lightcurve and can also spoil the achromaticity implicit in equation 1. In our survey, we measure flux differences that are created by subtracting a reference image. Since the flux from unlensed sources is subtracted from an image to form the difference image, blending is not a problem unless the unlensed sources are variable. Blending by variable sources does introduce variations in the baseline flux and adversely affects the fit.
For a difference image the microlensing lightcurve takes the form
$$\mathrm{\Delta }F(t)F(t)F_{\mathrm{ref}}=\mathrm{\Delta }F_{\mathrm{bl}}+F_0(A(t)1)$$
(4)
where $`F_{\mathrm{ref}}`$ is the reference image flux and $`\mathrm{\Delta }F_{\mathrm{bl}}F_0F_{\mathrm{ref}}`$. Thus, if in the reference image the source is not lensed, $`F_{\mathrm{ref}}=F_0`$ and therefore $`\mathrm{\Delta }F_{\mathrm{bl}}0`$. Only if the source is amplified in the reference image will $`\mathrm{\Delta }F_{\mathrm{bl}}`$ be non-zero and negative.
For unresolved sources, a situation known as pixel lensing (and the one most applicable to stars in M31), those microlensing events that can be detected typically have high amplification. In the high amplification limit, $`t_\mathrm{E}`$ and $`\beta `$ are highly degenerate (Gould, 1996; Baltz & Silk, 2000) and difficult to extract from the lightcurve. It is therefore advantageous to parameterise the event duration in terms of the half-maximum width of the peak,
$$t_{\mathrm{FWHM}}=t_\mathrm{E}w(\beta ),$$
(5)
where
$$w(\beta )=2\sqrt{2f(f(\beta ^2))\beta ^2}$$
(6)
and
$$f(x)=\frac{x+2}{\sqrt{x(x+4)}}1$$
(7)
(Gondolo, 1999). $`w(\beta )`$ has the limiting forms $`w(\beta 1)\beta \sqrt{3}`$ and $`w(\beta 1)\beta (\sqrt{2}1)^{1/2}`$.
### 3.2 Simulation parameters
The parameters that characterise microlensing events fall into two categories: “microlensing parameters” such as $`\beta `$, $`t_{\mathrm{max}}`$, and $`t_\mathrm{E}`$, and parameters that describe the source such as its brightness $`F_{0,\mathrm{r}}`$, its r-i colour $`𝒞`$, and its position. We survey many lines-of-sight across the face of M31. Furthermore, all types of stars can serve as a source for microlensing. Therefore, our artificial event catalogue must span a rather large parameter space. This parameter space is summarised in Table 2 and motivated by the following arguments:
$``$ Peak times and baseline fluxes
We demand that the portion of the lightcurve near peak amplitude is well-sampled and therefore restrict $`t_{\mathrm{max}}`$ to one of the four INT observing seasons. The reference images are constructed from exposures obtained during the first season. If a microlensing event occurs during the first season and if the source is amplified in one or more exposures during this season, the baseline in the difference image will be below the true baseline. For an actual event in season one, this off-set is absorbed in one of the fit parameters for the lightcurve. For artificial events, the baseline is corrected by hand.
$``$ Event durations
Limits on the duration of detectable events follow naturally from the setup of the survey and the requirement that events are sampled through their peaks. Since the INT exposures are combined nightly, events with $`t_{\mathrm{FWHM}}<1\mathrm{day}`$ are practically undetectable except for very high amplifications. On the other hand, events with $`t_{\mathrm{FWHM}}`$ approaching the six-month length of the observing season are also difficult to detect with the selection probability decreasing linearly with $`t_{\mathrm{FWHM}}`$. Because gaps in the time coverage of our survey will affect our sensitivity to short events more strongly than to long events, sampling should be denser at shorter timescales. To limit computing time and ensure statistically significant results spread over a wide range of event durations, we simulate events at six discrete values of $`t_{\mathrm{FWHM}}`$: 1, 3, 5, 10, 20 and 50 days.
$``$ Source fluxes and colours
Faint stars are more abundant than bright ones. On the other hand, microlensing events are more difficult to detect when the source is a faint star. The competition between these two effects means that there is a specific range of the source luminosity function that is responsible for most of the detectable microlensing events.
The maximum flux difference during a microlensing event is
$$\mathrm{\Delta }F_{\mathrm{max}}=F_0\left(\frac{\beta ^2+2}{\beta \sqrt{\beta ^2+4}}1\right)$$
(8)
where we are ignoring the $`\mathrm{\Delta }F_{\mathrm{bl}}`$ term in equation 4. Let $`\mathrm{\Delta }F_{\mathrm{det}}`$ be the detection threshold for $`\mathrm{\Delta }F_{\mathrm{max}}`$. A lower bound on $`\mathrm{\Delta }F_{\mathrm{max}}`$ implies an upper bound on $`\beta `$ which, through equation 8, is a function of the ratio $`F_0/F_{det}`$: $`\beta _\mathrm{u}=\beta _\mathrm{u}\left(F_0/F_{det}\right)`$. The probability that a given source is amplified to a detectable level scales as $`\beta _\mathrm{u}^2`$. In Fig. 2 we show both the R-band luminosity function, $`N_{}`$, from Mamon & Soneira (1982) and the product of this luminosity function with $`\beta _\mathrm{u}^2`$ assuming a detection threshold of $`F_{\mathrm{det}}=1\mathrm{ADU}\mathrm{s}^1`$. The latter provides a qualitative picture of the distribution of detectable microlensing events. This distribution peaks at an absolute R-band magnitude of approximately 0 indicating that most of the sources for detectable microlensing events are Red Giant Branch (RGB) stars.
Since there is no point in simulating events we cannot detect we let the impact parameter $`\beta `$ vary randomly between 0 and $`\beta _\mathrm{u}`$. Table 2 summarises the fluxes and values for $`\beta _\mathrm{u}`$ used in the simulations.
For the artificial event catalogue, we use source stars with a r fluxes at several discrete values between 0.01 and 10 ADU s<sup>-1</sup>. Typically the r$``$i colours of RGB stars range between $`𝒞=0.5`$ and $`2.0`$. We assume $`𝒞=0.75`$ for our artificial events. As a check of the dependence of the detection efficiency with colour, we also simulate events with $`𝒞=1.25`$.
$``$ Position in M31
Lightcurve quality and detection efficiency vary with position in M31 for several reasons. The photometric sensitivity and therefore the detection efficiency depend on the amount of background light from M31 and are lowest in the the bright central areas of the bulge. Difference images from these areas are also highly crowded with variable-star residuals which influence the photometry and add noise to the microlensing lightcurves. To account for the position-dependence of the detection efficiency, artificial events are generated across the INT fields. To be precise, the artificial event catalogue is constructed in a series of runs. For each run, artificial events are placed on a regular grid with spacing of a 45 pixels ($`15\mathrm{}`$) so that there are 3916 artificial events per chip. The grid is shifted randomly between runs by a maximum of 10 pixels.
To summarise, artificial events are characterised by the parameters $`t_{\mathrm{FWHM}}`$, $`F_0`$, $`𝒞`$, $`t_{\mathrm{max}}`$, $`\beta `$, and their angular position. These events are added as residuals to the difference images using the PSF in the subregion of the event. The residuals also include photon noise. The new difference images are analysed as in Sect. 2 and lightcurves are built for all artificial events detected as variable objects.
## 4 Microlensing event selection
The vast majority of variable sources in our data set are variable stars. In this section we describe an automated algorithm that selects candidate microlensing lightcurves from this rather formidable background. Our selection criteria pick out lightcurves that have a flat baseline and a single peak with the “correct” shape. The criteria take the form of conditions on the $`\chi ^2`$ statistic that measures the goodness-of-fit of an observed lightcurve to equation 4. The fit involves seven free parameters: $`t_{\mathrm{max}}`$, $`\beta `$, $`t_\mathrm{E}`$ , $`F_{0,\mathrm{r}}`$, $`F_{0,\mathrm{i}}`$, $`\mathrm{\Delta }F_{\mathrm{bl},\mathrm{r}}`$, and $`\mathrm{\Delta }F_{\mathrm{bl},\mathrm{i}}`$. To increase computation speed we first obtain rough estimates for $`t_{\mathrm{max}}`$ and $`t_\mathrm{E}`$ from the r lightcurve and then perform the full 7-parameter fit using both r and i lightcurves.
Gravitational lensing is achromatic and therefore the observed colour of a star undergoing microlensing remains constant in contrast with the colour of certain variables. While we do not impose an explicit achromaticity condition, changes in the colour of a variable source show up as a poor simultaneous r and i fit. Because many red variable stars vary little in colour, as defined by measurable differences in flux ratios, the lightcurve shape and baseline flatness are better suited for distinguishing microlensing events from long period variable stars (LPVs) than a condition on achromaticity.
Lightcurves must contain enough information to fit adequately both the peak of the microlensing event and the baseline. We therefore impose the following conditions: (1) The r and i lightcurves must contain at least 100 data points. (2) The peak must be sampled by several points well-above the baseline. (3) The upper half of the peak, as defined in the difference-image lightcurve, must lie completely within a well-sampled observing period. The second condition can be made more precise. We allow for one of the following two possibilities: (a) 4 or more data points in the r-lightcurve are $`3\sigma `$ above the baseline or (b) 2 or more points in r and 1 or more points in i are $`3\sigma `$ above the baseline. (The r data is weighted more heavily than the i data because it is generally of higher quality and because i was not sampled as well during the first season.) The third condition insures that we sample both rising and falling sides of the peak. We note that there are periods during the last two seasons where we do not have data due to bad weather. The periods we use are the following: 01/08/1999-13/12/1999, 04/08/2000-23/01/2001, 13/08/2001-16/10/2001, 01/08/2002-10/10/2002, and 23/12/2002-31/12/2002.
The selection of candidate microlensing events is based on the $`\chi ^2`$-statistic for the fit of the observed lightcurve to equation 4 as well as $`\mathrm{\Delta }\chi ^2\chi _{\mathrm{flat}}^2\chi ^2`$ where $`\chi _{\mathrm{flat}}^2`$ is the $`\chi ^2`$-statistic for the fit of the observed lightcurve to a flat line. Our $`\chi ^2`$-cuts are motivated by simulations of artificial microlensing events. In Fig. 3 we show the distribution of artificial events with $`t_{\mathrm{FWHM}}=50,\mathrm{\hspace{0.17em}10},`$ and 1 days (panels a, b, and c respectively) and for all variable sources in one of the CCDs (panel d). In Fig. 4, we show the variable sources from all CCDs that satisfy conditions 1-3. The plots are presented in terms of $`\chi ^2/N`$ and $`\mathrm{\Delta }\chi ^2/N`$ where $`N`$ is the number of data points in an event. We choose the following cuts:
$$\mathrm{\Delta }\chi ^2>1.5N$$
(9)
and
$$\chi ^2<\left(N7\right)f\left(\mathrm{\Delta }\chi ^2\right)+3\left(2\left(N7\right)\right)^{1/2}$$
(10)
where $`f\left(\mathrm{\Delta }\chi ^2\right)=\mathrm{\Delta }\chi ^2/100+1`$. The first criterion is meant to filter out peaks due to noise or variable stars. The second criterion corresponds to a $`3\sigma `$-cut in $`\chi ^2`$ for low signal-to-noise events. The $`\chi ^2`$ threshold increases with increasing $`\mathrm{\Delta }\chi ^2`$. Panels a-c of Fig. 3 show a trend where $`\chi ^2`$ increases systematically with $`\mathrm{\Delta }\chi ^2`$. This effect is due to the photometry routine in DIFIMPHOT which underestimates the error in flux measurements for high flux values. The function $`f`$ is meant to compensate for this effect.
The selection criteria appear as lines in Figs. 3 and 4. (To draw these lines, we take $`N=309`$ though in practice $`N`$ is different for individual lightcurves.)
## 5 Candidate events
Of the 105 477 variable sources 28 667 satisfy conditions 1-3. Of these, 14 meet the criteria set by equations 9 and 10. The positions of 12 of these events in the $`\chi ^2/N\mathrm{\Delta }\chi ^2/N`$ plane are shown in Fig. 4.
### 5.1 Sample description
In Table 3 we summarise the properties and fit parameters of the 14 candidate microlensing events. The first column gives the assigned names of the events using the nomenclature from Paper I. The numbering reflects the fact that candidates 4, 5, 6, and 12 from Paper I are evidently variable stars since they peaked a second time in the fourth season. The other 10 events from Paper I are “rediscovered” in the current more robust analysis. Four additional candidates, events 15, 16, 17, and 18, are presented. Event 16 is the same as PA-99-N1 from Paulin-Henriksson et al. (2003) and was not selected in our previous analysis because the baseline was too noisy due to a nearby bright variable star. It now passes our selection criteria thanks to the smaller aperture used for the photometry (see discussion below). The three other events all peaked in the fourth observing season and are reported here for the first time.
The coordinates of the events are given in columns 2 and 3 of Table 3; their positions within the INT fields are shown in Fig. 5. The fit parameters, $`\chi ^2`$, and $`\mathrm{\Delta }\chi ^2`$ are given in the remaining columns. In Appendix A we show the r and i lightcurves, thumbnails from the difference images for a number of epochs, and a comparison of $`\mathrm{\Delta }`$r and $`\mathrm{\Delta }`$i for points near the peak. The latter provides an indication of the achromaticity of the event. The lightcurves include data points from observations at the 4m Mayall telescope on Kitt Peak (KP4m) though the fits use only INT data.
We have already seen that variable stars can mimic microlensing events. Blending of variable stars is also a problem since it leads to noisy baselines. This problem was rather severe in Paper I causing us to miss event PA-99-N1 found by the POINT-AGAPE collaboration (Paulin-Henriksson et al., 2003). In an effort to reduce the effects of blending by variable stars, we use a smaller aperture when fitting the PSF to residuals in the difference images. Nevertheless, some variable star blending is unavoidable, especially in the crowded regions close to the centre of M31. Event 3 provides an example of this effect. A faint positive residual is visible in the 1997 KP4m difference image as shown in Fig. 6. The residual is located one pixel ($`0.21^{\prime \prime }`$) from the event and is likely due to a variable star. It corresponds to the data point in the lightcurve $``$1000 days before the event and well-above the baseline (see Fig. 23). The KP4m data point from 2004 is also above the baseline but in this and other difference images, no residual is visible. The implication is that variable stars can influence the photometry even when they are too faint to be detected directly from the difference images.
Good simultaneous fits are obtained in both r and i for all candidate events. Event 7 has a high $`\chi ^2/N`$ of 1.98, but since $`\mathrm{\Delta }\chi ^2/N`$ is very high, the event easily satisfies our selection criteria. In high S/N events, secondary effects from parallax or close caustic approaches can cause measurable deviations from the standard microlensing fit. In addition, as discussed above, we tend to underestimate the photometric errors at high flux levels. An et al. (2004b) studied this event in detail and found that the deviations from the standard microlensing shape of the POINT-AGAPE lightcurve are best explained by a binary lens. The somewhat high $`\chi ^2`$ for events 10 and 15 are probably because they are located in regions of high surface brightness.
All of the candidate events are consistent with achromaticity, though for events with low S/N, it is difficult to draw firm conclusions directly from the lightcurves or $`\mathrm{\Delta }r^{}`$ vs. $`\mathrm{\Delta }i^{}`$ plots. The values for $`F_{0,\mathrm{r}}`$ and $`𝒞`$ for the events give some indication of the properties of the source stars. The unlensed fluxes are consistent with the expected range of $`0.110\mathrm{ADU}\mathrm{s}^1`$ and the colours for most of the events are typical of RGB stars. Note however that for many of the events, the uncertainties for $`F_0`$, $`\beta `$, and $`t_{\mathrm{FWHM}}`$ are quite large. These uncertainties reflect degeneracies among the lightcurve fit parameters.
The number of candidate events varies considerably from season to season. We find 7 events in the first season, 4 in the second season, none in the third season and 3 in the fourth season. The paucity of events during the third and fourth seasons is not surprising given that we have fewer epochs for those seasons (see table 1). In particular, the gaps in time coverage during those seasons conspired against the detection of short duration events.
### 5.2 Comparison with other surveys
The POINT-AGAPE collaboration published several analyses of the INT observations. In Paulin-Henriksson et al. (2003) they presented four convincing microlensing events from the first two observing seasons using stringent selection criteria. In particular, they restricted their search to events with high S/N and $`t_{\mathrm{FWHM}}<25\mathrm{days}`$. They argued that one of these events (PA-00-S3) is probably due to a stellar lens in the M31 bulge. This event lies in the region of the bulge excluded from our analysis (see Fig. 1). The other three events, PA-99-N1, PA-99-N2, and PA-00-S4, correspond respectively to our events 15, 7, and 11. Evidently, the remaining eight events from our analysis of the first two INT seasons did not satisfy their rather severe selection criteria.
In Belokurov et al. (2005), the POINT-AGAPE collaboration analysed data from the first three INT observing seasons without any restrictions on the event duration. Using different selection criteria from their previous analyses, they found three high quality candidates. Two of these events were already known (PA-00-S4 or MEGA-ML-11 and PA-00-S3). The one new event is present in our survey but does not pass our selection criteria because of a high $`\chi ^2`$. The lightcurve for this event, along with our best-fit model, is shown Fig. 7. The model does not do a good job of reproducing the observed lightcurve behaviour. In particular, the observed lightcurve appears to be asymmetric about the peak time $`t_{\mathrm{max}}`$. The observed r-lightcurve is systematically below the model 15-20 days prior to $`t_{\mathrm{max}}`$. Both r and i lightcurves are above the model 10-15 days after $`t_{\mathrm{max}}`$. Since there are no data available on the rising part of the peak, $`t_{\mathrm{max}}`$ is poorly constrained and may in fact be less than the $`770\mathrm{days}`$ used in the fit. The shape of our r lightcurve is similar to the one presented in Belokurov et al. (2005) (NB. They removed one epoch close to the peak centre that is present in our lightcurve.) In i the peak shapes are somewhat different.
Peak asymmetries can be caused by secondary effects such as parallax. In our opinion, a more likely explanation for this case is that the event is a nova-like eruptive variable. Granted, the event appears to be achromatic. But classical novae can be achromatic on the declining part of the lightcurve (see, for example, Darnley et al. (2004)), precisely where there is data. If this is a classical nova, it would be a very fast one, with a decline rate corresponding to $``$0.6 mag per day.
Calchi Novati et al. (2005) found six candidate microlensing events in an analysis of the three-year INT data set. Of these events, four are the same as reported by Paulin-Henriksson et al. (2003) and two are new events: PA-00-N6 and PA-99-S7. The latter of these is located in the bright part of the southern field excluded in our analysis (Fig. 1). Candidate event PA-00-N6 is present in our data, but was only detected in one epoch in our automatic SExtractor residual detection step and therefore did not make it into the catalogue of variable sources. Calchi Novati et al. (2005) do not detect our events 1, 2, 3, 8, 9, 10, 13, and 14, which all peak in the first two observing seasons. Evidently, these events do not satisfy their S/N constraints.
## 6 Detection efficiency
We determine the detection efficiency for microlensing events by applying the selection criteria from Sect. 4 to the catalogue of artificial events from Sect. 3. As discussed above, simulated lightcurves are generated by adding artificial events to the difference images and then passing the images through the photometry analysis routine designed for the actual data. Those lightcurves that satisfy the selection criteria for microlensing form a catalogue of simulated detectable microlensing events. The detection efficiency is the ratio of the number of these events to the original number of artificial events.
We first check that our artificial event catalogue includes the portion of the source luminosity function responsible for most of the detectable events. The function $`N_{}\beta _\mathrm{u}^2`$ in Fig. 2 is meant to give a qualitative picture of the detectability of microlensing as a function of source luminosity. Here we consider the function $`P_{\mathrm{det}}N_{}\beta _\mathrm{u}^2ϵ`$ where $`ϵ`$ is detection efficiency as a function of $`F_{0,r}`$ integrated over $`\beta `$, $`t_{\mathrm{FWHM}}`$ and position. $`P_{\mathrm{det}}`$ gives the relative probability for detection of a microlensing event as a function of the source luminosity. As shown in Fig. 8, the range $`0.01`$ to $`10\mathrm{ADU}\mathrm{s}^1`$ adequately covers the peak of this probability distribution.
Our goal is to represent the detection efficiency in terms of a simple portable function of a few key parameters. We adopt a strategy whereby the detection efficiency is modelled as functions of $`t_{\mathrm{FWHM}}`$ and $`\mathrm{\Delta }F_{\mathrm{max}}`$ for individual subregions of the two fields. The parameters $`\beta `$ and $`t_{\mathrm{max}}`$ are “integrated out” and $`𝒞`$ is fixed to the value $`0.75`$. This strategy is motivated by the following considerations.
In Fig. 9 we plot the detection efficiencies as a function of $`\beta `$ for four different values of $`t_{\mathrm{FWHM}}`$. In each of the panels, the efficiencies are integrated over position within a single chip of the INT fields. The top (bottom) panels are for the south-east chip of the north (south) field. The right (left) panels are for bright (faint) source stars. The general trend is for the detection efficiency to increase with increasing $`t_{\mathrm{FWHM}}`$ and decreasing $`\beta `$. This trend is expected since longer duration events are more likely to be observed near the peak and smaller values of $`\beta `$ imply larger amplification factors. For $`F(r)=10\mathrm{ADU}\mathrm{s}^1`$, $`t_{\mathrm{FWHM}}10\mathrm{days}`$ and small $`\beta 0.7`$, the detection efficiencies decrease with decreasing $`\beta `$. The decrease is more severe for the $`t_{\mathrm{FWHM}}=50`$ day events where the detection efficiency actually drops below that for the $`t_{\mathrm{FWHM}}=10`$ day events. The problem may be that we underestimate the photometric error at high fluxes therefore causing $`\chi ^2`$ to be systematically high. Moreover, $`50`$ days is a substantial fraction of the observing season and therefore some long duration events may not meet the requirement that the peak be entirely within a single season.
Since the shape of the microlensing lightcurve does not depend strongly on $`\beta `$ we expect no significant dependence of the detection efficiency on the intrinsic source brightness. This point is illustrated in Fig. 10 where we plot the detection efficiencies as a function of $`1/\mathrm{\Delta }F_{\mathrm{max}}`$ for events with $`t_{\mathrm{FWHM}}=50`$ days. We integrate the efficiencies over positions within single CCDs and show the results for four of the eight CCDs in our fields. The curves vary by at most 30% over three orders of magnitude in $`F(r)`$. The implication is that an explicit $`F(r)`$ dependence in the detection efficiency will not change the results significantly.
We next test whether the detection efficiency depends on the colour $`𝒞`$ of the source. In addition to the main artificial event catalogue, we generate artificial events with $`𝒞=1.25`$ and r unchanged for a part of the north field. Fig. 11 compares the detection efficiencies for the two colours and shows that there is no significant difference, except for the very highest signal to noise events. The discrepancy at high S/N reflects the problem discussed above with our estimates of the photometric errors at high flux. This problem is worse for redder sources which have a higher i-band flux.
Motivated by the shapes of the curves in Fig. 11, we choose a Gaussian in $`1/\mathrm{\Delta }F_{\mathrm{max}}`$ where the position of the peak depends on $`t_{\mathrm{FWHM}}`$. The explicit functional form is taken to be:
$$ϵ=c_1\left(1t_{\mathrm{FWHM}}/112\right)e^{c_2(1/\mathrm{\Delta }F_{\mathrm{max}}c_3)^2}$$
(11)
where
$$c_3=d_1\mathrm{ln}(t_{\mathrm{FWHM}})+d_2.$$
(12)
The factor multiplying the Gaussian takes into account the sharp decrease in detection efficiency for events with duration comparable to or longer than the observing season. The parameters $`c_1`$, $`c_2`$, $`d_1`$ and $`d_2`$ are determined by fitting simultaneously the detection efficiencies for all values of $`t_{\mathrm{FWHM}}`$ to equation 11. Fig. 12 shows an example of these fitting formulae to the detection efficiencies.
Fig. 10 illustrates the dependence of the detection efficiencies on location in the INT fields. This dependence is due mainly to variations in galaxy surface brightness but also to the presence of bad pixels and saturated-star defects. As discussed above, we account for the spatial dependence by fitting the detection efficiency separately for subregions of the fields. To be precise, we divide each chip into 32 subregions, $``$3′$`\times `$3′ in size. For each of these regions we average 14 640 simulated events (2 440 per choice of $`t_{\mathrm{FWHM}}`$).
## 7 Extinction
Microlensing surveys such as MEGA and POINT-AGAPE are motivated, to a large extent, by the argument that a MACHO population in M31 would induce a near-far asymmetry in the microlensing event distribution. In the absence of either extinction or significant intrinsic asymmetries in the galaxy, the distribution of self-lensing events and variable stars masquerading as microlensing events would be near-far symmetric. The detection of a near-far asymmetry would then provide compelling evidence in favour of a significant MACHO population.
Recently, An et al. (2004a) found a near-far asymmetry in the distribution of variables which they attribute to differential extinction across the M31 disk. That differential extinction is significant is also witnessed by several dust features including two prominent dust lanes on the near side of the disk.
We construct a simple model for differential extinction in M31 and test it to against the distribution of LPVs. In the next section, we incorporate this extinction model into our calculations for the theoretical event rate.
Following Walterbos & Kennicutt (1988) we assume that the dust is located in a thin layer in the mid-plane of the disk. Along a given line-of-sight, only light from behind the dust layer is absorbed. Because of the galaxy’s high inclination, the fraction of stars located behind the dust layer is higher for lines-of-sight on the near side of the disk than for those on the far side, as illustrated in Fig. 13. Therefore, even if the distribution of dust is intrinsically symmetric, extinction will have a greater effect on the near side of the disk.
Based on these assumptions the observed intensity along a particular line-of-sight is
$$_{\mathrm{obs}}=_{\mathrm{front}}+_{\mathrm{back}}e^\tau $$
(13)
where $`_{\mathrm{front}}`$ ($`_{\mathrm{back}}`$) is the intensity of light originating from in front of (behind) the dust layer and $`\tau `$ is the optical depth. This equation can be rewritten in terms of the total intrinsic intensity, $`_{\mathrm{intr}}`$, and the fraction $`x`$ of light that originates from in front of the dust layer:
$$_{\mathrm{obs}}=x_{\mathrm{intr}}+(1x)_{\mathrm{intr}}e^\tau .$$
(14)
The three unknowns in this equation, $`_{\mathrm{intr}}`$, $`x`$, and $`e^\tau `$, depend on wavelength. Rewriting equation 14 for the B-band we have
$$e^{\tau _B}=\frac{_{\mathrm{obs}}(B)/_{\mathrm{intr}}(B)x_B}{1x_B}.$$
(15)
As a first approximation we assume that $`_{\mathrm{obs}}(I)=_{\mathrm{intr}}(I)`$ so that
$$e^{\tau _B}=\frac{_{\mathrm{obs}}(B)/(𝒞_{BI}_{\mathrm{obs}}(I))x}{1x}$$
(16)
where $`𝒞_{BI}_{\mathrm{intr}}(B)/_{\mathrm{intr}}(I)`$ is the intrinsic $`IB`$ colour of the stellar population. An improved estimate of $`_{\mathrm{intr}}(I)`$ is obtained by transforming the extinction factor from B to I via the standard reddening law (Savage & Mathis, 1979). The calculation is repeated several times
We approximate $`x_B`$ and $`x_I`$ from a simple model of the galaxy wherein the intrinsic (i.e., three-dimensional) light distribution $`\eta \left(𝐱\right)`$ for the disk and bulge are taken to be double exponentials. In cylindrical coordinates for M31, we have
$$\eta ^i\left(𝐱\right)=\eta _0e^{r/h_R^i}e^{z/h_z^i}$$
(17)
where the superscript $`i`$ denotes either the disk or bulge, $`\eta _0`$ is a normalisation constant, and $`h_R`$ and $`h_z`$ are the radial and vertical scale lengths, respectively. Different scale lengths are used for B and I because the two bands have different sensitivities to young and old populations of stars. Young stars tend to lie closer to the disk mid-plane than old ones. Our choices for the parameters are given in Table 4. The values of the disk scale lengths and the bulge-to-disk-ratios are taken from Walterbos & Kennicutt (1988). The scale lengths for bulge are adapted from their de Vaucouleurs fit while the disk scale heights are based on the distribution of different stellar populations in the Milky Way disk. The observables $`_{\mathrm{obs}}(I)`$ and $`_{\mathrm{obs}}(B)`$ are from Guhathakurta et al. (2005) who cover a 1.7$`\mathrm{°}\times `$5$`\mathrm{°}`$ field centred on M31. We derive colour profiles from their mosaics which are found to be similar to the profiles in Walterbos & Kennicutt (1988). The colour is approximately constant within 30$`\mathrm{}`$ and becomes bluer at larger radii.
Our I-band extinction map for M31 is shown in Fig. 14. The major dust lanes are clearly visible in the northern field and, as expected, the derived extinction is much larger on the near side of the galaxy than on the far side. The I-band attenuation is $`<40\%`$ and reaches a maximum in the innermost dust lane and a few smaller complexes.
Our model almost certainly underestimates the effect of extinction across the M31 disk. The approximation $`_{\mathrm{obs}}(I)_{\mathrm{intr}}(I)`$ is a poor starting point in the limit of large optical depths. For $`\tau 1`$, most of the light in both B and I from behind the dust layer is absorbed and therefore $`_{\mathrm{obs}}(B)/_{\mathrm{obs}}(I)𝒞_{BI}`$. However substituting this result into equation 15 gives $`\mathrm{exp}\left(\tau \right)1`$, an obvious contradiction. By the same token, if the dust is distributed in high-$`\tau `$ clumps, then $`I`$ and $`B`$ wavelengths will be absorbed by equal amounts given essentially by the geometric cross section of the clumps. Moreover, the thin-layer approximation tends to yield an underestimate of the extinction factor (Walterbos & Kennicutt, 1988). Finally, scattering increases the flux observed towards the dust lanes and therefore also leads one to underestimate the extinction factor. Some of these problems can be solved by using infrared data in the construction of the extinction map. In a future paper we plan to use 2MASS data in order to derive a more accurate model for differential extinction in M31.
We can use the distribution of variable stars in our survey to test and refine the extinction model. The underlying assumption of this exercise is that the intrinsic distribution of variables is the same on the near and far sides of the disk. We begin by determining the periods of the variable stars using a multi-harmonic periodogram (Schwarzenberg-Czerny, 1996) suitably modified to allow for unevenly sampled data. A six-term Fourier series is then fit to each lightcurve yielding additional information such as the amplitude of the flux variations. Only variables with lightcurves that are well-fit by the Fourier series are used.
We will use LPVs to test the extinction model because they generally belong to quite old stellar populations. This is an advantage because the majority of the microlensing source stars also belong to older populations which are more smoothly distributed over the galaxy than younger variables such as Cepheids. We select LPVs with periods between 150 and 650 days and focus on two regions of our INT fields. One of these is located on the near-side of the disk where extinction is expected to be high while the other is located symmetrically about the M31 centre on the far side. Fig. 15 shows the spatial distribution of the LPVs. Since extinction reduces the amplitude of the flux variations and the average flux by the same factor we can study extinction by comparing the distributions in $`\mathrm{\Delta }F`$ for the near and far sides. These flux variation distributions are shown in Fig. 16. For low $`\mathrm{\Delta }F`$, where the shapes of the distributions are dominated by the detection efficiency, results for the near and far side agree. For high $`\mathrm{\Delta }F`$, where the detection efficiency for variables approaches 100%, one finds a large discrepancy between the near and far-side distributions.
To test whether this discrepancy is indeed due to extinction we transform the coordinates of LPVs on the far side to their mirror image on the near side. The amplitude of the flux variation is then reduced by the model extinction factor suitably transformed from I to r (Savage & Mathis, 1979). The new distribution, shown in Fig. 16, is still significantly above the near-side distribution at large $`\mathrm{\Delta }F`$ though it does provide a better match than the original far-side distribution. The implication is that our model underestimates extinction. To explore this point further we consider models in which $`\tau `$ is replaced by $`c\tau `$ where $`c>1`$. In Fig. 16, we show the distributions of the far side LPVs for $`\tau 2\tau `$ (long-dashed line) and $`\tau 2.5\tau `$ (dot-dashed line). Apparently, the bright end of the (mirror) far-side distribution with $`\tau `$ increased by a factor of 2.5 agrees with the bright end of the near-side distribution. We therefore conclude that our original model does indeed underestimate the effects of extinction. In some places this will be stronger than in others, but over the probed region the model underestimates extinction effectively by perhaps a factor of 2.5 in $`\tau `$.
## 8 Theoretical predictions
The detection efficiencies found in Sect. 6 allow us to predict the number and distribution of events given a specific model for the galaxy. Though M31 is one of the best studied galaxies, a number of the parameters crucial for microlensing calculations, are not well-known. Chief among these are the mass-to-light ratios of the disk and bulge, $`\left(M/L\right)_\mathrm{d}`$ and $`\left(M/L\right)_\mathrm{b}`$, respectively. The light distributions for these components are constrained by the surface brightness profile while the mass distributions of the disk, bulge, and halo are constrained by the rotation curve and line-of-sight velocity dispersion profile. However, the mass-to-light ratios are poorly constrained primarily because the shapes of the disk and halo contributions to the rotation curve are similar (e.g. van Albada et al., 1985). One can compensate for an increase in $`\left(M/L\right)_\mathrm{d}`$ by decreasing the overall density of the halo. Stellar synthesis models (Bell & de Jong, 2001), combined with observations of the colour profile of M31, can be used to constrain the mass-to-light ratios though these models come with their own internal scatter and assumptions. Another poorly constrained parameter is the thickness of the disk which affects the disk-disk self-lensing rate.
In this section we describe theoretical calculations for the expected number of events in the MEGA-INT survey. We consider a suite of M31 models which span a wide range of values in $`\left(M/L\right)_\mathrm{d}`$ and $`\left(M/L\right)_\mathrm{b}`$. The dependence of the microlensing rate on other parameters is also explored.
### 8.1 Self-consistent models of M31
The standard practice for modelling disk galaxies is to choose simple functional forms for the space density of the disk, bulge, and halo tuned to fit observational data. For microlensing calculations, velocity distributions are also required. Typically, one assumes that the velocity distribution for each of the components is isotropic, isothermal, and Maxwellian with a dispersion given by the depth of the gravitational potential or, in the case of the bulge, the observed line-of-sight velocity dispersion. (But see Kerins et al. (2001) where the effects of velocity anisotropy are discussed.) This approach can lead to a variety of problems. First, these “mass models” do not necessarily represent equilibrium configurations, that is, self-consistent solutions to the collisionless Boltzmann and Poisson equations. A system initially specified by the model may well relax to a very different state. Another issue concerns dynamical instability. Self-gravitating rotationally supported disks form strong bars. This instability may be weaker or absent altogether if the disk is supported, at least in part, by the bulge and/or halo. Therefore, models with very high $`\left(M/L\right)_\mathrm{d}`$ are the most susceptible to bar formation and can be ruled out.
In order to overcome these difficulties we use new, multi-component models for disk galaxies developed by Widrow & Dubinski (2005). The models assume axisymmetry and incorporate an exponential disk, a Hernquist model bulge (Hernquist, 1990), and an NFW halo (Navarro et al., 1996). They represent self-consistent equilibrium solutions to the coupled Poisson and collisionless Boltzmann equations and are generated using the approach described in Kuijken & Dubinski (1995).
The phase-space distribution functions (DFs) for the disk, bulge, and halo ($`f_{\mathrm{disk}},f_{\mathrm{bulge}},`$ and $`f_{\mathrm{halo}}`$ respectively) are chosen analytic functions of the integrals of motion. For the axisymmetric and time-independent system considered here, the angular momentum about the symmetry axis, $`J_z`$, and the energy, $`E`$, are integrals of motion. Widrow & Dubinski (2005) assume that $`f_{\mathrm{halo}}`$ depends only on the energy while $`f_{\mathrm{bulge}}`$ incorporates a $`J_z`$-dependence into the Hernquist model DF to allow for rotation. For both halo and bulge, the DFs are “lowered” as with the King model (King, 1966) so that the density goes to zero at a finite “truncation” radius. The disk DF is a function of $`E`$, $`J_z`$, and an approximate third integral of motion, $`E_z`$, which corresponds to the energy associated with vertical motions of stars in the disk (Kuijken & Dubinski, 1995).
Self-consistency requires that the space density, $`\rho `$, and gravitational potential, $`\psi `$, satisfy the following two equations:
$$\rho =d^3v\left(f_{\mathrm{disk}}+f_{\mathrm{bulge}}+f_{\mathrm{halo}}\right)$$
(18)
and
$$^2\psi =4\pi G\rho .$$
(19)
Self-consistency is achieved through an iterative scheme and spherical harmonic expansion of $`\rho `$ and $`\psi `$. Straightforward techniques allow one to generate an N-body representation suitable for pseudo-observations of the type described below. The N-body representations also provide very clean initial conditions for numerical simulations of bar formation and disk warping and heating.
The DFs are described by 15 parameters which can be tuned to fit a wide range of observations. In addition, one must specify mass-to-light ratios if photometric data is used. Our strategy is to compare pseudo-observations of M31 with actual observational data to yield a $`\chi ^2`$-statistic. Minimisation of $`\chi ^2`$ over the model parameter space – performed in Widrow & Dubinski (2005) by the downhill simplex method (see e.g. Press et al., 1992) – leads to a best-fit model.
Following Widrow & Dubinski (2005) (see, also Widrow et al. (2003) who carried out a similar exercise with the original Kuijken & Dubinski (1995) models) we utilise measurements of the surface brightness profile, rotation curve, and inner (that is, bulge region) velocity profiles. We use R-band surface brightness profiles for the major and minor axes from Walterbos & Kennicutt (1988). (Widrow & Dubinski (2005) used the global surface brightness profile from Walterbos & Kennicutt (1988) which was obtained by averaging the light distribution in elliptical rings. The use here of both major and minor axis profiles should yield a more faithful bulge-disk decomposition.) The theoretical profiles are corrected for internal extinction using the model described in the previous section. In addition, a correction for Galactic extinction is included. We assume photometric errors of $`0.2\mathrm{mag}`$. We use a composite rotation curve constructed from observations by Kent (1989) and Braun (1991) that run from 2 to 25 kpc in galactocentric radius. Values and error bars for the circular speed are obtained at intervals of $`10\mathrm{arcmin}2.2\mathrm{kpc}`$ using kernal smoothing (Widrow et al., 2003). Finally, we use kinematic measurements from McElroy (1983) to constrain the dynamics in the innermost part of the galaxy. We smooth his data along the minor axis to give values for the line-of-sight stellar rotation and velocity dispersion at $`0.5\mathrm{kpc}`$ and $`1.0\mathrm{kpc}`$. The values at these radii are insensitive to the effects of a central supermassive object and reflect the dynamics of the bulge stars with little disk contamination (McElroy, 1983). An overall $`\chi ^2`$ for the model is calculated by combining results from the three types of data. Photometric and kinematic data are given equal weight; the circular rotation curve measurements are weighted more heavily than the bulge velocity and dispersion measurements. To be precise, we use
$$\chi ^2=\frac{1}{\sqrt{2}}\left(\chi _{\mathrm{sbp}}^2+\frac{1}{3}\chi _{\mathrm{disp}}^2+\frac{2}{3}\chi _{\mathrm{rc}}^2\right)$$
(20)
where $`\chi _{\mathrm{sbp}}^2`$, $`\chi _{\mathrm{bulge}}^2`$, and $`\chi _{\mathrm{rc}}^2`$ are the individual $`\chi ^2`$-statistics for the photometric, bulge kinematics, and rotation curve measurements.
Our reference model (model A1) is constructed with $`\left(M/L\right)_\mathrm{d}=2.4`$ and $`\left(M/L\right)_\mathrm{b}=3.6`$. These values are motivated by the stellar population synthesis models of Bell & de Jong (2001). Along the far side of the minor axis, where the surface brightness profile is relatively free of extinction, the $`BR`$ colour is $`1.8`$ in the bulge region and $`1.6`$ in the disk region Walterbos & Kennicutt (1988). A correction for Galactic extinction brings these numbers down by $`0.18`$. Substituting into the appropriate formula from Table 1 of Bell & de Jong (2001) yield the mass-to-light ratios chosen for this model. In Fig. 17 we compare predictions for model A1 with observations. Shown are the surface brightness profiles along major and minor axes and the circular rotation curve. Not shown is the excellent agreement between model and observations for the stellar rotation and dispersion measurements in the bulge region. The reduced $`\chi ^2`$ statistic for this model is $`1.06`$ (see Table 5).
In model A1, the scale height of the disk was fixed to a value of $`1.0\mathrm{kpc}`$. Note that our model uses a $`\mathrm{sech}^2`$-law for the vertical structure of the disk. A $`\mathrm{sech}^2`$-scale height of $`1\mathrm{kpc}`$ is roughly equivalent to an exponential scale height of $`0.50.7\mathrm{kpc}`$. The observations used in this study do not provide a tight constraint on the scale height of the disk and so we appeal to observations of edge-on disk galaxies. Kregel et al. (2002) studied correlations between the (exponential) vertical scale height and other structural parameters such as the radial scale height and asymptotic circular speed in a sample of 34 edge-on spirals. Using these correlations we arrive at an exponential scale height for $`M31`$ of $`0.6\mathrm{kpc}`$ with a fairly large scatter.
We also fix the disk truncation radius for this model to $`28\mathrm{kpc}`$ which is at the high end of the range favoured in Kregel et al. (2002). Lower values appear to be inconsistent with the measured surface brightness profile. The remaining parameters for the disk, bulge, and halo DFs are varied in order to minimise $`\chi ^2`$.
Table 5 outlines other models considered in this paper. Models B1-E1 explore the $`\left(M/L\right)_\mathrm{b}\left(M/L\right)_\mathrm{d}`$ plane. The $`\chi ^2`$ for these models are generally quite low, a reflection of the model degeneracy mentioned above. In these models, disk and bulge “mass” are traded off against halo mass. Previous investigations (Widrow & Dubinski, 2005) suggest that model E1 is unstable to the formation of a strong bar while the other models are stable against bar formation or perhaps allow for a weak bar.
The aforementioned models used values for the extinction factor derived in Sect. 7. As discussed in that section, there are a number of reasons to expect that this model underestimates the amount of extinction in M31. Indeed, our analysis of the near-far asymmetry in LPVs favours a higher optical depth by a factor of $`2.5`$, that is, the substitution $`e^\tau e^{2.5\tau }`$. For this reason, we consider a parallel sequence of models, A2-F2, with high extinction. Note that the $`\chi ^2`$ for these models are as good as if not better than those for the corresponding low-extinction models.
### 8.2 Event rate calculation
The event rate is calculated by performing integrals over the lens and source distribution functions. The rate for lenses to enter the lensing tube of a single source is
$$d^5R=\frac{f_l(l_l,𝐯_l)}{_l}2R_Ev_{}dl_ld𝐯_ld\beta $$
(21)
where $`f_l`$ is the DF for the lens population, $`l_l`$ is the observer-lens distance ($`D_{\mathrm{OL}}`$ in the notation of equation 2), $`v_{}`$ is the transverse velocity of the lens with respect to the observer-source line-of-sight, and $`_l`$ is the mass of the lens. In writing this equation, we assume all lenses have the same mass.
For a distribution of sources described by the DF $`f_s`$, equation 21 is replaced by the following expression for the rate per unit solid angle
$`{\displaystyle \frac{dR}{d\mathrm{\Omega }}}={\displaystyle \frac{f_l(l_l,𝐯_l)}{_l}\frac{f_s(l_s,𝐯_s)}{\left(M/L\right)_sL_s}2R_Ev_{}}`$
$`\times dl_ld𝐯_ll_s^2dl_sd𝐯_sd\beta `$ (22)
where $`l_s`$ is the observer-source distance, $`\left(M/L\right)_s`$ is the mass-to-light ratio of the source and $`L_s`$ is the source luminosity. (For the moment, we treat all sources as being identical.)
We perform the integrals using a Monte Carlo method. The DFs are sampled at discrete points:
$$f_p(l_p,𝐯_p)=\frac{\mathrm{\Sigma }_p}{N_p}\underset{i=1}{\overset{N_p}{}}\delta (l_pl_i)\delta (𝐯_p𝐯_i)$$
(23)
where $`p\{l,s\}`$, $`\mathrm{\Sigma }_p`$ is the surface density of either lens or source population, and $`N_p`$ is the number of points used to Monte Carlo either lens or source populations. The nine-dimensional integral in equation 22 is replaced by a double sum and an integral over $`\beta `$:
$$\frac{dR}{d\mathrm{\Omega }}=𝒮_{sl}\underset{i,j}{}_0^{\beta _u}𝑑\beta _{ij}$$
(24)
where
$$𝒮_{sl}=\frac{\mathrm{\Sigma }_l\mathrm{\Sigma }_s}{N_l_lN_sL_s\left(M/L\right)_s}$$
(25)
and
$$_{ij}(2R_Ev_{})_{ij}l_j^2.$$
(26)
Note that $`𝒮`$ depends on the line of sight densities of the lens and source distributions along with characteristics of the two populations. $`_{ij}`$ depends on the coordinates and velocities of the lens and source (hence the $`ij`$ subscripts). The sum is restricted to lens-source pairs with $`l_l<l_s`$. For each lens-source pair, the Einstein crossing time, $`t_{\mathrm{E},ij}`$ is easily calculated. The differential event rate is then
$$\frac{d^2R}{d\mathrm{\Omega }dt_E}=𝒮_{sl}\underset{i,j}{}_0^{\beta _u}𝑑\beta _{ij}\delta (t_{\mathrm{E},ij}t_\mathrm{E}).$$
(27)
### 8.3 Stellar and MACHO populations
The formulae in the previous section apply to the six lens-source combinations in our model: disk-disk, disk-bulge, bulge-disk, bulge-bulge, halo-disk, and halo-bulge. As written the formulae assume homogeneous populations. For the disk and bulge populations, we modify equation 27 to include integrals over the mass and luminosity functions as appropriate. We write the luminosity function (LF) as
$$\frac{dN}{dM_R}=Ag(M_R)$$
(28)
and the mass function as
$$\frac{dN}{d}=Bh(,_0)$$
(29)
where $`A`$ and $`B`$ are normalisation constants and $`_0`$ is the lower bound for the mass function (MF). We take the function $`g`$ from Mamon & Soneira (1982) and the function $`h`$ from Binney & Merrifield (1998) (their equation 5.16) with the power-law form $`dN/d^{1.8}`$ extended to $`_0`$. $`A`$ and $`B`$ are evaluated separately for the disk and bulge populations. In the case of the disk, we assume that $`30\%`$ of the mass is in the form of gas. The LF is normalised to give $`\overline{L}=L_{}`$ with the proviso that $`L_s`$ in equation 22 is given in solar units. To determine the normalisation constant $`B`$ of the mass function, we write
$$Bh(_{},_0)=\left(\frac{dN}{dM_V}\frac{dM_V}{d}\right)|_=_{}$$
(30)
where the V-band LF is again from Mamon & Soneira (1982) and $`dM_V/d`$ is from Kroupa et al. (1993). Equation 30 is evaluated at solar values for convenience. The relation
$$\left(\frac{M}{L}\right)_R=\frac{Bh(,_0)𝑑}{Ag(M_R)L(M_R)𝑑M_R}$$
(31)
can then be solved for $`_0`$. Thus, a disk with high $`M/L`$ contains more low-mass stars than a disk with low $`M/L`$.
For simplicity, and because we lack a model for what MACHOs actually are, we assume all MACHOs have the same mass, $`_\mathrm{M}`$ that is
$$\frac{dN}{d}=\delta \left(_M\right).$$
(32)
The value of $`_\mathrm{M}`$ will directly determine the number density of MACHOs for a given halo mass density. Since the MACHOs only provide lenses and no sources for microlensing, a higher value of $`_\mathrm{M}`$ and thus a lower number density, will result in a lower number of microlensing events. A given value of $`_\mathrm{M}`$ can be considered as the average mass of a more elaborate MACHO mass function.
### 8.4 Theoretical prediction for the number of events
Recall that the efficiency $`ϵ`$ is written as a function of $`t_{\mathrm{FWHM}}`$ and $`\mathrm{\Delta }F_{\mathrm{max}}`$. (The efficiency also depends on the line of sight.) These quantities are explicit functions of $`\beta `$, $`F_r`$, and $`t_\mathrm{E}`$. Thus, the expected number of events per unit solid angle is
$`{\displaystyle \frac{d}{d\mathrm{\Omega }}}=EAB𝒮_{ls}{\displaystyle \underset{i,j}{}}{\displaystyle _0^{\beta _u}}𝑑\beta {\displaystyle 𝑑M_Rg\left(M_R\right)}`$
$`\times {\displaystyle }d_lh(,_0)_{ij}ϵ(t_{\mathrm{FWHM}},\mathrm{\Delta }F)`$ (33)
where $`E`$ is the overall duration of the experiment. Our survey covers four half-year seasons and so, with our choice of units for $`ϵ`$ and $`dR/d\mathrm{\Omega }`$, we have $`E=2`$.
The number of events expected in each of the 250 bins used for the extinction calculation and labelled by “k” is
$$_k=\mathrm{\Delta }\mathrm{\Omega }\left(\frac{dE}{d\mathrm{\Omega }}\right)_k$$
(34)
where $`\mathrm{\Delta }\mathrm{\Omega }=9\mathrm{arcmin}^2`$ is the angular area of a bin. $`_k`$ carries an additional label (suppressed for notational simplicity) which denotes the lens-source combination. The total number of events is $`=_k`$.
### 8.5 Binary lenses
Our microlensing selection criteria are based on the assumption that the lenses are single point-mass objects. However, at least half of all stars are members of multiple star systems. Microlensing lightcurves for a lens composed of two or more point masses can deviate significantly from the standard lightcurve (Schneider & Weiss, 1986) and may therefore escape detection. The deviations are strongest when the source crosses or comes close to the so-called caustics, positions in the source plane where the magnification factor is formally infinite. (The actual magnification factor is finite due to the finite size of the source.) The size of the caustic region is largest when the separation of the components of the lens is comparable to the Einstein radius corresponding to the total mass (equation 2). Mao & Paczynski (1991) estimated that $``$10% of microlensing events towards the bulge of the Milky Way (mainly self-lensing events) should show strong binary characteristics such as caustic crossings. Since the Einstein radius for bulge-bulge self-lensing toward the Milky Way and M31 are comparable, we can expect a similar 10% effect in our survey. Baltz & Gondolo (2001) perform a similar analysis for pixel-lensing surveys and estimate that in the order of 6% of self-lensing events from normal stellar populations will exhibit caustic crossings. Since the majority of detected events will have low signal-to-noise, we can assume that deviations other than caustic crossings in most cases will not strongly affect our detection efficiency. Therefore, to account for binary lenses, the calculated theoretical predictions for self-lensing are revised downward by $`10\%`$.
### 8.6 Results
Table 5 presents the theoretical predictions for the total number of events expected in the MEGA-INT four-year survey. The results are given for both self-lensing ($`_{\mathrm{self}}`$) and halo lensing ($`_{\mathrm{halo}}`$). The values quoted for $`_{\mathrm{halo}}`$ assume 100% of the halo is in the form of MACHOs. In other words, these values should be multiplied by the MACHO halo fraction in order to get the expected number of events for a MACHO component. We note that lensing by the Milky Way halo is not included in these results. This possible contribution is expected to be small, since the number of microlensing events from a 100% MW halo is a few times lower than for a 100% M31 halo (Gyuk & Crotts, 2000; Baillon et al., 1993) for MACHO masses around 0.5M.
We also consider the near-far asymmetry for self and halo lensing. In Fig. 18, we show the cumulative distribution of events for self and halo lensing as a function of the distance from the major axis, $`s`$. We take $`s`$ to be positive on the far side of the disk. For this plot, we choose model A1 but since the distributions are normalised to give 14 total events, the difference between the models is rather inconsequential. We see that both self and halo lensing models do a good job of describing the event distribution in the inner $`0.2^{}`$. The halo distribution does a somewhat better job of modelling the three events between $`s=0.2^{}`$ and $`s=0.3^{}`$. Neither halo nor self lensing models predict anywhere near two events for $`s>0.35^{}`$.
To further explore the distribution, we define the asymmetry parameter $`𝒜`$:
$$𝒜=\frac{_ks_k}{}.$$
(35)
In Table 5 we give values for $`𝒜_{\mathrm{self}}`$ and $`𝒜_{\mathrm{halo}}`$. We also provide an average $`𝒜_{\mathrm{ave}}`$ which assumes that MACHOs make up the shortfall between the expected number of events and the observed value of 14. In cases where the expected number of events is greater than 14, we set $`𝒜_{\mathrm{ave}}=𝒜_{\mathrm{self}}`$. The asymmetry parameter for the 14 candidate events is $`𝒜_{\mathrm{data}}=0.125`$.
The general trend, in terms of total expected number of events, is that as the mass-to-light ratios are increased, $`_{\mathrm{self}}`$ increases and $`_{\mathrm{halo}}`$ decreases. There are counter examples. In model C1, the $`\left(M/L\right)_\mathrm{b}`$ (as compared with model A1) leads to a less massive disk and lower $`_{\mathrm{self}}`$. Recall that for each choice of mass-to-light ratios, the remaining parameters are adjusted to minimise $`\chi ^2`$. The process can lead to rather complicated interdependencies between the model parameters. The self-lensing rate decreases with decreasing $`h_z`$ as illustrated with model F1. The self-lensing rate is generally reduced in the high extinction models relative to the low extinction ones. Finally we see that the halo event rate decreases with increasing MACHO mass. Models G and H illustrate this point and span the range in $`_M`$ identified by Alcock et al. (2000) as the most probable mass range for Milky Way MACHOs.
The timescale distribution is easily calculated using the method outlined in the previous section. Essentially, one calculates $`t_{\mathrm{FWHM}}`$ for each lens-source pair in the Monte Carlo sum. In Fig. 19 we show the cumulative timescale distribution of our candidate microlensing event sample and model A1. In constructing the curves for self and halo lensing, we have scaled the distributions to give a total of 14 events.
## 9 Discussion
The numbers expected for events due to self-lensing across the models probed in Table 5 fall within the narrow range of 10-16. The relative insensitivity of $`_{\mathrm{self}}`$ to changes in the mass-to-light ratios is a result of our approach to constructing models; changes in $`\left(M/L\right)_\mathrm{b}`$ and $`\left(M/L\right)_\mathrm{d}`$ are compensated by changes in the structural parameters of the disk, bulge, and halo so as to minimise $`\chi ^2`$ for the fit to the rotation curve and surface brightness data. Consider models D1 and E1. The mass-to-light ratios differ by a factor of $`2`$ while $`_{\mathrm{self}}`$ differs by only a factor of 1.4; with the low $`M/L`$ values in model D1, the rotation curve data drive up the disk and bulge luminosity distributions at the expense of a poorer fit for the photometric data. A balance is struck and the net result is that the change in $`_{\mathrm{self}}`$ is significantly smaller than what one might expect.
The consistency of the number of candidate events with the number of predicted self-lensing events is contrary to the results of the analysis of the first three seasons of INT data by the POINT-AGAPE collaboration. Calchi Novati et al. (2005) present six high quality, short duration microlensing candidates with one of these events attributed to M32-M31 lensing. They also model the detection efficiency and calculate number of expected self- and halo-lensing events for a variety of M31 models. In all of their models, the number of events for self-lensing is predicted to be less than $`1.5`$. Since this number is significantly less than the observed number, they conclude that some of the events are due to MACHOs and estimate that the MACHO halo fraction is at least 20%.
Calchi Novati et al. (2005) use the model from Kerins et al. (2001) which features a bulge following Kent (1989), an exponential $`sech^2`$ disk and a spherical, nearly isothermal halo. They use the same structural parameters for the three components as Kerins et al. (2001) but take $`(M/L_B)_b=3`$ and $`(M/L)_d=4`$. This model for the stellar mass distribution in M31 predicts an inner rotation curve that is significantly lower than the observed one, and so an extra ‘dark bulge’ component is required as well as the isothermal halo. Calchi Novati et al. (2005) do not consider microlensing by this dark bulge in their model, but instead attribute all surplus microlensing to the halo. In our model the stellar bulge is more massive, with $`M/L`$ that is sufficient to reproduce the inner rotation curve, and there is no non-lensing dark bulge component. It appears to provide sufficient microlensing events to explain the observations.
Furthermore, the choice of $`0.3\mathrm{kpc}`$ for the $`sech^2`$ scale height is small by perhaps a factor of 3 if M31 is a typical spiral galaxy as represented in the survey by Kregel et al. (2002). Thickening the disk increases the disk-disk self-lensing rate.
For our models, the number of events due to self-lensing is consistent with the total number of events observed but not inconsistent with a significant MACHO fraction for the halo of M31. We can make this statement more quantitative by treating halo events as a Poisson process with background due to self-lensing and employing the approach of Feldman & Cousins (1998). We let $`n`$ be the number of observed events consisting of MACHO events with mean $`f_{\mathrm{halo}}`$, where $`f`$ is the MACHO fraction, and a background due to self-lensing with known mean $`_{\mathrm{self}}`$. For this analysis, we ignore the background due to variables and background supernovae. The probability distribution function is
$$P\left(n|f\right)=\left(f_{\mathrm{halo}}+_{\mathrm{self}}\right)^n\mathrm{exp}\left[\left(f_{\mathrm{halo}}+_{\mathrm{self}}\right)\right]/n!.$$
(36)
To obtain confidence intervals for $`f`$:
1. Calculate $`P\left(n|f\right)`$ for $`N`$ values of $`f\{0,1\}`$ and sort from high to low. The maximum of $`P`$ defines the most probable value of $`f`$. The values of $`P`$ are normalised so that the sum of all sampled values of $`P`$ is 1.
2. Accept values of $`f`$ starting from the highest value of $`P`$ until the sum of $`P`$ exceeds the desired confidence level. The largest and smallest values of accepted $`f`$ define the confidence interval.
In Table 6 we provide most probable values of $`f`$ and 95% confidence intervals for all of the models in Table 5. We provide these values both for the case of the full sample of 14 observed candidate events ($`n`$=14), as well as for the case of 11 observed events ($`n`$=11), for reasons discussed below.
We next turn to the distribution of events across the M31 disk as represented by the asymmetry parameters. From Table 5 we see that $`𝒜_{\mathrm{self}}<𝒜_{\mathrm{halo}}<𝒜_{\mathrm{data}}`$. The (weak) asymmetry in the self-lensing distribution is due to extinction. Note that the values are significantly below $`𝒜_{\mathrm{data}}`$ even for the high extinction models.
The asymmetry parameter for the halo is significantly higher than that for self-lensing events and close to, though still below, $`𝒜_{\mathrm{data}}`$. However, the asymmetry parameter for combinations of self and halo lensing are well below $`𝒜_{\mathrm{data}}`$. Evidently, the distribution of candidate events is difficult to explain with any reasonable combination of self and halo lensing.
The large asymmetry in the data is due, for the most part, to events 11, 13, and 14 (see Table 7). It is therefore worth considering alternative explanations for these events. As argued in Paulin-Henriksson et al. (2002), the lens for event 11 likely resides in M32 and since we have not included M32 in our model, this event should be removed from the analysis. Doing so leads to a modest reduction in $`𝒜_{\mathrm{data}}`$.
Events 13 and 14 may be more difficult to explain. For model A1, the predicted number of self-lensing events with $`s>s\left(\mathrm{event}\mathrm{\hspace{0.17em}18}\right)`$ is $`0.005`$ while the predicted number of MACHO events in the same range in $`s`$ is $`0.14f`$. Thus, the probability of having two events either from self or halo lensing is exceedingly small, unless the halo fraction is very large. However, since some contamination by variable stars of our sample can not be excluded, one or both of these events may be a variable star. We note, for example, that event 13 has the lowest S/N in our sample. The probability of having one event for MACHO lensing with $`f=0.20`$ is $`3\%`$, small, but not vanishingly so.
A closer inspection of the model is also warranted. Recall that our models assume axisymmetry whereas M31 exhibits a variety of non-axisymmetric features such as disk warping. This point is illustrated in the isophotal map by Hodge & Kennicutt (1982). From the map, one finds that event 13 lies on the B=24 (R=22.6) contour while model A1 predicts R=23.5. Thus, the model may in fact underestimate the surface brightness of the disk by a factor of 2, and hence the disk-disk self-lensing rate by a factor of 4. (The reason for the discrepancy is not completely clear. The contours on the far side do appear to be “boxier” than those predicted by the model.)
It is interesting to note that events 13 and 14 are coincident with the location of the giant stellar stream discovered by Ibata et al. (2001). This stream runs across the southern INT field, approximately perpendicular to the major axis and over M32. Indeed, M32 may be the progenitor of the stream (Merrett et al., 2003). The average V-band surface brightness of the stream is $`\mathrm{\Sigma }_V30\pm 0.5`$ mag arcsec<sup>-2</sup> (Ibata et al., 2001) but this is measured far from the projected positions of events 13 and 14. The surface brightness of the stream might be significantly higher near the position of M32. Perhaps the most conservative statement one can make about the stream is that it is not bright enough to distort the contours near events 13 and 14, that is, it cannot be brighter than the disk at these radii. The microlensing event rate due to stars in the stream is of course enhanced relative to the rate for self-lensing by the ratio of the distance from the stream to the disk and the thickness of the disk, that is, by a factor of $`20`$. The stream-disk lensing rate might be further enhanced if the stars in the stream have a large proper motion relative to the disk. These arguments suggest that the number of stream-disk events in the vicinity of M32 might be $`0.030.1`$; perhaps high enough to explain one event.
Fig. 20 provides a summary of our results with respect to the expected number of events and the asymmetry parameter. The points with error bars represent the data for the 4 cases considered in Table 7. The solid circles and lines correspond to the high extinction case; the open circles and dotted lines correspond to the low extinction case. The circles assume pure self lensing while the lines trace out the values for increasing MACHO fraction with the tick-mark indicating the position of $`f=0.2`$. Once again, we see that the asymmetry parameter for the data is higher than that for any of the models. Removing events 13 and 14 does improve the situation as does increasing the optical depth $`\tau `$; the asymmetry remains a little higher but consistent with the models.
## 10 Conclusions
This paper presents the analysis of four seasons of M31 observations at the INT, a subset of the MEGA survey of M31. The observations were carried out to search for MACHOs in the halo of M31. Our fully automated search algorithm identified 14 candidate microlensing events from over $`10^5`$ variable sources. Three of the candidates were previously unpublished. The spatial and timescale distributions are consistent with microlensing.
The core of this paper is the comparison of this candidate event sample with a calculation of the expected number of events from self and halo lensing. This calculation breaks into three parts: a model for the extinction across the M31 disk; a model for the detection efficiency; and a suite of self-consistent disk-bulge-halo models for M31.
The results with regard to the fundamental question of whether there is a significant MACHO fraction in the halo are inconclusive. Based on the total number of events, we find that the most probable MACHO halo fraction $`f`$ varies between $`0`$ and $`0.1`$ depending on the model. Our event rate analysis is consistent with a total absence of MACHOs as the confidence intervals for all of our models include $`f=0`$. On the other hand we can not exclude some MACHO component, since the confidence intervals extend typically up to $`f=0.25`$ and even up to $`f=0.4`$ for a few models.
The spatial distribution of the candidate events is highly asymmetric and does seem to favour a MACHO component. However, for different reasons it is questionable whether the 3 candidate events that largely determine the asymmetry signal should be used in this analysis. Thus, we conclude that both from the observed number of events, and from their spatial distribution we find no compelling evidence for the presence of MACHOs in the halo of M31.
###### Acknowledgements.
We would like to thank Raja Guhathakurta and Phil Choi for making their M31 map available to us and Stephane Courteau for useful conversations. We also thank all observers who have performed the observations for this survey and the ING staff. Support is acknowledged from STScI (GO 10273) and NSF (grants 0406970 and 0070882). JdJ and KK thank the LKBF for travel support.
## Appendix A Candidate event lightcurves
On the following pages, for each of the 14 candidate microlensing events in our sample, the r and i lightcurves and thumbnails taken from the difference centred on the event positions are shown, together with a short discussion. Apart from the INT r and i data, KP4m R and I data points are also plotted in the lightcurves. The fits shown are however the fits done to only the INT data.
### MEGA-ML-1
Located close to the centre of M31, this event has a rather noisy baseline. Apart from the background of very faint variables there are some variable sources clearly visible in the difference images. As can be seen in the thumbnails in figure 21(b) a bright variable is located just a few pixels from the position of the candidate event. Another, fainter variable is seen at a similar distance above and to the left. The other variable sources are further away and should have no influence on the photometry.
### MEGA-ML-2
This candidate event is located very close to MEGA-ML-1 and therefore has the same problems connected to being close the centre of M31. In the thumbnails of days 94, 754, and 1208 we see a variable source a few pixels to the left of the event position. This variable is brighter in r than in i, which causes the r baseline to be the most noisy.
### MEGA-ML-3
This candidate event is also located close to the M31 centre. In figure 6 we already demonstrated that a very faint variable source is positioned $``$0.25″away from this candidate event. In the i thumbnails another variable is visible just above and to the right of the event. This variable has a bright episode between days 440 and 480, causing the bump in the baseline in the i lightcurve.
### MEGA-ML-7
By far the brightest event in our sample, the thumbnails of MEGA-ML-7 show a very bright residual close to the peak centre. Since the peak occurs during the first season, some of the exposures used for creating the reference image contained a significant amount of the magnified flux, so that the baseline lies at a negative difference flux. There are some variables nearby, but none of them are close or bright enough to significantly influence the photometry. The distance to the centre of M31 is also quite large ($``$22′), reducing the background of faint variable sources. As pointed out by Paulin-Henriksson et al. (2003), there are some systematic deviations from the best fit microlensing model. An et al. (2004c) find that this anomaly can be explained by a binary lens.
### MEGA-ML-8
This near side event is located $``$23′ from the centre of M31. A variable that is particularly bright in i is situated about 2.4″ NW of the candidate event, but should not have much of an effect on the photometry. The baselines of the lightcurves indeed look stable and well-behaved.
### MEGA-ML-9
Peak coverage is poor for this candidate event, but the baselines are stable. The thumbnails show quite a lot of faint variables, two of which are located very close, approximately 1″ to the left of the event position, accounting for the noise in the i baseline that is higher than in the r lightcurve.
### MEGA-ML-10
This event is a beautiful example of a combined lightcurve with KP4m and INT data. Peak coverage in INT i is poor, but the KP4m I data points follow the fit (derived only from INT data) very well. A fairly bright variable is situated slightly above and to the right of the event position and there is a hint of a very faint variable about 1″ to the left. Although the INT baseline in i is noisy, the r and both KP4m R and I lightcurves show an very stable and well-behaved baseline.
### MEGA-ML-11
A high signal-to-noise event with a good fit and stable baseline. There is some noise in the i baseline, caused by the variable source that is visible in the thumbnails of days 6 and 756 at $``$1.3″ above the event position. During the fourth observing season a few bad columns were lying exactly on top of the event position, so that there is only 1 INT data point available. However, the KP4m data show that the baseline remains flat everywhere.
### MEGA-ML-13
This candidate event has the lowest signal-to-noise of our sample. It is situated far out in the far side of the disk at $``$31′ from the centre of the galaxy and the relatively low galaxy background makes it possible to detect these kind of faint events. Due to the y-axis scale the i the baseline looks quite noisy, but it is in fact not significantly more so than for other candidate events. The thumbnails of days 398 and 520 show that the closest variable source is located $``$1.4″ below and to the left of the event, which explains the scatter in the i baseline.
### MEGA-ML-14
At $``$35.5′ from the M31 centre, this candidate event is the most far out in the disk of all events in our sample. The i photometry of this candidate event is compromised by the variable source at $``$1.3″. From the i thumbnails one can also see that the event lies at the edge of a fringe, making the background in the lower half of the thumbnails brighter than in the upper half. This can also cause some extra scatter in the photometry. Overall, however, the microlensing fit is very good and both INT and KP4m lightcurves show a stable baseline.
### MEGA-ML-15
This event is again located close to the centre of M31 and presumably has a strong background of faint variable sources. In the thumbnails also several variables are visible very close to the event position, both in r and in i. The lightcurve baselines are rather noisy because of this, but show no coherent secondary bumps and the KP4m baselines are very stable.
### MEGA-ML-16
Not selected in our first analysis of the first two seasons of INT data (de Jong et al. 2004) due to baseline variability, the i lightcurve of this event is strongly influenced by a bright variable situated just 1.1″ to the north. Using a smaller extraction aperture for the photometry in the present analysis, the i baseline is still very noisy and the same is true for the KP4m I-band data. The INT r and KP4m R data are much better behaved and the r peak is fit very well by the microlensing fit.
### MEGA-ML-17
The i baseline is slightly noisy, but the r and both KP4m lightcurves are well-behaved. In the thumbnails no very close variables are visible.
### MEGA-ML-18
This candidate event shows quite large scatter in the baseline and also in the peak. Faint variables might be the culprits, although the event is not located very close to the galaxy centre ($``$15.1′). The thumbnails show no variable sources very close to the event position, however they do show that this event is situated on the edge of a fringe running diagonally across the thumbnails. This fringe and the fact that it can change position slightly between frames is the most probable cause for the noisy i photometry. |
warning/0507/astro-ph0507394.html | ar5iv | text | # Identification of the Red Supergiant Progenitor of Supernova 2005cs: Do the Progenitors of Type II-P Supernovae Have Low Mass?
## 1 Introduction
Direct identification of the progenitors of supernovae (SNe) provides vital information on their explosion mechanisms, a key issue in SN studies. The white dwarfs thought to give rise to thermonuclear Type Ia SNe have such exceedingly low luminosities that they cannot at present be detected in external galaxies. Core-collapse SNe (of Type II, Type Ib, and Ic), on the other hand, arise from far more luminous, massive stars. Unfortunately, even these progenitors are so faint that their detection (ground-based or space-based) is confined to nearby galaxies (distances $`<`$10 Mpc), in which SN discoveries are relatively rare.
Until now, out of more than 3,200 SNe discovered since 1885, only 5 genuine SNe have had their progenitors identified: SN 1987A (a peculiar, subluminous SN II) in the Large Magellanic Cloud (e.g., Gilmozzi et al. 1987; Sonneborn et al. 1987), SN 1993J (an unusual SN IIb) in NGC 3031 (M81; Aldering et al. 1994; Van Dyk et al. 2002), SN 1999ev (Type II) in NGC 4274 (Maund & Smartt 2005), SN 2003gd (Type II) in NGC 628 (M74; Van Dyk et al. 2003a; Smartt et al. 2004), and SN 2004et (Type II) in NGC 6946 (Li et al. 2005a, 2005b). In addition, identifications of precursors that give rise to the super-outbursts of luminous blue variable stars in other galaxies, occasionally misclassified as SNe, have been made for several other objects (Zwicky 1964; Ryder et al. 1993; Van Dyk et al. 1999b; Van Dyk et al. 2000).
Nature provided us with a rare opportunity to increase the sample of directly identified SN progenitors with the discovery of SN 2005cs in the Whirlpool Galaxy (hereafter, M51). The supernova was discovered by Kloehr (2005) at about magnitude 14 on 2005 June 28.9 (UT times are used throughout this paper), and confirmed in images obtained with the 0.76-m Katzman Automatic Imaging Telescope (KAIT; Li et al. 2000; Filippenko et al. 2001) at Lick Observatory on 2005 June 30.25 (Li 2005). An optical spectrum taken with the F. L. Whipple Observatory 1.5-m Tillinghast telescope on 2005 June 30.23 showed SN 2005cs to be a young SN II (Modjaz et al. 2005), with P-Cygni-like line profiles of the hydrogen Balmer series and helium superimposed on a blue continuum.
Just five months before the discovery of SN 2005cs, M51 was fortuitously observed by the Hubble Heritage team (GO/DD program 10452; PI: S. Beckwith) with the Wide Field Channel (WFC) of the Advanced Camera for Surveys (ACS) on-board the Hubble Space Telescope (HST). A large four-color (F435W, F555W, F658N, and F814W) mosaic image of the nearly face-on spiral galaxy NGC 5194 (M51a, the SN 2005cs host) and its interacting companion, NGC 5195 (M51b), was obtained in six ACS pointings (with four dithered exposures at each pointing), and the resulting color composite image was released to the community on 2005 April 25 to celebrate the fifteenth anniversary of the successful operation of HST (Mutchler et al. 2005). These high-resolution ($`0.^{\prime \prime }05`$ pixel<sup>-1</sup>) images are also the deepest ever obtained of M51, reaching limiting magnitudes of 27.3, 26.5, and 25.8 in the combined F435W ($`B`$), F555W ($`V`$), and F814W ($`I`$) images, respectively. M51 had also been observed in several bands with the Wide Field and Planetary Camera 2 (WPFC2), but not to this depth or, generally, at such high spatial resolution.
M51 was also observed by the Near Infrared Camera and Multi-Object Spectrometer (NICMOS) on-board HST in GTO program 7237 (PI: N. Scoville) in Cycle 7; the SN site was imaged in five bands (F110W, F160W, F187N, F190N, and F222M) on 1998 June 28. Together, the ACS and NICMOS data therefore provide images of unprecedented quantity and quality for possibly identifying and studying the progenitor star of SN 2005cs.
In § 2 we report on our direct identification of the progenitor of SN 2005cs, and in § 3 we describe the progenitor’s nature inferred from analysis of these pre-SN HST data. Further discussion is in § 4, and we summarize our conclusions in § 5. We note that M51 was also host to SN 1945A (in NGC 5195; Type I) and SN 1994I (in NGC 5194; Type Ic).
## 2 Identification of the Progenitor
In Table 1 we list the HST ACS/WFC and NICMOS data that we analyzed here. Many additional pre-SN observations of M51 exist, obtained with other instruments on-board HST, including WFPC2 images of SN 1994I. However, the ACS data provide the deepest and the highest-resolution optical images currently available of the galaxy prior to the SN, supplemented by the deep NICMOS data in the near-infrared.
### 2.1 Registration of the Ground-based Observations
To initially locate the SN 2005cs progenitor in these images, we utilized the KAIT SN confirmation observations. We identified six to eight stars in common between the ACS/WFC and KAIT images and measured their pixel coordinates. Using the task IRAF<sup>1</sup><sup>1</sup>1IRAF (Image Reduction and Analysis Facility) is distributed by the National Optical Astronomy Observatories, which are operated by the Association of Universities for Research in Astronomy, Inc., under cooperative agreement with the National Science Foundation. /GEOMAP, we performed a geometrical transformation between the two sets of coordinates and were able to match them to $`<3.0`$ ACS pixels root-mean-square (rms $`<0.^{\prime \prime }15`$). \[This rather large uncertainty arises from the relatively poor spatial resolution in the KAIT images ($`0.^{\prime \prime }8`$ pixel<sup>-1</sup>) and the mediocre seeing under which the images were obtained ($``$2$`\mathrm{}.`$5 FWHM).\] To within this positional uncertainty, we were able to identify a single object at the SN site in the F814W image. However, this object is not detected in the images in any of the other ACS passbands (Li et al. 2005c).
To better isolate the putative progenitor star, we obtained higher-resolution ($`0.^{\prime \prime }187`$ pixel<sup>-1</sup>) SN images with the 3.6-m Canada-France-Hawaii Telescope (CFHT, + MegaCam) on 2005 July 2.28, under good seeing conditions ($`0.^{\prime \prime }7`$ FWHM). We obtained 3$`\times `$10 s and 2$`\times `$90 s exposures in the Sloan $`i`$ band. (Only the images from one chip containing the SN site, out of a total of 36 chips, were analyzed.) As a result, we are able to detect many more faint objects in the CFHT images than in the KAIT images. We identified 20–30 stars (or compact star clusters) around the SN position that are present in both the CFHT and ACS images. The CFHT images were geometrically transformed to the ACS images using IRAF/GEOMAP, with uncertainties of $`<0.6`$–0.8 ACS pixel, and the SN location transformed onto the ACS images is consistent in all five CFHT exposures (to within this uncertainty). Increasing the number of stars included in the transformation does not appreciably reduce the uncertainty, indicating that we have approached the limit of the transformation accuracy. An independent astrometric tie conducted using 90 stars or compact clusters in both sets of images, fit to a third-order polynomial function with full cross-terms for the transformation, resulted in a consistent localization with a similar precision.
The final error-weighted mean pixel location for all measurements from the CFHT images of the SN is $`X=4177.85\pm 0.8`$ and $`Y=3421.17\pm 0.8`$ in the ACS mosaic image “h\_m51\_i\_s05\_drz\_sciḟits” (Table 1), which corresponds to $`\alpha (\mathrm{J2000})=13^h29^m52.^\mathrm{s}764\pm 0.^\mathrm{s}004`$, $`\delta (\mathrm{J2000})=+47^{}10\mathrm{}36.^{\prime \prime }09\pm 0.^{\prime \prime }04`$ from the ACS image world coordinate system (WCS). Direct inspection of the F814W image confirmed our initial identification of the candidate progenitor: only one object is within the 0.8 ACS pixel uncertainty, at position $`\alpha (\mathrm{J2000})=13^h29^m52.^\mathrm{s}760`$, $`\delta (\mathrm{J2000})=+47^{}10\mathrm{}36.^{\prime \prime }11`$ (Li et al. 2005d), and we consider this object most likely to be the progenitor star. The difference between the transformed and measured position for this star is $`\mathrm{\Delta }X=0.71`$ ACS pixel ($`0.^{\prime \prime }036`$) and $`\mathrm{\Delta }Y=0.43`$ ACS pixel ($`0.^{\prime \prime }022`$), both of which are within 1$`\sigma `$ of the positional uncertainties.
A comparison between the ACS/WFC color composite mosaic image and the combined CFHT image is shown in Figure 1.
### 2.2 Registration of the HST/ACS SN Observations
Subsequent images of M51 showing SN 2005cs were obtained with HST/ACS on 2005 July 11, and with NICMOS on 2005 July 13, as part of program GO-10182 (PI: A. V. Filippenko). The details of the observations are listed in Table 2. The ACS images were observed with the High Resolution Channel (HRC) of ACS. The HRC provides a finer spatial resolution ($`0.^{\prime \prime }025`$ pixel<sup>-1</sup>) than the WFC ($`0.^{\prime \prime }05`$ pixel<sup>-1</sup>), but with a smaller field-of-view, $`29\mathrm{}\times 25\mathrm{}`$ (compared to $`202\mathrm{}\times 202\mathrm{}`$ for WFC).
Using six point sources within 10$`\mathrm{}`$ of the SN site, we performed a geometrical transformation between the F250W ACS/HRC SN image and the ACS/WFC F814W pre-SN image. The transformation has an overall error of 0.1 ACS pixel ($`0.^{\prime \prime }005`$). When the SN position measured from the ACS/HRC image is transformed to the ACS/WFC image, it is coincident with our progenitor candidate, with a difference of $`\mathrm{\Delta }X=0.12`$ ACS pixel ($`0.^{\prime \prime }006`$) and $`\mathrm{\Delta }Y=0.10`$ ACS pixel ($`0.^{\prime \prime }005`$). In Figure 2 we show a comparison between the ACS/HRC F250W image and the ACS/WFC F814W image after image registration, while in Figure 3 we show the SN site in all the four ACS/WFC passbands. We now conclude to a high degree of certainty that we have identified in the F814W image the progenitor of SN 2005cs.
### 2.3 Registration of the HST/NICMOS Observations
To determine whether the progenitor star is also seen in the pre-SN HST/NICMOS observations, we attempted to match geometrically the NICMOS images showing SN 2005cs to the pre-SN NICMOS images. Out of the five NICMOS passbands in which M51 was imaged, we consider only the F110W ($`J`$) and F160W ($`H`$) exposures to be deep enough to be useful for our purpose.
We first identified five bright sources in common between the new and the pre-SN NICMOS images. We performed a first-order geometrical transformation between the images which achieved an rms accuracy of 0.08 NICMOS pixel ($`0.^{\prime \prime }016`$). After the images were registered in common, we identified an additional 8–10 fainter objects in both sets of images. A broader geometrical transformation then resulted in a registration with rms accuracy 0.15 NICMOS pixel ($`0.^{\prime \prime }03`$). This larger uncertainty in the transformation for the larger number of fiducials is due mainly to the measurement uncertainty in the centroids for stars in the undersampled and relatively shallow NICMOS images.
In Figure 4 we show a comparison of the registered NICMOS F110W images before and after SN 2005cs. The position of SN 2005cs is marked as a white circle in the left panel. The right panel shows the pre-SN NICMOS F110W image, with the position of SN 2005cs marked with a $`0.^{\prime \prime }09`$ radius circle, which represents the 3$`\sigma `$ uncertainty in the image registration. We show a 2$`\mathrm{}\times `$2$`\mathrm{}`$ close-up of the SN 2005cs progenitor environment on the pre-SN NICMOS images in Figure 5. For comparison, the ACS F814W image is shown in the left panel, with the progenitor of SN 2005cs marked with an illustrative circle. The pre-SN NICMOS F110W and F160W images are shown in the middle and right panels, with a $`0.^{\prime \prime }09`$ radius circle that represents the 3$`\sigma `$ uncertainty in the image registration. Just outside the 3$`\sigma `$ error circle is a relatively bright source in the NICMOS images, which we tentatively identify as the counterpart of the star immediately to the northwest of the SN progenitor seen in the F814W image. Note that the bright object immediately to the southwest of the progenitor seen in the F814W image, which we suggest is a blue compact cluster, has only a faint or undetected counterpart in the near-infrared bands.
We marginally detect a source near the SN progenitor position in the F110W image and, somewhat more suggestively, in the F160W image. However, this source seems to be offset slightly from the exact progenitor position, suggesting it may be a different, very red star. However, the NICMOS source is within 2–3$`\sigma `$ of the progenitor position, and it is possible this source is the same as the star detected in the ACS F814W image. To be conservative, we consider the brightness of this source in the NICMOS images to be an upper limit to the flux of the SN progenitor.
## 3 Photometry of the Progenitor
We have almost certainly identified the progenitor for SN 2005cs in deep pre-SN HST images (see §4). We attempt now to measure the brightnesses of this star in the various bands. From Figure 3 it is apparent that SN 2005cs occurred in a crowded field; in particular, the bright, blue object (likely a compact star cluster) immediately to the southwest of the fainter star will complicate these measurements in the ACS bands.
### 3.1 ACS photometry
We perform photometry on the ACS images following the prescription outlined by Sirianni et al. (2005). Using the IRAF/DAOPHOT package (Stetson 1987), we combined the four individually “drizzled” images produced by the HST ACS calibration pipeline in each of the four observed passbands. We have confirmed that the photometric scale is consistent between our combined images and the mosaic image produced by the Hubble Heritage team.
For the combined F814W and F555W images, we manually selected 377 and 384 (respectively) bright, isolated stars to construct a spatially varying model point-spread function (PSF) across the field. Since photometric accuracy for the progenitor would be improved if we could cleanly remove the contribution of the neighboring bright compact cluster, we carefully examined the characteristics of the cluster (i.e., its light profile) and identified about 20–30 similar objects in each image. Using our model stellar PSF, we first removed all the stars in the neighborhood of our selected clusters. Then, we constructed a spatially constant model cluster “point”-spread function (there are not enough clusters to produce a spatially varying model), which was then fit to the bright contaminating cluster. Inspection of images in which the cluster was subtracted revealed that the model fit did a good job in removing its contribution in the F814W image, but some small subtraction residuals remained in the F555W band.
Nonetheless, with the effect of the cluster minimized as much as possible, we determined the brightness of objects in the SN environment (within a radius of $`0.^{\prime \prime }5`$) using a $`0.^{\prime \prime }15`$-radius aperture PSF (to maximize the signal-to-noise ratio) fit in an iterative process, in which the brightest 2–3 stars were measured first and subtracted away, then the next brightest stars, and so on. By performing the photometry iteratively, we avoided potential errors which could result from fitting all the bright and faint stars simultaneously. The progenitor’s brightness was measured in the F814W image when all other sources had been cleanly removed.
For the F555W image, since the progenitor is not detected, we derived an upper limit to its brightness. We subtracted stars of various magnitudes near the location of the progenitor (determined from the F814W image) and visually inspected the subtracted images. While this process involves some subjectivity and thus may have large uncertainties, we nevertheless find that when stars brighter than 25.8 mag are subtracted from the F555W image, an apparent residual is left in the subtracted images. We consider this magnitude (25.8) the detection limit in the SN environment. This limit is brighter than the global limiting magnitude, $`>26.5`$ mag (Mutchler et al. 2005), due to the presence of the neighboring cluster and its imperfect subtraction.
We determined photometric corrections from the $`0.^{\prime \prime }15`$ PSF fitting radius to a standard $`0.^{\prime \prime }5`$-radius aperture, using several isolated, bright stars. These are $`0.19\pm 0.02`$ mag in F814W and $`0.14\pm 0.02`$ mag in F555W. We then employed the tabulated correction to infinite aperture of $`0.092\pm 0.001`$ mag for the F555W image (Sirianni et al. 2005). Correcting the F814W photometry to infinite aperture was more involved: because both the PSF and aperture correction at F814W depends on a star’s color, we first estimated an effective wavelength $``$8200 Å for the peak of the star’s spectral energy distribution, based on the spectral type suggested by the limit on the F555W$``$F814W color for the progenitor. Following Sirianni et al. (2005), we derived a value for the correction, based on this wavelength; however, it is quite similar to the tabulated synthetic correction in Sirianni et al. (2005) of $`0.087\pm 0.001`$ mag, which we adopt for all the stars in the region. Because of the fortuitous location of the SN pixel position near a readout amplifier for the chip, the ACS charge transfer efficiency correction is negligible for that position.
We then corrected the photometry for interstellar extinction. The Galactic reddening along the line of sight to M51 is only $`E(BV)`$ = 0.035 mag (Schlegel et al. 1998); however, the presence of Na I D line absorption in the SN spectrum at the redshift of M51 suggests some additional host-galaxy extinction toward the SN (Modjaz 2005). We therefore adopted a total reddening $`E(BV)`$ = 0.10$`\pm `$0.05 mag (Richmond 2005) toward the SN progenitor. For the other stars in the SN environment, we lack constraints on the host extinction and corrected only for the Galactic component. Finally, we transformed the F814W and F555W magnitudes to Johnson-Cousins $`V`$ and $`I`$ (Sirianni et al. 2005) and obtained $`I=24.15\pm 0.20`$, $`V>25.5`$ mag for the progenitor. Adopting a distance modulus $`\mu =29.6\pm 0.3`$ mag (Richmond et al. 1996), the progenitor has $`M_I^05.5`$ mag and $`(VI)^0>1.3`$ mag, consistent with a red supergiant of spectral type K0 or later.
Since the F435W and F658N images are not of sufficiently high signal-to-noise ratio to derive meaningful limits on the progenitor brightness (the progenitor is not detected in either image) and enhance our knowledge of the progenitor beyond what we have now learned from the F555W and F814W images, we do not consider them further.
### 3.2 NICMOS Photometry
We performed photometry on the subsampled drizzled NICMOS F110W and F160W pre-SN images, with resolution $`0.^{\prime \prime }1`$ pixel<sup>-1</sup>. The photometric zero points appropriate for Cycle 7, when these NICMOS observations were obtained, were adopted. Since the SN progenitor is, at best, marginally detected, and it occurred in a crowded region, we employed a procedure for the photometry similar to that described above, in §3.1, for the ACS F555W image: stars were subtracted iteratively in order of decreasing brightness near the SN location, using the appropriate TinyTim model PSFs (Krist & Hook 2003) for the NICMOS filters, until the visually most satisfactory subtracted image was achieved. As a result, we derived $`J`$ = 22.5$`\pm `$0.5 and $`H`$ = 21.5$`\pm `$0.5 mag for the marginally detected emission close to the SN site, and adopted these magnitudes for the limiting brightness of the progenitor in each band.
## 4 Discussion
### 4.1 Is SN 2005cs of Type II-Plateau?
Type II SNe can be further divided into several subclasses; see the review of SN types in Filippenko (1997) for more details and references. The two classical subtypes are the Type II-plateau (SNe II-P), with a pronounced plateau phase seen in their optical light curves, and the Type II-linear (SNe II-L), with a linear decline after their maximum brightness. Additional subtypes include the Type II-narrow (SNe IIn), with narrow emission lines in their spectra (often, but not always, superimposed on a broader emission component), and the SN 1987A-like (peculiar, subluminous, with a unique photometric behavior). Moreover, the Type IIb SNe, such as SN 1993J in M81 (Filippenko et al. 1993; Nomoto et al. 1993), manifest themselves as SNe II at early times, but then experience a metamorphosis into a Type Ib SN at late times. To some extent the shapes of the optical light curves correlate with the spectral subtype for SNe II.
We have been able to follow SN 2005cs in $`UBVRI`$ with KAIT roughly every other night since its discovery, sampling the light curves for $`>`$20 d since explosion, so we should already be getting an indication of the SN subtype. Unfortunately, the absolute calibration for SN 2005cs has not yet been established, but in Figure 5 we show its $`UBVRI`$ light curves relative to the recent prototypical SNe II-P 1999em (Leonard et al. 2002b; Hamuy et al. 2001) and 2004et (Li et al. 2005a). The SN 2005cs light curves were obtained by performing differential aperture photometry between the SN and two or more bright stars in the field. In the figure we show the curves adjusted in magnitude to match the light curves for SN 1999em and SN 2004et; the number of days since explosion has not been adjusted. The final photometry for SN 2005cs will be published in a forthcoming paper, when proper absolute calibration for the field has been obtained, along with analysis of available optical spectra.
Even from these relative light curves, it can be clearly seen that SN 2005cs is currently undergoing a plateau phase, especially in the $`V`$, $`R`$, and $`I`$ bands. SN 2005cs also follows very similar photometric evolution to that of SN 1999em and SN 2004et in these three bands, but seems to be evolving somewhat faster in the $`U`$ and $`B`$ bands. At the distance of M51, SN 2005cs had a peak $`V`$-band magnitude on the plateau of only $`15.6`$, so it is rather subluminous. SN 2005cs also had a relatively small early expansion velocity, $``$7500 km s<sup>-1</sup>, derived from the absorption minimum of H$`\beta `$ (Modjaz et al. 2005). We believe there should be little doubt that SN 2005cs is a SN II-P, albeit possibly somewhat unusual and subluminous, relative to the prototypes.
### 4.2 An X-ray Flash from SN 2005cs?
An X-ray flash was detected in a 2136 s Swift/XRT observation of M51 on 2005 July 6.231, 8 days after the discovery of SN 2005cs (Immler et al. 2005). This generated interest as to whether or not the flash was related to SN 2005cs. The X-ray flash is a 6$`\sigma `$ detection, and is located about 10$`\mathrm{}`$ from the SN 2005cs position. No X-ray source was detected 96 minutes later in a 2031 s exposure, as well as in all other Swift/XRT observations of M51 on June 30, July 3, 5, 6 and 7. Given the timing of the X-ray flash ($`>`$8 days after the SN explosion), and the relative offset of the detection from the SN position, speculation emerged that the new source was unrelated to SN 2005cs. However, since the XRT PSF has a half-power diameter of 18$`\mathrm{}`$ at 1.5 keV, based on position alone the possibility that the emission is from SN 2005cs cannot be excluded.
From the photometry obtained with the Palomar Observatory robotic 60-in telescope, Gal-Yam (2005) reports that no photometric anomaly was detected for SN 2005cs around the time of the detection of the X-ray flash in M51, and yet no additional variable sources were detected within 20$`\mathrm{}`$ of SN 2005cs. We obtained with KAIT a sequence of $`UBVRI`$ images for SN 2005cs, starting on July 6.246, just 20 minutes after the Swift/XRT detection of the X-ray flash. The photometry from and timing of these images are indicated with an arrow in Figure 5. Our data are consistent with what Gal-Yam (2005) found: the SN evolved photometrically as expected, following the trend before and after the X-ray flash detection, particularly in $`U`$.
We further compare the ACS/HRC images of the SN 2005cs field in the UV bands taken on July 11.501, 5.27 days after the X-ray flash detection, to the pre-SN ACS/WFC F435W ($`B`$) images. We do not detect any apparent new sources within 10$`\mathrm{}`$ of SN 2005cs, to limiting magnitudes $``$22.5.
Thus, either the X-ray flash was unrelated to SN 2005cs, or it did not produce any detectable anomalous behavior at optical wavelengths shortly thereafter. The nature of the X-ray flash remains a mystery.
### 4.3 The Spectral Type and Mass of the Progenitor
We can further constrain the spectral type of the progenitor of SN II-P 2005cs, a red supergiant, from the HST photometry. In Figure 6 we show the star’s implied spectral energy distribution (SED); obviously, the star was detected only in $`I`$, so we can place only upper limits on its $`VJH`$ brightness in this diagram. The $`I`$ brightness and the upper limits have been corrected for reddening. Also shown for comparison are SEDs of some late-type supergiants (Drilling & Landolt 2003; Tokunaga 2003), reddened by the assumed $`E(BV)`$ to the SN and all normalized to the progenitor’s $`I`$ magnitude. We see that the $`V`$ upper limit constrains the spectral type to later than about K0, while the $`JH`$ upper limits constrain it to earlier than M5.
We can further venture to estimate the initial mass of this star by comparing its intrinsic brightness and color limit to theoretical massive stellar evolutionary tracks. Lejeune & Schaerer (2001) have generated tracks for a range of zero-age main sequence (ZAMS) masses for several different metallicities and assuming enhanced mass loss for the most massive stars. From metallicity measurements in M51 (Zaritsky et al. 1990), we estimate that at the SN 2005cs site, $`12+\mathrm{log}`$ \[O/H\] = 9.06 $`\pm `$ 0.04 dex. This is 0.26 dex higher than the solar abundance (Grevesse & Sauval 1998), implying that the metallicity in the SN environment could be somewhat higher than solar. We therefore consider the tracks from Lejeune & Schaerer for both $`Z=0.02`$ (i.e., solar) and $`Z=0.04`$ (the next higher metallicity).
In Figure 7 we show the intrinsic color-magnitude diagram for the SN environment (the progenitor and other stars within a $`0.^{\prime \prime }5`$ radius of the SN) and overlay the tracks for the two possible metallicities for a variety of stellar masses. The total photometric uncertainties shown are the measurement uncertainties and the uncertainty in the distance modulus added in quadrature. It is also clear from the diagram that the SN environment is abundant in red supergiants with masses of 7–15 $`M_{}`$, although from the astrometric arguments above, we have eliminated these stars as potential candidates for the progenitor.
The location in the diagram of the progenitor itself suggests that it had $`M_{\mathrm{ZAMS}}7`$–9 $`M_{}`$. This estimated mass is right at the $`8M_{}`$ theoretical lower limit for core collapse in massive stars (Woosley & Weaver 1986). Furthermore, we consider it highly unlikely that we have not correctly identified the SN progenitor, if it is indeed a core-collapse SN. The low reddening to the SN and the limiting magnitude in the F814W image, $`M_I^04`$ mag, imply that any star which has escaped detection, yet somehow remains a possible progenitor candidate, has a ZAMS mass ($`<7M_{}`$) formally below the core-collapse limit.
### 4.4 Implications from the SN 2005cs Progenitor
We consider our identification of the progenitor of SN 2005cs quite secure and the estimates for its inferred spectral type and ZAMS mass compelling. SN 2005cs is only the third SN II-P to have its progenitor directly identified on pre-SN images. The other two are SN 2003gd (Van Dyk et al. 2003; Smartt et al. 2004; Hendry et al. 2005), which also had a progenitor with initial mass (7–9 $`M_{}`$) very near the theoretical limit for core collapse, and SN 2004et, with an initial mass in the range 13–20 $`M_{}`$ (Li et al. 2005a). The former was a subluminous SN, also with a red supergiant progenitor, and the latter may have been unusual, with apparently a yellow supergiant progenitor.
Upper limits on the initial masses of other, well-studied SNe II-P based on pre-SN images have also been established: $`M_{\mathrm{ZAMS}}15M_{}`$ for SN 1999em (Smartt et al. 2002), $`15M_{}`$ for SN 1999gi (Leonard et al. 2002a), $`13M_{}`$ for SN 2001du (Van Dyk et al. 2003b; Smartt et al. 2003), and $`12M_{}`$ for SN 1999br (Maund & Smartt 2005). From these detections and upper limits, a trend is emerging for SNe II-P, the most common core-collapse SNe, that the majority (if not all) of them arise from stars with masses in the range $``$8–15 $`M_{}`$. We note that a progenitor in the mass range 20–40 $`M_{}`$ has yet to be found for a normal SN II-P, and thus the fate of these very massive stars still needs to be observationally verified. However, because of the steeply declining mass function, it may be quite some time before a dearth of very massive progenitors presents a significant challenge to theory.
It is also unclear whether stars more massive than 20 $`M_{}`$ actually become normal SNe II-P, or give rise to SNe II-L, such as SN 1979C (Branch et al. 1981; Van Dyk et al. 1999a), or SNe IIn, such as SN 1988Z (Stathakis & Sadler 1991; Chugai & Danziger 1994) and SN 1995N (Fransson et al. 2002). Massive blue variable stars ($`>30`$–40 $`M_{}`$) may undergo super-outbursts (sometimes misclassified as SNe) while in the luminous blue variable stage (e.g., Van Dyk et al. 2000), en route to becoming Wolf-Rayet stars, SNe Ic, and perhaps even collapsars (which appear to be responsible for long-duration gamma-ray bursts).
We have already discussed that SN 2005cs could be subluminous, as was SN 2003gd (Van Dyk et al. 2003; Hendry et al. 2005). Its post-plateau photometric behavior will corroborate this (we predict it should decline markedly in brightness in all bands after the plateau phase, much more so than for prototypical SNe II-P). Zampieri et al. (2003) proposed that these subluminous SNe II-P with low expansion velocities and low <sup>56</sup>Ni yields, such as SN 1999br, SN 2003gd, and now possibly SN 2005cs, originate from high-mass ($`20M_{}`$) progenitors in which the rate of early infall of stellar material on the collapsed core is large. They further postulated that events of this type could form a black hole remnant, giving rise to significant fallback and late-time accretion. It is apparent from the direct mass estimates for both the SN 2003gd and SN 2005cs progenitors, as well as from the upper mass limit for SN 1999br, that these subluminous SNe are produced by relatively low-mass red supergiants, which may end their lives in less energetic explosions.
We also note that the low energy output and the low synthesized yields for <sup>56</sup>Ni suggest that, even though the rate of such events is higher than SNe from more massive stars, the impact on the energy input and chemical evolution of their host galaxies per event is relatively small. Certainly, it is thought that at very early times in galactic history ($`z>6`$), the initial mass function was heavily skewed to very massive stars of essentially zero metallicity, and these stars contributed entirely to the early galactic evolution (e.g., Matteucci & Calura 2005; Umeda & Nomoto 2005). The contribution of lower-mass stars has increased steadily over cosmic history, although the explosions of the more massive stars could still have a large impact on galaxies today. A careful SN rate calculation such as that being conducted by Leaman et al. (2004), and detailed models for nucleosynthesis in the SN II explosions (Woosley & Weaver 1995), will provide useful information on the relative contribution to the chemical evolution of galaxies from SNe of various progenitor masses.
## 5 Conclusions
In this paper we have analyzed HST/ACS and NICMOS data for M51 before the discovery of SN 2005cs, and we have identified the progenitor of SN 2005cs. The secure identification of the progenitor is achieved by geometrical transformations, using new images of SN 2005cs from KAIT, CFHT, and HST/ACS. We measure $`I`$ = 24.15$`\pm `$0.20, $`V>25.5`$, $`J>22.5`$, and $`H>21.5`$ mag for the SN progenitor. Together, this information suggests that SN 2005cs originated from a red supergiant of spectral type K0–M3, with an initial mass of 7–9 $`M_{}`$. A significant trend appears to be emerging, that SNe II-P arise from massive stars with initial masses $`8`$–15 $`M_{}`$.
SN 2005cs is a very important addition to those SNe whose progenitors have been directly identified, a very small (but growing) number of objects. This identification would have been very difficult without the deep images of superior spatial resolution obtained with HST/ACS; the progenitor star would almost certainly have been lost in the light of the neighboring bright star cluster had only inferior images been available. We will continue to monitor SN 2005cs with KAIT and other telescopes to better define the nature of the SN itself; it remains to be seen how this SN will ultimately behave.
The work of A.V.F.’s group at U.C. Berkeley is supported by National Science Foundation grant AST-0307894, and by NASA/HST grant GO-10182 from the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555. He is also grateful for a Miller Research Professorship at U.C.B., during which part of this work was completed. We thank Stefan Immler for his private communications on the X-ray flash in M51. KAIT was made possible by generous donations from Sun Microsystems, Inc., the Hewlett-Packard Company, AutoScope Corporation, Lick Observatory, the NSF, the University of California, and the Sylvia & Jim Katzman Foundation. |
warning/0507/math0507052.html | ar5iv | text | # Zariski pairs on sextics II
## 1. Introduction
We continue to study Zariski pairs in sextics. In this paper, we study Zariski pairs of sextics which are not irreducible. The idea of the construction of Zariski partner sextic for reducible cases is quit different from the irreducible case. It is crucial to take the geometry of the components and their mutual intersection data into account. When there is a line component, flex geometry (i.e., linear geometry) is concerned to the geometry of sextics of torus type and non-torus type. When there is no linear components, the geometry is more difficult to distinguish sextics of torus type. For this reason, we introduce the notion of conical flexes.
We have observed in that the case $`\rho (C,5)=6`$ is critical in the sense that the Alexander polynomial $`\mathrm{\Delta }_C(t)`$ can be either trivial or non-trivial for sextics. If $`\rho (C,5)>6`$ (resp. $`\rho (C,5)<6`$), the Alexander polynomial is not trivial (resp. trivial) (). For the definition of $`\rho (C,5)`$-invariant, see . Thus we concentrate ourselves in this paper the case $`\rho (C,5)=6`$. In , we have classified the possible configurations for reducible sextics of torus type. In particular, the configurations with $`\rho (C,5)=6`$ are given as in Theorem 1 below. Hereafter we use the same notations as for denoting component types. For example, $`C=B_1+B_5`$ implies that $`C`$ has a linear component $`B_1`$ and a quintic component $`B_5`$. We denote the configuration of the singularities of $`C`$ by $`\mathrm{\Sigma }(C)`$.
###### Theorem 1.
() Assume that $`C`$ is a reducible sextic of torus type with $`\rho (C,5)=6`$ and only simple singularities. Let $`\mathrm{\Sigma }_{in}`$ be the inner singularities. Then the possible configurations of simple singularities are as follows.
1. $`\mathrm{\Sigma }_{in}=[A_5,4A_2]:`$ $`C=B_5+B_1`$ and $`\mathrm{\Sigma }(C)=[A_5,4A_2,2A_1],[A_5,4A_2,3A_1]`$, $`[A_5,4A_2,4A_1]`$.
2. $`\mathrm{\Sigma }_{in}=[2A_5,2A_2]`$:
1. $`C=B_1+B_5`$: $`\mathrm{\Sigma }(C)=[2A_5,2A_2,2A_1],[2A_5,2A_2,3A_1]`$.
2. $`C=B_1+B_1^{}+B_4`$: $`\mathrm{\Sigma }(C)=[2A_5,2A_2,3A_1],[2A_5,2A_2,4A_1]`$.
3. $`C=B_2+B_4`$: $`\mathrm{\Sigma }(C)=[2A_5,2A_2,2A_1]`$, $`[2A_5,2A_2,3A_1]`$.
4. $`C=B_3+B_3^{}`$: $`\mathrm{\Sigma }(C)=[2A_5,2A_2,3A_1]`$.
3. $`\mathrm{\Sigma }_{in}=[E_6,A_5,2A_2]`$: $`C=B_1+B_5`$, $`\mathrm{\Sigma }(C)=[E_6,A_5,2A_2,2A_1]`$, $`[E_6,A_5,2A_2,3A_1]`$.
4. $`\mathrm{\Sigma }_{in}=[3A_5]`$:
1. $`C=B_1+B_5`$: $`\mathrm{\Sigma }(C)=[3A_5,2A_1]`$.
2. $`C=B_2+B_4`$: $`\mathrm{\Sigma }(C)=[3A_5,2A_1]`$.
3. $`C=B_1+B_1^{}+B_4`$: $`\mathrm{\Sigma }(C)=[3A_5,3A_1]`$.
4. $`C=B_3+B_3^{}`$: $`\mathrm{\Sigma }(C)=[3A_5]`$, $`[3A_5,A_1]`$, $`[3A_5,2A_1]`$.
5. $`C=B_1+B_2+B_3`$: $`\mathrm{\Sigma }(C)=[3A_5,2A_1],[3A_5,3A_1]`$.
6. $`C=B_1+B_1^{}+B_1^{\prime \prime }+B_3`$: $`\mathrm{\Sigma }(C)=[3A_5,3A_1]`$, $`[3A_5,4A_1]`$,
7. $`C=B_2+B_2^{}+B_2^{\prime \prime }`$: $`\mathrm{\Sigma }(C)=[3A_5,3A_1]`$.
5. $`\mathrm{\Sigma }_{in}=[2A_5,E_6]`$:
1. $`C=B_1+B_5`$: $`\mathrm{\Sigma }(C)=[E_6,2A_5,2A_1]`$.
2. $`C=B_2+B_4`$: $`\mathrm{\Sigma }(C)=[E_6,2A_5,2A_1]`$.
3. $`C=B_1+B_1^{}+B_4`$: $`\mathrm{\Sigma }(C)=[E_6,2A_5,3A_1]`$.
6. $`\mathrm{\Sigma }_{in}=[A_8,A_5,A_2]`$: $`C=B_1+B_5`$, $`\mathrm{\Sigma }(C)=[A_8,A_5,A_2,2A_1]`$, $`[A_8,A_5,A_2,3A_1]`$.
7. $`\mathrm{\Sigma }_{in}=[A_{11},2A_2]`$:
1. $`C=B_2+B_4`$: $`\mathrm{\Sigma }(C)=[A_{11},2A_2,2A_1]`$, $`[A_{11},2A_2,3A_1]`$.
2. $`C=B_3+B_3^{}`$: $`\mathrm{\Sigma }(C)=[A_{11},2A_2,3A_1]`$.
8. $`\mathrm{\Sigma }_{in}=[A_{11},A_5]`$:
1. $`C=B_1+B_5`$: $`\mathrm{\Sigma }(C)=[A_{11},A_5,2A_1]`$.
2. $`C=B_2+B_4`$: $`\mathrm{\Sigma }(C)=[A_{11},A_5,2A_1]`$.
3. $`C=B_3+B_3^{}`$: $`\mathrm{\Sigma }(C)=[A_{11},A_5]`$, $`[A_{11},A_5,A_1]`$,$`[A_{11},A_5,2A_1]`$,
4. $`C=B_1+B_2+B_3`$: $`\mathrm{\Sigma }(C)=[A_{11},A_5,2A_1]`$, $`[A_{11},A_5,3A_1]`$,
9. $`\mathrm{\Sigma }_{in}=[A_{17}]`$: $`C=B_3+B_3^{}`$, $`\mathrm{\Sigma }(C)=[A_{17}]`$, $`[A_{17},A_1]`$, $`[A_{17},2A_1]`$.
Our main result in this paper is:
###### Theorem 2.
There are Zariski partner sextics with the above configurations with the following exceptions:
1. $`\mathrm{\Sigma }(C)=[A_5,4A_2,4A_1]`$ with $`C=B_5+B_1`$.
2. $`\mathrm{\Sigma }(C)=[2A_5,2A_2,4A_1]`$ with $`C=B_1+B_1^{}+B_4`$.
3. $`\mathrm{\Sigma }(C)=[E_6,A_5,2A_2,3A_1]`$ with $`C=B_5+B_1`$.
4. $`\mathrm{\Sigma }(C)=[3A_5,4A_1]`$ with $`C=B_3+B_1+B_1^{}+B_1^{\prime \prime }`$.
5. $`\mathrm{\Sigma }(C)=[E_6,2A_5,3A_1]`$ with $`C=B_4+B_1+B_1^{}`$.
The non-existence of sextics of non-torus type with the above exceptional configurations will be explained by flex geometry. The existence will be also explained by the flex geometry for those which has a line components and by conical flex geometry for the component type $`B_4+B_2,B_2+B_2^{}+B_2^{\prime \prime }`$.
###### Remark 3.
(1) In the list of Theorem 2, there are certainly several cases which are already known. For example, the configuration $`C=B_3+B_3^{}`$ with one singularity $`A_{17}`$ is given by Artal .
(2) In this paper, we only studied possible Zariski pairs of reducible sextics $`(C,C^{})`$ where $`C`$ is of torus type and $`C^{}`$ is not of torus type. On the other hand, the possibility of Zariski pairs among reducible sextics of the same class is not discussed here. Several examples are known among reducible sextics of non-torus type. For such cases, Alexander polynomials can not distinguish the differnece. See papers
## 2. Reducible sextics of non-torus type
To compute explicit polynomials defining reducible sextics, it is not usually easy to look for special degenerations into several irreducible components starting from the generic sextics $`_{i+j6}a_{ij}x^iy^j`$. Recall that we have classified all possible reducible simple configurations in and it is easier to start from a fixed reducible decomposition. In fact, the geometry of the configuration of a reducible sextic depends very much on the geometry of each components. A smooth point $`PC`$ is called a flex point if the intersection multiplicity of the tangent line and $`C`$ at $`P`$ is strictly greater than 2. First we recall the following fact for flex points ().
###### Lemma 4.
Let $`C:F(X,Y,Z)=0`$ be an irreducible plane curve of degree $`n`$ with singularities $`\{P_1,\mathrm{},P_k\}`$. Then the number of flexes $`\iota (C)`$ is given by
$$\iota (C)=3n(n2)\underset{i=1}{\overset{k}{}}\epsilon (P_i;C)$$
where the second term $`\epsilon (P_i;C)`$ is the flex defect and given by the local intersection number of $`C`$ and the hessian curve of $`C`$ at $`P_i`$.
Generic flex defect of simple singularities we use are
(1) $`\epsilon (A_1)=6,\epsilon (A_2)=8,\epsilon (A_{3\iota 1})=9\iota ,(\iota 2),\epsilon (E_6)=22`$
Recall that flex points of a curve are described by the hessian of the defining homogeneous equation. When we have an affine equation $`C:f(x,y)=0`$, flex points in $`𝐂^\mathrm{𝟐}`$ are described by $`f(x,y)=flex_f(x,y)=0`$ () where
$$flex_f(x,y):=f_{xx}f_y^22f_{xy}f_xf_y+f_{yy}f_x^2$$
This is an easy way to check flex points from the affine equation.
A sextic $`C`$ is of (2,3)-torus type if we can take a defining polynomial of the form $`f_2(x,y)^3+f_3(x,y)^2=0`$ where $`\text{degree}f_j=j`$. The intersections $`f_2=f_3=0`$ are singular points of $`C`$ and we call them inner singularities. For a given sextic $`C`$ of torus type whose singularities are simple, the possible inner singularities are
$`(\mathrm{}):\{A_2,A_5,A_8,A_{11},A_{14},A_{17},E_6\}`$.
A convenient criterion for $`C`$ to be of torus type is the existence a certain conic $`C_2`$ such that $`C_2C\mathrm{\Sigma }(C)`$ (Tokunaga’s criterion , Lemma 3, ).
A sextic of torus type $`C`$ is called of linear torus type if the conic polynomial $`f_2`$ can be written as $`f_2(x,y)=\mathrm{}(x,y)^2`$ for some linear form $`\mathrm{}(x,y)`$ (). A sextic of linear torus type can have only $`A_5,A_{11},A_{17}`$ as inner singularties and the location of these singularities are colinear.
The proof of Theorem 2 is done by giving explicit examples. For the better understanding of the existence or non-existence of the Zariski pairs, we divide the above configurations into the following classes.
1. $`C`$ has a quintic component. The corresponding component type is $`B_5+B_1`$ and $`B_1`$ is a flex tangent line.
2. $`C`$ has a quartic component. There are two subcases.
1. $`C=B_4+B_1+B_1^{}`$. In this case, two line components are flex tangent lines.
2. $`C=B_4+B_2`$.
3. $`C`$ has a cubic component. There are two subcases.
1. Sextics of linear torus type.
2. Sextics, not of linear torus type.
4. $`C=B_2+B_2^{}+B_2^{\prime \prime }`$.
## 3. Configuration coming from quintic flex geometry
Let $`B_5`$ be an irreducible quintic and let $`P`$ be a flex point of $`B_5`$. We denote the tangent line at $`P`$ by $`L_P`$. We say that $`P`$ is a flex of torus type (respectively a flex of non-torus type) if $`B_5L_P`$ is a sextic of torus type (resp. of non-torus type). The following configurations are mainly related to the flex geometry of certain quintics. (By ’flex geometry’, we mean the geometry of the tangent lines at the flex points and the curve.) Recall that $`\mathrm{\Sigma }(B_5)`$ is the configuration of the singularities of $`B_5`$. Let $`\iota `$ be the number of flex points on $`B_5`$.
1. $`C=B_5+B_1`$ with $`\mathrm{\Sigma }(C)=[A_5,4A_2,kA_1],k=2,\mathrm{\hspace{0.17em}3},\mathrm{\hspace{0.17em}4}`$. Then $`\mathrm{\Sigma }(B_5)=[4A_2,(k2)A_1]`$ for $`k=2,\mathrm{\hspace{0.17em}3},\mathrm{\hspace{0.17em}4}`$ and $`\iota =13,\mathrm{\hspace{0.17em}7},\mathrm{\hspace{0.17em}1}`$ respectively.
2. $`C=B_5+B_1`$ with $`\mathrm{\Sigma }(C)=[2A_5,2A_2,kA_1],k=2,3`$. Then $`\mathrm{\Sigma }(B_5)=[A_5,2A_2,(k2)A_1]`$ and $`\iota =11,\mathrm{\hspace{0.17em}5}`$ respectively.
3. $`C=B_5+B_1`$ with $`\mathrm{\Sigma }(C)=[E_6,A_5,2A_2,kA_1],k=2,\mathrm{\hspace{0.17em}3}`$. Then $`\mathrm{\Sigma }(B_5)=[E_6,2A_2,(k2)A_1]`$ and $`\iota =7,\mathrm{\hspace{0.17em}1}`$ respectively for $`k=2,3`$. The case $`k=3`$ corresponds to sextics of torus type.
4. $`C=B_5+B_1`$ with $`\mathrm{\Sigma }(C)=[E_6,2A_5,2A_1]`$. The quintic $`B_5`$ has $`\mathrm{\Sigma }(C)=[E_6,A_5]`$ and $`\iota =5`$.
5. $`C=B_5+B_1`$ with $`\mathrm{\Sigma }(C)=[3A_5,2A_1]`$. Then $`\mathrm{\Sigma }(B_5)=[2A_5]`$ and $`\iota =9`$.
6. $`C=B_5+B_1`$ with $`\mathrm{\Sigma }(C)=[A_8,A_5,A_2,kA_1],k=2,3`$. Then $`\mathrm{\Sigma }(B_5)=[A_8,A_2,(k2)A_1]`$ and $`\iota =10,\mathrm{\hspace{0.17em}4}`$ for $`k=2,3`$.
7. $`C=B_5+B_1`$ with $`\mathrm{\Sigma }(C)=[A_{11},A_5,2A_1]`$. Then $`\mathrm{\Sigma }(B_5)=[A_{11}]`$ and $`\iota =9`$.
We are going to show the stronger assertion for the above configurations: the Zariski partner sextic of non-torus type are simply given by replacing the flex line components $`B_1`$ for the above cases, if $`B_5`$ has at least two flex points.
Let $`\mathrm{\Xi }`$ be a configuration of singularities on $`B_5`$, which is one of the above list. Let $`(\mathrm{\Xi };5)`$ be the configuration space of quintics $`B_5`$ such that $`\mathrm{\Sigma }(B_5)=\mathrm{\Xi }`$. We considere it as a topological subspace of the space of quintics. For our purpose, it is enough to consider the marked configuration subspace $`(\mathrm{\Xi };5)^{}`$ which consists of the pair $`(B_5,P)`$, where $`B_5(\mathrm{\Xi };5)`$ and $`P`$ is a flex point of torus type. The following describes the existence of sextics of non-torus type with the above configurations.
###### Theorem 5.
Let $`\mathrm{\Xi }`$ be a configuration of singularities on $`B_5`$, which is one of the above list. 1. The configuration subspace $`(\mathrm{\Xi };5)^{}`$ is connected for each $`\mathrm{\Xi }`$.
2. For each $`\mathrm{\Xi }[4A_2,2A_1],[E_6,2A_2,A_1]`$ and $`B_5(\mathrm{\Xi };5)^{}`$, a Zariski pair sextics are given as $`\{B_5L_P,B_5L_Q\}`$ where $`P`$ and $`Q`$ are flex points of torus-type and of non-torus type respectively.
3. For these two exceptional cases, we have the equality $`(\mathrm{\Xi };5)^{}=(\mathrm{\Xi };5)`$ and a quintic $`B_5(\mathrm{\Xi };5)`$ does not contain any flexes of non-torus type.
###### Remark 6.
Let $`\iota _t,\iota _{nt}`$ be the respective number of flex points of torus type and of non-torus type on a generic $`B_5(\mathrm{\Xi };5)^{}`$. We do not need the precise number $`\iota _t,\iota _{nt}`$ for our purpose. The sum $`\iota =\iota _t+\iota _{nt}`$ is described by Lemma 3. By an explicit computation, we have the next table which describes the distributions of number of flex points. The second line is the configuration of singularity and the last line is the pair of flex numbers $`(\iota _t,\iota _{nt})`$.
| 1 | 2 | 3 | 4 | 5 | 6 | 7 |
| --- | --- | --- | --- | --- | --- | --- |
| $`4A_2`$ $`4A_2+A_1`$ $`4A_2+2A_1`$ | $`A_5+2A_2`$ $`A_5+2A_2+A_1`$ | $`E_6+2A_2`$ $`E_6+2A_2+A_1`$ | $`E_6+A_5`$ | $`2A_5`$ | $`A_8+A_2`$ $`A_8+A_2+A_1`$ | $`A_{11}`$ |
| (1,12) (1,6) (1,0) | (1,10) (1,4) | (1,6) (1,0) | (1,4) | (1,8) | (1,9) (1,3) | (1,8) |
Proof. First recall that the topology of the complement of the sextics $`B_5L_P`$ for a flex point $`P`$ of torus type and non-torus type are different. They can be distinguished by Alexander polynomial (). Therefore to show the assertion about the positivity $`\iota _{nt}>0`$, it is enough to check the assertion by some quintic $`B_5`$. Examples will be given in the next subsection. Secondly, the irreducibility of the configuration space $`(\mathrm{\Xi };5)`$ of quintics $`f_5(x,y)=0`$ with singularities $`\mathrm{\Xi }=[4A_2,\mathrm{\hspace{0.17em}2}A_1],[E_6,\mathrm{\hspace{0.17em}2}A_2,A_1]`$ are easily proves as follows. For $`\mathrm{\Xi }=[4A_2,\mathrm{\hspace{0.17em}2}A_1]`$, the dual curves of quintics in this configuration space are quartics with configuration $`[A_2,2A_1]`$. As the irreducibility of the configuration space $`([A_2,2A_1];4)`$ is easy to be checked, the irreducibility of $`([4A_2,2A_1];5)`$ follows. Take $`\mathrm{\Xi }=[E_6,2A_2,A_1]`$. For a quintic $`B_5`$ with $`\mathrm{\Sigma }(B_5)=\mathrm{\Xi }`$, the dual curve $`B_5^{}`$ is again a quartic with $`[A_2,\mathrm{\hspace{0.17em}2}A_1]`$ (thus mapped into the same configuration space with the dual of quintics with $`[4A_2,2A_1]`$). Hoever we can not apply the same argument. The reason is that the dual curve $`B_5^{}`$ is not generic in the configuration space $`([A_2,2A_1];4)`$: the quartic $`B_5^{}`$ has not 4 flexes but three flexes, one flex of flex order 4 (=dual of $`E_6`$) and 2 flexes of flex order 3 (i.e., dual of $`2A_2`$). Thus we need another argument. Note that any three singular points can not be colinear on $`B_5`$ by Bézout theorem. We can consider the slice condition:
$`()`$: $`E_6`$ is at $`(1,0)`$ and two $`A_2`$ are at $`(0,1),(0,1)`$ and one $`A_1`$ at $`(1,0)`$.
It is easy to compute that a Zariski open subset of this slice has the normal form:
$$\begin{array}{c}h:=e_{\mathit{1}}^{}{}_{}{}^{3}4y^3x^22e_{\mathit{1}}^{}{}_{}{}^{3}x^3+x^4e_{\mathit{1}}^{}{}_{}{}^{3}+xe_{\mathit{1}}^{}{}_{}{}^{3}+3y^5e_{\mathit{1}}^{}{}_{}{}^{2}+e_{\mathit{1}}^{}{}_{}{}^{3}y^4+3ye_{\mathit{1}}^{}{}_{}{}^{2}+10e_{\mathit{1}}^{}{}_{}{}^{3}y^2x^3\hfill \\ \hfill +18e_{\mathit{1}}^{}{}_{}{}^{3}y^2x^212e_\mathit{1}y^2x^2+6e_{\mathit{1}}^{}{}_{}{}^{3}y^2x12e_\mathit{1}y^2x9yx^4e_{\mathit{1}}^{}{}_{}{}^{2}12ye_{\mathit{1}}^{}{}_{}{}^{2}x^3\\ \hfill +12yxe_{\mathit{1}}^{}{}_{}{}^{2}+6yx^2e_{\mathit{1}}^{}{}_{}{}^{2}6y^3x^2e_{\mathit{1}}^{}{}_{}{}^{2}6y^3e_{\mathit{1}}^{}{}_{}{}^{2}2e_{\mathit{1}}^{}{}_{}{}^{3}y^2+e_{\mathit{1}}^{}{}_{}{}^{3}x^57e_{\mathit{1}}^{}{}_{}{}^{3}y^4x\\ \hfill +12e_\mathit{1}y^4x12y^3xe_{\mathit{1}}^{}{}_{}{}^{2}2x^2e_{\mathit{1}}^{}{}_{}{}^{3}+8y^34y4y^58yx+8yx^3+4yx^4\\ \hfill +8y^3x,e_10,\pm 2/\sqrt{3}\end{array}$$
Thus the irreducibility of $`([E_6,2A_2,A_1];5)`$ follows. For each of them we know that the number of flex points is 1 and $`(\mathrm{\Xi };5)^{}\mathrm{}`$. On the other hand, the topology of sextics $`B_5L_P`$, of torus type and non-torus type, are distinguished by the Alexander polynomials $`(t^2t+1)(t1)`$ and $`(t1)`$ respectively. This implies $`(\mathrm{\Xi };5)^{}=(\mathrm{\Xi };5)`$. ∎
### 3.1. Example for sextics with quintic components
We gives examples of sextics with a quintic components.
1. $`C=B_5+B_1,\mathrm{\Sigma }=[A_5,4A_2,kA_1],k=2,3`$: First we consider the case $`k=2`$. The quintic has $`4A_2`$ and 13 flex points.
$$\begin{array}{c}C:((\frac{28}{153}x+\frac{8}{153})y^4+(\frac{52}{51}x^3\frac{20}{51}x^2\frac{152}{153}x\frac{16}{153})y^2+x^5\frac{2}{17}x^4\frac{193}{51}x^3\hfill \\ \hfill +\frac{116}{51}x^2+\frac{124}{153}x+\frac{8}{153}\left)\right(\frac{80}{17}x+\frac{880}{51}\frac{320}{51}y)\end{array}$$
The sextics of torus type is obtained by replacing the line component by $`\frac{28}{153}x+\frac{8}{153}`$.
Next, we consider the case $`k=3`$. Th equintic $`B_5`$ has $`4A_2+A_1`$. The flex which gives a sextic of torus type is $`(1,0)`$.
$$\begin{array}{c}f_5:=\frac{385}{16}x^4y+\frac{3885}{16}xy^2+\frac{1897}{128}xy^4+\frac{345}{16}yx^3\frac{441}{4}y^3x+\frac{529}{4}yx^2+73y+72y^3\hfill \\ \hfill +\frac{403}{128}x^3y^2\frac{16783}{128}x^2y^2\frac{811}{4}yx\frac{869}{64}y^4+\frac{7087}{128}x\frac{3675}{32}y^2+\frac{3201}{512}x^5\\ \hfill +\frac{601}{32}x^4y^5+\frac{313}{8}x^2y^3\frac{11511}{256}x^2\frac{997}{64}\frac{10167}{512}x^3\end{array}$$
$`B_5`$ has two obvious flex points: $`P:=(1,0)`$ and $`Q:=(1520/293,287/293)`$, where $`P`$ is a flex of torus type and $`Q`$ is a flex of non-torus type. There are 5 other flex points whose $`x`$-coordinates are the solution of
$$\begin{array}{c}R_1:=5926214587003x^532698277751050x^4+69779834665700x^3\hfill \\ \hfill 72918583611000x^2+37638730560000x7728486400000=0\end{array}$$
We can check that the roots of $`R_1=0`$ corresponds to flexes of non-torus type as follows. (The same argument applies to other cases.) Note that any conics which is passing through 4 $`A_2`$ of $`B_5`$ are given by
$$h_2:=y^2\frac{1}{2}d_{01}yx+d_{01}y\frac{1}{2}x^2+\frac{1}{19}x^2d_{01}\frac{5}{2}x\frac{10}{19}d_{01}x+3+\frac{16}{19}d_{01}$$
Thus if there is a flex $`P(a,b)`$ of torus type (so $`R_1(a)=0`$), there is a cubic form $`h_3(x,y)`$ such that the sextic $`C=B_5L_P`$ is described as $`C:=\{h_3^2+h_2^3=0\}`$. On the other hand, put $`S_2(x,d_{01})`$ be the polynomial of degree 2 in $`x`$ defined by $`S_2(x,d_{01})=R(h_2,f_5,y)/P(x)^2`$ where $`R(h_2,f_5,y)`$ is the resultant of $`h_2`$ and $`f_5`$ in $`y`$ and $`P(x)=0`$ is the defining polynomial for the x-coordinates of 4 $`A_2`$. Then $`S_2`$ must be $`c(xa)^2`$ for some $`c0`$. Let $`b_1(d_{01})`$ be the discriminant polynomial of $`S_2`$ in $`x`$ and let $`b_2(d_{01})`$ be the resultant of $`S_2(x,d_{01})`$ and $`R_1(x)`$ in $`x`$. Thus we obtain two polynomials $`b_1(d_{01}),b_2(d_{01})`$ of the parameter $`d_{01}`$ which must have a common root: We can check that $`b_1(d_{01})=b_2(d_{01})=0`$ has no common root in $`d_{01}`$.
2. $`C=B_5+B_1`$, $`\mathrm{\Sigma }=[2A_5,2A_2,jA_1],j=2,3`$. The quintic $`B_5`$ has $`A_5+2A_2`$.
$$\begin{array}{c}j=2,[2A_5,2A_2,2A_1]:(\frac{16145}{1024}y^3\frac{93}{64}y^5\frac{877727}{8192}y^2\frac{110055}{1024}y\frac{329525}{16384}x\hfill \\ \hfill +\frac{11025}{16384}x^3\frac{543975}{16384}x^2+\frac{1751733}{16384}y^2x^2+\frac{79625}{4096}x^5\frac{235529}{16384}y^4+\frac{1999101}{16384}y^2x^3\\ \hfill \frac{100809}{8192}y^3x+\frac{1199495}{8192}yx^3+\frac{289995}{4096}yx^2\frac{1199495}{8192}yx+\frac{150225}{4096}yx^4\\ \hfill \frac{498845}{4096}y^2x\frac{275703}{16384}y^4x+\frac{23889}{8192}y^3x^2+\frac{18625}{512}x^4\frac{52025}{16384})(y+1)\end{array}$$
A sextic of torus type is give by replacing the line component by $`x1=0`$.
$$\begin{array}{c}j=3,[2A_5,2A_2,3A_1]:(\frac{2}{7}+x^5\frac{4}{7}x^22x^3+\frac{2}{7}x^4+x\frac{2}{7}y^2x^3+\frac{12}{7}x^2y^2\frac{1}{7}xy^4\hfill \\ \hfill \frac{6}{7}y^2x\frac{4}{7}y^2+\frac{2}{7}y^4\left)\right(\frac{44064}{34157767}y\sqrt{963+1182\sqrt{6}}+\frac{16521840}{34157767}\frac{7198560}{34157767}\sqrt{6}\\ \hfill \frac{28320}{34157767}y\sqrt{963+1182\sqrt{6}}\sqrt{6}+\frac{7328592}{34157767}\sqrt{6}x\frac{2704104}{4879681}x)\end{array}$$
The quintic $`B_5`$ has $`A_5+2A_2+A_1`$ and 5 flex points and among them, there exists a unique flex of torus type. The tangent line at this flex of torus type is given by $`2x=0`$.
3. A sextic $`C=B_5+B_1`$ with $`\mathrm{\Sigma }(C)=[E_6,A_5,2A_2,2A_1]`$ is given by
$$\begin{array}{c}f:=(\frac{53}{141}x+\frac{3}{47}y+y^5\frac{50}{47}y^3+\frac{4}{47}y^2x\frac{769}{141}y^4x\frac{614}{141}y^2x^3+2y^2+\frac{53}{141}x^5+\frac{56}{141}yx^2\hfill \\ \hfill \frac{10}{3}yx+\frac{1174}{141}y^3x+\frac{1256}{141}y^3x^2+\frac{10}{3}x^3y\frac{69}{47}y^4\frac{25}{47}x^4+\frac{50}{47}x^2\frac{106}{141}x^3\\ \hfill \frac{1462}{141}y^2x^2\frac{65}{141}yx^4\frac{25}{47}\left)\right(y+1\frac{8}{3}x)\end{array}$$
The quintic has 7 flex points and there is a unique one among them which is of torus type at $`(\frac{2400}{1357},\frac{357}{1357})`$.
4. $`C=B_5+B_1`$ with $`[E_6,2A_5,2A_1]`$. The quintic $`B_5`$ has $`E_6+A_5`$ and it has 5 flexes. Among them, there is a unique flex of torus type. A sextic of non-torus type:
$$\begin{array}{c}f:=(4451+9742y^2x+4639y^4x9501y14381x423y^5\sqrt{33}351\sqrt{33}\hfill \\ \hfill 16343x^3+6546y^38005y^4+3554y^2+19373x^219373y^2x^2+9836x^4\\ \hfill +2955y^5+10266yx^32936x^519020y^3x+19020yx14661yx^2\\ \hfill +14661y^3x^2+1521\sqrt{33}y1098y^3\sqrt{33}+1593y^4\sqrt{33}+756x^4\sqrt{33}\\ \hfill +1215x^2\sqrt{33}1242y^2\sqrt{33}1917x^3\sqrt{33}+297x\sqrt{33}+999y^2x^3\sqrt{33}\\ \hfill +36yx^4\sqrt{33}+1458y^2x\sqrt{33}+918yx^3\sqrt{33}+2052y^3x\sqrt{33}423yx^2\sqrt{33}\\ \hfill 2052yx\sqrt{33}+423y^3x^2\sqrt{33}1215y^2x^2\sqrt{33}1755y^4x\sqrt{33}+6077y^2x^3\\ \hfill 5124yx^4)(y+1)\end{array}$$
A sextic torus type is given by replacing $`y+1`$ by the flex tangent at $`(\alpha ,\beta )`$ where
$$\alpha :=\frac{1476423}{6805087}+\frac{176748}{6805087}\sqrt{33},\beta :=\frac{1469468}{6805087}\frac{931392}{6805087}\sqrt{33}$$
5. $`C=B_5+B_1`$ with $`[3A_5,2A_1]`$.
$$\begin{array}{c}\frac{12}{6279955}(2516+27\sqrt{69})(251610064x^445xy^4\sqrt{69}+90xy^2\sqrt{69}\hfill \\ \hfill +108\sqrt{69}y^3x^2108\sqrt{69}yx^2+180x^5\sqrt{69}+54y^2\sqrt{69}45x\sqrt{69}27y^4\sqrt{69}\\ \hfill 108x^4\sqrt{69}27\sqrt{69}+10060yx^45030y^310064yx^22516y^4\\ \hfill +10064y^3x^2+2515y^5+10060x^2840xy^4+1680xy^2+5032y^210060y^2x^2\\ \hfill +3360x^5+2515y840x)\\ \hfill \left((\frac{31104}{2497}\frac{1620}{2497}\sqrt{69})y+\frac{108}{2497}(52\sqrt{69}499)(x1)\right)\end{array}$$
The flex point which gives a sextic of torus type is
$$(\alpha ,\beta ),\alpha :=\frac{957138004}{22902646825}+\frac{2339358408}{22902646825}\sqrt{69},\beta :=\frac{540908244}{916105873}\frac{52210443}{916105873}\sqrt{69}$$
6. a. A sextic of torus type with $`[A_8,A_5,A_2,2A_1]`$ with line component is given by:
$$\begin{array}{c}f:=\left(60y^2+60yx^2\right)^3\hfill \\ \hfill +\left(\frac{81}{25}y^3+(\frac{6849}{25}x\frac{9639}{25})y^2+(\frac{162}{25}x^2+\frac{6849}{25}x+\frac{1944}{5})y+x^3\frac{162}{25}x^2\right)^2\end{array}$$
and the line component is defined by $`y1=0`$. It has 10 flexes and the flex at $`(0,1)`$ gives a flex of torus type. Other flex tangent lines give a sextic of non-torus type. For example
$$\begin{array}{c}f_6:=(\frac{324}{25}x^5\frac{26244}{625}x^4\frac{2223126}{625}y^2x^3+\frac{40132557}{625}y^3x^2\frac{428706}{625}yx^4\hfill \\ \hfill \frac{43281837}{625}y^2x^2\frac{134993439}{625}y^5+\frac{271568079}{625}y^4\frac{8419248}{125}y^3\frac{3779136}{25}y^2\\ \hfill +\frac{629856}{125}yx^2+\frac{1733076}{625}yx^3\frac{26628912}{125}y^2x+\frac{132035022}{625}y^3x+\frac{1109538}{625}y^4x)\\ \hfill (\frac{64039734}{25}y+\frac{18974736}{5}+\frac{33205788}{25}x)\end{array}$$
b. A sextic of torus type with $`[A_8,A_5,A_2,3A_1]`$ and with component type $`B_5+B_1`$ is given by
$$\begin{array}{c}f:=(\frac{15}{2}y^2\frac{15}{2}y16x^2)^3\hfill \\ \hfill +(\frac{455}{24}y^3+(\frac{80}{3}x+\frac{245}{6})y^2+(\frac{140}{3}x^2+\frac{80}{3}x\frac{175}{8})y+64x^3\frac{140}{3}x^2)^2\end{array}$$
It has the line component $`y1=0`$ and the quintic has 4 flex points among which only the flex (0,1) is of torus type. An example of sextic of non-torus type is given by
$$\begin{array}{c}f_6:=(\frac{17920}{3}x^5\frac{19600}{9}x^4+\frac{3500}{3}y^2x\frac{6125}{3}yx^2+\frac{47600}{9}x^3y\frac{30625}{64}y^2+\frac{332125}{192}y^3\hfill \\ \hfill \frac{11275}{3}y^3x^2+5800y^2x^2\frac{19600}{9}xy^3+\frac{9100}{9}xy^4\frac{44240}{9}x^3y^2+\frac{40720}{9}yx^4\\ \hfill \frac{1170775}{576}y^4+\frac{450025}{576}y^5)(\frac{2222000000}{255584169}y+\frac{2156000000}{255584169}+\frac{492800000}{85194723}x)\end{array}$$
7. A quintic with $`A_{11}`$ has 9 flex points, among which there exists a unique flex of torus type. In the following example, our quintic has a flex of torus type at $`(0,1)`$ so the the sextic of torus type is given by
$$\begin{array}{c}f:=\left(\frac{28}{25}y^2+\frac{28}{25}yx^2\right)^3\hfill \\ \hfill +\left(\frac{511}{100}y^3+(\frac{28}{25}x\frac{7}{500})y^2+(\frac{91}{20}x^2+\frac{28}{25}x\frac{637}{125})yx^3+\frac{91}{20}x^2\right)^2\end{array}$$
and the line component is $`y1=0`$. $`B_5`$ has 8 flex of non-torus type. We can take one at $`(1,1)`$ so that a sextic of non-torus type is given by
$$\begin{array}{c}f_6:=(\frac{6176793}{250000}y^5\frac{7154}{625}y^4x+\frac{7194719}{250000}y^4\frac{245049}{5000}y^3x^2\frac{859901}{31250}y^3+\frac{98}{3125}y^3x\hfill \\ \hfill +\frac{35672}{3125}y^2x\frac{405769}{15625}y^2+\frac{13181}{5000}y^2x^2\frac{7}{250}y^2x^3\frac{2548}{125}yx^3+\frac{57967}{1250}yx^2\\ \hfill +\frac{7833}{400}yx^4\frac{8281}{400}x^4+\frac{91}{10}x^5)(\frac{3306744}{15625}y+\frac{1928934}{15625}+\frac{275562}{3125}x)\end{array}$$
## 4. Configuration coming from quartic geometry
### 4.1. Configuration coming from quartic flex geometry
We consider the sextics $`C`$ with component type $`B_4+B_1+B_1^{}`$. The corresponding possible configurations are
(a) $`\mathrm{\Sigma }(C)=[3A_5+3A_1]`$ and $`\mathrm{\Sigma }(B_4)=[A_5]`$ or (b) $`\mathrm{\Sigma }(C)=[2A_5+2A_2+kA_1],k=3,4`$ and $`\mathrm{\Sigma }(B_4)=[2A_2+(k3)A_1]`$ or (c) $`\mathrm{\Sigma }(C)=[E_6+2A_5+3A_1]`$ and $`\mathrm{\Sigma }(B_4)=[E_6]`$.
Let $`P,Q`$ be two flex points on $`B_4`$ and let $`L_P,L_Q`$ be the flex tangents. We say that a pair of flex points $`\{P,Q\}`$ are a flex pair of torus type if the sextic $`B_4L_PL_Q`$ is a sextic of torus type.
###### Theorem 7.
Case (a) $`\mathrm{\Sigma }(C)=[3A_5+3A_1]`$. The quartic $`B_4`$ has one $`A_5`$ and 6 flex points and two line components are flex tangent lines. There exist two flex pairs of torus type. The other choices give sextics of non-torus type.
Case (b) $`\mathrm{\Sigma }(C)=[2A_5+2A_2+kA_1],k=3,4`$ . The quartic $`B_4`$ has $`2A_2`$ or $`2A_2+A_1`$ according to $`k=3`$ or $`4`$ and $`B_4`$ has 8 or 2 flex points respectively. For the case, $`k=3`$, there are both flex pairs of torus type and of non-torus type. For $`k=4`$, the choice of $`\{P,Q\}`$ is unique and it is a pair of torus type.
Case (c) $`\mathrm{\Sigma }(C)=[E_6+2A_5+3A_1]`$ . $`B_4`$ has two flexes and they gives a pair of torus type. Thus there is no sextic of non-torus type with $`E_6+2A_5+3A_2`$ with component type $`B_4+B_1+B_1^{}`$.
Proof. As the configuration spaces of quartics with one $`A_5`$, or $`2A_2`$ or $`2A_2+A_1`$ or $`E_6`$ are connected, it is enough to check the assertion by an example.
For the non-existence, note that a quartic with $`\mathrm{\Sigma }(B_4)=2A_2+A_1`$ or $`\mathrm{\Sigma }(B_4)=[E_6]`$ has exactly 2 flexes. Thus the existence of sextic of torus type $`B_4+B_1+B_1^{}`$ with $`[2A_5+2A_2+4A_1]`$ or $`[E_6+2A_5+3A_1]`$ implies that there does not exist sextic of non-torus type with these two configurations∎
Example I. We consider the quartic $`B_4:=\{g_4=0\}`$ with one $`A_5`$:
$$\begin{array}{c}g_4:=6391350x^2+351x^4+468x^3108x+288y+1608Iy^2x^2\sqrt{3}+1452Iyx\sqrt{3}\hfill \\ \hfill 676Iyx^3\sqrt{3}1032Iy^2x\sqrt{3}1452Iy^3x\sqrt{3}288y^3918y^2+279y^4\\ \hfill 648y^3x+1350y^2x^2+108y^2x936yx^3+648yx432Iy^2\sqrt{3}+776Iy^3\sqrt{3}\\ \hfill 1608Ix^2\sqrt{3}152I\sqrt{3}776Iy\sqrt{3}+1032Ix\sqrt{3}+584I\sqrt{3}y^4+728Ix^3\sqrt{3}\end{array}$$
$`B_4`$ has an $`A_5`$ singularity at $`(1,0)`$ and 6 flexes at
$$\begin{array}{c}P_1=(0,1),P_2=(0,1),P3=(\frac{130}{1069}+\frac{370}{1069}I\sqrt{3},\frac{263}{1069}+\frac{156}{1069}I\sqrt{3})\hfill \\ \hfill P4:=(\frac{2190}{13333}\frac{4790}{39999}I\sqrt{3},\frac{12671}{13333}\frac{4492}{39999}I\sqrt{3}),\\ \hfill P_5=(\frac{2116}{6841}I\sqrt{3}+\frac{632}{20160427}\sqrt{9098887465462588I\sqrt{3}}+\frac{4086}{6841}\\ \hfill +\frac{498}{20160427}I\sqrt{3}\sqrt{9098887465462588I\sqrt{3}},\\ \hfill \frac{9056}{47887}\frac{6133}{47887}I\sqrt{3}+\frac{4}{47887}\sqrt{9098887465462588I\sqrt{3}})\\ \hfill P_6=(\frac{632}{20160427}\sqrt{9098887465462588I\sqrt{3}}+\frac{4086}{6841}\\ \hfill \frac{498}{20160427}I\sqrt{3}\sqrt{9098887465462588I\sqrt{3}}+\frac{2116}{6841}I\sqrt{3},\\ \hfill \frac{9056}{47887}\frac{6133}{47887}I\sqrt{3}\frac{4}{47887}\sqrt{9098887465462588I\sqrt{3}})\end{array}$$
It is easy to check that $`\{P_1,P_3\},\{P_2,P_4\},\{P_5,P_6\}`$ give sextic of torus type. The other cases give sextics of non-torus type. For example, a nice sextic of non-torus type is given by taking the tangent lines at $`P_1`$ and $`P_2`$: $`B_1+B_1^{}:(y1)(y+1)=0`$.
II. Now we consider the quartic with $`2A_2`$ (case (b) with $`k=3`$).
$$\begin{array}{c}f_4(x,y):=\frac{254143}{4096}x^4\frac{251}{16}x^3+\frac{11}{32}yx^3\frac{5893}{2048}y^2x^2\frac{2761}{1024}yx^2\frac{126093}{2048}x^2\frac{11}{32}yx+\frac{1}{32}y^3x\hfill \\ \hfill +\frac{251}{16}x+\frac{5893}{2048}y^2\frac{1957}{4096}+\frac{211}{4096}y^4+\frac{2761}{1024}y\frac{251}{1024}y^3\end{array}$$
It has 8 flex points and four flexes are explicitly written as
$$P_1=(1,0),P_2=(1,0),P_3=(0,1),P_4=(\frac{16064}{64025},\frac{61977}{64025})$$
Sextics of torus type are given by taking tangent lines at $`\{P_1,P_2\},\{P_3,P_4\}`$. As a sextic of non-torus type, we can take the tangent lines at $`P_1,P_3`$ so that the sextics is given by adding two lines $`(x1)(4y4+16x)=0`$. The configuration space of quartic with $`2A_2+A_1`$ is connected. Each quartic has two flex points and with two tangent lines $`B_1,B_1^{}`$, $`B_4B_1B_1^{}`$ gives a sextic of torus type with $`[2A_5,2A_2,4A_1]`$. Thus there is no sextic of non-torus type $`C=B_4+B_1+B_1^{}`$ with configuration $`[2A_5,2A_2,4A_1]`$.
### 4.2. Conical geometry of quartic
Now we consider the configuration with component type $`B_4+B_2`$ will be considered here. The corresponding configurations are
1. $`\mathrm{\Sigma }(C)=[2A_5,2A_2,2A_1],[2A_5,2A_2,3A_1]`$.
2. $`\mathrm{\Sigma }(C)=[3A_5,2A_1]`$.
3. $`\mathrm{\Sigma }(C)=[E_6,2A_5,2A_1]`$.
4. $`\mathrm{\Sigma }(C)=[A_{11},2A_2,2A_1],[A_{11},2A_2,3A_1]`$.
5. $`\mathrm{\Sigma }(C)=[A_{11},A_5,2A_1]`$.
Zariski pairs with the above configurations with fixed component type $`B_4+B_2`$ can not be explained by the flex geometry.
We have to generalize the notion of flex points. Let $`B`$ be a given irreducible plane curve of degree $`d`$. Let $`\mathrm{\Phi }`$ be a linear system of conics and let $`\alpha `$ be the dimension of $`\mathrm{\Phi }`$. For a general smooth point $`PB`$, the maximal intersection number of $`I(B,B_2;P)`$ for $`B_2\mathrm{\Phi }`$ is $`\alpha `$. We say $`P`$ is a conical flex point with respect to $`\mathrm{\Phi }`$ if the intersection number $`I(B,B_2;P)\alpha +1`$. If $`dim\mathrm{\Phi }=5`$ (so $`\mathrm{\Phi }`$ is the family of all conics), we say simply that $`P`$ is a conical flex point.
1. Let us consider the case $`\mathrm{\Sigma }(C)=[2A_5,2A_2,2A_1],[2A_5,2A_2,3A_1]`$. We consider first a sextic of torus type $`C=\{f_2^3+f_3^2=0\}`$ which decomposes into a quartic $`B_4`$ and a quartic $`B_2`$:
$$\begin{array}{c}f(x,y):=\left(y^22+2x^2\right)^3+\hfill \\ \hfill \left(25y^3+(13x23)y^2+(26x^2+26)y+13x^323x^213x+23\right)^2\\ \hfill =(y^2+x^21)(177x^4598x^3676yx^3+344x^2+849x^2y^2+1196yx^2650y^3x\\ \hfill +598x+676yx598xy^2521151y^2+1150y^31196y+626y^4)\end{array}$$
Our quartic is defined by
$$\begin{array}{c}g_4(x,y)=(177x^4598x^3676yx^3+344x^2+849x^2y^2+1196yx^2650y^3x\hfill \\ \hfill +598x+676yx598xy^2521151y^2+1150y^31196y+626y^4)=0\end{array}$$
Note that the singularities $`B_4`$ are $`2A_2`$. The intersection $`B_2B_4`$ makes two $`A_5`$ at $`P:=(1,0)`$ and $`Q:=(1,0)`$. We consider the linear system $`\mathrm{\Phi }`$ of conics of dimension 2 which are defined by the conics $`C_2:=\{h_2(x,y)=0\}`$ such that $`I(C_2,B_4;P)3`$. Then we consider the conical flex points $`R=(a,b)B_4`$ with respect to $`\mathrm{\Phi }`$, which is described by the condition $`h_2\mathrm{\Phi }`$ such that $`I(h_2,B_4;R)3`$. We found that there are $`11`$ conical flex points. Two of them can be explicitly given as $`S_1=Q`$ and $`S_2:=(0,1)`$. The corresponding conics are given as
$`h_2(x,y)=x^2+y^21,\{h_2=0\}B_4\{P,Q\}`$
$`g_2(x,y)=(y^2+2yx2yx^2+4x3),\{g_2=0\}B_4\{P,S_2\}`$
We can easily check that the sextic
$$\begin{array}{c}(y^2+2yx2yx^2+4x3)(177x^4598x^3676yx^3+344x^2+849x^2y^2\hfill \\ \hfill +1196yx^2650xy^3+598x+676yx598xy^2521151y^2+1150y^3\\ \hfill 1196y+626y^4)=0\end{array}$$
is not of torus type. Thus $`B_4\{h_2(x,y)=0\}`$ and $`B_4\{g_2(x,y)=0\}`$ is a Zariski pair.
Similarly the case $`[2A_5,2A_2,3A_1]`$ can be treated in the same way. We start from a sextic of torus type:
$$\begin{array}{c}f(x,y)=\left(\frac{49}{64}y^2\frac{15}{64}+\frac{15}{64}x^2\right)^3+\hfill \\ \hfill \left(\frac{131}{256}y^3+(\frac{729}{512}x\frac{297}{256})y^2+(\frac{387}{256}x^2+\frac{387}{256})y+\frac{729}{512}x^3\frac{297}{256}x^2\frac{729}{512}x+\frac{297}{256}\right)^2\\ \hfill =\frac{27}{262144}(y^2+x^21)(19808x^441796yx^332076x^3+40521y^2x^26865x^2+34056x^2y\\ \hfill 32076xy^2+32076x+41796xy14148xy^37770y^234056y\\ \hfill 1815y^4+11528y^312943)\end{array}$$
The quartic $`B_4`$ has two $`2A_2+A_1`$ and we consider the linear system $`\mathrm{\Phi }`$ of conics intersecting $`B_4`$ at $`P=(1,0)`$ with intersection number 3. We find that there exist 5 conical flex points with respect to $`\mathrm{\Phi }`$, and among them we have two explicit ones: $`Q=(1,0)`$ and $`(0,1)`$. We see that the conic corresponding to $`(0,1)`$ gives a Zariski partner sextic $`f_6=0`$ to $`C=\{f=0\}`$.
$$\begin{array}{c}f_\mathit{6}:=(5y^264xy+64y+69x^2128x+59)(19808x^441796yx^332076x^3+40521y^2x^2\hfill \\ \hfill 6865x^2+34056x^2y32076xy^2+32076x+41796xy14148xy^37770y^234056y\\ \hfill 1815y^4+11528y^312943)\end{array}$$
###### Remark 8.
The calculation of conical flex points are usually very heavy. We used maple 7 to compute in the following recipe. a. First compute the normal form of $`h_2\mathrm{\Phi }`$. It contains two parameters. b. Assume $`(u,v)B_4`$. Put $`gg_4(x,y):=(x+u,y+v)`$ and $`hh_2(x,y):=h_2(x+u,y+v)`$. Consider the maximal contact coordinate at $`(u,v)`$: $`\mathrm{\Phi }(x)=a_1x+a_2x^2+a_3x^3`$ and put $`GG4(x):=gg_4(x,\mathrm{\Phi }(x))`$ and $`HH_2(x):=hh_2(x,\mathrm{\Phi }(x))`$. Our assumption implies that $`\text{Coeff}(GG4,x,1)=\text{Coeff}(GG2,x,2)=0`$ and $`\text{Coeff}(HH_2,x,0)=\text{Coeff}(HH_2,x,1)=\text{Coeff}(HH_2,x,2)=0`$. Solve the equations $`\text{Coeff}(GG4,x,1)=\text{Coeff}(GG2,x,2)=0`$ in $`a_1,a_2`$. Then solve the equations $`\text{Coeff}(HH_2,x,0)=\text{Coeff}(HH_2,x,1)=0`$ in the remaining parameters of the linear system. Then we get two equations in $`u,v`$:
$$g_4(u,v)=\text{Coeff}(HH_2,x,2)=0$$
c. Use the resultant computation to solve the above equations to obtain the possibility of conical flex points.
2. $`\mathrm{\Sigma }(C)=[3A_5,2A_1]`$ with $`B_4+B_2`$:
$$\begin{array}{c}B_4:g_4(x,y)=(6yx+\frac{1710}{91}yx^2\frac{1466}{91}x^3y\frac{1992}{91}y^2x6y^3x+\frac{790}{91}y^3+y^4\frac{4904}{91}x^3+\frac{1992}{91}x\hfill \\ \hfill \frac{790}{91}y+\frac{939}{91}y^2x^2+\frac{1161}{91}y^2+\frac{1968}{91}x^2\frac{1252}{91}+\frac{2196}{91}x^4)\end{array}$$
It has an $`A_5`$ at $`P:=(1,0)`$. We consider the linear system $`\mathrm{\Phi }`$ of conics of dimension 2 whose conic are intersecting with $`B_4`$ at $`(0,1)`$ with intersection number 3. We find 14 conical flex points with respect to $`\mathrm{\Phi }`$ in which two are explicit: $`R=(\frac{498727}{500817},\frac{266266}{500817})`$ and $`Q=(0,1)`$. The corresponding conics $`f_2=0,k_2=0`$ intersecting $`B_4`$ at $`P,Q`$ oe $`P,R`$ are given by the following and they gives sextics of non-torus type and of torus type respectively.
$`f_2(x,y):=y^21+{\displaystyle \frac{171}{79}}x^2`$
$`k_2(x,y)=(y^2+({\displaystyle \frac{268}{759}}x+{\displaystyle \frac{4424}{2277}})y+{\displaystyle \frac{85291}{19987}}x^2{\displaystyle \frac{268}{759}}x{\displaystyle \frac{6701}{2277}})`$
3. $`\mathrm{\Sigma }(C)=[E_6,2A_5,2A_1]`$: We start the next quartic
$$\begin{array}{c}B_4:y^4+(\frac{195}{64}x+\frac{169}{64})y^3+(\frac{105}{32}x^2\frac{33}{8}x+\frac{27}{32})y^2+(\frac{143}{64}x^3+\frac{117}{64}x^2+\frac{195}{64}x\frac{169}{64})y\hfill \\ \hfill +\frac{45}{32}x^4\frac{19}{8}x^3\frac{21}{16}x^2+\frac{33}{8}x\frac{59}{32}=0\end{array}$$
Sextics of non-torus type and of torus type are given by the conics $`B_2,B_2^{}`$:
$`B_2:y^21+{\displaystyle \frac{9}{13}}x^2=0`$
$`B_2^{}:y^2+({\displaystyle \frac{770}{1147}}x+{\displaystyle \frac{156}{1147}})y+{\displaystyle \frac{11025}{14911}}x^2+{\displaystyle \frac{770}{1147}}x{\displaystyle \frac{1303}{1147}}=0`$
4. We consider the configurations $`[A_{11},2A_2,2A_1],[A_{11},2A_2,3A_1]`$. First we consider two cuspidal quartics with $`[2A_2]`$:
$`f_\mathit{4}:=5805x^42916Ix^3\sqrt{2}+3888Ix^3y\sqrt{2}1269x^2y^2729x^23834x^2y`$
$`3888Ixy^2\sqrt{2}+108Ix\sqrt{2}y^3+2916Ixy\sqrt{2}+1323y^31971y^281y^4`$
$`+729y=0`$
It can makes sextics of torus type and non-torus type with configuration $`[A_{11},2A_2,2A_1]`$ with respective conics:
$`f_2(x,y)=yx^2`$
$`h_2(x,y)=({\displaystyle \frac{1}{8}}I\sqrt{2}y^2+xy{\displaystyle \frac{3}{4}}Iy\sqrt{2}{\displaystyle \frac{5}{4}}I\sqrt{2}x^23x+{\displaystyle \frac{5}{8}}I\sqrt{2})`$
They correspond to the conical flex points $`(0,0)`$ and (0,1). The other conical flexes are very heavy to be computed.
Next we consider the configuration $`[A_{11},2A_2,3A_1]`$ which is produced by a quartic $`B_4`$ with $`2A_2+A_1`$ and a conic $`B_2`$ with a single tangent at a conical flex.
$$B_4:\frac{1}{16}y^4+\frac{3}{4}xy^3+\frac{59}{8}y^2x^2y^2x\frac{1}{8}y^2+\frac{27}{4}yx^36yx^2\frac{3}{4}yx+\frac{17}{16}x^4x^3\frac{9}{8}x^2+x+\frac{1}{16}=0$$
We find three conical flex points
$$P1:=(1/19,12/19),\mathit{P2}:=(\frac{4}{13},\frac{15}{13}),\mathit{P3}:=(1,\mathrm{\hspace{0.17em}0})$$
The corresponding conic which are tangent at the respective conical flex point $`P_i,i=1,2,3`$ are given by
$`g_{21}:=y^2+({\displaystyle \frac{1270}{141}}x+{\displaystyle \frac{10}{141}})y+{\displaystyle \frac{38711}{423}}x^2{\displaystyle \frac{5486}{423}}x{\displaystyle \frac{457}{423}}`$
$`g_{22}:=y^2+({\displaystyle \frac{7462}{2517}}x{\displaystyle \frac{5462}{2517}})y+{\displaystyle \frac{13007}{7551}}x^2{\displaystyle \frac{21902}{7551}}x+{\displaystyle \frac{8831}{7551}}`$
$`g_{23}:={\displaystyle \frac{3}{22}}y^2{\displaystyle \frac{61}{11}}yx+y{\displaystyle \frac{73}{66}}x^2{\displaystyle \frac{1}{33}}x+{\displaystyle \frac{71}{66}}`$
The corresponding sextic $`f_4(x,y)g_{2latexConicallRi}(x,y)=0`$ is of torus type for $`i=1`$ and of non-torus type for i=2,3. The torus decomposition of $`f_4g_{21}`$ is given by $`z_{21}^3+z_{31}^2=0`$ where
$`z_{\mathit{21}}:={\displaystyle \frac{1}{423}}\mathrm{\hspace{0.17em}423}^{(2/3)}(3y^2+24yx+293x^234x3)`$
$`z_{\mathit{31}}:={\displaystyle \frac{1}{6768}}I(5189x3477x^2+20045x^33y54yx`$
$`+2361yx^25y^2+247y^2x+3y^3)\sqrt{6768}`$
5. Lastly, we consider the configuration $`[A_{11},A_5,2A_1]`$ which is associated to a quartic $`B_4`$ with an $`A_5`$ singularity and a conic $`B_2`$ tangent at a conical flex point with intersection number 6. As a quartic, we take:
$$\begin{array}{c}\mathit{f4}:=y^4+xy^3+\frac{7}{15}x^2y^22xy^23y^2+\frac{2}{15}x^3y\frac{2}{15}x^2y+xy+2y+\frac{4}{75}x^4\hfill \\ \hfill \frac{2}{15}x^3\frac{1}{3}x^2\end{array}$$
$`B_4`$ has apparently 26 conical flex points. (The calculation is very heavy.) We take four explicit conical flex points:
$$P_1:=(0,0),P_2:=(\frac{540}{493},\frac{250}{493}),P_3:=(\frac{270}{301},\frac{58}{301}),P_4:=(\frac{270}{193},\frac{50}{193})$$
After an easy computation, the respective conics are given as
$`n_{\mathit{21}}:=y^2+({\displaystyle \frac{5}{59}}x{\displaystyle \frac{50}{59}})y+{\displaystyle \frac{25}{177}}x^2`$
$`n_{\mathit{22}}:=y^2+({\displaystyle \frac{10845}{262699}}x{\displaystyle \frac{273650}{262699}})y+{\displaystyle \frac{135675}{262699}}x^2+{\displaystyle \frac{153000}{262699}}x+{\displaystyle \frac{70000}{262699}}`$
$`n_{\mathit{23}}:=y^2+({\displaystyle \frac{16681}{32607}}x{\displaystyle \frac{17318}{10869}})y+{\displaystyle \frac{30251}{163035}}x^2{\displaystyle \frac{1544}{32607}}x{\displaystyle \frac{112}{10869}}`$
$`n_{\mathit{24}}:=y^2+({\displaystyle \frac{5225}{1633}}x+{\displaystyle \frac{350}{71}})y+{\displaystyle \frac{4225}{1633}}x^2+{\displaystyle \frac{13000}{1633}}x+{\displaystyle \frac{10000}{1633}}`$
Put $`f_{6j}(x,y):=f_4(x,y)n_{2j}(x,y)`$ and $`C^{(j)}=\{f_{6j}=0\}`$. It is also easy to see that $`C^{(1)},C^{(2)}`$ are of non-torus type and $`C^{(3)},C^{(4)}`$ are of torus type.
## 5. Flex geometry of cubic curves
### 5.1. Configurations coming from cubic flex geometry: a cubic component and a line component.
Let us consider first configurations which occurs in sextics which have at least a cubic component $`B_3`$. We divides into the following cases.
1. $`C=B_3+B_3^{}`$.
1. $`\mathrm{\Sigma }(C)=[A_{17}],[A_{17},A_1],[A_{17},2A_1]`$.
2. $`\mathrm{\Sigma }(C)=[A_{11},A_5],[A_{11},A_5,A_1],[A_{11},A_5,2A_1],`$.
3. $`\mathrm{\Sigma }(C)=[A_{11},2A_2,3A_1]`$
4. $`\mathrm{\Sigma }(C)=[3A_5],[3A_5+A_1],[3A_5,2A_1]`$.
2. $`C=B_3+B_2+B_1`$.
1. $`\mathrm{\Sigma }(C)=[A_{11}+A_5+2A_1],[A_{11}+A_5+3A_1]`$.
2. $`\mathrm{\Sigma }(C)=[3A_5+2A_1],[3A_5+3A_1]`$ .
3. $`C=B_3+B_1+B_1^{}+B_1^{\prime \prime }`$ with configuration $`[3A_5+3A_1],[3A_5+4A_1]`$.
We first consider the cases (2) and (3). In these cases, there are one cubic component $`B_3`$ and at least one line component $`B_1`$. Recall that the configurations in $`(2)`$ and $`(3)`$ occurs as sextics of linear torus type. For a reducible sextic $`C`$ which is classified in either (2) or (3), the necessary and sufficient condition for $`C`$ to be of torus type is there exists a line $`L`$ containing inner singularities. In the case of $`\mathrm{\Sigma }(C)=[A_{11},A_5]`$, $`L`$ is also tangent to the tangent cone of $`A_{11}`$. We first recall the following basic geometry for cubic curves.
###### Proposition 9.
1. A smooth cubic $`C`$ has 9 flex points. Among $`84`$ choices of three flex points, 12 colinear triples of flexes.
2. A nodal cubic has 3 flex points, and they are colinear.
3. A cuspidal cubic has one flex point.
For the proof of the assertion 1, see Example below.
###### Corollary 10.
The configuration $`[3A_5+4A_1]`$ with components type $`B_3+B_1+B_1^{}+B_1^{\prime \prime }`$ does not exist as a sextic of non-torus type.
Proof. The cubic has a node and three line components are flex tangent lines at three flex points. We know that such configuration exists as a sextics of linear torus type . As the configuration space of one nodal cubics is connected, every sextics $`B_3+B_1+B_1^{}+B_1^{\prime \prime }`$ is of torus type.∎
(2) $`C=B_3+B_2+B_1`$ with $`\mathrm{\Sigma }(C)=[3A_5,kA_1],k=2,3`$. In this case, $`B_3`$ is either smooth or nodal and two intersection points $`B_3B_2`$ generates $`2A_5`$. The third $`A_5`$ is generated by a flex tangent line $`B_1`$.
###### Proposition 11.
Assume that a cubic $`B_3`$ and a conic $`B_2`$ are intersecting at two points $`P,Q`$ with respective intersection number 3, producing $`2A_5`$-singularities. Then the line passing through $`P,Q`$ intersects $`B_3`$ at another point, say $`R`$, and $`R`$ is a flex point of $`B_3`$.
This Proposition describes sextics of torus type and non-torus type with configuration $`[3A_5,jA_1],j=0,1,2`$ ((2-b)).
Proof of Proposition 11. Assume that $`P=(0,1)`$ and $`Q=(0,1)`$ with the tanget lines $`y=\pm 1`$ respectively. Then by an easy computation, the cubic $`B_3`$ is defined by a polynomial
$$f_3:=y^3+y^2a_{12}xy^2a_{00}+ya_{21}x^2y+a_{30}x^3a_{00}a_{21}x^2a_{12}x+a_{00}=0$$
and the conic is given as $`y^2+a_{21}x^21=0`$. Then $`R`$ is given as $`(0,a_{00})`$ and we can easily see that $`R`$ is a flex of $`B_3`$.∎
By the same calculation, we see that
###### Proposition 12.
Assume that a cubic $`B_3`$ and a conic $`B_2`$ are intersecting at one points $`P`$ with intersection multiplicity $`6`$ producing an $`A_{11}`$-singularity Then the tangent line passing at $`P`$ intersects another point $`RB_3`$ and $`R`$ is a flex point of $`B_3`$.
Proof is similar. Putting $`P=(0,0)`$ and assuming $`y=0`$ as the tangent line, the cubic is written as
$$\begin{array}{c}f_3:=(y^3t_{2}^{}{}_{}{}^{4}y^2xt_3a_{01}t_4t_2+y^2xa_{21}t_3t_{2}^{}{}_{}{}^{2}+2y^2xa_{11}t_{3}^{}{}_{}{}^{2}t_2+2y^2xa_{01}t_{3}^{}{}_{}{}^{3}\hfill \\ \hfill y^2xa_{11}t_4t_{2}^{}{}_{}{}^{2}y^2t_{2}^{}{}_{}{}^{2}a_{01}t_4y^2t_{2}^{}{}_{}{}^{3}a_{21}y^2t_{2}^{}{}_{}{}^{2}a_{11}t_3+yt_{2}^{}{}_{}{}^{4}a_{21}x^2\\ \hfill +yt_{2}^{}{}_{}{}^{4}a_{11}x+yt_{2}^{}{}_{}{}^{4}a_{01}x^3t_{2}^{}{}_{}{}^{4}a_{01}t_3x^3t_{2}^{}{}_{}{}^{5}a_{11}a_{01}t_{2}^{}{}_{}{}^{5}x^2)/t_2^4\end{array}$$
and $`R=(0,xt_2a_{11}+a_{01}t_2+xa_{01}t_3)`$.
###### Lemma 13.
Assume that a conic $`B_2`$ is tangent to an irreducible curve $`C`$ of degree $`d3`$ at a smooth point $`PC`$ so that $`I(B_2,C;P)3`$. Then $`P`$ is not a flex point of $`C`$.
Proof. Let $`h_2(x,y)=0`$ be a conic equation which defines $`B_2`$. In fact, if $`P=(a,b)`$ is a flex point of $`C`$, $`I(C,B_2;P)3`$ implies that $`C`$ is locally parametrized as $`y_1(x)=t_1x_1+t_3x_1^3+\text{(higher terms)}`$ where $`(x_1,y_1)=(xa,yb)`$ assuming the tangent line is not $`xa=0`$. As $`B_2`$ does not have any flex, the equation $`h_2(x,y)=0`$ is solved as $`y_1=s_1x_1+s_2x_1^2+\text{(higher terms)}`$ with $`s_20`$. Thus $`I(C,B_2;P)=1`$ or $`2`$ according to $`s_1t_1`$ or $`s_1=t_1`$.∎
First we consider the case (2-b). Then the cubic is either smooth or have a node. Thus it has at least 3 flexes. As the intersection $`B_3B_2`$ are not flex points, we can find another flex point $`SB_3`$, $`SR`$. Taking flex tangent $`L:=T_SB_3`$ as the line component, the corresponding sextic is not of torus type. As a cuspidal cubic has a unique flex points, we see that a sextic $`C=B_3+B_2+B_1`$ with $`[3A_5+A_2+2A_1]`$ does not appear as a sextic of non-torus type.
Now we consider (3). Assume that the cubic is smooth. Then there are 9 flex points and $`B_1,B_1^{},B_1^{\prime \prime }`$ are flex tangents. The sextic is of torus type if and only if three flexes are colinear . This case, the configuration is $`[3A_5+3A_1]`$.
Finally we consider the case (2-a). In this case, $`B_3B_2`$ is a single point $`P`$ and $`I(B_3,B_2;P)=6`$ and the intersection singularity is $`A_{11}`$. $`B_3`$ has at most a node and so it has at least 3 flex points. Taking a line component which is the flex tangent at $`S`$ other than $`R`$, we get a sextic $`C=B_3+B_2+B_1`$ with $`[A_{11}+A_5+2A_1]`$ or $`[A_{11}+A_5+3A_1]`$. ∎
We omit explicit examples for (2-a) and (2-b) as they can be easily obtained from sextic of torus type with the same configuration () and replacing the flex line $`B_1`$. We only gives an example of (3).
Example $`B_3+B_1+B_1^{}+B_1^{\prime \prime }`$, with configuration $`[3A_5+3A_1]`$. In this case, the cubic $`B_3`$ is smooth and has $`9`$ flexes and three lines $`B_1,B_1^{},B_1^{\prime \prime }`$ are the tangent lines at flexes $`P_1,P_2,P_3`$ (see below) of $`B_3`$. Three $`A_1`$ are the intersections of lines. We know that $`B_3L_PL_QL_R`$ is of torus type if and only if $`P,Q,R`$ are colinear, where $`L_P`$ is the tangent line at the flex point $`P`$. Let $`B_3:f_3(x,y)=0`$. An Example of such a sextic of non-torus type is given by
$`f_3(x,y)=y^3+((3+{\displaystyle \frac{1}{2}}I\sqrt{3})x2)y^2y+(1{\displaystyle \frac{1}{2}}I\sqrt{3})x^3+(3+{\displaystyle \frac{1}{2}}I\sqrt{3})x+2`$
$`f(x,y)=(y^3+((3+{\displaystyle \frac{1}{2}}I\sqrt{3})x2)y^2y+(1{\displaystyle \frac{1}{2}}I\sqrt{3})x^3+(3+{\displaystyle \frac{1}{2}}I\sqrt{3})x+2)`$
$`(y1)(y+1)(I\sqrt{3}(x1)y)`$
We can moreover explicitly compute 9 flex points $`P_1,\mathrm{},P_9`$ as follows.
$$\begin{array}{c}P_1:=(1,\mathrm{\hspace{0.17em}0},\mathrm{\hspace{0.17em}1}),P_2:=(0,\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}1}),P_3:=(0,1,\mathrm{\hspace{0.17em}1}),P_4:=(0,\mathrm{\hspace{0.17em}2},\mathrm{\hspace{0.17em}1}),P_5:=(\frac{3}{5},\frac{4}{5},\mathrm{\hspace{0.17em}1}),\hfill \\ \hfill P_6:=(\frac{1}{2}+\frac{1}{6}I\sqrt{3},\frac{1}{2}\frac{1}{6}I\sqrt{3},\mathrm{\hspace{0.17em}1}),P_7:=(\frac{15}{14}+\frac{3}{14}I\sqrt{3},\frac{1}{14}+\frac{3}{14}I\sqrt{3},\mathrm{\hspace{0.17em}1}),\\ \hfill P_8:=(\frac{21}{38}+\frac{9}{38}I\sqrt{3},\frac{31}{38}\frac{3}{38}I\sqrt{3},\mathrm{\hspace{0.17em}1}),P_9:=(\frac{33}{62}+\frac{3}{62}I\sqrt{3},\frac{37}{62}+\frac{9}{62}I\sqrt{3},\mathrm{\hspace{0.17em}1})\end{array}$$
Thus by a direct checking, we find the following 12 triples which are colinear.
$`𝒞_1:=[P_1,P_2,P_6],𝒞_2:=[P_1,P_3,P_7],𝒞_3:=[P_1,P_4,P_5]`$
$`𝒞_4:=[P_1,P_8,P_9],𝒞_5:=[P_2,P_3,P_4],𝒞_6:=[P_2,P_5,P_8]`$
$`𝒞_7:=[P_3,P_5,P_9],𝒞_8:=[P_3,P_6,P_8],𝒞_9:=[P_4,P_6,P_9]`$
$`𝒞_{10}:=[P_4,P_7,P_8],𝒞_{11}:=[P_5,P_6,P_7],𝒞_{12}:=[P_2,P_7,P_9]`$
### 5.2. Sextics with two cubic components: $`C=B_3+B_3^{}`$.
Now we consider sextics with two cubic curves $`B_3,B_3^{}`$. The possible configurations are
1. $`C=B_3+B_3^{}`$.
1. $`\mathrm{\Sigma }(C)=[A_{17}],[A_{17},A_1],[A_{17},2A_1]`$.
2. $`\mathrm{\Sigma }(C)=[A_{11},A_5],[A_{11},A_5,A_1],[A_{11},A_5,2A_1],`$.
3. $`\mathrm{\Sigma }(C)=[A_{11},2A_2,3A_1]`$
4. $`\mathrm{\Sigma }(C)=[3A_5],[3A_5+A_1],[3A_5,2A_1]`$.
First we consider two cubics $`B_3,B_3^{}`$ which are tangent at the origin with intersection number $`9`$. Let $`f(x,y)=0`$ and $`f^{}(x,y)=0`$ be the defining polynomials of $`B_3`$ and $`B_3^{}`$ respectively and we may assume that the tangent line of $`B_3`$ is given by $`y=0`$. Let $`y=_{i=2}^{\mathrm{}}t_ix^i`$ be the solution of $`f(x,y)=0`$ at $`O`$. Then by the assumption $`I(B_3,B_3^{};O)=9`$, we must have $`\text{ord}_xf^{}(x,_{i=2}^{\mathrm{}}t_ix^i)=9.`$
###### Lemma 14.
The sextic $`C:=B_3B_3^{}`$ is of torus type if and only if $`t_2=0`$. This implies that $`O`$ is a flex of $`B_3`$.
Proof. This assertion is given by Artal in . Our proof is computational. In fact, if $`t_2=0`$, $`O`$ is a flex for both $`B_3,B_3^{}`$ and we see that $`y=0`$ is a flex tangent line for $`B_3,B_3^{}`$. Thus by Tokunaga’s criterion, $`y^2=0`$ is the conic which gives a linear torus decomposition. For the detail about linear torus decomposition, we refer .
Assume that $`t_20`$ and we prove that any such $`C`$ is of non-torus type. In fact, supposing $`C`$ to be a sextic of torus type, take a torus decomposition $`f(x,y)f^{}(x,y)=f_2(x,y)^3+f_3(x,y)^2`$. Put $`y_1:=y_{i=2}^{\mathrm{}}t_ix^i`$. So the assumption implies that
$$f(x,y_1+\varphi (x))\times f^{}(x,y_1+\varphi (x))=y_1(y_1+cx^9)+\text{(higher terms)},c0$$
and thus $`y_1^{}:=y_{i=2}^8t_ix^i`$ is the maximal contact coordinate and it is also the solution of $`f_3(x,y)=0`$ in $`y`$ mod $`x^9`$ and $`f_2(x,_{i=2}^8t_ix^i)0\text{modulo}(x^6)`$. For the existence of a non-trivial conic $`f_2(x,y)`$, we see that the coefficient must satisfy:
(2) $`J_0:=3t_4t_3t_2+2t_3^3+t_5t_2^2=0`$
(conics are five dimensional but we have 6 equation $`coeff(f_2(x,_{i=2}^8t_ix^i),x,j)=0`$ for $`j=0,\mathrm{},5`$). Then we examine the other equations
$$(\mathrm{})\text{Coeff}(g(x,\underset{i=2}{\overset{8}{}}t_ix^i),x,j)=0,j=0,\mathrm{},8$$
where $`g(x,y)`$ is a cubic which corresponds to either $`f(x,y)`$ or $`f^{}(x,y)`$. Write $`g(x,y)`$ as a generic cubic form with 10 coefficients (but by scalar multiplication, one coefficient can be normalized to be 1, say that of $`x^6`$ is 1, and we have 9 free coefficients), we solve the equations ($`\mathrm{}`$) from $`j=0`$ to $`j=0`$ to $`j=8`$ to express the coefficients in rational functions of $`t_2,\mathrm{},t_7`$. At the last step, we have one free coefficient undetermined, say the coefficient $`c`$ of $`x^ay^b`$ and a linear equation $`\text{Coeff}(g(x,_{i=2}^8t_ix^i),x,8)=0`$. This is written as $`K_1c+K_0=0`$ where $`K_1,K_0`$ are rational functions of $`t_2,\mathrm{},t_7`$. Thus to have two non-trivial cubic solutions, we need to have $`K_1=K_0=0`$. Now we can easily check that there are no solutions if we assume that $`J_0(t_2,\mathrm{},t_5)=0`$. The other assertion will be checked in the explicit construction of examples.∎
If $`t_20`$, there is no line $`L`$ such that $`L`$ intersects only at $`O`$. Thus the sextic can not be of torus type by the classification in . However the above argument is useful for the explicit computation of non-torus sextics.
(1-a) Let us consider the case $`[A_{17}],[A_{17},A_1],[A_{17},2A_1]`$. The sextics of non-torus type with above configurations are obtained using above computation ($`t_20`$). Their Zariski partners are cubics intersecting at flex points.
(a-1) $`[A_{17}],C=B_3+B_3^{}`$, $`C:=\{f_3(x,y)g_3(x,y)=0\}`$ where
$`f_3:=y^3y^2+(x^2+x)yx^3+x`$
$`g_3:=x{\displaystyle \frac{10}{9}}y^3x^2y+{\displaystyle \frac{10}{9}}xyy^2{\displaystyle \frac{10}{9}}x^3`$
(a-2) $`[A_{17}+A_1],C=B_3+B_3^{}`$ with $`B_3`$ has a node: $`C:=\{f_3(x,y)g_3(x,y)=0\}`$ where
$`f_3:=x3xy+3x^2y2x^2y^2+2y^3+4y^2x+x^3`$
$`g_3:=x+xy+11x^2y10x^2y^22y^3+8y^2x+5x^3`$
(a-3) $`[A_{17}+2A_1],C=B_3+B_3^{}`$ $`C:=\{f_3(x,y)g_3(x,y)=0\}`$, where cubics are nodal and
$$\begin{array}{c}f_3:=\frac{1}{48}I(48x96x^2264xy+124y^348y^2+411y^2x+264x^2y+48x^3\hfill \\ \hfill 104Ixy\sqrt{3}+181Iy^2x\sqrt{3}+104Ix^2y\sqrt{3}+16Ix\sqrt{3}32Ix^2\sqrt{3}\\ \hfill +16Ix^3\sqrt{3}16Iy^2\sqrt{3}+52Iy^3\sqrt{3})\sqrt{3}/(1+I\sqrt{3})\\ \hfill g_3:=\frac{1}{8}(48x72x^2+10Ix^2y\sqrt{3}25Iy^2x\sqrt{3}+56Ixy\sqrt{3}216xy+68y^348y^2\\ \hfill +231y^2x+138x^2y+23x^3+7Ix^3\sqrt{3}+8Ix^2\sqrt{3}16Ix\sqrt{3}12Iy^3\sqrt{3}\\ \hfill +16Iy^2\sqrt{3})/((3+I\sqrt{3})(1+I\sqrt{3}))\end{array}$$
### 5.3. Exceptional configuration: $`[A_{11},2A_2,3A_1]`$ with two cubic components.
In this case, we do the similar computation. We compute sextics $`C=B_3B_3^{}`$ such that $`B_3`$ and $`B_3^{}`$ have two $`A_2`$ at $`(0,1)`$ and $`(1,0)`$ respectively and they intersect at $`(0,1)`$ with intersection number $`6`$ to make $`A_{11}`$. We have the following sextics of non-torus type.
$`f(x,y):=f_3(x,y)g_3(x,y)`$
$`f_3:={\displaystyle \frac{1}{44652}}(47\sqrt{3}+168+195I+74I\sqrt{3})(195y+168Iy^3168Iy^2168Iy`$
$`156Ix^2+74\sqrt{3}y^374\sqrt{3}y^274\sqrt{3}y60\sqrt{3}x^2+60\sqrt{3}yx^247I\sqrt{3}y^3`$
$`+47I\sqrt{3}y^2+156Iyx^2+48Ix^2\sqrt{3}+168x^2+47Iy\sqrt{3}48Iyx^2\sqrt{3}`$
$`168yx^2+195y^2195y^3195+168I+74\sqrt{3}47I\sqrt{3}+48x^3)`$
$`g_3(x,y)={\displaystyle \frac{1}{35636460}}(30331989I+1313I\sqrt{3}+1361\sqrt{3})(90720x+14898y15336yx`$
$`+3021\sqrt{3}52722x^211749\sqrt{3}y^38305\sqrt{3}y^2+6465\sqrt{3}y+3021\sqrt{3}x^2`$
$`1785\sqrt{3}yx^2+438yx^287573y^224417y^3+75384y^2x+65388x^3`$
$`6042x\sqrt{3}7839Iy^3+36333Iy^2+60705Iy+16533Ix^2+5736I\sqrt{3}`$
$`33066Ix4680yx\sqrt{3}+1362y^2x\sqrt{3}6383I\sqrt{3}y^3+22753I\sqrt{3}y^2`$
$`+231Iyx^2+5736Ix^2\sqrt{3}+34872Iy\sqrt{3}60936Iyx11472Ix\sqrt{3}`$
$`27870Iy^2x5532Iyx^2\sqrt{3}29340Iyx\sqrt{3}17868Iy^2x\sqrt{3}`$
$`+78054+16533I)`$
### 5.4. Examples of (b) and (d).
As the corresponding sextics of torus type are linear, we only need to check the singularities are not colinear.
(b) $`\mathrm{\Sigma }(C)\{A_{11},A_5\}`$. We put $`A_{11}`$ at (0,0) with tangent line $`x=0`$ and $`A_5`$ at (1,0).
(b-1) $`\mathrm{\Sigma }(C)=[A_{11},A_5]:`$
$`f(x,y)=`$ $`(y^3+(9x1)y^2+7x^38x^2+x)\times `$
$`(2y^3+(5x1)y^2+(x^2+x)y+4x^35x^2+x)`$
(b-2) $`\mathrm{\Sigma }(C)=[A_{11},A_5,A_1]:`$
$$\begin{array}{c}f(x,y)=\frac{1}{55}(16xy2y^2+16x^2y+14y^2x8x^2+5x^3+3x)\hfill \\ \hfill (11xy+4y^211x^2y+98y^2x5x^2+11x^3+14y^36x)\end{array}$$
(b-3) $`\mathrm{\Sigma }(C)=[A_{11},A_5,2A_1]:`$
$$\begin{array}{c}f(x,y):=(175y^3+11x^2\sqrt{3}88x^2y^3\sqrt{3}30y^2x\sqrt{3}18yx^2\sqrt{3}+27x+61x^3+48y^2\hfill \\ \hfill 94yx+83yx^2126y^2x+Ix^2\sqrt{3}+5Iy^2\sqrt{3}17Iy^2xIx\sqrt{3}3Iyx\\ \hfill +11I\sqrt{3}y^36y^2\sqrt{3}11x\sqrt{3}8Ix^2+21Iy^2+8Ix27Iy^3+25Iy^2x\sqrt{3}\\ \hfill +8Iyx\sqrt{3}+15Iyx^2\sqrt{3}+2Iyx^2+27yx\sqrt{3})(175y^3+11x^2\sqrt{3}88x^2\\ \hfill +y^3\sqrt{3}30y^2x\sqrt{3}+18yx^2\sqrt{3}+27x+61x^3+48y^2+94yx83yx^2\\ \hfill 126y^2x3Iyx+11I\sqrt{3}y^36y^2\sqrt{3}11x\sqrt{3}27Iy^3+8Iyx\sqrt{3}\\ \hfill +15Iyx^2\sqrt{3}+2Iyx^227yx\sqrt{3}25Iy^2x\sqrt{3}+8Ix^221Iy^28Ix\\ \hfill +Ix\sqrt{3}Ix^2\sqrt{3}5Iy^2\sqrt{3}+17Iy^2x)\end{array}$$
(d) $`\mathrm{\Sigma }(C)\{3A_5\}`$, $`C=B_3+B_3^{}`$. We put three $`A_5`$ at $`(0,1),(0,1)`$ and $`(1,0)`$.
$$\begin{array}{c}[3A_5]:f(x,y):=(\frac{3}{7}+xy+y^3\frac{3}{7}x^3x^2\frac{1}{7}yx^2y^2x\frac{3}{7}y^2)\hfill \\ \hfill (\frac{4}{5}xy+1+\frac{4}{5}y^2x\frac{3}{5}x^2+\frac{3}{5}yx^2+y^3y^2+\frac{2}{5}x^3)\end{array}$$
$$\begin{array}{c}[3A_5,A_1]:f(x,y):=(y^3+\frac{57}{4}y^2x+\frac{1}{4}y^2yyx+\frac{1}{4}x^3+\frac{1}{4}x^2\frac{1}{4}x\frac{1}{4})\hfill \\ \hfill (y^3+\frac{71}{5}y^2xy^2+\frac{49}{5}yx^220yxy+\frac{67}{5}x^3\frac{101}{5}x^2+\frac{29}{5}x+1)\end{array}$$
$$\begin{array}{c}[3A_5+2A_1]:f(x,y):=(21y^2x12I\sqrt{3}y^2x+12yx^2I\sqrt{3}yx^2+12yx2I\sqrt{3}yxIy\sqrt{3}\hfill \\ \hfill +Iy^3\sqrt{3}3+3y^2+3x^33x+3x^2)(1+49x^317y^2x+84yx^263x^2\\ \hfill 12yx+y^2+15x7I\sqrt{3}yx^26I\sqrt{3}yx4I\sqrt{3}y^2x+Iy\sqrt{3}Iy^3\sqrt{3})\end{array}$$
## 6. Three conics
In this section, we study the last case $`C=B_2+B_2^{}+B_2^{\prime \prime }`$ with the configuration of the singularities $`[3A_5,3A_1]`$. Such a sextic is given when each pair of conics are intersecting at two points: at one point, with intersection multiplicity 3 and at another point, transversely. We can understand Zariski pairs in this situation using conical flex points. Assume that the respective defining polynomials of $`B_2,B_2^{},B_2^{\prime \prime }`$ are $`f_2(x,y),g_2(x,y),h_2(x,y)`$ and the location of two $`A_5`$’s are $`P_1=(0,1),P_2=(0,1)`$ with respective tangent cones are $`y1=0`$. We assume further $`P_1B_2B_2^{}`$ and $`P_2B_2B_2^{\prime \prime }`$. We fix $`B_2,B_2^{}`$ generically and consider a linear system $`\mathrm{\Phi }`$ of conics $`B_2^{\prime \prime }`$ of dimension 2 such that $`B_2^{\prime \prime }`$ and $`B_2`$ are tangent at $`P_2`$. Under this situation we assert that
###### Proposition 15.
There exist 5 conical flex points $`Q_i,i=1,\mathrm{},5`$ on $`B_2^{}`$ with respect to $`\mathrm{\Phi }`$ so that $`Q_1`$ is a conical flex of torus type and the other are of non-torus type.
Proof. To avoid the complexity of the equation, we choose a generic $`B_2,B_2^{}`$ so that
$`f_2(x,y)=(y^21+x^2)`$
$`g_2(x,y)=(y^2+({\displaystyle \frac{2}{15}}\sqrt{130}x{\displaystyle \frac{2}{3}})y+{\displaystyle \frac{2}{3}}x^2+{\displaystyle \frac{2}{15}}\sqrt{130}x{\displaystyle \frac{1}{3}})`$
We find 5 conical flex points on $`B_2^{}`$:
$$\begin{array}{c}Q_1=(\frac{1}{9}\sqrt{130},\frac{17}{9})\hfill \\ \hfill Q_2=(\frac{7}{3}I\sqrt{10}\sqrt{3}\frac{2}{3}I\sqrt{13}\sqrt{10}\sqrt{3},4+\sqrt{13}+\frac{5}{3}I\sqrt{13}\sqrt{3}\frac{19}{3}I\sqrt{3})\\ \hfill Q3=(\frac{7}{3}I\sqrt{10}\sqrt{3}+\frac{2}{3}I\sqrt{13}\sqrt{10}\sqrt{3},4+\frac{19}{3}I\sqrt{3}+\sqrt{13}\frac{5}{3}I\sqrt{13}\sqrt{3})\\ \hfill Q4=(\frac{7}{3}I\sqrt{10}\sqrt{3}\frac{2}{3}I\sqrt{13}\sqrt{10}\sqrt{3},4\frac{5}{3}I\sqrt{13}\sqrt{3}\frac{19}{3}I\sqrt{3}\sqrt{13}),\\ \hfill Q5=(\frac{7}{3}I\sqrt{10}\sqrt{3}+\frac{2}{3}I\sqrt{13}\sqrt{10}\sqrt{3},4+\frac{19}{3}I\sqrt{3}\sqrt{13}+\frac{5}{3}I\sqrt{13}\sqrt{3})\end{array}$$
$`Q_1`$ gives a sextic of torus type so that $`B_2^{\prime \prime }`$ is given by
$$\begin{array}{c}h_2(x,y)=(\frac{338}{201}y\frac{104}{1005}y\sqrt{130}x\frac{104}{1005}\sqrt{130}x+\frac{137}{201}+y^2+\frac{32}{201}x^2)\hfill \end{array}$$
The other conical flex points give sextics of non-torus type. For example, $`Q_2`$ gives $`B_2^{\prime \prime }`$ described as:
$$\begin{array}{c}h_2(x,y)=75+150I\sqrt{3}+50\sqrt{13}45I\sqrt{13}\sqrt{3}104\sqrt{13}\sqrt{10}x80\sqrt{10}x\hfill \\ \hfill +72I\sqrt{10}x\sqrt{3}+300Iy\sqrt{3}+790y90Iy\sqrt{13}\sqrt{3}+100y\sqrt{13}\\ \hfill 104y\sqrt{13}\sqrt{10}x80xy\sqrt{10}+72Ixy\sqrt{10}\sqrt{3}45Iy^2\sqrt{13}\sqrt{3}+50y^2\sqrt{13}\\ \hfill +150Iy^2\sqrt{3}+715y^2+320x^2\end{array}$$ |
warning/0507/hep-ph0507084.html | ar5iv | text | # Heavy quark bound states above 𝑇_𝑐
## I Introduction
The anomalous suppression of the $`J/\psi `$ production in heavy ion collisions, which has been experimentally observed Na50 ; Na50bis in the depletion of the dilepton multiplicity in the region of invariant mass corresponding to the $`J/\psi `$ meson, was proposed long time ago as a possibly unambiguous signal of the onset of deconfinement satz . Indeed in Ref. satz it is argued that charmonium states can only be produced in the first instants after the nucleus-nucleus collision, before the formation of a thermalized QGP. Then, in their path through the deconfined medium, the original $`c\overline{c}`$ bound states tend to melt, since the binding (colour) Coulomb potential is screened by the large number of colour charges. This, in turn, produces an anomalous (with respect to normal nuclear absorption) drop in the $`J/\psi `$ yields.
In this picture it is implicitly assumed that, once the charmonium dissociates, the heavy quarks hadronize by combining with light quarks only (recombination leading to a secondary $`J/\psi `$ production is neglected). This assumption is certainly justified at the SPS conditions, due to the very small number of $`c\overline{c}`$ pairs produced per collision ($`N_{c\overline{c}}0.2`$ in a central collision), but at RHIC ($`N_{c\overline{c}}10`$) and LHC ($`N_{c\overline{c}}200`$) energies it is no longer warranted tew .
Moreover in a hadronic collisions only about $`60\%`$ of the observed $`J/\psi `$’s are directly produced, the remaining stemming from the decays of excited charmonium states (notably the $`\chi _c`$ and the $`\psi ^{}`$). Since each $`c\overline{c}`$ bound state dissociates at a different temperature, a model of sequential suppression was developed, with the aim of reproducing the $`J/\psi `$ suppression pattern as a function of the energy density reached in the heavy ion collision (the highest temperatures and energy densities being reached in the most central collisions) satz2 ; diga1 ; diga2 ; kar ; Kar97 . SPS experimental data for Pb-Pb collisions at different centralities seem indeed to support the dissociation pattern predicted by this model Na50 ; Na50bis .
Alternative mechanisms for the $`J/\psi `$ production, like the Statistical Coalescence Model (SCM), have also been proposed scm . In the SCM one assumes again that all the $`c\overline{c}`$ pairs are produced in hard processes at the initial stage of the collision. Any heavy quark bound state, if present, is assumed to melt in the QGP phase and the number of $`c\overline{c}`$ pairs in the fireball is considered fixed. At the chemical freeze-out open and hidden charmed hadrons are then produced with multiplicity ratios<sup>1</sup><sup>1</sup>1Due to the low rate of inelastic reactions full chemical equilibrium cannot be reached by charmed hadrons: their total multiplicity measured at SPS stays well above the thermal value. fixed by their masses, according to the laws of statistical mechanics. Hence, in such a scheme, the measured $`J/\psi `$ multiplicity is not related to the presence of charmonium bound states in the plasma phase, but to the statistical hadronization of the initially produced $`c`$($`\overline{c}`$) (anti-)quarks. This may lead, at LHC energies, to a completely different picture characterized by an enhanced charmonium production even if all the $`c\overline{c}`$ bound states dissociate during the plasma phase.
In any case the hypothesis that all the primary produced $`J/\psi `$’s melt during the QGP lifetime is hardly realized at SPS conditions. Hence models have been developed gran1 ; gran2 attempting to account both for the initial state production, eventually subject to in-medium dissociation, and for the thermal production at the hadronization.
Concerning the heavy quarkonia in the QGP phase, recent lattice data dat ; asa1 ; asa2 in quenched approximation (hence neglecting effects arising from virtual processes involving dynamical fermions), which display narrow peaks for the charmonium spectral functions in the pseudoscalar and vector channels (even up to $`T2T_c`$), seem to point to the existence of heavy quark bound states up to temperatures above $`T_c`$<sup>2</sup><sup>2</sup>2Actually, in Ref. asa1 the analysis has been done for $`s\overline{s}`$ mesons.. Clearly these results, if confirmed, would entail striking experimental consequences.
Actually, the meson spectral functions cannot be measured directly on the lattice. From the numerical simulations one gets the current-current correlation function along the (imaginary) temporal direction on a finite number of points. Such a correlator corresponds to the convolution of the meson spectral function with a thermal kernel. The spectral function can then be obtained only indirectly. With this aim in Refs. dat ; asa1 ; asa2 a procedure called Maximum Entropy Method (MEM) has been adopted. Clearly, an independent check of the results obtained with MEM appears desirable. This indeed is our scope in the present paper.
For this purpose, we first extract from lattice data a heavy quark potential accounting for thermal effects and then we solve the Schrödinger equation for the charmonium (and bottomonium). As shown in Refs. mc ; Kac1 ; Kac2 ; Kac3 ; Kac4 , from the Polyakov loop correlation function it is possible to extract the free energy (in the different color channels) of a heavy quark-antiquark pair placed at a distance $`r`$ in a thermal bath of gluons and light dynamical fermions. Once a good parameterization of the color singlet free energy is obtained, the entropy and internal energy contributions can be disentangled. Since the quarks acting as static sources of the color field are considered infinitely heavy, the internal energy coincides with the potential. The latter is then inserted into the Schrödinger equation, from which the binding energy of the different stable states — if there are any — and their evolution with the temperature are obtained.
Indeed a clear distinction between the $`Q\overline{Q}`$ free and potential energies is necessary in order to get a reliable estimate of the quarkonium dissociation temperature $`T_d`$ in the different spin-parity channels.
In Refs. diga1 ; diga2 ; wong1 ; wong2 , where the colour singlet free energy was directly inserted into the Schrödinger equation, the dissociation temperatures $`T_d=1.10T_c`$ diga2 and $`T_d=0.99T_c`$ wong1 were found for the $`J/\psi `$, all the other charmonium states melting well below $`T_c`$.
On the other hand, in Ref. wong3 , where a parameterization of the lattice color singlet potential (in the quenched approximation) was used, the temperature $`T_d`$ for the spontaneous dissociation of the $`J/\psi `$ was estimated to occur at about $`2T_c`$, a value even larger than the one obtained from the spectral analysis performed in Refs. dat ; asa1 ; asa2 . Also results for different charmonium and bottomonium states have been reported in Ref. wong3 . The case $`N_f=2`$ was addressed in Ref. shury , where the $`J/\psi `$ meson was found to be bound till $`T2.7T_c`$.
Actually, even if the potential supports the existence of bound states, other physical processes may lead to the dissociation of the quarkonium. First, if the $`Q\overline{Q}`$ binding energy is lower than the temperature — and assuming that the quarkonia have reached the thermal equilibrium with the plasma — a certain fraction of their total number will be thermally excited to resonant states according to a Bose-Einstein distribution: such a process is referred to as thermal dissociation wong1 ; wong3 . Furthermore, the collisions with the gluons and the light quarks of the plasma may lead to the collisional dissociation of the quarkonium. In this connection, the reaction $`g+J/\psi c+\overline{c}`$ was studied in detail in Ref. wong3 . Hence, in spite of the presence of a bound state solution of the Schrödinger equation till $`T2T_c`$, the $`J/\psi `$ turned out to be really stable only up to temperatures lower than $`2t_c`$ wong3 . Clearly the two processes above mentioned are not encoded in the Polyakov loop correlation function, where the heavy quarks act as unthermalized static sources of the color field, and have to be accounted for a posteriori.
Of course, if the process $`g+J/\psi c+\overline{c}`$ can lead to the dissociation of the charmonium, the same reaction can also occur in the opposite direction. Hence a consistent calculation of $`J/\psi `$ multiplicity implies the solution of a kinetic rate equation integrated over the lifetime of the QGP phase in which both processes (dissociation and recombination) enter wong3 ; raf . To carry out this detailed balance calculation the knowledge of the $`J/\psi `$ binding energy and wave function in the thermal bath turns out to be an important input. This is of relevance because, as mentioned above, the usual assumption in considering the $`J/\psi `$ suppression as a signature of deconfinement is that its production can occur only in the very initial stage of the collision. Really, if at SPS the role played by recombination is numerically negligible, this is no longer true at RHIC as pointed out in Ref. raf .
In any case in this paper we limit ourselves to check the existence of bound state solutions of the Schrödinger equation. A quantitative study of the dissociation and recombination processes is left for future work.
Here we take advantage of all the available lattice data, obtained not only in quenched QCD ($`N_f=0`$), but also including two and, more recently, three light flavors. We are then in a position to study also the flavor dependence of the dissociation process, a perspective not yet achieved by the parallel studies of the spectral functions, which are, as already mentioned, only available in quenched QCD.
The present paper is organized as follows. In Sec. II we present a parameterization of the color singlet free energy lattice data for the cases $`N_f=0`$ Kac1 , $`N_f=2`$ Kac4 and $`N_f=3`$ Petr , from which the heavy quark potential is obtained. In Sec. III we solve numerically the associated Schrödinger equation at different temperatures for the charmonium and bottomonium states, thus determining their dissociation temperature. Finally, in Sec. IV we present our conclusions.
## II Parameterization of the lattice data
In this section we provide a unified parameterization of the lattice data for the *color singlet* $`Q\overline{Q}`$ *free energy* $`F_1`$ in the case of quenched Kac1 , 2-flavor Kac4 and 3-flavor Petr QCD. The lattice findings are shown (in dimensionless units) in Fig. 1.
For what concerns the critical temperature we assume the values $`T_c=270`$ MeV ($`N_f=0`$), $`T_c=202`$ MeV ($`N_f=2`$) and $`T_c=193`$ MeV ($`N_f=3`$) given in Ref. Kac4 .
The free energy on the lattice is defined up to an additive normalization constant, which has to be fixed using some physical constraint. In Refs. Kac1 ; Kac4 ; Petr it has been normalized to match, at the shortest available distance for each temperature, the $`T=0`$ heavy quark potential. Such a normalization amounts to make the (reasonable) assumption that thermal effects become negligible at very short distances.
In Refs. Kac1 ; Kac4 the zero temperature heavy quark potential has been determined through a best fit procedure of the available $`T=0`$ lattice data with a Cornell-like parameterization<sup>3</sup><sup>3</sup>3Note that the $`1/r`$ term accounts for two different physical processes: the perturbative one-gluon-exchange at short distances and the transverse string fluctuations at large distances.:
$$\frac{V(r)}{\sqrt{\sigma }}=\frac{4}{3}\frac{\alpha }{r\sqrt{\sigma }}+\sqrt{\sigma }r,$$
(1)
$`\sigma `$ representing the string tension. The values $`\alpha =0.195(1)`$ for the case $`N_f=0`$ and $`\alpha =0.212(3)`$ for the case $`N_f=2`$ are given in Ref. Kac4 , where the value $`\sqrt{\sigma }=420`$ MeV is employed to translate the lattice results into physical units.
In Ref. Petr a similar parameterization (Cornell potential plus a $`1/r^2`$ term to mimic the effects of asymptotic freedom at the shortest distances reachable on the lattice) is employed for the case $`N_f=3`$. Note that in that work the free energy is provided directly in physical units: however, for the sake of comparison, we show it in dimensionless units, using for the string tension the value $`\sqrt{\sigma }=460`$ MeV extracted from the parameterization of the $`T=0`$ potential given in Ref. Petr .
In past calculations diga1 ; diga2 ; wong1 ; wong2 the free energy has been often identified with the heavy quark potential and inserted directly into the Schrödinger equation. However, a better candidate for a more appropriate finite temperature potential is given by the internal energy of the $`Q\overline{Q}`$ system, defined by the well known relation
$$F=UTS,$$
(2)
where
$$U=T^2\frac{(F/T)}{T}$$
(3)
is the internal energy and
$$S=\frac{F}{T}$$
(4)
is the entropy. One can see from the data in Fig. 1 that the role played by the entropy is more relevant at large distances.
Getting the internal energy from the free energy involves a derivative of the latter with respect to the temperature: it is thus clear that one needs an accurate parameterization of the temperature dependence of $`F`$.
In order to establish a suitable form for this parameterization of the lattice data, we first consider the two limits in which the underlying physics is supposed to be known.
At very short distances ($`r1/T`$) thermal effects are negligible and the colour singlet free energy is dominated by the perturbative one-gluon exchange with the typical behaviour:
$$F_1(r,T)\underset{rT1}{}\frac{4}{3}\frac{\alpha (r)}{r},$$
(5)
the coupling $`\alpha `$ depending only upon the $`Q\overline{Q}`$ separation.
On the other hand, for $`TT_c`$, the large distances ($`r1/T`$) behaviour of the free energy is expected to be dominated by the exchange of a resummed electrostatic gluon leading to the expression nad :
$$F_1(r,T)\underset{rT1}{}\frac{4}{3}\frac{\alpha (T)}{r}e^{m_D(T)r}+F_1(r=\mathrm{},T).$$
(6)
In this limit, the coupling $`\alpha `$ is a function of the temperature and $`m_D(T)`$ is the Debye screening mass arising from the dressing of the electrostatic gluon.
In the two above limits, the running of the coupling is determined by a Renormalization Group Equation (RGE), allowing to express $`\alpha =g^2/4\pi `$ as a known (at least in the weak coupling regime) function of an energy scale $`\mu `$. In the short distance limit the relevant energy scale is given by the inverse of the distance ($`\mu 1/r`$), while for large separations the major role in setting the scale is expected to be played by the temperature ($`\mu T`$).
Actually, in order to solve the Schrödinger equation for heavy quarkonia, one really needs a parameterization of the free energy covering the whole range of distances. For this purpose, on the basis of Eqs. (5) and (6) it appears convenient to cast the dependence of $`F_1`$ on $`r`$ and $`T`$ into the following functional form<sup>4</sup><sup>4</sup>4Note that the mass $`M(T)`$ appearing in this exponential will not necessarily coincide with the Debye screening mass of Eq. (6), the latter being determined by fitting only the large distance data.:
$$F_1(r,T)=\frac{4}{3}\frac{\alpha (r,T)}{r}e^{M(T)r}+C(T).$$
(7)
In order to recover from the above the short and large distance limits given by Eqs. (5) and (6) one can, e.g., assume $`\alpha `$ to depend on the following combination of $`r`$ and $`T`$:
$$\alpha (r,T)=\alpha (\mu =c_r/r+c_tT),$$
(8)
where $`\alpha (\mu )`$ is obtained by solving the RGE, while $`c_r`$ and $`c_t`$ are numerical coefficients to be fixed through a best fit of the data.
If supported by the data, Eq. (8) would allow to interpolate between the short distance regime ($`r1/T`$), where $`\mu 1/r`$, and the long distance one ($`r1/T`$), where on the contrary $`\mu T`$.
We now describe in detail the fitting procedure we have employed. The data for the colour singlet free energy, taken from Ref. Kac1 ($`N_f=0`$), Ref. Kac4 ($`N_f=2`$) and Ref. Petr ($`N_f=3`$), are displayed in Fig. 1 in dimensionless units, namely $`y=F_1/\sqrt{\sigma }`$ and $`x=r\sqrt{\sigma }`$, for different values of the temperature both below and above $`T_c`$. For each value of $`T>T_c`$ the data have been parameterized with the following (four parameter) fitting function:
$$y=\frac{4}{3}\frac{\alpha (\stackrel{~}{\mu })}{x}e^{a_3x}+a_0\text{with}\stackrel{~}{\mu }=\frac{a_1}{x}+a_2$$
(9)
where for $`\alpha (\stackrel{~}{\mu })`$ we use the RGE result obtained with the two-loop QCD beta-function quoted in Appendix A, the dimensionless variable $`\stackrel{~}{\mu }`$ being identified with the ratio $`\mu /\mathrm{\Lambda }_{\text{QCD}}`$.
The fitting procedure yields a very mild dependence on $`T`$ for the parameters $`a_1`$ and $`a_2`$. Actually, on the basis of the above discussion, one would have expected the coefficient $`a_2`$ to scale linearly with the temperature, but our finding might just signal that the range of temperatures spanned here is still not in the asymptotic regime $`TT_c`$. On the other hand, a constant value for $`a_1`$ is in agreement with the ansatz of Eq. (8). This has suggested the following procedure: a weighted average of the values obtained for $`a_1`$ has been performed, yielding $`a_1=0.2719(2)`$ for $`N_f=0`$, $`a_1=0.2687(7)`$ for $`N_f=2`$ and $`a_1=0.2354(17)`$ for $`N_f=3`$. We have then fitted again the data keeping $`a_1`$ fixed and using only $`a_0`$, $`a_2`$ and $`a_3`$ as free parameters. This parameterization works well — yielding values of $`\chi ^2`$ per degree of freedom of the order of 1 at all the temperatures — and it is compared to the data in Fig. 2.
The finite temperature lattice data are limited to distances $`r0.1÷0.2`$ fm. Since the data examined in this work have been normalized assuming that at short distances thermal effects are negligible, we should check that our parameterization does not introduce any sizable (and spurious) temperature dependence for small values of $`r`$, remaining in this region close to the $`T=0`$ perturbative potential (we remind that at $`T=0`$ free energy and internal energy coincide). In this respect our choice of refitting the data keeping the coefficient $`a_1`$ fixed fulfills this requirement.
The behaviour of our parameterization of the free energy for $`r<0.1`$ fm is displayed, at three different temperatures, in Fig. 3 where we indeed see that spurious short distance thermal effects appears negligible. It is gratifying (and somewhat surprising) that our curves seems to interpolate smoothly between the $`T=0`$ Cornell potential and the short distance perturbative potential. The Cornell curve reported in the figure is the one employed in Ref. Kac4 to fix the normalization of the free energy at the different temperatures. As already discussed at the beginning of this section it was obtained by fitting zero temperature lattice data which cover distances $`r0.05`$ fm. Of course for shorter distances the Cornell parameterization cannot account for running coupling effects (asymptotic freedom) and the perturbative calculation pet ; mel ; sch should provide more reliable results. More details on the one- and two-loop perturbative potential used in Fig. 3 are given in Appendix A.
The curves $`V_{\text{string}}`$ reported in Fig. 3 refer to the $`Q\overline{Q}`$ potential
$$V(r)=\frac{\pi }{12}\frac{1}{r}+\sigma r$$
(10)
obtained in the Nambu string model, the term $`1/r`$ arising from the quantum fluctuations of the flux tube in the transverse directions.
As a next step we have provided a parameterization of the $`T`$-dependence of $`a_0`$, $`a_2`$ and $`a_3`$, which allows one to perform analytically the derivative of $`F_1`$ with respect to the temperature (actually, for the calculation of the $`Q\overline{Q}`$ binding energies the knowledge of $`a_0`$ is unnecessary: we show also this contribution to the potential for completeness). Lacking compelling physical hints to the form of the $`T`$-dependence, we have sought for the simplest and smoothest expressions yielding a satisfactory interpolation of the values calculated from the fit to the lattice data (see Appendix B for details). The parameters and their interpolations are displayed in Fig. 4.
The value of $`a_0(T)`$ at each temperature is essentially fixed by the normalization of the data, which has been determined using similar procedures in all the lattice calculations. Notably this parameter, once plotted as a function of $`T/T_c`$, turns out to be fairly close for all values of $`N_f`$. On the other hand, the remaining two parameters display some quantitative flavor dependence, but within the same qualitative pattern. An exception to this behavior is represented by $`a_2(T)`$ for $`N_f=3`$, although one should note that the three-flavor case is the one where the parameterization (9) has the largest uncertainties. Moreover, the variation of $`a_2(T)`$ with $`T`$ is actually magnified by the scale of the figure, being generally around 20% for the range of temperatures considered here. As we shall see the resulting effective potential is rather robust with respect to these variations.
Note also that the strongest variation of the parameters with the temperature is confined around $`T`$ very close to $`T_c`$. Since extracting the internal energy from $`F_1`$ involves a derivative with respect to $`T`$, this is the temperature domain where the sensitivity to the details of the parameterization might be high. For this reason we felt it safer to use the resulting potentials at temperatures larger than $`T_c`$, say for $`T1.05T_c`$, where they turn out to be stable with respect to changes in the parameterization.
Finally, we recall that the quenched QCD data Kac1 are actually available also at temperatures much larger than those for the unquenched cases: here we limit our analysis up to about twice the critical temperature, since this is the range where lattice calculations with two and three flavors are available and, moreover, since this is the range of interest for the problem of the $`J/\psi `$ dissociation. Here we just remark that an analysis of the $`N_f=0`$ data at larger temperatures (till $`T5T_c`$) with the parameterization of Eq. (9) shows that the fitting parameters maintain, also at these high temperatures, the trend displayed in Fig. 4.
## III $`Q\overline{Q}`$ bound states in quark-gluon plasma
The $`Q\overline{Q}`$ free energy obtained in lattice calculations has been used in the past as an input for the $`Q\overline{Q}`$ potential energy in the non-relativistic Schrödinger equation diga1 ; diga2 ; wong1 ; wong2 . More recently Kac1 ; wong3 ; shury it has been recognized that the $`Q\overline{Q}`$ potential energy can be more appropriately identified with the $`Q\overline{Q}`$ internal energy (see Eq. (3))<sup>5</sup><sup>5</sup>5Actually, the author of Ref. wong3 tries to disentangle, in the total internal energy $`U_1(r,T)`$, the gluon and the $`Q\overline{Q}`$ contributions. This leads to a lower value for the $`J/\psi `$ dissociation temperature with respect to the result $`T_d2T_c`$ obtained with the full internal energy. .
Once the temperature dependence of the color-singlet free energy $`F_1`$ has been parameterized, the corresponding color-singlet internal energy $`U_1`$ is easily obtained and one can define an effective potential
$$V_1(r,T)=U_1(r,T)U_1(r\mathrm{},T)$$
(11)
to be used in the Schrödinger equation
$$\left[\frac{\mathbf{}^2}{2\mu }+V_1(r,T)\right]\psi (𝒓,T)=\epsilon (T)\psi (𝒓,T),$$
(12)
where $`\mu `$ is the reduced mass of the $`Q\overline{Q}`$ system.
In Fig. 5 we display, at two different temperatures, the effective $`Q\overline{Q}`$ potentials that we have obtained by using the parameterizations discussed in the previous section. At high temperature the form of $`V_1`$ is the one typical of a screened Coulomb potential, with the $`N_f=3`$ potential providing a stronger attraction; at temperatures close to $`T_c`$ the shape of $`V_1`$ appears somewhat distorted because of the strong temperature dependence of $`a_2`$ and $`a_3`$ (see Fig. 4), the $`N_f=3`$ potential being still more attractive. Remarkably, the behaviour of the color singlet potential energies obtained with our procedure appears qualitatively in agreement with the one given in Refs. Kac2 ; Kac5 , where the internal energy has been directly calculated on the lattice for $`N_f=0`$ and $`N_f=2`$, respectively.
We now employ the effective potential previously derived for the study of the charmonium and bottomonium spectroscopy above the critical temperature.
For what concerns the charmonium states we plot the values of the binding energies (Fig. 6) and of the mean square radii (Fig. 8) of the different bound state solutions of Eq. (12) as functions of $`T/T_c`$. For the sake of comparison the results obtained for a different number of dynamical fermions ($`N_f=0,2,3`$) are plotted in the same panel.
In our analysis we have chosen $`m_c=1.3`$ GeV for the charm quark mass; however, we have also studied the sensitivity of our results to variations of $`m_c`$ in the range given in the Particle Data Group (PDG) listings PDG04 , namely 1.15 GeV$`<m_c<`$1.35 GeV. For reference, in Fig. 8 the arrows point to the $`T=0`$ mean square radii obtained in a potential model calculation Eic80 . Note that potential models at $`T=0`$ typically follow a different philosophy: they fix the quark mass and the confining potential parameters to reproduce the mass of the lowest states, whereas here the potential is provided by lattice calculations and the quark mass is the running mass in the $`\overline{MS}`$ scheme.
The $`1P`$ and $`2S`$ states turn out to melt at temperatures $`1.1T_cT_d1.15T_c`$. On the contrary the $`J/\psi `$ stays bound up to temperatures $`1.7T_cT_d2.3T_c`$, the precise limits depending upon $`N_f`$. The lower bound ($`T_d1.7T_c`$) refers to the quenched case and appears in striking agreement with the limiting value obtained in Ref. asa2 through the study of the $`J/\psi `$ and $`\eta _c`$ spectral functions.
Note that the free energy measured on the lattice provides a spin averaged result of the singlet and triplet channels: hence, in Figs. 6 and 8 the $`J/\psi `$ and $`\eta _c`$ mesons appear degenerate. This degeneracy is expected to be removed by a short range spin-spin force, whose effect, at $`T=0`$, is often treated perturbatively assuming a contact interaction Bra04 :
$$H_{ss}=\frac{8\pi }{9}\frac{\alpha _s(\mu )}{m_qm_{\overline{q}}}𝝈_q𝝈_{\overline{q}}\delta (𝒓).$$
(13)
Again, the short range nature of this force makes plausible the assumption that it is not affected by thermal effects and that it can be employed also at finite temperatures. Then, for the $`J/\psi `$ and $`\eta _c`$ energy shifts one gets:
$`\mathrm{\Delta }E_{J/\psi }`$ $`=`$ $`{\displaystyle \frac{8\pi }{9}}{\displaystyle \frac{\alpha _s(\mu )}{m_cm_{\overline{c}}}}|\psi (0,T)|^2`$ (14a)
$`\mathrm{\Delta }E_{\eta _c}`$ $`=`$ $`{\displaystyle \frac{24\pi }{9}}{\displaystyle \frac{\alpha _s(\mu )}{m_cm_{\overline{c}}}}|\psi (0,T)|^2,`$ (14b)
where for $`\alpha _s(\mu )`$ we have used the two-loop expression of Eq. (20), evaluated at $`\mu =m_c`$. As one can see, the temperature dependance of this term stems entirely from the value of the $`c\overline{c}`$ wave function at the origin.
We display in Fig. 8 the binding energies of the $`J/\psi `$ and $`\eta _c`$ mesons above the deconfinement temperature. The spin-spin contribution gives rise to a mild reduction of the dissociation temperature, without altering the qualitative features of the results. For instance, for $`N_f=2`$ at $`T=1.05T_c`$ one gets a $`J/\psi \eta _c`$ splitting of 145 MeV (the experimental value at $`T=0`$ is about 117 MeV), whereas at the $`J/\psi `$ melting temperature, $`T1.8T_c`$, one gets a splitting of 13 MeV (although, strictly speaking, in this case the perturbative expansion is no longer justified).
In Figs. 9 and 10 we plot the binding energies and the mean square radii of the bottomonium bound states above $`T_c`$. We have chosen $`m_b=4.3`$ GeV for the bottom quark mass, but again we have studied the sensitivity to variations of $`m_b`$ in the range of the PDG listings, 4.1 GeV$`<m_b<`$4.4 GeV. The arrows in Fig. 10 point to the results of $`T=0`$ potential model calculations Eic80 . Of course, for heavier quark masses more states than in the charmonium case survive above the deconfinement temperature.
Both for the charmonium and for the bottomonium states the binding energy gets smaller and the mean square radius gets larger as the temperature grows. This should reflect into a huge increase of the elastic cross section $`\sigma (Q\overline{Q}Q\overline{Q})`$ at low energy. In fact, for a small relative momentum $`k`$ of the two particles, when in the partial wave expansion the contribution of the $`s`$-wave is dominant, the presence of a $`l=0`$ state near zero binding energy leads to a dramatic increase of the cross section as $`k0`$, according to the formula:
$$\sigma _{l=0}\underset{k0}{}\frac{4\pi }{k^2+2\mu |ϵ|},$$
(15)
where $`\mu `$ is reduced mass of the system and $`ϵ`$ the energy of the state near zero binding, no matter whether it is positive (virtual level) or negative (bound state). If something analogous happened also for the heavy-light states (for which lattice data are not available yet) this would clearly accelerate the thermalization of the heavy quarks.
From Figs. 56 and 9 the effective potential obtained with $`N_f=3`$ appears more attractive. Indeed the critical temperature $`T_c`$ in the three cases $`N_f=0,2,3`$ assumes different values, hence a naive comparison might not be so meaningful. Furthermore the finite temperature $`N_f=3`$ lattice data available so far are affected by larger errors, so that our parameterization presents larger uncertainties with respect to the other two cases.
## IV Conclusions
The picture of the deconfined phase of QCD for values of the temperature slightly exceeding $`T_c`$ — that is the ones accessible in the heavy-ion experiments presently performed at RHIC — has substantially evolved in the last few years.
For large values of the temperature ($`T3T_c`$) a description of the QGP in terms of a gas of weakly interacting quasiparticles bla appears reliable (even in a regime where the coupling $`g`$ is not so weak) and supported by the lattice data for the QGP thermodynamics thermlat .
On the other hand, for temperatures up to $`T2T_c`$ (namely the ones currently accessible in the experiments), the matter resulting from the heavy-ion collisions is nowadays often described shury as a strongly interacting QGP (sQGP). In particular, a striking feature of the QGP matter obtained at RHIC is its hydrodynamical behaviour teaney ; kolb (characterized by a very low viscosity), which manifests itself in particular in the elliptic flow olli observed in non-central collisions: the plasma obtained at RHIC seems to behave as a nearly ideal fluid whose expansion is driven by pressure gradients.
In order to explain the very small mean free path of the plasma particles required by the hydrodynamical scenario, a picture of the matter obtained at RHIC in terms of a system of hundreds of loosely bound states of *quasiparticles* ($`q\overline{q}`$, $`qg`$, $`gg`$…) has been recently proposed shury . In such a framework one has to resort to some assumptions. Gluons and light quarks are treated as quasiparticles endowed with quite heavy thermal masses obtained from lattice calculations and the potential felt by them in the different color channels is got from the lattice $`Q\overline{Q}`$ free energy under the hypothesis of *Casimir scaling*.
The presence of such a pattern of bound states in the range of temperatures $`T_cT2T_c`$ has been recently questioned on the basis of the analysis of the correlation between baryon number and strangeness koch . The evaluation of the correlation coefficient $`C_{BS}=3BS/S^2`$ should allow to discriminate between a scenario in which the relevant degrees of freedom in the QGP are weakly interacting quark and gluon *quasiparticles* or loosely bound states of the latters. The lattice data on the off-diagonal quark-number susceptibilities available so far seem to favour the first hypothesis.
Here we followed a different approach, starting from the case of a heavy $`Q\overline{Q}`$ pair placed in a thermalized QGP and extracting their interaction from the available lattice calculations. Indeed, we have exploited the lattice data for the heavy quark free energies to get information on the existence of $`c\overline{c}`$ and $`b\overline{b}`$ bound states above the deconfinement transition. We have examined the cases $`N_f=0`$ Kac1 , $`N_f=2`$ Kac4 and also $`N_f=3`$ Petr , where lattice data are getting available. For the color singlet free energy we have adopted a parameterization which accounts for the effects of asymptotic freedom at short distances and displays an exponential screening at large distances. From the free energy we have then extracted the heavy quark potential, to be inserted into the Schrödinger equation. The latter has been solved numerically for the two interesting cases of charmonium and bottomonium.
For what concerns the charmonium, we have found a dissociation temperature $`T_d2T_c`$ for the 1S states ($`\eta _c`$ and $`J/\psi `$). Also the excited states 1P and 2S appear to melt, but at temperatures slightly exceeding $`T_c`$.
On the other hand, the bottomonium spectrum displays a much larger number of bound states above $`T_c`$. In particular its ground state turns out to remain bound in the whole range of temperatures covered by our parameterization. By extrapolating the latter at larger $`T`$ we get a dissociation temperature $`T_d4÷6T_c`$, depending upon the number of dynamical fermions. At fixed $`T/T_c`$, both for the charmonium and for the bottomonium, the system turns out to be more bound when the number of light flavors is increased.
Remarkably, in the range of temperatures of experimental interest covered in this paper a number of loosely bound $`Q\overline{Q}`$ states exists. Since the existence of states near zero binding energy entails a huge increase of the elastic cross sections at low relative momenta, this can help the approach to thermal equilibrium also of the heavy quarks.
An analogous situation occurs in the cross-over from the BCS theory of superconductivity to the Bose-Einstein condensation, much explored lately in ensembles of alkaline fermionic atoms confined in a magnetic trap. Indeed, by smoothly changing the external magnetic field $`B`$, one induces a smooth change of the energy of the two interacting atoms: if this energy is close to the one of a Feshbach resonance, a huge increase in the cross-section occurs Tim99 ; Leg01 ; Dui04 ; OHa02 . The striking analogy between $`B`$ and $`T`$ as control parameters is self-imposing.
In connection with the previous discussion on heavy quark thermalization, it would be highly desirable to have lattice data also for the $`Q\overline{q}`$ states, which are not yet available.
A phenomenological model has been recently proposed assuming the existence of *resonant* (not bound) D- and B- meson states above $`T_c`$ van . In this scheme the transverse momentum distributions of the charmed quarks (anti-quarks) would approach their thermal equilibrium value much faster, due to isotropic resonant scattering on light anti-quarks (quarks). Indeed, recent PHENIX results for azimuthally averaged transverse momentum spectra of single electrons arising, in Au-Au collisions, from the decay of D- and B-mesons seem to be compatible with a thermalization scenario flow .
Furthermore, preliminary results from PHENIX and STAR on the elliptic flow of the above electrons tend to support a picture in which, due to strong rescattering, charmed quarks reach thermal equilibrium and follow the flow of the fireball PHE1 ; PHE2 ; STAR1 ; greco . Hence, an extension of the present approach to the $`Q\overline{q}`$ ($`q\overline{Q}`$) states might shed light on the mechanism of thermalization of the heavy quarks.
Coming back to the present results, they offer the relevant possibility of evaluating the charmonium (bottomonium) multiplicity produced in heavy ion collisions at the experimental conditions of SPS, RHIC and LHC.
Indeed, a reliable estimate of the dissociation temperature of the different quarkonia is essential in order to predict how many of them, after being produced in the hard initial processes, survive in the QGP phase thus contributing to the final measured yields. The latter of course also contain, e. g., the charmonia that might be thermally produced during the hadronization process by $`c\overline{c}`$ recombination.
In particular, the knowledge of the in-medium quarkonium wave function and binding energy would allow a hopefully reliable estimate of the gluon-dissociation cross section of the $`J/\psi `$ (and hence, via a detailed balance analysis, also of the cross section for the inverse process). Thus a kinetic rate equation accounting for both the dissociation and recombination processes can be tackled: its solution, depending on how many $`c\overline{c}`$ pairs are produced in the initial state, should provide the number of $`J/\psi `$ present at the end of the QGP phase. These issues will be addressed in future work.
## V Acknowledgments
We are grateful to O. Kaczmarek and P. Petreczky for providing us their lattice data. We wish to thank also C.Y. Wong for sending us a revised version of his paper.
## Appendix A Running coupling and perturbative potentials at short distances
In this appendix we collect the formulas resulting from the perturbative calculations of Refs. pet ; mel ; sch of the heavy quark potential at short distances, where asymptotic freedom guarantees that the perturbative approach is justified.
The static QCD potential turns out to be given by:
$$V(r)=C_F\frac{\alpha _{\overline{MS}}(\mu )}{r}\left(1+v_1(r,\mu )\frac{\alpha _{\overline{MS}}(\mu )}{\pi }+v_2(r,\mu )\frac{\alpha _{\overline{MS}}^2(\mu )}{\pi ^2}+\mathrm{}\right),$$
(16)
where mel
$`v_1(r,\mu )`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left[{\displaystyle \frac{31}{9}}C_A{\displaystyle \frac{20}{9}}T_Fn_f+2\beta _0\mathrm{log}(\mu r^{})\right],`$ (17)
$`v_2(r,\mu )`$ $`=`$ $`{\displaystyle \frac{1}{16}}[({\displaystyle \frac{4343}{162}}+4\pi {\displaystyle \frac{\pi ^2}{4}}+{\displaystyle \frac{22}{3}}\zeta _3)C_A^2({\displaystyle \frac{1798}{81}}+{\displaystyle \frac{56}{3}}\zeta _3)C_AT_Fn_f`$ (18)
$`\left({\displaystyle \frac{55}{3}}16\zeta _3\right)C_FT_Fn_f+\left({\displaystyle \frac{20}{9}}T_Fn_f\right)^2+\beta _0^2\left(4\mathrm{log}^2(\mu r^{})+{\displaystyle \frac{\pi ^2}{3}}\right)`$
$`+2(\beta _1+2\beta _0({\displaystyle \frac{31}{9}}C_A{\displaystyle \frac{20}{9}}T_Fn_f))\mathrm{log}(\mu r^{})].`$
In the above $`r^{}r\mathrm{exp}(\gamma _E)`$ ($`\gamma _E`$ is the Euler-Mascheroni constant), $`C_A=3`$, $`C_F=4/3`$, $`T_F=1/2`$ and $`\zeta _3=\zeta _R(3)`$ (the Riemann function of argument 3); moreover, $`\alpha _{\overline{MS}}(\mu )`$ is the QCD running coupling in the $`\overline{MS}`$ renormalization scheme coming from the solution of the RGE:
$$\beta (\alpha _s(\mu ^2))=\frac{1}{\alpha _s(\mu ^2)}\frac{\alpha _s(\mu ^2)}{\mathrm{log}\mu ^2}\underset{n=0}{\overset{\mathrm{}}{}}\beta _n\left(\frac{\alpha _s(\mu ^2)}{4\pi }\right)^{n+1}.$$
(19)
The one-loop perturbative potential is by obtained keeping only the first two terms of Eq. (16) and employing for $`\alpha _{\overline{MS}}(\mu )`$ the result
$$\alpha _{\overline{MS}}(\mu )=\frac{4\pi }{\beta _0\mathrm{ln}\frac{\mu ^2}{\mathrm{\Lambda }_{\text{QCD}}^2}}\left(1\frac{\beta _1}{\beta _0^2}\frac{\mathrm{ln}\left(\mathrm{ln}\frac{\mu ^2}{\mathrm{\Lambda }_{\text{QCD}}^2}\right)}{\mathrm{ln}\frac{\mu ^2}{\mathrm{\Lambda }_{\text{QCD}}^2}}\right),$$
(20)
which arises from the solution of Eq. (19) with the two-loop beta function (i.e. considering only the first two terms of the series). Indeed the coefficients $`\beta _0`$ and $`\beta _1`$ do not depend on the renormalization scheme and are given by:
$$\beta _0=11\frac{2}{3}N_f,\beta _1=102\frac{38}{3}N_f.$$
(21)
On the other hand, for the two-loop perturbative potential one has to keep all the terms up to order $`\alpha _s^3`$ displayed in Eq. (16). The evaluation of the running coupling should include also the three-loop coefficient $`\beta _2`$ of the beta-function, which is no more renormalization scheme independent. In evaluating the two-loop curve in Fig. 3 we have used Eq. (13) of Ref. gock , where the $`\overline{MS}`$ scheme is employed.
Clearly, Eq. (16) still depends on the parameters $`\mu `$ and $`\mathrm{\Lambda }_{\text{QCD}}`$. The scale $`\mu r^1`$ is quite arbitrary: we chose $`\mu =[r\mathrm{exp}(\gamma _E)]^1`$ as usually done in the literature. For what concerns $`\mathrm{\Lambda }_{\text{QCD}}`$ we employed the value suggested by a recent lattice collaboration gock $`\mathrm{\Lambda }_{\text{QCD}}=261`$ MeV for all the three cases ($`N_f=0,2,3`$).
## Appendix B $`T`$-dependence of the fit parameters
In this appendix we show the functions employed in fitting the temperature dependence of the parameters $`a_0(T)`$, $`a_2(T)`$ and $`a_3(T)`$ entering into the functional form we have adopted for the $`Q\overline{Q}`$ free energy (see Sect. II).
As already mentioned in the text, since we have no phenomenological or theoretical hints to the functional form of the $`T`$-dependence of these parameters, we have tried to use the simplest expressions yielding a smooth fit and a “reasonable” $`\chi ^2`$ (say, of the order of a few units).
For $`a_0(T)`$ and $`a_2(T)`$ we have used the following forms:
$`a_0`$ $`=`$ $`{\displaystyle \frac{A_1^{(0)}x^{A_3^{(0)}}\mathrm{exp}[A_0^{(0)}x]}{A_2^{(0)}x^21}},`$ (22)
$`a_2`$ $`=`$ $`A_0^{(2)}{\displaystyle \frac{A_1^{(2)}+A_2^{(2)}x^2+x^4}{A_3^{(2)}+A_4^{(2)}x+x^3}},`$ (23)
where $`xT/T_c`$ and the $`A_i^{(0)}`$’s and $`A_i^{(2)}`$’s are fit parameters.
Only in the case of $`a_3(T)`$ we have used an expression inspired by perturbative QCD, namely the one for the Debye mass (although, rigorously this parameter does not represent the Debye mass):
$$a_3=\left(1+\frac{N_F}{6}\right)^{1/2}x\left(A_0^{(3)}+\frac{A_1^{(3)}}{x^2}+\frac{A_2^{(3)}}{x^4}+\frac{A_3^{(3)}}{x^6}+\frac{A_4^{(3)}}{x^8}\right)g_{2\text{loop}}(x),$$
(24)
where $`g_{2\text{loop}}^2/4\pi \alpha _{\overline{MS}}`$ is the 2-loop coupling constant obtained by replacing $`\mu /\mathrm{\Lambda }_{\text{QCD}}`$ with $`4.8826x(2\pi T_c/\mathrm{\Lambda }_{\text{QCD}})x`$, having used $`T_c=202`$ MeV and $`\mathrm{\Lambda }_{\text{QCD}}=261`$ MeV. In the case $`N_f=3`$, because of the larger errors affecting the parameterization of the free energy, we made the fit of $`a_3`$ stiffer by setting $`A_4^{(3)}=0`$. |
warning/0507/cond-mat0507348.html | ar5iv | text | # Large atom number Bose-Einstein Condensate machines
## I Introduction
It has been a decade since Bose Einstein condensation (BEC) in atomic vapors was first observed Cornell and Wieman (2002); Ketterle (2002). The transition from a classical thermal gas to the quantum degenerate Bose-Einstein condensate occurs when the occupation number of the lowest energy states (phase space density) $`\rho =n\lambda _{dB}^3`$ is increased to $`1`$, where $`n`$ is the number density and $`\lambda _{dB}`$ is the thermal de Broglie wavelength of the atoms. So in principle, getting a BEC is easy: you simply cool down the gas until the critical phase space density is reached. In practice, the procedure is more complicated. A variety of different techniques are needed to increase the phase space density in several stages (Table 1). Furthermore, each atom has different properties and requires modifications to the cooling techniques. Major work by many groups around the world has now extended the cooling techniques to an impressive number of atomic species: <sup>87</sup>Rb Anderson et al. (1995), <sup>23</sup>Na Davis et al. (1995), <sup>7</sup>Li Bradley et al. (1995, 1997), <sup>1</sup>H Fried et al. (1998), <sup>85</sup>Rb Cornish et al. (2000), <sup>4</sup>He\* Roberts et al. (2001); Santos et al. (2001), <sup>41</sup>K Modugno et al. (2001), <sup>133</sup>Cs Weber et al. (2003), <sup>174</sup>Yb Takasu et al. (2003), and <sup>52</sup>Cr Griesmaier et al. (2005). Still, <sup>23</sup>Na and <sup>87</sup>Rb are the two atoms which appear to have the most favorable properties for laser and evaporative cooling and are the two work horses in the field.
A major difference between the various experiments is the way in which atoms are laser cooled and then loaded into a magnetic trap (or now sometimes into an optical trap) for evaporative cooling. Our approach at MIT employs atomic ovens and Zeeman slowing. Other approaches use vapor cell magneto-optical traps, often in a double MOT configuration and more recently surface MOTs. An important figure of merit of a BEC setup is the number of atoms in the condensate. Large atom number allows better signal-to-noise ratios, greater tolerance against misalignments, and greater robustness in day-to-day operation. Since 1996, the MIT sodium BEC setups have featured the largest alkali condensates. Our three setups routinely produce condensates with atom numbers between 20 and 100 million. Since diode lasers for rubidium are less expensive than the dye lasers for sodium, most new groups have chosen to work in rubidium. The most popular laser cooling setups for rubidium involve vapor cell MOTs, which do not obtain the condensate size of the MIT sodium experiments. There has been a widespread perception in the field that Zeeman slowing is the technique of choice for sodium and vapor cell traps for rubidium. The construction of vapor cell MOT rubidium condensate machines is extensively detailed in the complementary work of Ref. Lewandowski et al. (2003).
When the Center of Ultracold Atoms was created at MIT and Harvard, one major funded project was the translation of the techniques we had developed for sodium to rubidium and to create a rubidium BEC experiment with enormous condensates. The successful accomplishment of these goal is described in this paper. Since we are in the unique situation to have two similar experimental setups for Rb and Na, we are able to discuss similarities and differences between the optimization for the two species. Our conclusion is that the technique of using an atomic beam and Zeeman slower works as well for rubidium as for sodium, and we present the technical details of how to build an intense slow beam for both atomic species. Other figures-of-merit besides atom number include simplicity and reliability. In our experience, the Zeeman slowing technique is by far the simplest technique to generate an intense slow beam since it requires only a single laser beam of modest power. The length of the slower and therefore the overall size of the vacuum setup may look intimidating, but once it is built, it provides simple and reliable day-to-day operation without need for realignment.
The third-generation sodium experiment and the rubidium experiment were both designed with an additional vacuum chamber (“science chamber”) into which the BEC or evaporatively cooled atoms can be moved using optical tweezers. The multi-chamber design allows us to rapidly reconfigure the experimental setup in the science chambers while keeping the BEC production chambers under vacuum. This has allowed us to perform very different experiments in rapid succession Gustavson et al. (2002); Leanhardt et al. (2002); Chikkatur et al. (2002); Leanhardt et al. (2003a, b); Shin et al. (2004a); Pasquini et al. (2004); Shin et al. (2004b, c); Schneble et al. (2003, 2004); Campbell et al. (2005).
## II System Overview
Fig. 1 illustrates the layout of our system. A thermal atomic beam emanates from the oven and is decelerated with the Zeeman slower. In the main chamber, the slowed atoms are captured and cooled with a six-beam magneto-optical trap (MOT) Raab et al. (1987). To load the magnetic trap, the atoms are optically pumped into the F=1 hyperfine ground state. Atoms in the F=1, m<sub>F</sub>=-1 state are held by their attraction to the field minimum of the Ioffe-Pritchard magnetic trap.
The trapped sample is evaporatively cooled by removing hotter atoms through radio frequency (RF) induced transitions into untrapped states. Reducing the RF frequency lowers the effective depth of the magnetic trap, allowing us to progressively cool to higher densities and lower temperatures until the atoms reach BEC. Magnetically trapped atoms in the F=2, m<sub>F</sub>=+2 state have also been evaporated down to BEC.
Ultracold atoms can be transported from the main chamber into an adjoining auxiliary “science chamber” by loading the atoms into the focus of an optical tweezer and then translating the focus. In this manner we have transported <sup>23</sup>Na BECs Gustavson et al. (2002). Vibrational issues during transport cited in Gustavson et al. (2002) were reduced by the use of Aerotech ABL2000 series air bearing translation stages.
Technical issues related to the greater mass and higher three body recombination rate in <sup>87</sup>Rb were overcome by transporting ultracold atoms just above the transition temperature $`T_c`$, and then evaporating to BEC at the destination. The oven and Zeeman slower are tilted by 57 from horizontal to allow a horizontal orientation for the weak trapping axes of both the optical tweezers and magnetic trap.
Trapping ultracold atoms requires that they be isolated from the surrounding environment. The laser and magnetic trapping techniques confine the atoms in the center of the chamber, out of contact with the room temperature chamber walls. The atoms are still exposed to thermal black body radiation from the chamber walls, but are transparent to most of the spectrum. The transitions which the black body radiation can couple to are the optical transitions used for laser cooling and the microwave hyperfine transitions. For optical transitions, which have energies much greater then $`k_BT`$ the excitation rate is $`\frac{3}{\tau _{optical}}\mathrm{exp}\left(\mathrm{}\omega _{optical}/k_BT\right)`$, where $`\omega _{optical}`$ is the frequency of the transition and $`\tau _{optical}`$ the lifetime of the excited state. For rubidium in a 25C chamber this gives a characteristic excitation lifetime of $``$56 billion years. Raising the chamber temperature to 680C increases the optical excitation rate into the experimentally relevant domain of once per minute. The hyperfine transitions are significantly lower in energy compared to $`k_BT`$ and have an excitation rate of $`\frac{3}{\tau _{hfs}}\frac{k_BT}{\mathrm{}\omega _{hfs}}`$, which is once per year at 25C in <sup>87</sup>Rb since the ground state hyperfine spontaneous decay lifetime $`\tau _{hfs}`$ in an alkali atom is thousands of years. Neither of these excitation rates are limitations on current experiments.
Collisions with background gas molecules result in loss from the trap, necessitating low vacuum pressure for long atom cloud lifetime. We can magnetically trap ultracold atoms with lifetimes of several minutes in the $`<10^{11}`$ torr ultrahigh vacuum (UHV) environment of the main production chamber. To achieve this vacuum performance we have followed the general guidelines set out in Ref. O’Hanlon (1989) for constructing vacuum systems. The main chamber body was constructed of nonmagnetic 304 stainless steel and then electropolished to reduce the surface roughness. The only component placed inside the chamber was the RF evaporation antenna coil (Fig. 2).
The cloverleaf-style Ioffe-Pritchard magnetic trap coils fit inside two re-entrant bucket windows <sup>1</sup><sup>1</sup>1Simon Hanks of UKAEA, D4/05 Culham Science Center, Abingdon,UK, allowing them to be outside the chamber with an inter coil spacing of 25 mm (Fig. 2). The Zeeman slower tube is mounted between the main chamber and the oven chamber. The Zeeman slower coils that are around the Zeeman slower tube are also outside of the vacuum system but cannot be removed without breaking vacuum.
After assembling the chamber, we pumped out the system and reached UHV conditions by heating the system to accelerate outgassing. We heated the main chamber to 230C and the Zeeman slower to 170C (limited by the coil epoxy). Using a residual gas analyzer to monitor the main chamber, we “baked” until the partial pressure of hydrogen was reduced to less than $`10^7`$ torr and was at least ten times greater than the partial pressure of other gases. A typical bakeout lasted between 3 and 9 days, with temperature changes limited to less than $`50^{}`$C/hour. While we acknowledge the merit of using dry pumps as recommended in Ref. Lewandowski et al. (2003), we have not had any detrimental experiences using oil sealed rotary vane roughing pumps to back our turbo pumps. The vacuum in the main chamber is preserved after bakeout with a 75 L/s ion pump and a titanium sublimation pump. Refer to Sec. 3.4 of Ref. Chikkatur (2002) for more details of our bakeout procedures.
## III Oven
We generate large fluxes of thermal atoms for Zeeman slowing from effusive atomic beam ovens. An effusive beam is created by atoms escaping through a small hole in a heated chamber Pauly and Scoles (1988). The higher vapor pressure of rubidium requires a more complicated design, but lower operating temperature (110-150C Rb, 260C Na.) At room temperature, the vapor pressure of sodium ($`2\times 10^{11}`$ torr Alcock et al. (1984)) is roughly compatible with our UHV main chamber environment, while that of rubidium ($`4\times 10^7`$ torr Alcock et al. (1984)) is not. This dictated that the design of the rubidium oven prevent contaminating the main chamber with rubidium. Because of its greater complexity, further discussion will focus on the rubidium oven (Fig. 3). We expect that the rubidium oven would work as well for sodium, but instead we used a simpler design described in Ref. Chikkatur (2002).
A combination of active pumping and passive geometrical techniques were used to reduce extraneous rubidium transfer to the main chamber. A cold cup (I) is used to reduce rubidium vapor in the oven chamber by almost completely surrounding the oven aperture (J) with a cold surface. After bakeout, the combination of cold cup and oven chamber ion pump has achieved pressures as low as $`10^9`$ torr, although we have successfully made BECs with pressures of up to $`10^6`$ torr in this region. A combination of a differential pumping tube, ion pump, and the Zeeman slower tube provides a pressure differential of over 3 orders of magnitude between the oven and main chamber. This is sufficient to isolate the UHV environment from an oven pressure dominated by rubidium vapor at room temperature. When the oven is opened to replace rubidium and clean the cold cup, the main chamber vacuum is isolated with a pneumatic gate valve. The second gate valve can be used in case of failure of the first. While not used in our system, designers may want to consider gate valves with an embedded window from VAT to allow optical access along Zeeman slower or tweezer beamlines during servicing.
The oven is loaded with a sealed glass ampoule containing 5 g of rubidium in an argon atmosphere. To add rubidium, the ampoule is cleaned, placed in the oven, and baked out under vacuum while still sealed. We then break the ampoule under vacuum and heat the oven to 110C to produce the atomic beam. During operation, the machine is run as a sealed system, without the turbo-mechanical pump, to prevent accidental loss of the main chamber vacuum. Oven temperatures from 150C down to 110C produced similar sized <sup>87</sup>Rb BECs. Reducing the oven temperature increased the time between rubidium changes to greater than $`1000`$ hrs of operating time. This long operating cycle precluded the need for more complex recycling oven designs Walkiewicz et al. (2000).
## IV Zeeman slower
The atomic beams are slowed from thermal velocities by nearly an order of magnitude by scattering photons from a resonant, counter-propagating laser beam. When a photon with momentum $`\mathrm{}k`$ ($`k=2\pi /\lambda `$) is absorbed or emitted by an atom with mass $`m`$, the atom will recoil with a velocity change of $`v_r=\mathrm{}k/m`$ to conserve momentum. Atoms can resonantly scatter photons up to a maximum rate of $`\mathrm{\Gamma }/2`$, where $`1/\mathrm{\Gamma }=\tau `$ is the excited-state lifetime. This results in a maximum acceleration $`a_{max}=\frac{\mathrm{}k\mathrm{\Gamma }}{2m}`$ ($`1.1\times 10^5\text{m/s}^2`$ Rb, $`9.3\times 10^5\text{m/s}^2`$ Na). As the atoms decelerate, the reduced Doppler shift is compensated by tuning of the Zeeman shift with a magnetic field Phillips and Metcalf (1982) to keep the optical transition on resonance. We designed our slower to decelerate the atoms at a reduced rate $`fa_{max}`$ where $`f50\%`$ is a safety factor to allow for magnetic field imperfections and finite slower laser intensity.
Our slowers are designed along the lines of Ref. Barrett et al. (1991), with an increasing magnetic field and $`\sigma ^{}`$ polarized light scattering off the F=2, m<sub>F</sub>=-2 $``$ F=3, m$`_F^{}`$=-3 cycling transition. Before the slowing begins, the atoms are optically pumped into the F=2,m<sub>F</sub>=-2 state. The large magnetic field at the end of the slower corresponds to a large detuning from the low velocity, low magnetic field resonance frequency. This large detuning allows the slowing light to pass through the MOT without distorting it due to radiation pressure. Within the slower, the quantization axis is well-defined by the longitudinal magnetic field and optical pumping out of the cycling transition is strongly suppressed by the combination of light polarization and Zeeman splitting.
We slow <sup>87</sup>Rb atoms from an initial velocity of $``$350 m/s with a tailored 271 G change in magnetic field (Fig. 4). An additional uniform $``$200 G bias field was applied along the length of the slower to ensure that neighboring hyperfine levels were not near resonance in either the slower or the MOT. The slower cycling transition light is detuned -687 MHz from the F=2 $``$ F=3 transition. The slowing laser intensity is $`I/I_{sat}8`$, giving a maximum theoretical deceleration of 89% of $`a_{max}`$. “Slower repumping” light copropagates with the cycling transition light and is detuned -420 MHz from the F=1 $``$ F=1 transition to match the Doppler shift of the unslowed thermal atoms from the oven. A flux of $`10^{11}`$ <sup>87</sup>Rb atoms/s with a peak velocity of 43 m/s was measured from our slower with an oven temperature of 150C. This is signifigantly greater flux then the $`8\times 10^8`$ Rb/sec vapor cell loading rate quoted by Lewandowski et al. (2003).
The higher temperature of the sodium oven, along with the atoms’ lower mass, results in a greater initial velocity of 950 m/s. This requires a slower with a much larger magnetic field change of 1150 Gauss. To reduce the maximum magnitude of the magnetic fields we use the spin flip variant of the increasing field design by shifting the zero crossing of the magnetic field from the beginning of the slower to the middle. The first segment then becomes a decreasing field slower, with current flowing in the opposite direction of the second, increasing field segment. In the low magnetic field region near the zero crossing, the atoms are theoretically vulnerable to optical pumping out of the cycling transition, but we did not find this to be a problem experimentally. Ref. Slowe et al. (2004) has demonstrated a high flux spin flip style slower for <sup>87</sup>Rb. The sodium slowing beam is detuned -1.0 GHz from the F=2 $``$ F=3 transition and has an intensity of $`I/I_{sat}4`$, giving a laser power limited maximum deceleration of 80% of $`a_{max}`$. Unlike the rubidium slower, light for optical pumping was generated by adding 1.75 GHz sidebands to the slowing light using an electro-optical modulator.
The sodium slower coils were broken up such that the first segment had an initial field of 440 G and a length of 52 cm and the second segment had a final field of 710 G and a length of 43 cm. The sodium slower was tested as depicted in Fig. 5, with a measured flux of $`3\times 10^{11}`$ <sup>23</sup>Na atoms/s with a peak velocity of 100 m/s.
## V Lasers
Resonant laser light is used to slow, cool, trap, and detect the atoms. All laser light is prepared on a separate optics table and delivered to the apparatus (Fig. 1) through single-mode optical fibers. Because stray resonant light can heat the atoms during evaporation, black cloth separates the two tables. All frequency shifting and attenuation of the light is done with acousto-optic modulators. Mechanical shutters are also placed in front of each fiber coupler to block any light which might leak through the modulators and disturb the atoms. Atomic energy levels and laser frequencies used are indicated in Fig. 6.
We use different techniques for generating laser light at the resonant wavelengths of <sup>87</sup>Rb (780 nm) and <sup>23</sup>Na (589 nm). Commercially available (Toptica DL100,TA100) external cavity diode lasers and semiconductor tapered amplifiers are used to create 350 mW and 35 mW of light resonant with the <sup>87</sup>Rb F=2$``$ F=3 and F=1$``$ F=1 transitions at 780 nm respectively. The lasers are stabilized with a polarization sensitive saturated absorption spectroscopy lock Yoshikawa et al. (2003); Pearman et al. (2002). This modulation-free technique optically creates a derivative signal of the absorption spectra that is locked with a proportional+integral gain servo loop. The locking signal fluctuation indicates a frequency jitter of $`<1`$MHz over several seconds, which is much less than the 6.1 MHz natural linewidth of <sup>87</sup>Rb.
The <sup>87</sup>Rb MOT uses a total of 60 mW of light near the F=2$``$F=3 cycling transition for trapping/cooling. The F=2$``$F=3 transition in the MOT is only approximately a closed cycle and atoms are often optically pumped into the F=1 ground state. To “repump” these atoms back into the F=2 state we use 10 mW of light on the F=1$``$ F=1 transition. Likewise, to deliberately transfer atoms from the F=2 to F=1 manifold we can introduce a few mW of “depumping” light resonant with the F=2$``$ F=2 transition. Powers are quoted after fiber coupling, measured as delivered to the apparatus table. After frequency shifting the slower cycling and repumping light to their desired detunings only a few mW of power are available. Each of these beams is then amplified to 35-40 mW by injection locking Siegman (1986) a free running Sanyo DL7140-201 laser diode. The two amplified beams are then overlapped and coupled into a fiber, which delivers 18 mW of slower cycling light and 6 mW of slower repumping light.
For <sup>23</sup>Na we use a Coherent 899 dye laser pumped by a Spectra Physics Millenia laser (532 nm, $`8.5\mathrm{W}`$). Typically 1.2 W of 589 nm light is generated by the dye laser. The laser frequency was referenced to an external saturation-absorption lock-in scheme and locked to a Fabry-Perot cavity. Stable operation was improved by using a precision dye nozzle (Radiant Dyes, Germany), high pressure dye circulator at 12 bars, and stabilized temperatures for the room and dye.
For more detailed information on the generation of the laser light for sodium MOTs, see Sec. $`3.4`$ of Ref. Stamper-Kurn (2000). Typical delivered laser powers are 80 mW for the MOT light, 20 mW for the repumping light, 40 mW for the slowing light and less than one mW for the imaging beam. Electro-optic modulators allow the addition of high frequency sidebands ($``$1.8 GHz) on the slowing and MOT light for repumping without the use of an additional laser beam. Recent advances in single frequency high power fiber and diode pumped solid state lasers <sup>2</sup><sup>2</sup>2IPG Photonics EAD and RLM series amplifiers. have made nonlinear techniques such as sum frequency generation Moosmüller and Vance (1997); Bienfang et al. (2003) and frequency doubling Thompson et al. (2003) interesting alternatives as resonant light sources.
## VI Magneto Optical Trap
The MOT Raab et al. (1987) is the workhorse of atomic physics for creating large samples of ultracold atoms. We use a six-beam MOT, which doubles as an optical molasses when the magnetic gradient field is off. Similar to Ref. Lewandowski et al. (2003) the <sup>87</sup>Rb apparatus uses a bright MOT. The <sup>87</sup>Rb MOT equilibrates to around $`4\times 10^{10}`$ atoms after $`2`$ s of loading, operating in a magnetic field gradient of 16.5 G/cm with cycling beams detuned -18 MHz from the F=2$``$ F=3 transition and a peak intensity 5.3 mW/cm<sup>2</sup>. To increase the efficiency of the transfer into the magnetic trap, we briefly compress the <sup>87</sup>Rb MOT and then switch off the magnetic field gradient to cool the atoms with optical molasses. The <sup>87</sup>Rb MOT is compressed by linearly ramping the gradient to 71 G/cm in 200 ms and simultaneously sweeping the detuning to -45 MHz in 400 ms. We use 5ms of “gray” molasses, where the repumper power is dropped by 95%, the optical trapping power is ramped down to 50%, and the detuning is swept from -18MHz to -26MHz. The molasses phase requires reduction of imbalances in intensity between beams and also to residual magnetic fields Shang et al. (1991). After the molasses phase, 0.5-1 ms of “depumping” light is applied to put all the <sup>87</sup>Rb atoms into the F=1 level before loading into the magnetic trap. Exact MOT and molasses parameters were found through empirical optimization, and all listed numbers should be considered as guides.
The <sup>23</sup>Na apparatus uses a dark-spot MOT Ketterle et al. (1993), with a detuning of -15MHz, peak beam intensity of 8.8 mW/cm<sup>2</sup> and a magnetic field gradient of 11G/cm. A 4 mm diameter opaque circle blocks light in the middle of a single repumper beam, creating a region at the center of the MOT where trapped atoms are optically pumped into the F=1 state. The <sup>23</sup>Na MOT equilibrates after a few seconds of loading. The effectiveness of the dark-spot in <sup>23</sup>Na has precluded the need for the compression and molasses phases as in <sup>87</sup>Rb. 99% of trapped atoms reside inside the dark spot, and the <sup>23</sup>Na atoms are sufficiently cold and dense to be directly loaded into the magnetic trap.
## VII Magnetic Trap
Atoms in weak magnetic field seeking states can be trapped in a magnetic field minimum. Our magnetic trap is a high current Ioffe-Pritchard (IP) trap with a cloverleaf style winding that can hold F=1, m$`{}_{F}{}^{}=1`$ or F=2, m$`{}_{F}{}^{}=+2`$ ground state atoms of <sup>87</sup>Rb and <sup>23</sup>Na with long lifetimes. An IP trap has an anisotropic, “cigar”-shaped, 3D harmonic trap for energies which are small compared to the trap minimum $`g_Fm_F\mu _BB_0`$ and a 2D linear/1D harmonic trap at higher energies (See Fig. 7 and Appendix C). This linear regime at higher energies (higher cloud temperatures) is more efficient for evaporatively cooling hot atoms Ketterle and van Druten (1996), while the finite bias field at the minimum prevents Majorana spin flip loss of colder atoms.
Fig. 8 shows an expanded view of the magnetic trap coils. The two sets of four cloverleaf coils create radial gradients $`B^{}`$ along $`\widehat{x}`$ and $`\widehat{y}`$, while the curvature coils produce a parabolic field curvature $`B^{\prime \prime }`$ in the $`\widehat{z}`$ direction. The curvature coils also produce a substantial bias field (Table 4, Appendix C) along $`\widehat{z}`$, which is balanced by a roughly homogeneous field from the antibias coils, resulting in a low residual bias field $`B_0`$ of $``$1 G at the center of the trap. The subtraction of the large magnetic fields from the curvature and antibias coils can make the residual bias field $`B_0`$ susceptible to jitter from current noise. To prevent this we drive current through both coils in series from the same power supply (Appendix C, Fig. 12), reducing the effect of current noise in the residual bias field $`B_0`$ by $``$30. When assembled the anti bias coils enclose the cloverleaf coils and the MOT coils surround the curvature coils.
When the magnetic trap is initially turned on, the strength and shape of the confinement are adjusted to match that of the laser cooled atoms to preserve phase space density. Additional current is applied to the curvature coils, increasing the residual bias field and decreasing the radial confinement to make a roughly spherical magnetic trap that match the spherical MOT. After loading atoms in the trap, the additional curvature coil current is reduced over one second to adiabatically change the trap geometry to the tightly confining cigar shape, favorable for evaporative cooling. Sec. 2.3.2 of Ref. Ketterle et al. (1999) has an extensive discussion of mode matching magnetic traps to MOTs. The adiabatic compression technique is reviewed in Ref. Ketterle and van Druten (1996).
## VIII Control and Imaging
Two computers run the apparatus; one controls the various parts of experiment and the other processes images from a camera which images the atoms. The control computer has custom built National Instruments (NI) LabWindows based software to drive analog (2 NI Model PCI 6713, 8 channels of 12 bit analog, 1MS/s update) and digital output (2 NI Model PCI-6533, 32 channels of binary TTL, 13.3 MS/s update) boards. The control computer also controls an Agilent 33250A 80MHz function generator through a GPIB interface, and triggers a Princeton Instruments NTE/CCD-1024-ED camera through a ST-133 controller to capture the absorption images.
BECs are typically imaged 10-40 ms after release from the trap. Ref. Castin and Dum (1996); Ketterle et al. (1999) provides the details of analyzing condensates after free expansion. Atoms are first optically pumped into the F=2 state in $`200\mu `$s and then an absorption image is taken using on resonance F=2 $``$ F=3 light. Detuning off resonance causes dispersion (lensing) as the light passes through the cloud of atoms and can distort the image. The intensity of the imaging probe is kept lower than the saturation intensity to prevent bleaching of the transition, which would lead to errors in atom number counting. Typical exposure times are 50-200 $`\mu `$s. Sec. 3 of Ref. Ketterle et al. (1999) discusses other imaging techniques that can also be used to probe BECs.
## IX Evaporation
Evaporative cooling works by selectively removing hot atoms from the trapped cloud, while the remaining atoms rethermalize to a lower temperature. The efficiency of cooling depends on $`\eta `$, the ratio of trap depth or energy of the escaping atoms to the temperature $`k_BT`$, and is reduced by the rate of heating. The speed of this process depends on how quickly the atoms rethermalize. In a magnetic trap evaporation is implemented through RF induced transitions between trapped and untrapped states. A given RF frequency corresponds to a shell of constant $`\mu _m\left|𝐁\right|`$ where the transitions occur. Atoms that pass through this shell enter untrapped states and are lost; thus RF provides a flexible mechanism to control the magnetic trap depth. Our RF antenna consists of two rectangular loops of wire, 10 cm x 2 cm, positioned 3 cm above and below the condensate as depicted in Fig. 2.
To evaporate thermal atoms to a BEC, we sweep the RF frequency over several seconds using an Agilent 33250A synthesizer amplified with a 5 W RF amplifier (Mini-Circuits ZHL-5W-1). Typical evaporation curves for <sup>87</sup>Rb would ramp from 60 MHz down to $``$0.8 MHz in 15 to 40 seconds. Forced RF evaporative cooling is very efficient, increasing phase space density by $`>10^6`$ (Table 1). Fig. 9 shows the drop in temperature as the trap depth (calculated from the RF frequency) is lowered during evaporation of <sup>87</sup>Rb. Evaporation curves are frequently adjusted in the interest of tuning evaporation speed, atom number, density, and/or reproducibility. For instance, the atom number can be increased by decompressing the magnetic trap near the end of the evaporation. This reduces the effects of three body recombination heating by lowering the final condensate density. Such decompression techniques have allowed us to create nearly pure condensates with $`N_c20\times 10^6`$ in both <sup>87</sup>Rb and <sup>23</sup>Na with lifetimes in excess of 5 seconds.
Decompressing the trap shifts its center due to gravitational sag and imperfections in the balance of magnetic fields between the coils. Such movements can excite oscillations in the cloud, which results in the condensation of BECs which are not rest. Even in the absence of excitations, the magnetic field gradients must exert a force on the atoms which is greater then gravity for them to remain trapped. This limits the extent to which magnetic traps can be decompressed. Specially designed gravito-magnetic traps have been decompressed down to 1 Hz Leanhardt et al. (2003a) to investigate very cold, dilute BECs.
## X Deep Trap Limitations
A major difference we have observed between <sup>87</sup>Rb and <sup>23</sup>Na condensates is the unexpectedly high decay rate of <sup>87</sup>Rb in tightly confining deep traps, such as those used for transport in an optical trap Gustavson et al. (2002). At the typical densities of condensates, the lifetime and heating are usually determined by three-body recombination decay. However, the factor of four difference in the three-body rate coefficients (Table 2) was insufficient to explain this major discrepancy in behavior.
We investigated this issue in a magnetic trap instead of an optical trap. While it is easier experimentally to create tight trapping and hence high densities in an optical trap, both the trap frequencies and trap depth are functions of the optical power. This makes it difficult to separate density dependent effects, which are strongly affected by the trap frequency, from trap depth effects in an optical trap. In contrast, in a magnetic trap the trap depth can be controlled independently of the trap frequencies by adjusting the RF frequency which flips atoms to untrapped states.
There are two possible processes, both involving secondary collisions, which can greatly enhance the heating and losses due to the primary three-body collisions.
The first process is collisional avalanches, similar to a chain reaction, where the energetic products of three-body recombination collide with further atoms while they leave the condensate. This process would depend on the collisional opacity $`n\sigma l`$ (where $`\sigma =8\pi a^2`$ is the atom atom scattering cross section) and would increase dramatically when it exceeds the critical opacity of 0.693 Schuster et al. (2001). This process is independent of trap depth.
The second possible process can already happen at lower collisional opacities and relies on the retention of primary or secondary collision products by the trap in the so-called Oort cloud Ketterle et al. (1999); Burt et al. (1997). Subsequently, when those atoms slosh back into the trapped condensate, they can cause heating and trap loss. The retention of collision products in the Oort cloud should depend on whether the trap depth is larger or smaller than their energies.
Fig. 10 shows the initial loss rates measured for a large and a small BEC as a function of the magnetic trap depth. At low trap depths (5 $`\mu `$K) both the large and small condensate decay rates are in agreement with established three body recombination rates Tolra et al. (2003). Therefore, the avalanche effect does not significantly contribute to the observed decay rate, although the calculated opacity for the larger condensate was 0.88 and may not be far away from the onset of avalanches. Evidence for avalanches was obtained at an opacity of 1.4 in Schuster et al. (2001).
In the larger condensate at higher trap depth, the decay rate strongly increases, supporting the second process involving the Oort cloud as the likely mechanism. For times longer then 500 ms, the large <sup>87</sup>Rb condensate in the deeper traps was heated away and only a few thermal atoms remained. In contrast, at low trap depths the large condensate was still fully condensed after 20 seconds, and had an atom number in agreement with the expected losses from three body decay.
We speculate that this enhancement of the three-body losses was not observed in <sup>23</sup>Na for the following reasons. Three body recombination results in a diatomic molecule and an atom which fly apart with a total kinetic energy (2/3 to the atom, 1/3 to the molecule to conserve momentum) equal to the binding energy of the diatomic molecule in the highest vibrational state. This binding energy can be estimated from the scattering length as $`E_0\mathrm{}^2/ma^2`$ Borca et al. (2003) ($`200\mu `$K in <sup>87</sup>Rb, $``$ 2.7 mK in <sup>23</sup>Na). The direct decay products will only be retained if the trap depth is greater then their kinetic energies (min $`70\mu `$K for <sup>87</sup>Rb, $`900\mu `$K for <sup>23</sup>Na).
In addition, <sup>23</sup>Na decay products are less likely to undergo secondary collision processes from either primary three body products or hot Oort cloud atoms due to an elastic scattering cross section $`\sigma `$ which is 3.6 times smaller than for <sup>87</sup>Rb. The individual products of secondary collisions can have a spectrum of energies, further lowering the energy threshold for their retention by the Oort cloud. The combination of these three factors, three body rate, scattering cross section, and binding energy result in an estimated increase in the loss rate for <sup>87</sup>Rb condensates over <sup>23</sup>Na by a factor of $``$200.
The optical trap depths needed for transporting condensates in our system are significant fraction of the primary <sup>87</sup>Rb decay product energy, but a small fraction of that for <sup>23</sup>Na. Therefore <sup>23</sup>Na condensates can be easily transported using optical tweezers. For <sup>87</sup>Rb the preferred method is to transport a cloud at temperatures just above condensation, where the density is lower, and evaporate to BEC after transport.
## XI Discussion
<sup>87</sup>Rb and <sup>23</sup>Na are the two most popular species for BEC research. We have constructed two machines with nearly identical designs and can discuss differences in performance and operation. Key properties of the two atoms are highlighted in Table 2. The four principal differences are in vapor pressure, resonant wavelength, recoil velocity, and collisional properties.
The high vapor pressure of rubidium allows the operation of the oven at lower temperatures, but requires a more elaborate design of cold plates to avoid deposition of rubidium on surfaces of the UHV chamber which are at room temperature. An optimized Zeeman slower for rubidium will be about twice as long than that for sodium at similar oven temperatures, the stopping length L for the most probable velocity in a beam of temperature T being $`L=\frac{3k_BT}{\mathrm{}k\mathrm{\Gamma }}`$, assuming the maximum spontaneous light force. In our systems the gain from the greater light force is balanced out by the higher operating temperatures required of the sodium oven to produce comparable flux, resulting in both the rubidium and sodium slower being about 1 m in length.
Due to the higher recoil velocity, the slow sodium beam has a larger divergence than the rubidium beam. By keeping the distance between the end of the slower and the MOT to a minimum, we estimate quantitative transfer of atoms from the slower to the MOT. Our setup for rubidium was almost identical, but we expect that the requirement of keeping the slower and the MOT so close could be more relaxed for rubidium. Although we have not tried it, we expect that our rubidium experiment would work for sodium without changes to the oven, vacuum or magnet designs.
On the laser side, a major difference is the availability of low cost high power laser diodes in the near infrared region around 780 nm. In our experience a well run dye laser system provides similar or even superior performance to a diode laser system with several master and slave lasers. Our dye lasers at 589 nm tend to be more stable than semiconductor lasers during a long run of the experiment and need less day-to-day tweaking. However, occasionally they require major maintenance in terms of dye changes or full optical realignment. Another advantage is the visibility of the laser light and the atomic fluorescence. The near infrared 780 nm light is only modestly visible, whereas the sodium line at 589 nm is near the peak of human eye sensitivity and allows fine alignments of the laser beams and the magneto-optical trap without cameras, IR cards, or IR viewers.
<sup>87</sup>Rb has favorable properties for laser cooling and atom interferometry because of its greater mass, lower recoil velocity, and larger excited state hyperfine structure. While greater mass and longer resonant wavelength give <sup>133</sup>Cs an even lower recoil velocity, its complicated collisional behavior at low magnetic fields makes it difficult to cool to BEC. The lowest molasses temperature in rubidium is a factor of ten lower than for sodium. However, in BEC experiments the laser cooling is optimized to for large atom numbers and high initial elastic collision rates in the magnetic trap, and not for the lowest temperature. For laser cooling sodium at high atom numbers, the Dark SPOT technique Ketterle et al. (1993) is crucial to avoid rescattering of light, whereas it is only used in some rubidium experiments. One possible reason is that the larger excited state hyperfine structure allows for larger detunings from the cycling transition without exciting other hyperfine states. At the end of the day, although with somewhat different techniques, the laser cooling part works equally well for both atoms.
Both atomic species have favorable collisional properties for evaporative cooling. The elastic scattering cross section of <sup>87</sup>Rb atoms at low temperature is four times higher than in <sup>23</sup>Na. However, elastic collision rates after laser cooling are comparable since <sup>23</sup>Na atoms are faster. A peculiarity of <sup>87</sup>Rb is that the two ground electronic state hyperfine levels have similar scattering lengths, which can be advantageous for studies on spinor condensates and atomic clock transitions. Related to that, spin relaxation between the two hyperfine levels is almost completely suppressed. Mixtures of F=1 and F=2 atoms can be kept for seconds Harber et al. (2002), whereas in <sup>23</sup>Na they decay on ms time scales Görlitz et al. (2003). Both atoms have several Feshbach resonances below 1100 G Inouye et al. (1998); Stenger et al. (1999); Marte et al. (2002). Here, <sup>87</sup>Rb has the disadvantage, that the widest known resonance is only 200 mG wide compared to 1 G for <sup>23</sup>Na and requires more stable magnetic fields. Another difference is the higher rate of three-body collisions for <sup>87</sup>Rb atoms. As we discussed in Sec. X, this imposes limitations on trapping and manipulating dense <sup>87</sup>Rb condensates.
## XII Conclusions
In this paper, we have presented details for designing BEC machines with high performance and flexibility, and we hope that this description is useful for designing new experiments. Given the recent developments in the field, there is more than enough room for new experiments to join in the exploration of atom optics and many-body physics with quantum-degenerate atomic gases.
Funding for the <sup>87</sup>Rb machine was provided by the NSF MIT-Harvard Center for Ultracold Atoms. Funding for the <sup>23</sup>Na machine came from NSF, the ARO MURI program, NASA, and the ONR.
We thank S. Gupta, A. Görlitz, and A. E. Leanhardt for their contributions in the construction of the <sup>23</sup>Na machine; J. C. Mun and P. Medley for ongoing contributions to the <sup>87</sup>Rb machine; and MIT UROP students P. Gorelik and X. Sun for various contributions to the <sup>87</sup>Rb machine. The authors would also like to thank M. Saba and D. Kielpinksi for critical reading of this manuscript.
## Appendix A Oven
To sustain a high flux atomic beam, the background vacuum pressure must be low enough that the mean free path between collisions is much greater than the length of the beam. To generate an effusive beam with a thermal distribution of velocities, the size of the hole through which the atoms escape must be smaller than the mean free path inside the oven. We observed in sodium that at higher pressures (e.g. temperatures) the flux of slowable atoms does not increase and that the velocity distribution narrows. This phenomena is well understood Pauly and Scoles (1988), and limits the diffusive flux from a single aperture oven.
During servicing, a clean ampoule is essential for rapid recovery of good vacuum pressure. The ampoule is cleaned by submerging it in a 50/50 mixture by volume of acetone and isopropanol for 20 minutes, air drying it, and then by placing it in the oven while still sealed. This removes most of the water from the glass surface, which would otherwise require more time to pump away. After installation the ampoule is baked for 24 hours under vacuum at 150-180C to remove the remaining contaminates before it is broken.
To prevent accumulation of metal at the aperture (Fig. 3, J), the oven nozzle temperature (Fig. 3, K) is kept hotter ($`10^{}`$C in rubidium and $`90^{}`$C in sodium) than the rest of the oven . The velocity distribution of the beam is determined by the nozzle temperature (Fig. 3, K). On the other hand, the vapor pressure in the oven, which controls the beam flux, is dominated by the coldest spot in the elbow and bellows. The factor of two discrepancy between the observed and calculated (Table 3) rubidium oven lifetimes at 110C can be accounted for by a spot $`10^{}`$C colder then the lowest measured oven temperature. The specifics of this cold spot depend on how the oven is insulated.
## Appendix B Zeeman Slower
Every photon which scatters off an atom to slow the atom is radiated in a random direction, increasing the atoms spread in transverse velocity. The beam emerging from the tube needs to have sufficient forward mean velocity to load the MOT efficiently. Because of the random direction of the emission recoil, N photon scatterings increase the transverse velocity by $`v_r\sqrt{N/3}`$, or $`\sqrt{v_r\mathrm{\Delta }v/3}`$. The <sup>23</sup>Na slower operates with a recoil induced transverse exit velocity of $`3`$m/s, with a final forward velocity of 30 m/s so that the spatial transverse spread in the slowed beam matches the MOT capture area. The smaller initial and recoil velocities in the <sup>87</sup>Rb slower reduce the transverse velocity to $`0.8`$m/s, making MOT capture matching less critical.
An additional concern in both slowers is the fate of atoms not captured by the MOT. In <sup>87</sup>Rb we were concerned with the potential adverse impact a deposited film may have on the vapor pressure, and installed a cold plate near the slower window on the main chamber to capture desorbed Rb. Vacuum pressure has not been an issue and we have never needed to chill this cold plate. The opposite problem arises in <sup>23</sup>Na, where metal deposition on the slower window reduces the transmission of slower light. We have found heating the slowing beam vacuum port window to 90C prevents long term buildup.
### B.1 Slower Construction
The vacuum portion of the <sup>87</sup>Rb slower is a 99 cm long nonmagnetic 304 stainless steel tube with a 19 mm OD and 0.9 mm wall. The rear end of the tube is connected to the main chamber by a DN 16 CF rotatable flange, while the oven end of the tube has a narrow, 50mm long flexible welded bellows ending in another DN 16 CF rotatable flange. The retaining ring on this flange was cut in half for removal, so that the premounted coil assembly could be slid over the vacuum tube.
As shown in Fig. 1 the slower tube enters the main chamber at an angle of $`33^{}`$ from the vertical to accommodate access for optical tweezers. The oven and the Zeeman slower are supported two meters above the the experimental table in order to preserve the best optical and mechanical access to the main chamber. Aluminum extrusion from 80/20 Inc. was used to create the support framework.
Our <sup>87</sup>Rb slower was fabricated with a single layer bias solenoid and three increasing field segments (Fig. 1 and 11). The optimum configuration of currents and solenoid winding shapes was found by computer simulated winding of the solenoids one loop at a time, starting at the high field end and tapering the last few loops to best match the desired field profile. An alternative fabrication technique would be to apply a large uniform bias field and subtract away unwanted field with counter current coils. Stray fields from the Zeeman slower can have a detrimental effect on the MOT, particularly during sudden turnoff. An additional bias coil around the main chamber along the axis of the slower can compensate for this effect.
## Appendix C Magnetic Trap Appendix
### C.1 Ioffe Pritchard Trapping Potential
The field near the minimum of a Ioffe-Pritchard trap is approximately
$$𝐁=B_0\left(\begin{array}{c}0\\ 0\\ 1\end{array}\right)+B^{}\left(\begin{array}{c}x\\ y\\ 0\end{array}\right)+\frac{B^{\prime \prime }}{2}\left(\begin{array}{c}xz\\ yz\\ z^2\frac{1}{2}\left(x^2+y^2\right)\end{array}\right)$$
(1)
which realizes trap frequencies
$`m\omega _{}^2=`$ $`\mu _m{\displaystyle \frac{B^2}{B_0}}`$ (2)
$`m\omega _z^2=`$ $`\mu _mB^{\prime \prime }`$ (3)
Typical trap parameters of $`B^{}`$=223 G/cm, $`B^{\prime \prime }`$=100 G/cm<sup>2</sup>, $`B_0`$=1G (Fig. 7, Table 4) have frequencies of $`(\omega _{},\omega _z)/2\pi `$ of $`(200,9)`$ Hz for <sup>87</sup>Rb and $`(390,18)`$ Hz for <sup>23</sup>Na. Further details of Ioffe-Pritchard magnetic traps are discussed in Sec 2.3.2 of Ketterle et al. (1999) and Ch. 5. of Durfee (1999).
### C.2 Circuitry
Fig. 12 is representative of the magnetic trap circuit. We drive the magnetic trap coils with Lambda EMI DC power supplies in fixed current mode. Current to the cloverleaf coils is supplied from a Model ESS 30-500 15kW power supply, while the axial currents are driven with two Model EMS 20-250 5kW power supplies. Each power supply is protected against damage from reverse current with an International Rectifier SD600N04PC high-current diode.
To switch the high currents we use PowerEx models CM1000HA-24H and CM600HA-24H Integrated Gate Bipolar Transistors (IGBTs) controlled with PowerEx BG1A-F IGBT driver kits. The IGBTs and high current diodes dissipate several hundred W during operation, and are cooled with chilled water. Efficient heat sinking is critical for reliable operation, as thermal dissipation limits the maximum DC current. Fast turnoff of current on an inductive load, such as a coil, results in a large voltage spike. We have added a “debounce” circuit to each of the coil systems (Fig. 12) to control this process and prevent damage. The circuit consists of two different elements: a varistor (V) and a diode (D) in series with a low impedance resistor (R). The varistor shorts the circuit at high voltages to prevent this spike, and the diode (D) in series with a 1$`\mathrm{\Omega }`$ resistor (R) dissipates the remaining current after varistor shutoff. Ref. Melton and Pollak (1996) contains a through analysis of the behavior of a similar circuit. All control signals are electrically isolated from the high-current circuits to prevent voltage spikes from damaging connected hardware. Rapid, controlled magnetic field shutoff is important for quantitative interpretation of images taken after ballistic expansion.
### C.3 Wire
Both the slower and the magnetic trap coils were fabricated using square hollow core (0.125 in./side, 0.032 in. wall) Alloy 101 soft temper copper tubing from Small Tube Products, Inc. of Altoona, PA, wrapped with double Dacron glass fuse insulation by Essex Group Inc., Magnet Wire & Insulation of Charlotte, NC. The coils are held together with Hysol Epoxi-Patch 1C White high temperature epoxy that is bakable to 170C. Chilled water is forced through the hollow core of the copper wires to dissipate the $``$10kW of power generated from resistive heating in the magnetic trap and Zeeman slower coils. 200 psi of differential pressure is required for sufficient coolant flow. We designed all our coils to increase the cooling water temperature by less then 50C. Ch. 3 of Ref. Montgomery (1969) has an extensive discussion of water cooling in continuously powered resistive magnets. For our wire the following empirical values were measured,
$`\rho \left[\mathrm{\Omega }/m\right]`$ $`=`$ $`2.65\times 10^3`$ (4)
$`Q\left[ml/sec\right]`$ $`=`$ $`2.07\sqrt{{\displaystyle \frac{\mathrm{\Delta }P\left[psi\right]}{L\left[m\right]}}}`$ (5)
$`\mathrm{\Delta }T\left[{}_{}{}^{}C\right]`$ $`=`$ $`259I^2\left[Amps\right]\rho \sqrt{{\displaystyle \frac{L^3\left[m\right]}{\mathrm{\Delta }P\left[psi\right]}}}`$ (6)
where Q is the water flow rate in ml/sec, I<sup>2</sup> $`\rho `$ L is the power dissipated by the coil, $`\mathrm{\Delta }`$P the pressure drop in psi (1 psi=6.89 kPa), and L the length of the coil in meters.
### C.4 Fabrication
All of the components for each half of the magnetic trap were epoxied together for stability. Each assembly was then mounted in the bucket windows with an aluminum mounting plate backed by four threaded Alloy 316 stainless steel rods. No ferromagnetic materials were used in the mounting because of concern for irreproducibility from hysteresis effects. Table 4 lists the windings and typical parameters for each coil. |
warning/0507/math0507603.html | ar5iv | text | # Rendezvous numbers in normed spaces
## 1. Introduction
It was shown by O. Gross that for a compact, connected metric space $`(X,d)`$ there exists a unique number $`r:=r(X)`$ such that for each finite point system $`P=\{x_1,x_2,\mathrm{},x_n\}X`$, $`n`$ one finds a point $`xX`$ with the average distance to $`P`$ being exactly $`r`$, that is
(1)
$$\frac{1}{n}\underset{i=1}{\overset{n}{}}d(x,x_i)=r.$$
This number $`r(X)`$ is called the *rendezvous number* of the space $`X`$. Using this strict definition, that is requiring the very existence of a point $`x`$ with exact equality under (1), it is well-known that both compactness and connectedness are crucial assumptions. However, one can relax on the requirements considering so-called *weak rendezvous numbers*, meaning that there exist two points $`y,zX`$ with their average distances to the points $`x_j`$ being less or equal than $`r`$ and greater or equal than $`r`$, respectively, . Clearly, for connected spaces this is equivalent to the existence of a strong rendezvous number. Hence it is not surprising that one can prove the existence and uniqueness of such weak rendezvous numbers under the hypothesis of compactness, see e.g., . However, dropping the compactness condition one can not expect uniqueness as, for example, the case of $`C(K)`$ spaces shows . Furthermore, sticking to connected spaces but relaxing on compactness is also insufficient to prove the existence of rendezvous numbers. For example, the unit sphere of $`\mathrm{}_p`$ spaces have no rendezvous number (unless $`p=2,+\mathrm{}`$) , , , .
In we employed a systematic potential theoretic approach to rendezvous numbers and introduced a modified definition of these numbers, considering also closure of the occurring average distance sets in the construction. In the classical case of compact sets and continuous kernels (e.g., distances on compact spaces), closure is superfluous as a continuous image of a compact set is also compact, hence closed. However, in the more general case of non-compact sets, like unit spheres of infinite dimensional Banach spaces, and also in case of more general, only lower semicontinuous kernels, this approach provides its yield. In particular, with the new definition we have found very general existence results, far beyond the setting of metric spaces. It turned out that abstract potential theory on locally compact spaces with a lower semicontinuous kernel is an appropriate framework for such investigations. In it was also indicated that the local compactness assumption on the space $`X`$ is not necessary and the results go through to metric, but not necessarily locally compact spaces as well.
We analysed further consequences of this approach in the context of metric spaces in , extending and explaining a good deal of previous knowledge. In the present paper we continue the study of rendezvous- and average numbers in normed spaces.
Let us fix some notation and introduce the necessary notions. In the abstract potential theory developed by Fuglede and Ohtsuka the usual assumptions are the following. $`X`$ is a locally compact, Hausdorff space and $`k:X\times X_+\{+\mathrm{}\}`$ is a lower semicontinuous, symmetric, positive *kernel* function. Nevertheless, we will consider possibly infinite dimensional Banach spaces, so the local compactness assumption needs to be relaxed. We will accomplish this task on the cost of allowing special kernel functions only, such as $`k(x,y):=xy`$, which is just the usual kernel appearing in connection with rendezvous numbers.
###### Definition 1.
For arbitrary $`H,LX`$ the *general $`n^{\text{th}}`$ Chebyshev constant* of $`L`$ with respect to $`H`$ is defined as
$$M_n(H,L):=\underset{w_1,\mathrm{},w_nH}{sup}\underset{xL}{inf}\frac{1}{n}\left(\underset{k=1}{\overset{n}{}}k(x,w_k)\right).$$
and the *$`n^{\text{th}}`$ general dual Chebyshev constant* of $`L`$ relative to $`H`$ is
$$\overline{M}_n(H,L):=\underset{w_1,\mathrm{},w_nH}{inf}\underset{xL}{sup}\frac{1}{n}\left(\underset{j=1}{\overset{n}{}}k(x,w_j)\right).$$
The first part of the definition is due to Ohtsuka . By standard considerations, just as in the case of classical Chebyshev constants, one sees that $`M_n(H,L)`$ and $`\overline{M}_n(H,L)`$ converge to some $`M(H,L)`$, $`\overline{M}(H,L)[0,+\mathrm{}]`$ (see, e.g., , or ). Furthermore
$$\underset{n}{sup}M_n(H,L)=\underset{n\mathrm{}}{lim}M_n(H,L)\text{and}\underset{n}{inf}\overline{M}_n(H,L)=\underset{n\mathrm{}}{lim}\overline{M}_n(H,L).$$
The limits $`M(H,L)`$, $`\overline{M}(H,L)`$ above are called the *Chebyshev constant* and the *dual Chebyshev constant* of $`L`$ with respect to $`H`$.
If $`X`$ is a Hausdorff topological space, let us denote by $`𝔐(X)`$ the set of positive, regular Borel measures on $`X`$ and by $`𝔐_1(X)`$ the subset of probability measures. The notation $`𝔐_1^\mathrm{\#}(X)`$ is used for probability measures with finite support. Given a set $`HX`$, for the subfamily of measures concentrated on $`H`$ (or supported on $`H`$, in case $`H`$ is closed, cf. \[11, pp. 144–146\]) we use the analogous notations $`𝔐(H)`$, $`𝔐_1(H)`$ and $`𝔐_1^\mathrm{\#}(H)`$, respectively. The *potential* of a measure $`\mu 𝔐(X)`$ is
$$U^\mu (x):=_Xk(x,y)d\mu (y).$$
In the classical potential theoretic literature various notions of *energies* appear. Already Fuglede and Ohtsuka introduced the following two-variate versions of energies (see also ).
###### Definition 2.
Let $`H,LX`$ be fixed, and $`\mu 𝔐_1(X)`$ be arbitrary. First put
(2)
$$Q(\mu ,H):=\underset{xH}{sup}U^\mu (x),\text{and also}\underset{¯}{Q}(\mu ,H):=\underset{xH}{inf}U^\mu (x).$$
Then the *quasi-uniform energy* and *dual quasi-uniform energy* of $`L`$ with respect to $`H`$ are
(3)
$$q(H,L):=\underset{\mu 𝔐_1(H)}{inf}Q(\mu ,L)\text{and}\underset{¯}{q}(H,L):=\underset{\nu 𝔐_1(H)}{sup}\underset{¯}{Q}(\nu ,L).$$
We use the notation $`M(H):=M(H,H)`$, $`\overline{M}(H):=\overline{M}(H,H)`$, $`\underset{¯}{q}(H):=\underset{¯}{q}(H,H)`$ and $`q(H):=q(H,H)`$ for the diagonal (classical) cases of the quantities given in Definitions 1 and 2.
###### Remark 3.
It is not surprising that in general the quantities $`M(H),q(H)`$ etc. do not posses any monotonicity properties as functions of the set $`H`$. The worst consequence of this is the lack of good “inner regularity” properties, for example $`q(H)=inf_{KH}q(K)`$ fails to hold (we use the abbreviation $`KH`$ to express that $`K`$ is a compact subset of $`H`$). However, fixing one variable the functions $`M(H,L)`$, $`\overline{M}(L,H)`$, $`\underset{¯}{q}(H,L)`$ and $`q(L,H)`$ are increasing with respect to $`H`$ and decreasing with respect to $`L`$, and the above mentioned problem disappears. This particularly explains the relevance and importance of the above two-variable definitions to our subject, see also .
###### Definition 4.
For arbitrary subsets $`H,LX`$ the $`n^{\text{th}}`$ *(extended) rendezvous set* of $`L`$ with respect to $`H`$ is
(4) $`R_n(H,L)`$ $`:={\displaystyle \underset{w_1,\mathrm{},w_nH}{}}\overline{\mathrm{conv}}\{p_n(x):={\displaystyle \frac{1}{n}}{\displaystyle \underset{j=1}{\overset{n}{}}}k(x,w_j):xL\},`$
$`R_n(H)`$ $`:=R_n(H,H).`$
Correspondingly, one defines
(5) $`R(H,L)`$ $`:={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}R_n(H,L),`$ $`R(H)`$ $`:=R(H,H).`$
Similarly, one defines the *(extended) average set* of $`L`$ with respect to $`H`$ as
(6) $`A(H,L)`$ $`:={\displaystyle \underset{\mu 𝔐_1(H)}{}}\overline{\mathrm{conv}}\{U^\mu (x):xL\},`$ $`A(H)`$ $`:=A(H,H).`$
###### Remark 5.
Denoting the interval
(7)
$$A(\mu ,L):=[\underset{¯}{Q}(\mu ,L),Q(\mu ,L)]=\overline{\mathrm{conv}}\{U^\mu (x):xL\},$$
we see that $`R_n(H,L)`$, $`R(H,L)`$ and $`A(H,L)`$ are all of the form $`_\mu A(\mu ,H)`$, with $`\mu `$ ranging over all averages of $`n`$ Dirac measures at points of $`H`$, over $`𝔐_1^\mathrm{\#}(H)`$ and over all of $`𝔐_1(H)`$, respectively.
###### Remark 6.
Let us explain how the above notions relate to the usual definitions of rendezvous numbers or average numbers. Suppose that $`(X,k)`$ is a metric space and that the set $`L`$ is compact. Then there is no need for the closure in the above definitions, since in this case the potential $`U^\mu `$ is continuous, so the set $`A(\mu ,L)`$, being the continuous image of the compact set $`L`$, is compact. This means that a number $`r_+`$ belongs to $`R(H,L)`$ if and only if for any finite system of (not necessarily distinct) points $`x_1,\mathrm{},x_nH`$ one finds points $`y,zL`$ satisfying
(8)
$$\frac{1}{n}\underset{j=1}{\overset{n}{}}k(y,x_j)r\text{and}\frac{1}{n}\underset{j=1}{\overset{n}{}}k(z,x_j)r.$$
This is the usual definition of weak rendezvous numbers in metric spaces (see ). In the next step, we can assume that $`L`$ is connected. In this case, (8) is further equivalent to the existence of a “rendezvous point” $`xL`$ with
(9)
$$\frac{1}{n}\underset{j=1}{\overset{n}{}}k(x,x_j)=r.$$
Of course in the above reasoning an arbitrary probability measure $`\mu `$ can replace the average of Dirac measures. To sum up, for compact and connected sets $`L`$ of metric spaces, $`R(L)`$ (and $`A(L)`$) is a single point, and it is the classical rendezvous (or average) number of $`L`$ (results of Gross , Elton and Stadje ). For further discussion and examples see .
From the above definitions it is easy to identify the lower and upper endpoints of the rendezvous and the average intervals (see ).
###### Proposition 7.
For arbitrary subsets $`H,LX`$ we have
(10) $`R_n(H,L)`$ $`=[M_n(H,L),\overline{M}_n(H,L)],`$ $`R_n(H)`$ $`=[M_n(H),\overline{M}_n(H)],`$
(11) $`R(H,L)`$ $`=[M(H,L),\overline{M}(H,L)],`$ $`R(H)`$ $`=[M(H),\overline{M}(H)],`$
(12) $`A(H,L)`$ $`=[\underset{¯}{q}(H,L),q(H,L)],`$ $`A(H)`$ $`=[\underset{¯}{q}(H),q(H)].`$
The questions of existence and uniqueness of rendezvous or average numbers, are two naturally posed problems, which were investigated in in the potential theoretic framework on locally compact spaces. In fact, non-emptyness of the rendezvous, respectively the average intervals means that in the above formulation (10) the formal lower endpoints of the intervals do not exceed the upper endpoints (we use the convention $`[a,b]=\mathrm{}`$ if $`a>b`$). While uniqueness is the same as that the respective interval reduces to one point. We recall the following two results from .
###### Theorem 8.
Let $`X`$ be a locally compact Hausdorff space, $`\mathrm{}HLX`$ be arbitrary, and let $`k`$ be any nonnegative, symmetric kernel on $`X`$. Then the intervals $`R_n(H,L)`$, $`R(H,L)`$ and $`A(H,L)`$ are nonempty.
###### Theorem 9.
Let $`X`$ be any locally compact Hausdorff topological space, $`k`$ be any l.s.c., nonnegative, symmetric kernel function, and $`\mathrm{}KX`$ compact. Then $`A(K)`$ consists of one single point. Furthermore, if $`k`$ is continuous, then even $`R(K)`$ consists of only one point.
When the rendezvous or the average interval $`R(K)`$ respectively $`A(K)`$ consists of one point only, this single point is denoted by $`r(K)`$ or $`a(K)`$, respectively.
Let us close this introduction with a few remarks to explain the idea of the present approach. Investigating the polarisation constant problem, it was found in that for certain cases the Chebyshev constants of the unit spheres $`S^2`$ and $`S^3`$ appear as polarisation constants. The arising questions led to the systematic analysis of Chebyshev constants and also transfinite diameters and minimal energies in the general potential theoretical framework . Meanwhile, the second author took part in working out a general approach which might be termed as appropriate averaging over $`S^n`$, to estimate the polarisation constant . However, it turned out that part of the results achieved through such a potential theory flavoured approach, were already obtained by García-Vázquez and Villa , who used Gross’ Theorem on the existence of rendezvous numbers successfully in the context. That suggested that perhaps there is a way to relate the two methods, or even the underlying theories, i.e., potential theory and rendezvous numbers. Our paper stems from this observation.
## 2. Rendezvous numbers for normed spaces
In the last decade many results were obtained regarding the numerical values of rendezvous numbers of concrete spaces and sets, see, e.g., . In this context, the following terminology was introduced.
###### Definition 10.
Let $`(X,)`$ be a normed space with unit closed ball $`B_X`$ and unit sphere $`S_X`$. Considering $`S_X`$ with the norm-distance, the rendezvous numbers of the arising metric space are called the rendezvous numbers of the normed space $`X`$. Accordingly, we use the script notation
(13)
$$_n(X):=R_n(S_X),(X):=R(S_X)\text{and}𝒜(X):=A(S_X).$$
It is clear that for finite dimensional normed spaces the above notion is a special case of the general notion described in Section 1. However, for infinite dimensional normed spaces the metric space $`S_X`$ will not be locally compact, as is usually assumed in the potential theoretic setup. There are two ways to tackle this, one being the extension of the theory to not necessarily locally compact but metric spaces, as is done in . There we assume that the topology arises from a metric, but relax on local compactness. Conversely, in a number of cases it is possible to consider a different topology, in which the metric is still lower semicontinuous, while $`S_X`$ becomes locally compact, hence Fuglede-type general potential theory applies. Note that in this case the topology is *not* the metric topology, which deserves some care when working with the theory. In particular, the average sets $`A(H,L)`$, referring to regular Borel measures of the space, may be different for different topologies. On the other hand, for a fixed kernel $`R(H,L)=[M(H,L),\overline{M}(H,L)]`$ is independent of any topology.
###### Proposition 11.
Let $`X`$ be any abstract set, $`k0`$ be a symmetric function from $`X\times X`$ to $`_+\{+\mathrm{}\}`$, and $`\mathrm{}HLX`$ be arbitrary subsets. Then $`R(H,L)\mathrm{}`$.
###### Proof.
The definition of rendezvous intervals, as well as the corresponding statement in Proposition 7, are independent of the topology of the underlying space. Therefore, we can just take the discrete topology of $`X`$, and note that the kernel $`k`$ becomes continuous, hence l.s.c., in this topology. Thus Theorem 8 applies and $`R(H,L)\mathrm{}`$. ∎
In view of Proposition 11, we trivially obtain the following.
###### Corollary 12.
Let $`X`$ be any normed space. Then the rendezvous set $`(X)`$ is non-empty.
This seemingly contradicts to some assertions on nonexistence in the literature: the reason is that we considered also the closure in the definition of the rendezvous and average intervals. The above proposition shows that in the context taking the closure is also helpful. Note that Baronti, Casini and Papini have already considered the closed version of the rendezvous sets, at least in normed spaces but they focused on the question of “attaining” the rendezvous numbers. In such investigations the geometry of Banach spaces plays an important role.
On the other hand, uniqueness, so nicely obtained for locally compact spaces, continuous kernels and compact sets, can not be concluded as already shown by a couple of examples in the literature (see , , ).
## 3. Average numbers for normed spaces
For compact sets $`K`$ and continuous kernels it is already known that $`A(K)=R(K)`$, and that there are counterexamples showing that in general compactness is needed (compare \[8, §6\]). Nevertheless, the assertion remains valid in normed spaces, too.
###### Theorem 13.
Let $`X`$ be any normed space. Then we have $`𝒜(X)=(X)\mathrm{}`$.
For not necessarily finite dimensional Banach spaces, we do not have the means to restrict considerations to compactly supported measures only. Instead, we prove the following result, whose easy consequence is the above theorem.
###### Theorem 14.
Let $`(Y,d)`$ be a metric space. Assume that the kernel $`k`$ is positive, symmetric and bounded and that $`\{k(,y):yY\}`$ is uniformly equicontinuous on $`(Y,d)`$. Then we have $`A(Y)=R(Y)\mathrm{}`$.
###### Lemma 15.
Assume that the kernel $`k`$ is positive, symmetric and bounded and that $`\{k(,y):yY\}`$ is uniformly equicontinuous on $`Y`$. Let $`\mu 𝔐_1(Y)`$ and $`\epsilon >0`$ be given arbitrarily. Then there exist $`m`$ and points $`x_jY`$ ($`j=1,\mathrm{},m`$) such that the potential $`U^\nu `$ of the measure $`\nu :=\frac{1}{m}_{j=1}^m\delta _{x_j}`$ approximates $`U^\mu `$ within $`\epsilon `$ uniformly on $`Y`$.
To prove this lemma we need the following elementary result.
###### Lemma 16.
For any $`\epsilon >0`$ and any finitely supported probability measure $`\nu `$, there exists a probability measure of the form $`\mu =\frac{1}{m}_{i=1}^m\delta _{z_i}`$ having the same support as $`\nu `$ and satisfying $`(1\epsilon )\nu \mu (1+\epsilon )\nu `$.
###### Proof of Lemma 15.
Without loss of generality we can assume that $`k1`$, and hence $`U^\sigma 1`$ for all probability measures $`\sigma `$. By the assumptions we find an $`r>0`$, such that $`|k(x^{},y)k(x^{\prime \prime },y)|<\epsilon /2`$, if $`d(x^{},x^{\prime \prime })<r`$. As $`\mu `$ is a *regular* Borel measure, for any given $`\eta >0`$ there exists $`Ksupp\mu Y`$ with $`\mu (K)1\eta `$. Take now $`\nu _K:=\mu _K/\mu _K`$ (with $`\mu _K`$ being the trace of $`\mu `$ on $`K`$) so that $`\nu _K𝔐_1(K)`$. Note that
$$|U^\mu (x)U^{\mu _K}(x)|\underset{yY}{sup}k(x,y)\mu \mu _K\eta ,$$
and
$$|U^{\nu _K}(x)U^{\mu _K}(x)||U^{\nu _K}(x)|(1\mu _K)\eta .$$
Consider now the covering of the compact set $`KY`$ by open balls $`B(y,r)`$, with $`yK`$. By compactness there exist a finite sub-covering, i.e., there exist $`n`$, $`y_jK`$, $`B_j:=B(y_j,r)`$ ($`j=1,\mathrm{},n`$) satisfying $`K_{j=1}^nB_j`$. Put $`D_1:=B_1`$, and for $`j=2,\mathrm{},n`$ put $`D_j:=B_j_{i=1}^{j1}B_i`$, $`\alpha _j:=\nu _K(D_j)0`$. Clearly $`_{j=1}^n\alpha _j=\nu _K(K)=1`$. Consider the finitely supported measure $`\sigma :=_{j=1}^n\alpha _j\delta _{y_j}𝔐_1(K)`$. Then we have
$`|U^{\nu _K}(x)U^\sigma (x)|`$ $`=|{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{D_j}}k(x,y)k(x,y_j)d\nu _K(y)|`$
$`{\displaystyle \underset{j=1}{\overset{n}{}}}\alpha _j\underset{yD_jB_j}{sup}|k(x,y)k(x,y_j)|\epsilon /2.`$
Finally, the application of Lemma 16 to $`\sigma `$ yields an approximating measure $`\nu :=\frac{1}{m}_{j=1}^m\delta _{x_j}`$ with $`m`$ and points $`x_jK`$ ($`j=1,\mathrm{},m`$) so that $`(1\eta )\sigma \nu (1+\eta )\sigma `$, and thus
$$|U^\sigma (x)U^\nu (x)|\eta |U^\sigma (x)|\eta .$$
Collecting all the above, we find
$$|U^\mu (x)U^\nu (x)|3\eta +\epsilon /2<\epsilon ,$$
if $`\eta <\epsilon /6`$. ∎
###### Proof of Theorem 14.
It is obvious that $`A(Y)R(Y)`$, hence it suffices to show the converse inclusion. Let $`\mu 𝔐_1(Y)`$ be arbitrary and consider $`A(\mu ,Y)=[a,b]`$, (where $`a=\underset{¯}{Q}(\mu ,Y)`$, $`b=Q(\mu ,Y)`$). Let us take $`\epsilon :=1/n`$ and look at the measure $`\nu :=\nu _n`$ provided by Lemma 15 to $`\mu `$ and $`\epsilon `$. Since the potential functions are uniformly close to each other on $`Y`$, their infima and suprema are also within $`\epsilon `$: that is, $`|\underset{¯}{Q}(\mu ,Y)\underset{¯}{Q}(\nu ,Y)|\epsilon `$, $`|Q(\mu ,Y)Q(\nu ,Y)|\epsilon `$. In other words, $`A(\nu _n,Y)[a\frac{1}{n},b+\frac{1}{n}]`$ and thus $`_{n=1}^{\mathrm{}}A(\nu _n,Y)[a,b]=A(\mu ,Y)`$. It follows that $`R(Y)=_{\nu 𝔐_1^\mathrm{\#}(Y)}A(\nu ,Y)A(\mu ,Y)`$ for all $`\mu 𝔐_1(Y)`$, hence $`R(Y)_{\mu 𝔐_1(Y)}A(\mu ,Y)=A(Y)`$ and the theorem is proved. ∎
###### Remark 17.
The assumptions of Theorem 14 are fulfilled, for instance, when $`(X,d)`$ is a metrisable topological vector space and $`k(x,y)=f(d(x,y))`$, where $`f`$ is a continuous function (see Section 4 below).
###### Lemma 18.
Let $`X`$ be a (not necessarily locally compact) Hausdorff topological vector space, $`k`$ a l.s.c., nonnegative, symmetric and convex kernel function on $`X\times X`$, and $`\mu 𝔐_1(X)`$. Then the potential function $`U^\mu `$ is convex.
###### Proof.
Take any $`x,y,zX`$ with $`z=\alpha x+(1\alpha )y`$, where $`0\alpha 1`$. We then have
$`U^\mu (z)`$ $`={\displaystyle _X}k(z,w)d\mu (w){\displaystyle _X}(\alpha k(x,w)+(1\alpha )k(y,w))d\mu (w)=`$
$`=\alpha U^\mu (x)+(1\alpha )U^\mu (y),`$
and that was to be proved. ∎
###### Lemma 19.
Let $`X`$ be a (not necessarily locally compact) Hausdorff topological vector space, $`k`$ a l.s.c., nonnegative, symmetric and convex kernel function on $`X\times X`$, $`HX`$ a bounded set, and $`H`$ be its boundary. Then for any $`\mu 𝔐_1(X)`$ the potential function $`U^\mu `$ satisfies $`sup_HU^\mu =sup_HU^\mu `$.
###### Proof.
We are to show $`sup_HU^\mu sup_HU^\mu `$, the other direction being obvious. Let now $`xH`$ be arbitrary: we show that $`U^\mu (x)sup_HU^\mu `$. Draw any straight line $`\mathrm{}`$ through $`x`$. Since $`H`$ is bounded, both closed half-lines of $`\mathrm{}`$, starting from $`x`$, contain boundary points of $`H`$; that is, if these points are $`y,z\mathrm{}`$, then $`x[y,z]`$ with $`y,zH`$. According to Lemma 18, the potential is convex, which immediately yields $`U^\mu (x)\mathrm{max}(U^\mu (y),U^\mu (z))sup_HU^\mu `$. ∎
###### Corollary 20.
If $`X`$ is a normed space, then $`q(S_X,S_X)=q(S_X,B_X)`$ and $`\overline{M}(S_X,S_X)=\overline{M}(S_X,B_X)`$.
###### Remark 21.
Note that for the other endpoints of the average intervals generally we may have strict inequality: $`\underset{¯}{q}(S_X,B_X)<\underset{¯}{q}(S_X,S_X)`$. As will be seen in Theorem 23, for any $`1<p<+\mathrm{}`$ $`(\mathrm{}_p)=2^{1/p}>1`$. However, for any measure $`\mu 𝔐_1(S_X)`$, it is clear that $`\mathrm{𝟎}B_X`$ satisfies $`U^\mu (\mathrm{𝟎})=_{S_X}1d\mu =1`$, hence $`\underset{¯}{q}(S_X,B_X)1`$. In fact, $`\underset{¯}{q}(S_X,B_X)=1`$ is also true, since for any set $`H`$ and for two points $`x,yH`$, the corresponding measure $`\nu :=\frac{1}{2}(\delta _x+\delta _y)`$ always provides, by the triangle inequality, $`\underset{¯}{Q}(\nu ,B_X)xy/2`$, which can be as large as half the diameter.
###### Remark 22.
It is straightforward to show that $`[M(H),\overline{M}(H)]=R(H)[\frac{1}{2}diam(H),diam(H)]`$ (see, e.g., ). Indeed, $`\overline{M}(H)diam(H)`$ is trivial, while the lower estimate $`\frac{1}{2}diam(H)M(H)`$ is essentially contained in the previous remark.
## 4. Rendezvous numbers for $`L_p`$ spaces
We already know that the rendezvous interval of a Banach space is not empty. Here we identify the rendezvous, hence the average intervals of the $`L_p`$ spaces.
Let $`(\mathrm{\Omega },,\mu )`$ be a measure space. To complement the whole scale $`1p<+\mathrm{}`$, we consider $`L_p:=L_p(\mathrm{\Omega },,\mu )`$ when $`0<p<1`$ as well. In this case $`L_p`$ will not be a Banach space, but if endowed with the metric
(14)
$$d(f,g)=_\mathrm{\Omega }|fg|^pd\mu ,$$
it is a complete, metrisable, topological vector space (of course we have to identify functions coinciding on $`\mu `$-null sets). First we calculate the rendezvous number with respect to the symmetric function $`fg:=d(f,g)^{1/p}`$ instead of the metric $`d`$, this fits well together with the case $`p1`$. Of course, now $`S_{L_p}`$ denotes the “unit sphere” with respect to $``$ for all $`0<p<+\mathrm{}`$.
###### Theorem 23.
Let $`0<p<+\mathrm{}`$ be arbitrary and consider $`L_p(\mathrm{\Omega },,\mu )`$ over either $``$ or $``$. If $`L_p`$ is infinite dimensional, we have $`a(S_{L_p})=r(S_{L_p})=2^{1/p}`$.
###### Proof.
The following applies for both the complex- and real valued cases, hence we do not mention the underlying number field any more. Further we write briefly $`L_p`$ instead of $`L_p(\mathrm{\Omega },,\mu )`$. We already know that $`R(S_{L_p})`$ is nonempty (see Proposition 11) and is a compact interval (indeed, in case $`p1`$ it is a subset of the interval $`[1,2]`$, see Remark 22). Moreover, we have $`R(S_{L_p})=[M(S_{L_p}),\overline{M}(S_{L_p})]`$, and by Theorem 13, Theorem 14 and Remark 17 we have $`A(S_{L_p})=R(S_{L_p})`$. Therefore we need only to show that $`M(S_{L_p})2^{1/p}`$ and that $`\overline{M}(S_{L_p})2^{1/p}`$.
Since $`L_p`$ is infinite dimensional, there exist $`w_jS_{L_p}`$, $`j`$ such that the sets $`A_j:=\{x:x\mathrm{\Omega },w_j(x)0\}`$ are pairwise disjoint. For any function $`gL_p`$ let us introduce the notation $`g_j:=\chi _{A_j}g`$.
*Part 1: $`M(S_{L_p})2^{1/p}`$.* With the functions $`w_j`$ and for any $`fS_{L_p}`$ we have
$`{\displaystyle \underset{j=1}{\overset{n}{}}}fw_j`$ $`={\displaystyle \underset{j=1}{\overset{n}{}}}\left({\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}fw_j_k^p\right)^{\frac{1}{p}}={\displaystyle \underset{j=1}{\overset{n}{}}}\left(fw_j_j^p+{\displaystyle \underset{k=1,kj}{\overset{\mathrm{}}{}}}f_k^p\right)^{\frac{1}{p}}=`$
(15) $`={\displaystyle \underset{j=1}{\overset{n}{}}}\left(w_jf_j^p+1f_j^p\right)^{\frac{1}{p}}.`$
Now we distinguish between the cases $`p<1`$ and $`p1`$. First let $`p<1`$, then using $`f_j1`$ and $`fS_{L_p}`$, we can continue (4)
$`{\displaystyle \underset{j=1}{\overset{n}{}}}fw_j={\displaystyle \underset{j=1}{\overset{n}{}}}\left(1f_j^p+w_jf_j^p\right)^{\frac{1}{p}}`$
$`{\displaystyle \underset{j=1}{\overset{n}{}}}\left(1f_j^p+\left|w_j_j^pf_j^p\right|\right)^{\frac{1}{p}}={\displaystyle \underset{j=1}{\overset{n}{}}}\left(1f_j^p+\left(1f_j^p\right)\right)^{\frac{1}{p}}=`$
$`=2^{1/p}{\displaystyle \underset{j=1}{\overset{n}{}}}\left(1f_j^p\right)^{\frac{1}{p}}2^{1/p}n{\displaystyle \frac{(n1)^{\frac{1}{p}}}{n^{\frac{1}{p}}}},`$
using again $`fS_{L_p}`$ and the convexity of the function $`xx^{1/p}`$ in the last step. Second, let $`p1`$, then we can write
$`{\displaystyle \underset{j=1}{\overset{n}{}}}fw_j`$ $`={\displaystyle \underset{j=1}{\overset{n}{}}}\left(1f_j^p+w_jf_j^p\right)^{\frac{1}{p}}`$
$`{\displaystyle \underset{j=1}{\overset{n}{}}}\left(1f_j^p+\left(1f_j\right)^p\right)^{\frac{1}{p}}=`$
(16) $`={\displaystyle \underset{j=1}{\overset{n}{}}}\left(1f_j^p\left(1f_j\right)^p+2\left(1f_j^p\right)^p\right)^{\frac{1}{p}}`$
$`{\displaystyle \underset{j=1}{\overset{n}{}}}\left(2\left(1f_j\right)^p\right)^{\frac{1}{p}}=2^{\frac{1}{p}}{\displaystyle \underset{j=1}{\overset{n}{}}}\left(1f_j\right)2^{\frac{1}{p}}(nn^{\frac{1}{q}}),`$
where $`\frac{1}{p}+\frac{1}{q}`$ and using again $`fS_{L_p}`$ and Hölder’s inequality in the last step. We see that in both cases for the given $`n`$-point distribution (concentrated on the $`w_j`$s, $`j=1,\mathrm{},n`$) the corresponding potential is minorised by the right hand sides divided by $`n`$, hence we find $`M_n(S_{L_p})2^{1/p}+o(1)`$ (as $`n\mathrm{}`$) and $`M(S_{L_p})2^{1/p}`$ follows. (Note that, e.g., calculates even the exact value of $`(\mathrm{}_1^n())`$, but here the task is a little bit different.)
*Part 2: $`\overline{M}(S_{L_p})2^{1/p}`$.* Let $`n`$ and consider the same functions $`w_j`$, $`j`$ as in the first part. Then for any point $`fS_{L_p}`$ and for any given parameter $`\eta >0`$ we have, by (4), for $`p1`$
(17) $`{\displaystyle \underset{j=1}{\overset{n}{}}}fw_j`$ $`{\displaystyle \underset{j=1}{\overset{n}{}}}\left(1+(1+f_j)^p\right)^{\frac{1}{p}}`$
$`{\displaystyle \underset{j:f_j>\eta }{}}\left(1+2^p\right)^{\frac{1}{p}}+{\displaystyle \underset{j:f_j\eta }{}}\left(1+(1+\eta )^p\right)^{\frac{1}{p}}`$
$`{\displaystyle \frac{1}{\eta ^p}}\left(1+2^p\right)^{\frac{1}{p}}+n2^{\frac{1}{p}}(1+\eta ),`$
and for $`p<1`$
$`{\displaystyle \underset{j=1}{\overset{n}{}}}fw_j`$ $`{\displaystyle \underset{j=1}{\overset{n}{}}}\left(1+(1+f_j^p)\right)^{\frac{1}{p}}`$
(18) $`{\displaystyle \underset{j:f_j>\eta }{}}3^{\frac{1}{p}}+{\displaystyle \underset{j:f_j\eta }{}}\left(1+(1+\eta ^p)\right)^{\frac{1}{p}}{\displaystyle \frac{3^{\frac{1}{p}}}{\eta ^p}}+n2^{\frac{1}{p}}(1+\eta ^p)^{\frac{1}{p}}.`$
Choosing, e.g., $`\eta :=n^{1/2p}`$, we obtain the estimate $`U^\nu (f)2^{1/p}+\mathrm{o}_p(1)`$ (as $`n+\mathrm{}`$, $`fS_{L_p}`$) for the measure $`\nu :=\frac{1}{n}_{j=1}^n\delta _{w_j}`$. It follows that for the given $`n`$-point distribution $`\nu `$ we have $`\overline{M}_n(S_{L_p})Q(\nu ,S_{L_p})=2^{1/p}+\mathrm{o}_p(1)2^{1/p}`$ ($`n\mathrm{}`$), and so even $`\overline{M}(S_{L_p})2^{1/p}`$. ∎
Note that for $`p1`$ already Lin showed that $`(\mathrm{}_p)\{2^{1/p}\}`$ for “strict” rendezvous numbers (actually by a similar argument). So this and the non-emptyness of the rendezvous interval (Corollary 12) give the above result for $`\mathrm{}_p`$ ($`1p<+\mathrm{}`$).
###### Corollary 24 (Wolf, Lin).
Let $``$ be an infinite dimensional Hilbert space over any of the number fields $``$ or $``$. Then we have $`𝒜()=()=\{\sqrt{2}\}`$.
###### Remark 25.
In the above proof we actually used the same point distribution in both parts of the proof, therefore we have proved the existence of $`\epsilon `$-*quasi-invariant measures*. We say that there exist *$`\epsilon `$-quasi-invariant measures* cf. \[9, Definition 5.9\] for the kernel $`k`$ on $`S`$, if for all $`\epsilon >0`$ there is $`\nu 𝔐_1(S)`$ satisfying $`Q(\nu ,S)\underset{¯}{Q}(\nu ,S)\epsilon `$. Generally, if $`S`$ is compact, by weak-compactness we obtain the existence of a true invariant measure $`\mu 𝔐_1(S)`$, i.e., whose potential $`U^\mu `$ is constant on $`S`$. Of course, for $`S=S_\mathrm{}_p`$ this is not the case. But as we saw above there exist $`\epsilon `$-quasi-invariant measures, and by \[9, Proposition 5.11\] it is already enough to conclude that the average interval reduces to one point.
###### Theorem 26.
Let $`0<p<1`$ be arbitrary and $`L_p(\mathrm{\Omega },,\mu )`$ be the vector space of $`p`$-integrable functions endowed with the metric $`d`$ defined in (14). For the rendezvous interval of the unit ball $`S_{L_p}`$ of $`L_p`$ we have $`A(S_{L_p})=R(S_{L_p})=2`$.
###### Proof.
By Remark 22 we know $`\overline{M}(S_{L_p})2`$. So we only need to estimate the lower endpoint of $`R(S_{L_p})`$ from below. This can be done analogously to (4) and (4) by considering the functions $`w_j`$ and sets $`A_j`$ used in the proof of Theorem 23. Let us further use the abbreviation $`|f|_j:=_{A_j}|f|^pd\mu `$. For an arbitrary $`fS_{L_p}`$ we can write
$`{\displaystyle \underset{j=1}{\overset{n}{}}}d(f,w_j)`$ $`={\displaystyle \underset{j=1}{\overset{n}{}}}\left({\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}|fw_j|_k\right)={\displaystyle \underset{j=1}{\overset{n}{}}}\left(|fw_j|_j+{\displaystyle \underset{k=1,kj}{\overset{\mathrm{}}{}}}|f|_k\right)=`$
$`={\displaystyle \underset{j=1}{\overset{n}{}}}\left(1|f|_j+|w_jf|_j\right)`$
$`{\displaystyle \underset{j=1}{\overset{n}{}}}\left(1|f|_j+1|f|_j\right)=2{\displaystyle \underset{j=1}{\overset{n}{}}}\left(1|f|_j\right)2(n1).`$
Therefore $`M_n(S_{L_p})21/n`$, hence $`M(S_{L_p})2`$. ∎
It was already pointed out by Lin in that $`lim_n\mathrm{}r(S_{\mathrm{}_p^n})=2^{1/p}`$, if $`1p<+\mathrm{}`$. The following result, inspired by an analogous argument of García-Vázquez and Villa , explains this phenomenon in view of Theorem 23 above.
###### Theorem 27.
Let $`X`$ be a normed space and $`X_n`$ an increasing sequence of subspaces such that $`_{n=1}^{\mathrm{}}X_n`$ is dense in $`X`$. Suppose that $`(X_n,_n)`$ is a normed space and $`\rho _nR(S_{(X_n,_n)})`$. Assume that
$$\underset{n\mathrm{}}{lim}\underset{xX_nS_X}{sup}|1x_n|=0.$$
Then any accumulation point $`\rho `$ of the sequence $`\rho _n`$ belongs to $`R(S_X)`$.
###### Proof.
Let $`\rho `$ be an accumulation point of the sequence $`\rho _n`$. Assume without loss of generality that $`\rho _n\rho `$. Let $`\epsilon >0`$ be given. Then for sufficiently large $`nn_0(\epsilon )`$ we have $`[\rho _n\epsilon ,\rho _n+\epsilon ][\rho 2\epsilon ,\rho +2\epsilon ]`$.
By definition $`\rho _n_m(X_n)`$ for all $`m`$. Let $`m`$ be fixed, and $`x_1,\mathrm{},x_mS_X`$ be arbitrary. Take any $`y_jX_nS_X`$ with $`y_jx_j\epsilon `$. Such $`y_j`$ exists in view of the denseness of $`_{n=1}^{\mathrm{}}X_n`$ in $`X`$. By assumption, we have $`|zz_n|\epsilon z`$ for all sufficiently large $`nn_1n_0`$ and all $`zX_n`$. In particular, $`|1y_j_n|\epsilon `$. By definition of the rendezvous interval $`_m(X_n)`$, we find $`z_nS_{(X_n,_n)}`$ satisfying
(19)
$$\frac{1}{m}\underset{j=1}{\overset{m}{}}\frac{y_j}{y_j_n}z_n_n[\rho _n\epsilon ,\rho _n+\epsilon ][\rho 2\epsilon ,\rho +2\epsilon ].$$
According to the above, for all $`nn_1`$ we have
$`z_n{\displaystyle \frac{z_n}{z_n}}_n`$ $`=\left|1{\displaystyle \frac{1}{z_n}}\right|z_n_n=\left|1{\displaystyle \frac{z_n_n}{z_n}}\right|\epsilon ,`$
and
$`y_j{\displaystyle \frac{y_j}{y_j_n}}_n`$ $`=\left|1{\displaystyle \frac{1}{y_j_n}}\right|y_j_n=\left|y_j_n1\right|\epsilon .`$
Using these two inequalities in (19) we obtain
$$\frac{1}{m}\underset{j=1}{\overset{m}{}}y_j\frac{z_n}{z_n}_n[\rho 4\epsilon ,\rho +4\epsilon ].$$
For $`n>n_1`$ we also know
$$\left|y_j\frac{z_n}{z_n}y_j\frac{z_n}{z_n}_n\right|\epsilon y_j\frac{z_n}{z_n}2\epsilon ,$$
therefore we can write
$$\frac{1}{m}\underset{j=1}{\overset{m}{}}y_j\frac{z_n}{z_n}[\rho 6\epsilon ,\rho +6\epsilon ]\text{and}\frac{1}{m}\underset{j=1}{\overset{m}{}}x_j\frac{z_n}{z_n}[\rho 7\epsilon ,\rho +7\epsilon ].$$
This shows $`\rho R_m(S_X)`$, which being valid for all $`m`$, gives $`\rho R(S_X)`$. ∎
This theorem immediately gives the following corollary.
###### Corollary 28.
Let $`X`$ be a normed space and $`X_n`$ an increasing sequence of finite dimensional subspaces such that $`_{n=1}^{\mathrm{}}X_n`$ is dense in $`X`$. Let $`\rho `$ be an accumulation point of the sequence $`r(S_{X_n})`$ ($`r(S_{X_n})`$ exists uniquely by the compactness of $`S_{X_n}`$). Then $`\rho (X)`$.
## 5. Chebyshev centres, entropy and rendezvous numbers
###### Definition 29.
Let $`KX`$ be a compact, convex subset of some normed space $`X`$, with $`d`$ being the metric induced by the norm. Then the Chebyshev centre $`c:=c(K)K`$ and the Chebyshev out-radius $`\rho :=\rho (K)`$ of the set $`K`$ are the centre and the radius, respectively, of the closed ball $`\overline{B}:=\overline{B}(c,\rho )`$ of minimal radius with $`cK\overline{B}(c,\rho )`$.
Clearly, for any compact $`K`$ such a minimal ball always exists, and for convex sets it is even unique. Note that it is important in the definition that $`c`$ be chosen within $`K`$; for discussion see and .
Quoting private communication from Esther and George Szekeres, in Cleary, Morris and Yost present the following beautiful result with a nice, direct elementary proof. Here we present our even shorter version relying on the potential theoretical background developed.
###### Theorem 30 (E. and G. Szekeres).
Let $`KX`$ be a compact, convex subset of some normed space $`X`$, with $`d`$ being the metric induced by the norm. Then the Chebyshev radius and the rendezvous number of the set $`K`$ equal: $`r(K)=\rho (K)`$.
###### Proof.
Existence and uniqueness of $`R(K)=\{r(K)\}`$ and also the equality $`R(K)=A(K)`$ are already known from Theorems 8 and 9. Further, if $`c`$ is a Chebyshev centre of $`K`$, then to all points $`xK`$ we have $`xc\rho `$. Hence for any probability measure $`\mu 𝔐_1(K)`$ also the potential satisfies $`U^\mu (c)\rho (K)`$. It follows that $`\underset{¯}{Q}(\mu ,K)\rho `$ for all $`\mu 𝔐_1(K)`$, hence also $`r(K)=\underset{¯}{q}(K)\rho (K)`$.
Conversely, for $`\epsilon >0`$ let $`y_jK`$ be arbitrary points ($`j=1,\mathrm{},n`$) satisfying $`Q(\nu ,K)<\overline{M}(K)+\epsilon `$ with $`\nu :=\frac{1}{n}_{j=1}^n\delta _{y_j}`$. As $`K`$ is convex, it contains the convex combination $`y:=\frac{1}{n}_{j=1}^ny_jK`$ of the given points. Thus by the convexity of the norm for arbitrary $`xK`$ the estimate $`yx\frac{1}{n}_{j=1}^ny_jx=U^\nu (x)\overline{M}(K)+\epsilon `$ holds. Hence $`\overline{B}(y,\overline{M}(K)+\epsilon )`$ covers $`K`$ and $`\rho (K)\overline{M}(K)+\epsilon `$, which implies also $`\rho (K)\overline{M}(K)=q(K)=r(K)`$. ∎
Recall that for a positive number $`t>0`$ and a set $`HX`$ of a metric space $`X`$ with metric $`d`$ the $`t`$-covering number $`N(t,H)`$ is defined as
$$N(t,H):=\mathrm{min}\{n:y_jH(j=1,\mathrm{},n)\text{such that}H\underset{j=1}{\overset{n}{}}B(y_j,t)\}.$$
If there is no finite set of balls of radius $`t`$ which can cover the set $`H`$, then we say that $`N(t,H)=+\mathrm{}`$. Similarly, if $`H,LX`$, then
$$N(t,H,L):=\mathrm{min}\{n:y_jH(j=1,\mathrm{},n)\text{such that}L\underset{j=1}{\overset{n}{}}B(y_j,t)\}$$
with $`\mathrm{min}\mathrm{}=+\mathrm{}`$ being in effect again. The next observation is almost obvious.
###### Proposition 31.
Let $`t>0`$ and $`H,LX`$. We have $`\overline{M}_n(H,L)t`$ for all $`n<N(t,H,L)`$. In particular, if $`N(t,H,L)=+\mathrm{}`$, then $`t\overline{M}(H,L)=supR(H,L)`$.
###### Proof.
Since $`n<N(t,H,L)`$, for any system of points $`y_jH(j=1,\mathrm{},n)`$, there exists some point $`xL`$ so that $`d(x,y_j)t`$ for all $`j=1,\mathrm{},n`$. Therefore, $`sup_{xL}_{j=1}^nd(x,y_j)nt`$ holds for all systems of $`n`$ points, whence $`\overline{M}_n(H,L)t`$, and (10) concludes the proof. ∎
Recall that a set $`HX`$ is called *totally bounded*, if $`N(t,H)<+\mathrm{}`$ for all $`t>0`$. In Banach spaces this is the same as the conditional compactness of $`H`$, i.e., that $`\overline{H}X`$ is compact set. The proposition shows that for subsets $`H`$ which are not totally bounded, there is always a positive lower bound of $`\overline{M}(H,H)`$. The above proposition, however easy, provides an essential help in describing some rendezvous numbers. For instance, there is an elegant interpretation of the following result.
###### Theorem 32.
Let $`K`$ be a compact Hausdorff topological space without isolated points, and consider $`C(K)`$ the Banach space of real- or complex-valued continuous functions over $`K`$. Then we have $`\overline{M}(S_{C(K)})=2`$.
###### Proof.
Denote by $`S`$ the unit sphere of $`C(K)`$. We show that $`N(t,S,S)=+\mathrm{}`$ for $`0<t<2`$, then by Proposition 31 $`\overline{M}(S)t`$ hence $`\overline{M}(S)2`$ will follow. Then, by Remark 22, we must actually have $`\overline{M}(S)=2`$.
So let $`0<t<2`$ and $`f_1,\mathrm{}f_mS`$. Further let $`x_jK`$ be one of the maximum points of $`|f_j|`$, i.e., $`|f_j(x_j)|=1`$ and take $`\epsilon >0`$ small such that $`t+\epsilon <2`$. By continuity, there exist neighbourhoods $`G_j`$ of $`x_j`$ with $`|f_j(x_j)f_j(y)|<\epsilon `$, $`(yG_j)`$, for all $`j=1,\mathrm{},m`$. Take $`y_jG_j`$ distinct points ($`x_j`$ is not an isolated point!). By Tietze’s Theorem there exists a continuous function $`gS`$, such that $`g(y_j)=f_j(x_j)`$, for all $`j=1,\mathrm{},m`$. But then $`S`$ can not be covered by the balls $`B(f_j,t)`$, because this particular $`g`$ is not covered. Thus we conclude $`N(t,S,S)=+\mathrm{}`$. ∎
The above result is already present in García-Vázquez, Villa and Lin , where the authors determine the rendezvous interval of $`C(K)`$. Their proofs follow the same line, we included it for the sake of illustration of the role of entropy.
The real-valued case in the following theorem is due to Wolf , see also Lin .
###### Theorem 33.
Let $`c_0`$ denote the Banach space of real or complex valued null-sequences. Then we have
(20)
$$𝒜(c_0)=(c_0)=[1,\sigma ],$$
where $`\sigma =3/2`$ or $`\sigma =\frac{1}{3}+\frac{2\sqrt{3}}{\pi }`$ in the case of $``$ respectively $``$-valued sequences.
###### Proof.
In the real case Wolf showed that $`r(S_\mathrm{}_{\mathrm{}}^n)=\sigma `$ (\[29, Proposition 1\]), while the corresponding equality in the complex case is due to García-Vázquez and Villa. So by Corollary 28 we see that $`\sigma (c_0)`$. Applying the same idea as in \[13, Theorem 5\], we can split the space as $`c_0=𝕂\times c_0`$, where $`𝕂`$ is either the complex or the real scalar field. Consider the measure $`\mu =\lambda \delta _0`$, where $`\lambda `$ is the normalised Haar measure on $`S_𝕂`$ and $`\delta _0`$ is the Dirac measure on $`c_0`$ at the constant $`0`$ sequence. Clearly $`\mu `$ is supported in $`S_{c_0}`$. In case of real-valued sequences, Wolf essentially showed $`Q(\mu ,S_{c_0})\sigma `$ (see proof of Proposition 1 in ). Moreover, one can repeat the arguments from to see that $`Q(\mu ,S_{c_0})\sigma `$ in the complex case, too. So in both cases we have $`q(S_{c_0})\sigma =\sigma (𝕂)`$ and as by Theorem 13 we know $`\overline{M}(S_{c_0})=q(S_{c_0})`$, we find $`\overline{M}(S_{c_0})\sigma `$. Because $`\sigma (c_0)`$, the only possibility is $`\overline{M}(S_{c_0})=\sigma `$.
To calculate the lower endpoint of the rendezvous interval, let now $`m`$ be fixed and $`x_1,\mathrm{},x_mS_{c_0}`$ be arbitrary. For $`\epsilon >0`$ take $`n_0`$ be so large that $`|x_j(n)|<\epsilon `$ whenever $`nn_0`$ for all $`j=1,\mathrm{},m`$. Now let $`z`$ be the element of $`c_0`$ being almost completely $`0`$ but $`1`$ at the $`n_0`$th coordinate. Then
$$\frac{1}{m}\underset{j=1}{\overset{m}{}}x_jz1+\epsilon ,$$
so $`M_m(S_{c_0})1`$, and therefore $`M(S_{c_0})1`$. But then by Remark 22 we have $`M(S_{c_0})=1`$. We have calculated the lower and upper endpoints of the rendezvous (and the average) intervals to arrive at the assertion. ∎ |
warning/0507/cond-mat0507246.html | ar5iv | text | # The magnetic susceptibility of exchange-disordered antiferromagnetic finite chains
## I Introduction
Low dimensional quantum antiferromagnets have been intensively studied over the years, and became a reference problem in Condensed Matter Physics, Strongly Correlated Systems, and Statistical Mechanics.general-1D Quantum spin chains may exhibit quantum phase transitions and, in the continuum limit, turned out to be an active area for application of quantum field theory in Condensed Matter Physics. Random quantum antiferromagnetic (AF) chains, which model the magnetic behavior of several quasi-one-dimensional compounds, have also attracted a great deal of attentionvillain . Theoretical studiesma1 ; daniel of the quantum spin-1/2 disordered Heisenberg antiferromagnet in 1D show that, for infinite chains, the magnetic susceptibility behavior is essentially insensitive to the specific disorder distribution. While the behavior in the thermodynamic limit is thus well established, there has not yet been a detailed study of how that limit is approached, and specifically how the behavior of the magnetic susceptibilities of finite chains might depend on the disorder distribution. This is an interesting question in light of current efforts towards nanoscale control and applications which, coupled to improved techniques of material growth, directed interest into magnetic nanostructures, in particular into finite magnetic chains.Gambardella
One of the most exciting potential applications of magnetism at the extreme microscopic level consists in utilizing the 2-level dynamics of the electron spin as a physical implementation of the quantum-bit (qubit) in a solid state quantum computer, where the required entanglement between qubits could be provided by the exchange coupling between electrons.LD For example, the quantum behavior of exchange-coupled electrons bound to an array of phosphorous donors in silicon is a key element in Kane’s proposal for a Si-based quantum computer.Kane Low-temperature magnetic susceptibility measurements of P doped Si have provided valuable information regarding the exchange distributions in randomly 3D-doped samples.Andres It is expected that similar measurements might also be relevant in characterizing linear arrays of donors. In Kane’s proposal, two-qubit operations are mediated by the exchange interaction between electrons bound to nearest-neighbor P atoms in the chain. Previous studies have shown that this exchange coupling $`J`$ is always AF, and that its strength is highly sensitive to the inter-donor relative positioning.Andres ; bk1 ; bk2 Indeed, changes of just one lattice parameter in the relative positioning of two P impurities may alter the magnitude of the coupling between them by orders of magnitude. Controllable exchange coupling compatible with Kane’s original proposal would be achieved if all donors in the chain could be positioned exactly along a single crystal axis.bk2 In this situation, the exchange coupling behavior is (as assumed by KaneKane ) similar to the hydrogenic case,hf i.e., it decays exponentially with increasing interdonor distance. Uncertainties in the interdonor distances along the chain would result in a narrow distribution of values for $`J`$ around a “target” value $`J_0`$. Assuming a substitutional donor positioning precision of about 1 nm in the Si lattice,fabrication1 the perfect alignment situation may be modeled by a trimodal exchange distribution $`P_{tri}(J)`$. However, if instead of perfect alignment the donor positions are randomly distributed among all substitutional sites within a small spherical region of 1 nm around the ideal (“target”) impurity sites, the peculiar band structure of Si leads to a wide distribution of exchange coupling, peaked at $`J=0`$,bk2 causing difficulties in the operation and control of the “exchange gates”. We approximate such distribution here by an exponential function, $`P_{exp}(J)`$ for $`J>0`$. It is therefore clear that the exchange disorder distribution is also highly sensitive to the positioning distribution of the P donors in Si.
The magnetic susceptibility behavior of infinite chains is essentially insensitive to the specific disorder distribution, as mentioned above, indicating that susceptibility measurements would not differentiate among the two cases mentioned above in this limit. Of course this may be different for finite chains, a situation which is also of practical interest in terms of guiding current fabrication efforts towards P donors positioning in Si.fabrication1 ; fabrication2 The aim of the present work is to shed light on this problem by investigating the relation between the magnetic response of linear chains of spins and the distributions of the exchange interaction within the chains. Our results indicate that for AF disordered chains with even and relatively small number $`N`$ of sites, the exchange distributions $`P_{tri}(J)`$ and $`P_{exp}(J)`$ lead to quite distinct low temperature behavior of the zero-frequency uniform magnetic susceptibility. Hence, the two distributions could be experimentally distinguished by sufficiently sensitive magnetic measurements, providing useful information regarding donor alignment.
This paper is organized as follows. In Sec. II we briefly review results available in the literature regarding infinite AF chains. In Sec. III we consider finite chains with relatively small number of spins, starting from spin pairs and trios for which analytical solutions are obtained, as well as 8-spin chains, which are solved numerically. In Sec. IV we analyze the crossover into the thermodynamic limit by solving for longer chains via a Quantum Monte-Carlo method. Our summary and conclusions are presented in Sec. V.
## II Infinite Antiferromagnetic Chains
The Hamiltonian describing an open chain with $`N`$ spins is
$$=\underset{i}{\overset{N1}{}}J_i\stackrel{}{S_i}\stackrel{}{S}_{i+1},$$
(1)
where the spin quantum number is $`S=1/2`$. Since we are interested in AF chains, we assume that $`J_i0`$ for every $`i`$. We briefly review in this section several pertinent results available in the literature for AF chains in the $`N\mathrm{}`$ limit.
The ground state of an ordered infinite chain $`(J_i=J`$ for every $`i`$) can be obtained from Bethe ansatz. Griffithsgriffiths has shown that the zero temperature susceptibility per spin of such system is finite and given by $`\chi (T=0)/\chi _0(J)=1/\pi ^2,`$ where $`\chi _0(J)=g^2(\mu _B)^2/J`$, where $`\mu _B`$ is the Bohr magneton and $`g(=2)`$ is the Landé factor. For general $`T`$, field theory methods \[$`k=1`$ Wess-Zumino-Witten (WZW) non-linear $`\sigma `$ model\] affleck give
$$\frac{\chi (T)}{\chi _0(J)}=\frac{1}{(\pi )^2}\left(1+\frac{1}{2\mathrm{ln}(T_0/T)}\right),$$
(2)
where $`T_0`$ is a temperature cutoff. Quantum Monte-Carlo (QMC) calculations of $`\chi (T)`$ have been carried out by Kim et al.Kim and their results for $`T_0/J=1.8`$ (we use units of energy for temperature, that is, the Boltzmann constant $`k_B`$ is set equal to one) are well fitted by the WZW expression (2).
According to real-space renormalization group theory,daniel the introduction of any amount of disorder drives the system into a random singlet phase, in which each spin forms a singlet pair with another spin; pairs with arbitrarily long distance also exist. Bonds among distant spins, however, correspond to very weak coupling. The low-temperature excitations basically involve breaking these weakest bonds, resulting in nearly-free spins giving rise to a Curie susceptibility modified by the statistics in the number of contributing spins. As a result the magnetic susceptibility at low $`T`$ diverges asdaniel
$$\chi (T0)1/\left[T(\mathrm{log}T)^2\right].$$
(3)
## III Finite Chains
The above results refer to chains in the thermodynamic limit ($`N\mathrm{}`$). For finite ordered rings (periodic boundary conditions), Bonner and Fisher bonner have calculated the susceptibility per spin $`\chi _N(T)`$ for up to N=11 spins based on direct diagonalization of $``$ in the presence of a magnetic field. They have found that $`\chi _N(T0)`$ exhibits distinct behavior depending on whether $`N`$ is even or odd. In the first case, pairs of neighboring spins tend to form singlets and $`\chi _{N=even}(T0)0`$, whereas in the second case, the occurrence of unpaired spins leads to a Curie-law behavior $`\chi _{N=odd}(T0)1/T\mathrm{}`$. These results immediately raise the question as to what extent the behavior of $`\chi _N(T)`$ changes by the introduction of disorder. Moreover, could we infer, based on the magnetic response, the type of exchange disorder distribution? The relevant distributions here are: (i) trimodal, with
$$P_{tri}(J)=(1/3)\{\delta (JJ_0)+\delta (J(1+W)J_0)+\delta (J(1W)J_0)\},$$
(4)
where $`J_0`$ and $`W`$ are both positive, with $`W<1`$, and (ii) exponential, with
$$P_{exp}(J)=\frac{1}{J_0}e^{J/J_0}\mathrm{\Theta }(J),$$
(5)
where $`\mathrm{\Theta }`$ is the step-function. Note that in both cases $`J=J_0`$, the exchange “target” value.
We consider initially the $`N=2`$ case, for which the susceptibility per spin is given byharaldsen
$$\frac{\chi _2(T,J)}{\chi _0(J)}=\frac{\beta J}{3+e^{\beta J}},$$
(6)
where $`\beta =1/T`$ and $`\chi _0(J)`$ is given above Eq. (2). Considering the average of $`\chi _2`$ over the above distributions, it is clear that $`\chi _2_{tri}=\underset{0}{\overset{\mathrm{}}{}}P_{tri}(J)\chi _2(T,J)𝑑J`$ vanishes as $`T0`$, as in the absence of disorder. For the exponential distribution, it is convenient to split the integral for $`\chi _2_{exp}`$ into two terms
$$\chi _2_{exp}=\frac{\beta }{J_0}_0^{\alpha /\beta }\frac{e^{J/J_0}}{3+e^{\beta J}}𝑑J+\frac{\beta }{J_0}_{\alpha /\beta }^{\mathrm{}}\frac{e^{J/J_0}}{3+e^{\beta J}}𝑑J,$$
(7)
where $`\alpha `$ is a constant, and $`\chi _2_{exp}`$ is given in units of $`(g\mu _B)^2`$. For sufficiently low temperatures, $`\alpha `$ can be chosen such that $`\alpha /\beta J_0`$ and $`e^\alpha <<1`$. Hence, in the first integral (corresponding to small values of $`J`$), we can approximate $`e^{J/J_0}1`$, while in the second one, the term $`e^{\beta J}`$ is always much greater than 1, leading to
$$J_0\chi _2_{exp}C_1+e^\alpha C_1,$$
(8)
where $`C_1=_0^\alpha 1/(3+e^x)𝑑x`$ $`_0^{\mathrm{}}1/(3+e^x)𝑑x=\mathrm{ln}(4)/3`$. Thus, as $`T0`$, $`\chi _2_{exp}`$ approaches a non-zero value, in contrast with the trimodal distribution result.
For an open $`N=3`$ chain, the susceptibility per spin isharaldsen
$$\chi _3(T,J_1,J_2)=(\beta /12)\frac{5+e^{\beta (J_1+J_2)/2}\mathrm{cosh}\left((\beta /2)\sqrt{J_1^2J_1J_2+J_2^2}\right)}{1+e^{\beta (J_1+J_2)/2}\mathrm{cosh}\left((\beta /2)\sqrt{J_1^2J_1J_2+J_2^2}\right)}.$$
(9)
It is clear that for non-negative values of $`J_1`$ and $`J_2`$, $`\chi _3`$ exhibits a Curie-like divergence as $`T0`$. As a consequence, the average of $`\chi _3`$ over the trimodal distribution, $`\chi _3_{tri}`$, also diverges as $`T`$ approaches 0. Regarding the exponential distribution, the change of variables $`J_1=J\mathrm{sin}\theta `$ and $`J_2=J\mathrm{cos}\theta `$ leads to
$$J_0^2\chi _3_{exp}=(\beta /12)_0^{\pi /2}𝑑\theta _0^{\mathrm{}}e^{(\mathrm{sin}\theta +\mathrm{cos}\theta )J/J_0}\frac{5+f(\beta J,\theta )}{1+f(\beta J,\theta )}J𝑑J,$$
(10)
where $`f(x,\theta )=e^{x(\mathrm{sin}\theta +\mathrm{cos}\theta )/2}\mathrm{cosh}\left((x/2)\sqrt{1\mathrm{sin}\theta \mathrm{cos}\theta }\right)`$. As for the $`N=2`$ case, we split the integral above into two, which we label $`I_1`$ and $`I_2`$, corresponding to $`0J\alpha /\beta `$ and $`\alpha /\beta J\mathrm{}`$, respectively. For sufficiently low $`T`$, $`\alpha `$ is again chosen such that $`\alpha /\beta J_0`$ and $`e^\alpha <<1`$. Since $`1\mathrm{sin}\theta +\mathrm{cos}\theta \sqrt{2}`$, the term $`e^{(\mathrm{sin}\theta +\mathrm{cos}\theta )J/J_0}`$ in the first integral (corresponding to small values of $`J`$) can be approximated by 1, giving
$$I_1\frac{1}{\beta ^2}_0^{\pi /2}𝑑\theta _0^\alpha \frac{5+f(x,\theta )}{1+f(x,\theta )}x𝑑x=C_1^{}/\beta ^2,$$
(11)
where $`C_1^{}`$ is $`T`$-independent. In the second integral, since $`\beta J\alpha `$ and $`\alpha `$ is large, the $`T`$-dependent term in both the numerator and denominator of the integrand is much larger than 1, so that
$$I_2_0^{\pi /2}𝑑\theta _{\alpha /\beta }^{\mathrm{}}e^{(\mathrm{sin}\theta +\mathrm{cos}\theta )J/J_0}J𝑑J_0^{\pi /2}𝑑\theta _0^{\mathrm{}}e^{(\mathrm{sin}\theta +\mathrm{cos}\theta )J/J_0}J𝑑J=C_2^{},$$
(12)
where $`C_2^{}`$ is $`T`$-independent. Thus, from Eq. (10), $`J_0\chi _3_{exp}=(\beta /12)(I_1+I_2)/J_0\beta C_2^{}/J_0`$, which diverges following a Curie law as $`T0`$.
For a better physical insight, it is instructive to compare the $`N=2`$ and $`N=3`$ cases. We note that in the low-$`T`$ regime, the $`T`$-dependent behavior of $`\chi _N`$ is dominated in the first case by the disorder distribution , while in the second by the odd parity of $`N`$. The Curie-type low-$`T`$ behavior of $`\chi _3(T)`$ is related to the occurrence of “unpaired” spins, independently of $`P(J)`$. It is interesting that the large-$`J`$ tail of the distribution gives the dominant contribution to the average behavior, namely the integral given by $`I_2`$ in Eq. (12), corresponding to the more strongly coupled 3-spin chains. On the other hand, for $`N=2`$, the behavior of $`\chi _2(T)`$ at low $`T`$ is dominated by the contribution from small values of $`J`$ (more precisely, from values of $`Jmin\{T,J_0\}`$), as shown in Eqs. (7) and (8). The occurrence of arbitrarily small values of $`J`$ in the exponential distribution weakens the tendency of neighboring spins to form singlets,bonner leading to a finite value of $`\chi _2(T=0)_{exp}`$, while $`\chi _2(T=0)_{tri}=0`$.
We expect the above considerations to apply to other finite values of $`N`$. As odd-parity $`N`$ chains always lead to unpaired spins, and thus to a Curie-like low-$`T`$ divergence of $`\chi _{N=odd}(T)`$, susceptibility measurements would not be useful for distinguishing among different exchange disorder distributions if $`N`$ is odd. Since analytical calculations become impractical as $`N`$ increases, in order to illustrate the even-parity case we have carried out numerical calculations for both $`\chi _8_{tri}`$ and $`\chi _8_{exp}`$ as a function of $`T`$. We obtain the susceptibility numerically by determining the spin-spin correlation function $`S_i^zS_j^z`$ from which, through the fluctuation-dissipation theorem, we obtain $`\chi =\beta _{i,j}S_i^zS_j^z`$. For each temperature we have averaged over 10,000 realizations of disordered configurations. For the trimodal distribution, two values for the width parameter have been considered, namely $`W=0.3`$ and $`0.5`$. Results are presented in Fig. 1, which shows susceptibility curves for chains with $`N=2`$ and 8, for trimodal and exponential disorder distributions, as well as for the ordered chain cases. We note that results for $`\chi _N(T)_{tri}`$ are qualitatively very similar, regardless of the width parameter $`W`$, even in the limit $`W=0`$, corresponding to the ordered chains. All curves reach a maximum and eventually decrease as $`T0`$, going to zero for $`T=0`$. We remark the obvious fact, also illustrated in Fig. 1, that as $`W`$ increases the maximum in $`\chi `$ and the sharp downturn toward zero value occur at lower temperatures. The results for the exponential distribution are markedly different, with an increasing $`\chi _N(T)_{exp}`$ for decreasing $`T`$. In principle the trimodal and exponential distributions might be identified and differentiated through low-temperature susceptibility measurements in such even-parity chains.
## IV Crossover regime and Thermodynamic Limit
Further increase in $`N`$ eventually leads to the thermodynamic limit behaviortodo given in Eq. (3). For even $`N`$, approach to such behavior, giving a divergent $`\chi `$ as $`T0`$, might be expected from the comparison between the results of $`\chi _2(T)_{exp}`$ and $`\chi _8(T)_{exp}`$ in Fig. 1, but it is not so clear for the trimodal distributions. Another puzzling point regards the sensitivity to even $`(N_{even})`$ or odd $`(N_{odd})`$ values of $`N`$. For small-$`N_{odd}`$, $`\chi `$ exhibits a Curie-like $`1/T`$ divergence at $`T0`$ independent of the disorder distribution, as discussed for $`N=3`$ in Sec. III, while small-$`N_{even}`$ chains are quite sensitive to $`P(J)`$. But odd- and even-chains results must approach each other, possibly with some sensitivity to the type of disorder (remanent from $`N_{even}`$), as $`N`$ increases.
In order to clarify these points we have calculated the susceptibility for larger values of $`N`$ using a stochastic series expansion (SSE) method. The SSE is a QMC method based on the Taylor expansion of the Boltzmann weight operator $`e^\beta `$ up to a very high order.anders The partition function and observables can then be evaluated via importance sampling of the different terms appearing in this series. Choosing a large enough order for the expansion, the systematic errors introduced by the truncation of the series are negligible. The SSE method also allows us to use importance sampling update schemes based on global changes of the system (cluster or loop updates) that are extremely efficient, especially for the highly symmetric Heisenberg model studied here. We have adjusted the precision of the data obtained through importance sampling so that the errors obtained on each individual realization are roughly one order if magnitude smaller than a typical difference between two realizations. We have considered open chains and averaged our results over 10,000 disorder realizations (except for the large size $`N=128`$ where we used periodic boundary conditions and averaged over only 100 realizations).
We present in Fig. 2 results for $`\chi _N(T)_{tri}`$ and $`\chi _N(T)_{exp}`$ for increasing $`N`$ up to $`N=17`$ for the exponential and trimodal $`(W=0.5)`$ distributions. We identify here the following trends towards the thermodynamic limit: (i) from the frames for $`N_{even}`$ (on the left) we see that the well differentiated disorder distribution results for $`N=8`$ in (a) approach each other as $`N_{even}`$ increases \[(c) and (e)\]; (ii) from the frames for $`N_{odd}`$ (on the right) we do not identify significative changes as the Curie behavior discussed in Sec. III for N=3 acquires logarithmic corrections \[eventually leading to the asymptotic behavior governed by Eq. (3)\] which are not easily captured on this scale and through this range of $`N_{odd}`$ values; (iii) following successive rows we see that the even-odd differences (namely going from $`N_{even}`$ to $`N_{even}+1`$) become less prominent as $`N_{even}`$ increases \[(a)-(b), (c)-(d), (e)-(f)\]. From the practical point of view, it is clear from (i) that, contrary to the $`N=8`$ results, for $`N=16`$ the two distributions lead to very similar qualitative types of behavior of the susceptibility, indicating that, for the particular distributions and temperature range considered here, they may not be clearly differentiated beyond this size of chains, regardless of the parity of $`N`$.
Approach to the thermodynamic limit, given in Eq. (3), is usually investigated by plotting $`(T\chi )^{1/2}`$ vs $`\mathrm{log}T/J_0`$, which leads to a linear behavior at low temperatures in this limit.todo In Fig. 3 we illustrate the crossover regime by presenting such scaled plots for $`N=16`$ and 128 for trimodal $`(W=0.5)`$ and exponential distributions. In all cases, straight lines (dashed for trimodal and solid for exponential) are drawn through the two lowest-$`T`$ calculated points. We note that for the trimodal distribution and $`N=16`$ such line does not include any other calculated point, whereas for $`N=128`$ it gives a good fitting for $`T/J_0`$ up to about $`0.25`$, indicating that the thermodynamic limit has been reached up to this temperature for this value of $`N`$. For the exponential distribution, the approach to the thermodynamic limit scaling with increasing $`N`$ is faster, although careful analysis of the data plotted in Fig. 3 shows that the points for $`N=16`$ and exponential distribution \[triangles in (a)\] down-turn from the solid line at small $`T`$, and the good agreement for $`T/J_0=1`$ in (a) is fortuitous, as it does not remain for $`N=128`$ \[triangles in (b)\]. It is interesting to note that the solid (dashed) line in Fig. 3(a) is nearly parallel to the solid (dashed) line in (b), indicating that some aspects of the $`N\mathrm{}`$ behavior are already captured at the smaller $`N`$ values at low-$`T`$. Fig. 3(b) shows that in the thermodynamic limit the susceptibility is quantitatively quite sensitive to the disorder distribution, although in practice this is probably not as valuable a tool to identify the disorder distribution as the differences encountered for the smaller even values of $`N`$ \[e.g. Fig. 2(a)\].
## V Summary and conclusions
We have investigated the low-$`T`$ behavior of the uniform magnetic susceptibility of exchange disordered spin-1/2 AF chains as a function of the number $`N`$ of spins in the chain and disorder, for trimodal and exponential disorder distributions. Formally, the key distinction among the two distributions considered here is that the exponential distribution is not bound, so that pair exchange coupling $`J`$ arbitrarily close to zero may occur, while the trimodal distribution is bound and $`J`$ does not become arbitrarily small.
For chains with even and relatively small number of spins ($`N_{even}8`$) the susceptibility displays distinct behaviors for the two exchange distributions which are expected to occur in nanochains of P donors in Si. According to our results in Fig. 1, it might be possible to identify the atomic-scale positioning of the P atoms in such chains, thus providing complementary information on whether sample preparation techniques meet the requirements compatible with Kane’s original proposal for a quantum computer hardware. For larger values of $`N`$ ($`N16`$) such differentiation would not be so straightforward, as illustrated in Fig. 2. We also note that $`\chi _N(T)`$ becomes less sensitive to $`N`$ being even or odd when $`N16`$. For completeness, we have also investigated the approach to the thermodynamic limit, and we have found that for $`N=128`$ the expected scaling of $`\chi (T0)`$ is already obtained, while for $`N=16`$ the system is still far from the thermodynamic limit, as discussed in Fig. 3.
We therefore identify three regimes as $`N`$ increases: (i) when $`N8`$ the even-$`N`$ disordered chains present quite distinct behaviors according to the exchange distribution; (ii) for intermediate $`N`$ values, illustrated here by $`N=16`$, the low temperature behavior of the magnetic susceptibility corresponding to different distributions is not easily differentiated, although the thermodynamic limit has not yet been reached; (iii) for larger $`N`$, illustrated here by $`N=128`$, the two distributions follow the thermodynamic limit scaling, which is a signature that the random singlet phase is formed at the lowest $`T`$.
###### Acknowledgements.
We thank S.L.A. de Queiroz, R.R. dos Santos and D.J. Priour for helpful suggestions. BK thanks the hospitality of the CMTC at the University of Maryland. This work has been partially supported in Brazil by CNPq, FAPERJ, the Millennium Institute of Nanoscience and FUJB. RTS acknowledges support from NSF-DMR-0312261 and NSF-INT-0203837. |
warning/0507/physics0507192.html | ar5iv | text | # Contents
## 1 Introduction
The strategy of simplifying a mechanical problem by exploiting a symmetry so as to reduce the number of variables is one of classical mechanics’ grand themes. It is theoretically deep, practically important, and recurrent in the history of the subject. Indeed, it occurs already in 1687, in Newton’s solution of the Kepler problem; (or more generally, the problem of two bodies exerting equal and opposite forces along the line between them). The symmetries are translations and rotations, and the corresponding conserved quantities are the linear and angular momenta.
This paper will expound one central aspect of this large subject. Namely, the relations between continuous symmetries and conserved quantities—in effect, Noether’s “first theorem”: which I expound in both the Lagrangian and Hamiltonian frameworks, though confining myself to finite-dimensional systems. As we shall see, this topic is underpinned by the theorems in elementary Lagrangian and Hamiltonian mechanics about cyclic (ignorable) coordinates and their corresponding conserved momenta. (Again, there is a glorious history: these theorems were of course clear to these subjects’ founders.) Broadly speaking, my discussion will make increasing use, as it proceeds, of the language of modern geometry. It will also emphasise Hamiltonian, rather than Lagrangian, mechanics: apart from mention of the Legendre transformation, the Lagrangian framework drops out wholly after Section 3.4.1.<sup>3</sup><sup>3</sup>3It is worth noting the point, though I shall not exploit it, that symplectic structure can be seen in the classical solution space of the Lagrangian framework; cf. (3) of Section 6.7.
There are several motivations for studying this topic. As regards physics, many of the ideas and results can be generalized to infinite-dimensional classical systems; and in either the original or the generalized form, they underpin developments in quantum theories. The topic also leads into another important subject, the modern theory of symplectic reduction: (for a philosopher’s introduction, cf. Butterfield (2006)). As regards philosophy, the topic is a central focus for the discussion of symmetry, which is both a long-established philosophical field and a currently active one: cf. Brading and Castellani (2003). (Some of the current interest relates to symplectic reduction, whose philosophical significance has been stressed recently, especially by Belot: Butterfield (2006) gives references.)
The plan of the paper is as follows. In Section 2, I review the elements of the Lagrangian framework, emphasising the elementary theorem that cyclic coordinates yield conserved momenta, and introducing the modern geometric language in which mechanics is often cast. Then I review Noether’s theorem in the Lagrangian framework (Section 3). I emphasise how the theorem depends on two others: the elementary theorem about cyclic coordinates, and the local existence and uniqueness of solutions of ordinary differential equations. Then I introduce Hamiltonian mechanics, again emphasising how cyclic coordinates yield conserved momenta; and approaching canonical transformations through the symplectic form (Section 4). This leads to Section 5’s discussion of Poisson brackets; and thereby, of the Hamiltonian version of Noether’s theorem. In particular, we see what it would take to prove that this version is more powerful than (encompasses) the Lagrangian version. By the end of the Section, it only remains to show that a vector field generates a one-parameter family of canonical transformations iff it is a Hamiltonian vector field. It turns out that we can show this without having to develop much of the theory of canonical transformations. We do so in the course of the final Section’s account of the geometric structure of Hamiltonian mechanics, especially the symplectic structure of a cotangent bundle (Section 6). Finally, we end the paper by mentioning a generalized framework for Hamiltonian mechanics which is crucial for symplectic reduction. This framework takes the Poisson bracket, rather than the symplectic form, as the basic notion; with the result that the state-space is, instead of a cotangent bundle, a generalization called a ‘Poisson manifold’.
## 2 Lagrangian mechanics
### 2.1 Lagrange’s equations
We consider a mechanical system with $`n`$ configurational degrees of freedom (for short: n freedoms), described by the usual Lagrange’s equations. These are $`n`$ second-order ordinary differential equations:
$$\frac{d}{dt}(\frac{L}{\dot{q}^i})\frac{L}{q^i}=0,i=1,\mathrm{},n;$$
(2.1)
where the Lagrangian $`L`$ is the difference of the kinetic and potential energies: $`L:=KV`$. (We use $`K`$ for the kinetic energy, not the traditional $`T`$; for in differential geometry, we will use $`T`$ a lot, both for ‘tangent space’ and ‘derivative map’.)
I should emphasise at the outset that several special assumptions are needed in order to deduce eq. 2.1 from Newton’s second law, as applied to the system’s component parts: (assumptions that tend to get forgotten in the geometric formulations that will dominate later Sections!) But I will not go into many details about this, since:
(i): there is no single set of assumptions of mimimum logical strength (nor a single “best package-deal” combining simplicity and mimimum logical strength);
(ii): full discussions are available in many textbooks (or, from a philosophical viewpoint, in Butterfield 2004a: Section 3).
I will just indicate a simple and commonly used sufficient set of assumptions. But owing to (i) and (ii), the details here will not be cited in later Sections.
Note first that if the system consists of $`N`$ point-particles (or bodies small enough to be treated as point-particles), so that a configuration is fixed by $`3N`$ cartesian coordinates, we may yet have $`n<3N`$. For the system may be subject to constraints and we will require the $`q^i`$ to be independently variable. More specifically, let us assume that any constraints on the system are holonomic; i.e. each is expressible as an equation $`f(r^1,\mathrm{},r^m)=0`$ among the coordinates $`r^k`$ of the system’s component parts; (here the $`r^k`$ could be the $`3N`$ cartesian coordinates of $`N`$ point-particles, in which case $`m:=3N`$). A set of $`c`$ such constraints can in principle be solved, defining a $`(mc)`$-dimensional hypersurface $`Q`$ in the $`m`$-dimensional space of the $`r`$s; so that on the configuration space $`Q`$ we can define $`n:=mc`$ independent coordinates $`q^i,i=1,\mathrm{},n`$.
Let us also assume that any constraints on the system are: (i) scleronomous, i.e. independent of time, so that $`Q`$ is identified once and for all; (ii) ideal, i.e. the forces that maintain the constraints would do no work in any possible displacement consistent with the constraints and applied forces (a ‘virtual displacement’). Let us also assume that the forces applied to the system are monogenic: i.e. the total work $`\delta w`$ done in an infinitesimal virtual displacement is integrable; its integral is the work function $`U`$. (The term ‘monogenic’ is due to Lanczos (1986, p. 30), but followed by others e.g. Goldstein et al. (2002, p. 34).) And let us assume that the system is conservative: i.e. the work function $`U`$ is independent of both the time and the generalized velocities $`\dot{q}_i`$, and depends only on the $`q^i`$: $`U=U(q^1,\mathrm{},q^n)`$.
So to sum up: let us assume that the constraints are holonomic, scleronomous and ideal, and that the system is monogenic with a velocity-independent work-function. Now let us define $`K`$ to be the kinetic energy; i.e. in cartesian coordinates, with $`k`$ now labelling particles, $`K:=\mathrm{\Sigma }_k\frac{1}{2}m_k𝐯_k^2`$. Let us also define $`V:=U`$ to be the potential energy, and set $`L:=KV`$. Then the above assumptions imply eq. 2.1.<sup>4</sup><sup>4</sup>4Though I shall not develop any details, there is of course a rich theory about these and related assumptions. One example, chosen with an eye to our later use of geometry, is that assuming scleronomous constraints, $`K`$ is readily shown to be a homogeneous quadratic form in the generalized velocities, i.e. of the form $`K=\mathrm{\Sigma }_{i,j}^na_{ij}\dot{q}^i\dot{q}^j`$; and so $`K`$ defines a metric on the configuration space.
To solve mechanical problems, we need to integrate Lagrange’s equations. Recall the idea from elementary calculus that $`n`$ second-order ordinary differential equations have a (locally) unique solution, once we are given $`2n`$ arbitrary constants. Broadly speaking, this idea holds good for Lagrange’s equations; and the $`2n`$ arbitrary constants can be given just as one would expect—as the initial configuration and generalized velocities $`q^i(t_0),\dot{q}^i(t_0)`$ at time $`t_0`$. More precisely: expanding the time derivatives in eq. 2.1, we get
$$\frac{^2L}{\dot{q}^j\dot{q}^i}\ddot{q}^j=\frac{^2L}{q^j\dot{q}^i}\dot{q}^j\frac{^2L}{t\dot{q}^i}+\frac{L}{\dot{q}^i}$$
(2.2)
so that the condition for being able to solve these equations to find the accelerations at some initial time $`t_0`$, $`\ddot{q}^i(t_0)`$, in terms of $`q^i(t_0),\dot{q}^i(t_0)`$ is that the Hessian matrix $`\frac{^2L}{\dot{q}^i\dot{q}^j}`$ be nonsingular. Writing the determinant as $``$, and partial derivatives as subscripts, the condition is that:
$$\frac{^2L}{\dot{q}^j\dot{q}^i}L_{\dot{q}^j\dot{q}^i}\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0}.$$
(2.3)
This Hessian condition holds in very many mechanical problems; and henceforth, we assume it. (If it fails, we enter the territory of constrained dynamics; for which cf. e.g. Henneaux and Teitelboim (1992, Chapters 1-5).) It underpins most of what follows: for it is needed to define the Legendre transformation, by which we pass from Lagrangian to Hamiltonian mechanics.
Of course, even with eq. 2.3, it is still in general hard in practice to solve for the $`\ddot{q}^i(t_0)`$: they are buried in the lhs of eq. 2.2. In (5) of Section 2.2.2, this will motivate the move to Hamiltonian mechanics.<sup>5</sup><sup>5</sup>5This is not to say that Hamiltonian mechanics makes all problems “explicitly soluble”: if only! For a philosophical discussion of the various meanings of ‘explicit solution’, cf. Butterfield (2004a: Section 2.1).
Given eq. 2.3, and so the accelerations at the initial time $`t_0`$, the basic theorem on the (local) existence and uniqueness of solutions of ordinary differential equations can be applied. (We will state this theorem in Section 3.4 in connection with Noether’s theorem.)
By way of indicating the rich theory that can be built from eq. 2.1 and 2.3, I mention one main aspect: the power of variational formulations. Eq. 2.1 are the Euler-Lagrange equations for the variational problem $`\delta L𝑑t=0`$; i.e. they are necessary and sufficient for the action integral $`I=L𝑑t`$ to be stationary. But variational principles will play no further role in this paper; (Butterfield 2004 is a philosophical discussion).
But our main concern, here and throughout this paper, is how symmetries yield conserved quantities, and thereby reduce the number of variables that need to be considered in solving a problem. In fact, we are already in a position to prove Noether’s theorem, to the effect that any (continuous) symmetry of the Lagrangian $`L`$ yields a conserved quantity. But we postpone this to Section 3, until we have developed some more notions, especially geometric ones.
We begin with the idea of generalized momenta, and the result that the generalized momentum of any cyclic coordinate is a constant of the motion: though very simple, this result is the basis of Noether’s theorem. Elementary examples prompt the definition of the generalized, or canonical, momentum, $`p_i`$, conjugate to a coordinate $`q^i`$ as: $`\frac{L}{\dot{q}^i}`$; (this was first done by Poisson in 1809). Note that $`p_i`$ need not have the dimensions of momentum: it will not if $`q^i`$ does not have the dimension length. So Lagrange’s equations can be written:
$$\frac{d}{dt}p_i=\frac{L}{q^i};$$
(2.4)
We say a coordinate $`q^i`$ is cyclic if $`L`$ does not depend on $`q^i`$. (The term comes from the example of an angular coordinate of a particle subject to a central force. Another term is: ignorable.) Then the Lagrange equation for a cyclic coordinate, $`q^n`$ say, becomes $`\dot{p}_n=0`$, implying
$$p_n=\text{constant, }c_n\text{ say}.$$
(2.5)
So: the generalized momentum conjugate to a cyclic coordinate is a constant of the motion.
It is straightforward to show that this simple result encompasses the elementary theorems of the conservation of momentum, angular momentum and energy: this last corresponding to time’s being a cyclic coordinate. As a simple example, consider the angular momentum of a free particle. The Lagrangian is, in spherical polar coordinates,
$$L=\frac{1}{2}m(\dot{r}^2+r^2\dot{\theta }^2+r^2\dot{\varphi }^2\mathrm{sin}^2\theta )$$
(2.6)
so that $`L/\varphi =0`$. So the conjugate momentum
$$\frac{L}{\dot{\varphi }}=mr^2\dot{\varphi }\mathrm{sin}^2\theta ,$$
(2.7)
which is the angular momentum about the $`z`$-axis, is conserved.
### 2.2 Geometrical perspective
#### 2.2.1 Some restrictions of scope
I turn to give a brief description of the elements of Lagrangian mechanics in terms of modern differential geometry. Here ‘brief’ indicates that:
(i): I will assume without explanation various geometric notions, in particular: manifold, vector, 1-form (covector), metric, Lie derivative and tangent bundle.
(ii): I will disregard issues about degrees of smoothness: all manifolds, scalars, vectors etc. will be assumed to be as smooth as needed for the context.
(iii): I will also simplify by speaking “globally, not locally”. I will speak as if the scalars, vector fields etc. are defined on a whole manifold; when in fact all that we can claim in application to most systems is a corresponding local statement—because for example, differential equations are guaranteed the existence and uniqueness only of a local solution.<sup>6</sup><sup>6</sup>6A note for afficionados. Of the three main pillars of elementary differential geometry—the implicit function theorem, the local existence and uniqueness of solutions of ordinary differential equations, and Frobenius’ theorem—this paper will use the first only implicitly (!), and the second explicitly in Sections 3 and 4. The third will not be used.
We begin by assuming that the configuration space (i.e. the constraint surface) $`Q`$ is a manifold. The physical state of the system, taken as a pair of configuration and generalized velocities, is represented by a point in the tangent bundle $`TQ`$ (also known as ‘velocity phase space’). That is, writing $`T_x`$ for the tangent space at $`xQ`$, $`TQ`$ has points $`(x,\tau ),xQ,\tau T_x`$. We will of course often work with the natural coordinate systems on $`TQ`$ induced by coordinate systems $`q`$ on $`Q`$; i.e. with the $`2n`$ coordinates $`(q,\dot{q})(q^i,\dot{q}^i)`$.
The main idea of the geometric perspective is that this tangent bundle is the arena for Lagrangian mechanics. So various previous notions and results are now expressed in terms of the tangent bundle. In particular, the Lagrangian is a scalar function $`L:TQ\mathrm{I}\mathrm{R}`$ which “determines everything”. And the conservation of the generalized momentum $`p_n`$ conjugate to a cyclic coordinate $`q_n`$, $`p_np_n(q,\dot{q})=c_n`$, means that the motion of the system is confined to a level set $`p_n^1(c_n)`$: where this level set is a $`(2n1)`$-dimensional sub-manifold of $`TQ`$.
But I must admit at the outset that working with $`TQ`$ involves limiting our discussion to (a) time-independent Lagrangians and (b) time-independent coordinate transformations.
(a): Recall Section 2.1’s assumptions that secured eq. 2.1. Velocity-dependent potentials and-or rheonomous constraints would prompt one to use what is often called the ‘extended configuration space’ $`Q\times \mathrm{I}\mathrm{R}`$, and-or the ‘extended velocity phase space’ $`TQ\times \mathrm{I}\mathrm{R}`$.
(b): So would time-dependent coordinate transformations. This is a considerable limitation from a philosophical viewpoint, since it excludes boosts, which are central to the philosophical discussion of spacetime symmetry groups, and especially of relativity principles. To give the simplest example: the Lagrangian of a free particle is just its kinetic energy, which can be made zero by transforming to the particle’s rest frame; i.e. it is not invariant under boosts.
#### 2.2.2 The tangent bundle
With these limitations admitted, we now describe Lagrangian mechanics on $`TQ`$, in five extended comments.
(1): $`2n`$ first-order equations; the Hessian again:—
The Lagrangian equations of motion are now $`2n`$ first-order equations for the functions $`q^i(t),\dot{q}^i(t)`$, falling in to two groups:
(a) the $`n`$ equations eq. 2.2, with the $`\ddot{q}^i`$ taken as the time derivatives of $`\dot{q}^i`$ with respect to $`t`$; i.e. we envisage using the Hessian condition eq. 2.3 to solve eq. 2.2 for the $`\ddot{q}^i`$, hard though this usually is to do in practice;
(b) the $`n`$ equations $`\dot{q}^i=\frac{dq^i}{dt}`$.
(2): Vector fields and solutions:—
(a): These $`2n`$ first-order equations are equivalent to a vector field on $`TQ`$: the ‘dynamical vector field’, or for short the ‘dynamics’. I write it as $`D`$ (to distinguish it from the generic vector field $`X,Y,\mathrm{}`$).
(b): In the natural coordinates $`(q^i,\dot{q}^i)`$, the vector field $`D`$ is expressed as
$$D=\dot{q}^i\frac{}{q^i}+\ddot{q}^i\frac{}{\dot{q}^i};$$
(2.8)
and the rate of change of any dynamical variable $`f`$, taken as a scalar function on $`TQ`$, $`f(q,\dot{q})\mathrm{I}\mathrm{R}`$ is given by
$$\frac{df}{dt}=\dot{q}^i\frac{f}{q^i}+\ddot{q}^i\frac{f}{\dot{q}^i}=D(f).$$
(2.9)
(c): So the Lagrangian $`L`$ determines the dynamical vector field $`D`$, and so (for given initial $`q,\dot{q}`$) a (locally unique) solution: an integral curve of $`D`$, $`2n`$ functions of time $`q(t),\dot{q}(t)`$ (with the first $`n`$ functions determining the latter). This separation of solutions/trajectories within $`TQ`$ is important for the visual and qualitative understanding of solutions.
(3): Canonical momenta are 1-forms:—
Any point transformation, or any coordinate transformation $`(q^i)(q^i)`$, in the configuration manifold $`Q`$, induces a basis-change in the tangent space $`T_q`$ at $`qQ`$. Consider any vector $`\tau T_q`$ with components $`\dot{q}^i`$ in coordinate system $`(q^i)`$ on $`Q`$, i.e. $`\tau =\frac{d}{dt}=\dot{q}^i\frac{}{q^i}`$; (think of a motion through configuration $`q`$ with generalized velocity $`\tau `$). Its components $`\dot{q}^{}_{}{}^{}i`$ in the coordinate system $`(q^i)`$ (i.e. $`\tau =\dot{q}^{}_{}{}^{}i\frac{}{q^{}_{}{}^{}i}`$) are given by applying the chain rule to $`q^{}_{}{}^{}i=q^{}_{}{}^{}i(q^k)`$:
$$\dot{q}^{}_{}{}^{}i\frac{q^i}{q^k}\dot{q}^k.$$
(2.10)
so that we can “drop the dots”:
$$\frac{\dot{q}^{}_{}{}^{}i}{\dot{q}^j}=\frac{q^{}_{}{}^{}i}{q^j}.$$
(2.11)
One easily checks, using eq. 2.11, that for any $`L`$, the canonical momenta $`p_i:=\frac{L}{\dot{q}^i}`$ form a 1-form on $`Q`$, transforming under $`(q^i)(q^i)`$ by:
$$p_i^{}:=\frac{L^{}}{\dot{q}^{}_{}{}^{}i}=\frac{q^k}{q^{}_{}{}^{}i}\frac{L}{\dot{q}^k}\frac{q^k}{q^{}_{}{}^{}i}p_k$$
(2.12)
That is, the canonical momenta defined by $`L`$ form a 1-form field on $`Q`$. (We will later describe this as a cross-section of the cotangent bundle.)
(4): Geometric formulation of Lagrange’s equations:—
We can formulate Lagrange’s equations in a coordinate-independent way, by using three ingredients, namely:
(i): $`L`$ itself (a scalar, so coordinate-independent);
(ii): the vector field $`D`$ that $`L`$ defines; and
(iii): the 1-form on $`TQ`$ defined locally, in terms of the natural coordinates $`(q^i,\dot{q}^i)`$, by
$$\theta _L:=\frac{L}{\dot{q}^i}dq^i.$$
(2.13)
(So the coefficients of $`\theta _L`$ for the other $`n`$ elements of the dual basis, the $`d\dot{q}^i`$ are defined to be zero.) This 1-form is called the canonical 1-form. We shall see that it plays a role in Noether’s theorem, and is centre-stage in Hamiltonian mechanics.
We combine these three ingredients using the idea of the Lie derivative of a 1-form along a vector field.
We will write the Lie derivative of $`\theta _L`$ along the vector field $`D`$ on $`TQ`$, as $`_D\theta _L`$. (It is sometimes written as $`L`$; but we need the symbol $`L`$ for the Lagrangian—and later on, for left translation.) By the Leibniz rule, $`_D\theta _L`$ is:
$$_D\theta _L=(_D\frac{L}{\dot{q}^i})dq^i+\frac{L}{\dot{q}^i}_D(dq^i).$$
(2.14)
But the Lie derivative of any scalar function $`f:TQ\mathrm{I}\mathrm{R}`$ along any vector field $`X`$ is just $`X(f)`$; and for the dynamical vector field $`D`$, this is just $`\dot{f}=\frac{f}{q^i}\dot{q}^i+\frac{f}{\dot{q}^i}\ddot{q}^i`$. So we have
$$_D\theta _L=(\frac{d}{dt}\frac{L}{\dot{q}^i})dq^i+\frac{L}{\dot{q}^i}d\dot{q}^i.$$
(2.15)
Rewriting the first term by the Lagrange equations, we get
$$_D\theta _L=(\frac{L}{q^i})dq^i+\frac{L}{\dot{q}^i}d\dot{q}^idL.$$
(2.16)
We can conversely deduce the familiar Lagrange equations from eq. 2.16, by taking coordinates. So we conclude that these equations’ coordinate-independent form is:
$$_D\theta _L=dL.$$
(2.17)
(5): Towards the Hamiltonian framework:—
Finally, a comment about the Lagrangian framework’s limitations as regards solving problems, and how they prompt the transition to Hamiltonian mechanics.
Recall the remark at the end of Section 2.1, that the $`n`$ equations eq. 2.2 are in general hard to solve for the $`\ddot{q}^i(t_0)`$: they lie buried in the left hand side of eq. 2.2. On the other hand, the $`n`$ equations $`\dot{q}^i=\frac{dq^i}{dt}`$ (the second group of $`n`$ equations in (1) above) are as simple as can be.
This makes it natural to seek another $`2n`$-dimensional space of variables, $`\xi ^\alpha `$ say ($`\alpha =1,\mathrm{},2n`$), in which:
(i): a motion is described by first-order equations, so that we have the same advantage as in $`TQ`$ that a unique trajectory passes through each point of the space; but in which
(ii): all $`2n`$ equations have the simple form $`\frac{d\xi ^\alpha }{dt}=f_\alpha (\xi ^1,\mathrm{}\xi ^{2n})`$ for some set of functions $`f_\alpha (\alpha =1,\mathrm{},2n)`$.
Indeed, Hamiltonian mechanics provides exactly such a space: it is usually the cotangent bundle of the configuration manifold, instead of its tangent bundle. But before turning to that, we expound Noether’s theorem in the current Lagrangian framework.
## 3 Noether’s theorem in Lagrangian mechanics
### 3.1 Preamble: a modest plan
Any discussion of symmetry in Lagrangian mechanics must include a treatment of “Noether’s theorem”. The scare quotes are to indicate that there is more than one Noether’s theorem. Quite apart from Noether’s work in other branches of mathematics, her paper (1918) on symmetries and conserved quantities in Lagrangian theories has several theorems. I will be concerned only with applying her first theorem to finite-dimensional systems. In short: it provides, for any continuous symmetry of a system’s Lagrangian, a conserved quantity called the ‘momentum conjugate to the symmetry’.
I stress at the outset that the great majority of subsequent applications and commentaries (also for her other theorems, besides her first) are concerned with versions of the theorems for infinite (i.e. continuous) systems. In fact, the context of Noether’s investigation was contemporary debate about how to understand conservation principles and symmetries in the “ultimate classical continuous system”, viz. gravitating matter as described by Einstein’s general relativity. This theory can be given a Lagrangian formulation: that is, the equations of motion, i.e. Einstein’s field equations, can be deduced from a Hamilton’s Principle with an appropriate Lagrangian. The contemporary debate was especially about the conservation of energy and the principle of general covariance (also known as: diffeomorphism invariance). General covariance prompts one to consider how a variational principle transforms under spacetime coordinate transformations that are arbitrary, in the sense of varying from point to point. This leads to the idea of “local” symmetries, which since Noether’s time has been immensely fruitful in both classical and quantum physics, and in both a Lagrangian and Hamiltonian framework.<sup>7</sup><sup>7</sup>7Cf. Brading and Castellani (2003). Apart from papers specifically about Noether’s theorem, this anthology’s papers by Wallace, Belot and Earman (all 2003) are closest to this paper’s concerns.
So I agree that from the perspective of Noether’s work, and its enormous later development, this Section’s application of the first theorem to finite-dimensional systems is, as they say, “trivial”. Furthermore, this application is easily understood, without having to adopt that perspective, or even having to consider infinite systems. In other words: its statement and proof are natural, and simple, enough that the nineteenth century masters of mechanics, like Hamilton, Jacobi and Poincaré, would certainly recognize it in their own work—allowing of course for adjustments to modern language. In fact, versions of it for the Galilei group of Newtonian mechanics and the Lorentz group of special relativity were published a few years before Noether’s paper; (Brading and Brown (2003, p. 90); for details, cf. Kastrup (1987)).<sup>8</sup><sup>8</sup>8Here again, ‘versions of it’ needs scare-quotes. For in what follows, I shall be more limited than these proofs, in two ways. (1): I limit myself, as I did in Section 2.2.1, both to time-independent Lagrangians and to time-independent transformations: so my discussion does not encompass boosts. (2): I will take a symmetry of $`L`$ to require that $`L`$ be the very same; whereas some treatments allow the addition to $`L`$ of the time-derivative of a function $`G(q)`$ of the coordinates $`q`$—since such a time-derivative makes no difference to the Lagrange equations.
Nevertheless, it is worth expounding the finite-system version of Noether’s first theorem. For:
(i): It generalizes Section 2.1’s result about cyclic coordinates, and thereby the elementary theorems of the conservation of momentum, angular momentum and energy which that result encompasses. The main generalization is that the theorem does not assume we have identified a cyclic coordinate. But on the other hand: every symmetry in the Noether sense will arise from a cyclic coordinate in some system $`q`$ of generalized coordinates. (As we will see, this follows from the local existence and uniqueness of solutions of ordinary differential equations.)
(ii): This exposition will also prepare the way for our discussion of symmetry and conserved quantities in Hamiltonian mechanics.<sup>9</sup><sup>9</sup>9Other expositions of Noether’s theorem for finite-dimensional Lagrangian mechanics include: Arnold (1989: 88-89), Desloge (1982: 581-586), Lanczos (1986: 401-405: emphasizing the variational perspective) and Johns (2005: Chapter 13). Butterfield (2004a, Section 4.7) is a more detailed version of this Section. Beware: though many textbooks of Hamiltonian mechanics cover the Hamiltonian version of Noether’s theorem (which, as we will see, is stronger), they often do not label it as such; and if they do label it, they often do not relate it clearly to the Lagrangian version.
In this exposition, I will also discuss en passant the distinction between:
(i) the notion of symmetry at work in Noether’s theorem, i.e. a symmetry of $`L`$, often called a variational symmetry; and
(ii) the notion of a symmetry of the set of solutions of a differential equation: often called a dynamical symmetry. This notion applies to all sorts of differential equations, and systems of them; not just to those with the form of Lagrange’s equations (i.e. derivable from an variational principle). In short, this sort of symmetry is a map that sends any solution of the given equation(s) (in effect: a dynamically possible history of the system—a curve in the state-space) to some other solution. Finding such symmetries, and groups of them, is a central part of the modern theory of integration of differential equations (both ordinary and partial).
Broadly speaking, this notion is more general than that of a symmetry of $`L`$. Not only does it apply to many other sorts of differential equation. Also, for Lagrange’s equations: a symmetry of $`L`$ is (with one caveat) a symmetry of the solutions, i.e. a dynamical symmetry—but the converse is false.<sup>10</sup><sup>10</sup>10An excellent account of this modern integration theory, covering both ordinary and partial differential equations, is given by Olver (2000). He also covers the Lagrangian case (Chapter 5 onwards), and gives many historical details especially about Lie’s pioneering contributions.
In this Section, the plan is as follows. We define:
(i): a (continuous) symmetry as a vector field (on the configuration manifold $`Q`$) that generates a family of transformations under which the Lagrangian is invariant; (Section 3.2);
(ii): the momentum conjugate to a vector field, as (roughly) the rate of change of the Lagrangian with respect to the $`\dot{q}`$s in the direction of the vector field; (Section 3.3).
These two definitions lead directly to Noether’s theorem (Section 3.4): after all the stage-setting, the proof will be a one-liner application of Lagrange’s equations.
### 3.2 Vector fields and symmetries—variational and dynamical
I need to expound three topics:
(1): the idea of a vector field on the configuration manifold $`Q`$; and how to lift it to $`TQ`$;
(2): the definition of a variational symmetry;
(3): the contrast between (2) and the idea of dynamical symmetry.
Note that, as in previous Sections, I will often speak, for simplicity, “globally, not locally”, i.e. as if the relevant scalar functions, vector fields etc. are defined on all of $`Q`$ or $`TQ`$. Of course, they need not be.
#### 3.2.1 Vector fields on $`TQ`$; lifting fields from $`Q`$ to $`TQ`$
We recall first that a differentiable vector field on $`Q`$ is represented in a coordinate system $`q=(q^1,\mathrm{},q^n)`$ by $`n`$ first-order ordinary differential equations
$$\frac{dq^i}{dϵ}=f^i(q^1,\mathrm{},q^n).$$
(3.18)
A vector field generates a one-parameter family of active transformations: viz. passage along the vector field’s integral curves, by a varying parameter-difference $`ϵ`$. The vector field is called the infinitesimal generator of the family. It is common to write the parameter as $`\tau `$, but in this Section we use $`ϵ`$ to avoid confusion with $`t`$, which often represents the time.
Similarly, a vector field defined on $`TQ`$ corresponds to a system of $`2n`$ ordinary differential equations, and generates an active transformation of $`TQ`$. But I will consider only vector fields on $`TQ`$ that mesh with the structure of $`TQ`$ as a tangent bundle, in the sense that they are induced by vector fields on $`Q`$, in the following natural way.
This induction has two ingredient ideas.
First, any curve in $`Q`$ (representing a possible state of motion) defines a corresponding curve in $`TQ`$, because the functions $`q^i(t)`$ define the functions $`\dot{q}^i(t)`$. (Here $`t`$ is the parameter of the curve.) More formally: given any curve in configuration space, $`\varphi :I\mathrm{I}\mathrm{R}Q`$, with coordinate expression in the $`q`$-system $`tIq(\varphi (t))q(t)=q^i(t)`$, we define its extension to $`TQ`$ to be the curve $`\mathrm{\Phi }:I\mathrm{I}\mathrm{R}TQ`$ given in the corresponding coordinates by $`q^i(t),\dot{q}^i(t)`$.
Second, any vector field $`X`$ on $`Q`$ generates displacements in any possible state of motion, represented by a curve in $`Q`$ with coordinate expression $`q^i=q^i(t)`$. Namely: for a given value of the parameter $`ϵ`$, the displaced state of motion is represented by the curve in Q
$$q^i(t)+ϵX^i(q^i(t)).$$
(3.19)
Putting these ingredients together: we first displace a curve within $`Q`$, and then extend the result to $`TQ`$. Namely, the extension to $`TQ`$ of the (curve representing) the displaced state of motion is given by the $`2n`$ functions, in two groups each of $`n`$ functions, for the $`(q,\dot{q})`$ coordinate system
$$q^i(t)+ϵX^i(q^i(t))\text{ and }\dot{q}^i(t)+ϵY^i(q^i(t),\dot{q}^i);$$
(3.20)
where $`Y`$ is defined to be the vector field on $`TQ`$ that is the derivative along the original state of motion of $`X`$. That is:
$$Y^i(q,\dot{q}):=\frac{dX^i}{dt}=\mathrm{\Sigma }_j\frac{X^i}{q^j}\dot{q}^j.$$
(3.21)
Thus displacements by a vector field within $`Q`$ are lifted to $`TQ`$. The vector field $`X`$ on $`Q`$ lifts to $`TQ`$ as $`(X,\frac{dX}{dt})`$; i.e. it lifts to the vector field that sends a point $`(q^i,\dot{q}^i)TQ`$ to $`(q^i+ϵX^i,\dot{q}^i+ϵ\frac{dX^i}{dt})`$.<sup>11</sup><sup>11</sup>11I have discussed this in terms of some system $`(q,\dot{q})`$ of coordinates. But the definitions of extensions and displacements are in fact coordinate-independent. Besides, one can show that the operations of displacing a curve within $`Q`$, and extending it to $`TQ`$, commute to first order in $`ϵ`$: the result is the same for either order of the operations.
#### 3.2.2 The definition of variational symmetry
To define variational symmetry, I begin with the integral notion and then give the differential notion. The idea is that the Lagrangian $`L`$, a scalar $`L:TQ\mathrm{I}\mathrm{R}`$, should be invariant under all the elements of a one-parameter family of active transformations $`\theta _ϵ:ϵI\mathrm{I}\mathrm{R}`$: at least in a neighbourhood of the identity map corresponding to $`ϵ=0`$, $`\theta _0id_U`$. (Here $`U`$ is some open subset of $`TQ`$, maybe not all of it.)
That is, we define the family $`\theta _ϵ:ϵI\mathrm{I}\mathrm{R}`$ to be a variational symmetry of $`L`$ if $`L`$ is invariant under the transformations: $`L=L\theta _ϵ`$, at least around $`ϵ=0`$. (We could use the correspondence between active and passive transformations to recast this definition, and what follows, in terms of a passive notion of symmetry as sameness of $`L`$’s functional form in different coordinate systems. I leave this as an exercise! Or cf. Butterfield (2004a: Section 4.7.2).)
For the differential notion of variational symmetry, we of course use the idea of a vector field. But we also impose Section 3.2.1’s restriction to vector fields on $`TQ`$ that are induced by vector fields on $`Q`$. So we define a vector field $`X`$ on $`Q`$ that generates a family of active transformations $`\theta _ϵ`$ on $`TQ`$ to be a variational symmetry of $`L`$ if the first derivative of $`L`$ with respect to $`ϵ`$ is zero, at least around $`ϵ=0`$. More precisely: writing
$$L\theta _ϵ=L(q^i+ϵX^i,\dot{q}^i+ϵY^i)\text{ with }Y^i=\mathrm{\Sigma }_j\frac{X^i}{q^j}\dot{q}^j,$$
(3.22)
we say $`X`$ is a variational symmetry iff the first derivative of $`L`$ with respect to $`ϵ`$ is zero (at least around $`ϵ=0`$). That is: $`X`$ is a variational symmetry iff
$$\mathrm{\Sigma }_iX^i\frac{L}{q^i}+\mathrm{\Sigma }_iY^i\frac{L}{\dot{q}^i}=0\text{ with }Y^i=\mathrm{\Sigma }_j\frac{X^i}{q^j}\dot{q}^j.$$
(3.23)
#### 3.2.3 A contrast with dynamical symmetries
The general notion of a dynamical symmetry, i.e. a symmetry of some equations of motion (whether Euler-Lagrange or not), is not needed for Section 3.4’s presentation of Noether’s theorem. But the notion is so important that I must mention it, though only to contrast it with variational symmetries.
The general definition is roughly as follows. Given any system of differential equations, $``$ say, a dynamical symmetry of the system is an active transformation $`\zeta `$ on the system $``$’s space of both independent variables, $`x_j`$ say, and dependent variables $`y^i`$ say, such that any solution of $``$, $`y^i=f^i(x_j)`$ say, is carried to another solution. For a precise definition, cf. Olver (2000: Def. 2.23, p. 93), and his ensuing discussion of the induced action (called ‘prolongation’) of the transformation $`\zeta `$ on the spaces of (in general, partial) derivatives of the $`y`$’s with respect to the $`x`$s (i.e. jet spaces).
As I said in Section 3.1, groups of symmetries in this sense play a central role in the modern theory of differential equations: not just in finding new solutions, once given a solution, but also in integrating the equations. For some main theorems stating criteria (in terms of prolongations) for groups of symmetries, cf. Olver (2000: Theorem 2.27, p. 100, Theorem 2.36, p. 110, Theorem 2.71, p. 161).
But for present purposes, it is enough to state the rough idea of a one-parameter group of dynamical symmetries (without details about prolongations!) for Lagrange’s equations in the familiar form, eq. 2.1.
In this simple case, there is just one independent variable $`x:=t`$, so that:
(a): we are considering ordinary, not partial, differential equations, with $`n`$ dependent variables $`y^i:=q^i(t)`$.
(b): prolongations correspond to lifts of maps on $`Q`$ to maps on $`TQ`$; cf. Section 3.2.1.
Furthermore, in line with the discussion following Lagrange’s equations eq. 2.1, the time-independence of the Lagrangian (time being a cyclic coordinate) means we can define dynamical symmetries $`\zeta `$ in terms of active transformations on the tangent bundle, $`\theta :TQTQ`$, that are lifted from active transformations on $`Q`$. In effect, we define such a map $`\zeta `$ by just adjoining to any such $`\theta :TQTQ`$ the identity map on the time variable $`id:t\mathrm{I}\mathrm{R}t`$. (More formally: $`\zeta :(q,\dot{q},t)TQ\times \mathrm{I}\mathrm{R}(\theta (q,\dot{q}),t)TQ\times \mathrm{I}\mathrm{R}`$.)
Then we define in the usual way what it is for a one-parameter family of such maps $`\zeta _s:sI\mathrm{I}\mathrm{R}`$ to be a (local) one-parameter group of dynamical symmetries (for Lagrange’s equations eq. 2.1): namely, if any solution curve $`q(t)`$ (equivalently: its extension $`q(t),\dot{q}(t)`$ to $`TQ`$) of the Lagrange equations is carried by each $`\zeta _s`$ to another solution curve, with the $`\zeta _s`$ for different $`s`$ composing in the obvious way, for $`s`$ close enough to $`0I`$.
And finally: we also define (in a manner corresponding to the discussion at the end of Section 3.2.2) a differential, as against integral, notion of dynamical symmetry. Namely, we say a vector field $`X`$ on $`Q`$ is a dynamical symmetry if its lift to $`TQ`$ (more precisely: its lift, with the identity map on the time variable adjoined) is the infinitesimal generator of such a one-parameter family $`\zeta _s`$.
For us, the important point is that this notion of a dynamical symmetry is different from Section 3.2.2’s notion of a variational symmetry.<sup>12</sup><sup>12</sup>12Since the Lagrangian $`L`$ is especially associated with variational principles, while the dynamics is given by equations of motion, calling Section 3.2.2’s notion ‘variational symmetry’, and this notion ‘dynamical symmetry’ is a good and widespread usage. But beware: it is not universal. As I announced in Section 3.1, a variational symmetry is (with one caveat) a dynamical symmetry—but the converse is false. Fortunately, the same simple example will serve both to show the subtlety about the first implication, and as a counterexample to the converse implication. This example is the two-dimensional harmonic oscillator.<sup>13</sup><sup>13</sup>13All the material to the end of this Subsection is drawn from Brown and Holland (2004a); cf. also their (2004). The present use of the harmonic oscillator example also occurs in Morandi et al (1990: 203-204).
The usual Lagrangian is, with cartesian coordinates written as $`q`$s, and the contravariant indices written for clarity as subscripts:
$$L_1=\frac{1}{2}\left[\dot{q}_{1}^{}{}_{}{}^{2}+\dot{q}_{2}^{}{}_{}{}^{2}\omega ^2(q_1^2+q_2^2)\right];$$
(3.24)
giving as Lagrange equations:
$$\ddot{q}_i+\omega ^2q_i=0,i=1,2.$$
(3.25)
But these Lagrange equations, i.e. the same dynamics, are also given by
$$L_2=\dot{q}_1\dot{q}_2\omega ^2q_1q_2.$$
(3.26)
The rotations in the plane are of course a variational symmetry of $`L_1`$, and a dynamical symmetry of eq. 3.25. But they are not a variational symmetry of $`L_2`$. So a dynamical symmetry need not be a variational one. Besides, these equations contain another example to the same effect. Namely, the “squeeze” transformations
$$q_1^{}:=e^\eta q_1,q_2^{}:=e^\eta q_2$$
(3.27)
are a dynamical symmetry of eq. 3.25, but not a variational symmetry of $`L_1`$. So again: a dynamical symmetry need not be a variational one.<sup>14</sup><sup>14</sup>14In the light of this, you might ask about a more restricted implication: viz. must every dynamical symmetry of a set of equations of motion be a variational symmetry of some or other Lagrangian that yields the given equations as the Euler-Lagrange equations of Hamilton’s Principle? Again, the answer is No for the simple reason that there are many (sets of) equations of motion that are not Euler-Lagrange equations of any Lagrangian, and yet have dynamical symmetries. Wigner (1954) gives an example. The general question of under what conditions is a set of ordinary differential equations the Euler-Lagrange equations of some Hamilton’s Principle is the inverse problem of Lagrangian mechanics. It is a large subject with a long history; cf. e.g. Santilli (1979), Lopuszanski (1999).
I turn to the first implication: that every variational symmetry is a dynamical symmetry. This is true: general and abstract proofs (applying also to continuous systems i.e. field theories) can be found in Olver (2000: theorem 4.14, p. 255; theorem 4.34, p. 278; theorem 5.53, p. 332).
But beware of a condition of the theorem. (This is the caveat mentioned at the end of Section 3.1.) The theorem requires that all the variables $`q`$ (for continuous systems: all the fields $`\varphi `$) be subject to Hamilton’s Principle. The need for this condition is shown by rotations in the plane, which are a variational symmetry of the familiar Lagrangian $`L_1`$ above. But it is easy to show that such a rotation is a dynamical symmetry of one of the Lagrange equations, say the equation for the variable $`q_1`$
$$\ddot{q}_1+\omega ^2q_1=0,$$
(3.28)
only if the corresponding Lagrange equation holds for $`q_2`$.
### 3.3 The conjugate momentum of a vector field
Now we define the momentum conjugate to a vector field $`X`$ to be the scalar function on $`TQ`$:
$$p_X:TQ\mathrm{I}\mathrm{R};p_X=\mathrm{\Sigma }_iX^i\frac{L}{\dot{q}^i}$$
(3.29)
(For a time-dependent Lagrangian, $`p_X`$ would be a scalar function on $`TQ\times \mathrm{I}\mathrm{R}`$, with $`\mathrm{I}\mathrm{R}`$ representing time.)
We shall see in the next Subsection’s examples that this definition generalizes in an appropriate way Section 2.1’s definition of the momentum conjugate to a coordinate $`q`$.
But first note that it is an improvement in the sense that, while the momentum conjugate to a coordinate $`q`$ depends on the choice made for the other coordinates, the momentum $`p_X`$ conjugate to a vector field $`X`$ is independent of the coordinates chosen. Though this point is not needed in order to prove Noether’s theorem, here is the proof.
We first apply the chain-rule to $`L=L(q^{}(q),\dot{q}^{}(q,\dot{q}))`$ and eq. 2.11 (“cancellation of the dots”), to get
$$\frac{L}{\dot{q}^i}=\mathrm{\Sigma }_j\frac{L}{\dot{q}^j}\frac{\dot{q}^j}{\dot{q}^i}=\mathrm{\Sigma }_j\frac{L}{\dot{q}^j}\frac{q^j}{q^i}.$$
(3.30)
Then using the transformation law for components of a vector field
$$X^i=\mathrm{\Sigma }_j\frac{q^i}{q^j}X^j.$$
(3.31)
and relabelling $`i`$ and $`j`$, we deduce:
$$p_X^{}=\mathrm{\Sigma }_iX^i\frac{L}{\dot{q}^i}=\mathrm{\Sigma }_{ij}X^j\frac{q^i}{q^j}\frac{L}{\dot{q}^i}=\mathrm{\Sigma }_{ij}X^i\frac{q^j}{q^i}\frac{L}{\dot{q}^j}=\mathrm{\Sigma }_iX^i\frac{L}{\dot{q}^i}p_X.$$
(3.32)
Finally, I remark incidentally that in the geometric formulation of Lagrangian mechanics (Section 2.2) , the coordinate-independence of $`p_X`$ becomes, unsurprisingly, a triviality. Namely: $`p_X`$ is obviously the contraction of $`X`$ as lifted to $`TQ`$ with the canonical 1-form on $`TQ`$ that we defined in eq. 2.13:
$$\theta _L:=\frac{L}{\dot{q}^i}dq^i.$$
(3.33)
We will return to this at the end of Section 3.4.1.
### 3.4 Noether’s theorem; and examples
Given just the definition of conjugate momentum, eq. 3.29, the proof of Noether’s theorem is immediate. (The interpretation and properties of this momentum, discussed in the last Subsection, are not needed.) The theorem says:
> Noether’s theorem for Lagrangian mechanics If $`X`$ is a (variational) symmetry of a system with Lagrangian $`L(q,\dot{q},t)`$, then $`X`$’s conjugate momentum is a constant of the motion.
Proof: We just calculate the derivative of the momentum eq. 3.29 along the solution curves in $`TQ`$, and apply Lagrange’s equations and the definitions of $`Y^i`$, and of symmetry eq. 3.23:
$`{\displaystyle \frac{dp}{dt}}=\mathrm{\Sigma }_i{\displaystyle \frac{dX^i}{dt}}{\displaystyle \frac{L}{\dot{q}^i}}+\mathrm{\Sigma }_iX^i{\displaystyle \frac{d}{dt}}\left({\displaystyle \frac{L}{\dot{q}^i}}\right)`$ (3.34)
$`=\mathrm{\Sigma }_iY^i{\displaystyle \frac{L}{\dot{q}^i}}+\mathrm{\Sigma }_iX^i{\displaystyle \frac{L}{q^i}}=0.\mathrm{QED}.`$
Examples:— This proof, though neat, is a bit abstract! So here are two examples, both of which return us to examples we have already seen.
(1): The first example is a shift in a cyclic coordinate $`q^n`$: i.e. the case with which our discussion of Noether’s theorem began at the end of Section 2.1. So suppose $`q^n`$ is cyclic, and define a vector field $`X`$ by
$$X^1=0,\mathrm{},X^{n1}=0,X^n=1.$$
(3.35)
So the displacements generated by $`X`$ are translations by an amount $`ϵ`$ in the $`q^n`$-direction. Then $`Y^i:=\frac{dX^i}{dt}`$ vanishes, and the definition of (variational) symmetry eq. 3.23 reduces to
$$\frac{L}{q^n}=0.$$
(3.36)
So since $`q^n`$ is assumed to be cyclic, $`X`$ is a symmetry. And the momentum conjugate to $`X`$, which Noether’s theorem tells us is a constant of the motion, is the familiar one:
$$p_X:=\mathrm{\Sigma }_iX^i\frac{L}{\dot{q}^i}=\frac{L}{\dot{q}^n}.$$
(3.37)
As mentioned in Section 3.1, this example is universal, in that every symmetry $`X`$ arises, around any point where $`X`$ is non-zero, from a cyclic coordinate in some local system of coordinates. This follows from the basic theorem about the local existence and uniqueness of solutions of ordinary differential equations. We can state the theorem as follows; (cf. e.g. Arnold (1973: 48-49, 77-78, 249-250), Olver (2000: Prop 1.29)).
Consider a system of $`n`$ first-order ordinary differential equations on an open subset $`U`$ of an $`n`$-dimensional manifold
$$\dot{q^i}=X^i(q)X^i(q^1,\mathrm{},q^n),qU;$$
(3.38)
equivalently, a vector field $`X`$ on $`U`$. Let $`q_0`$ be a non-singular point of the vector field, i.e. $`X(q_0)0`$. Then in a sufficiently small neighbourhood $`V`$ of $`q_0`$, there is a coordinate system (formally, a diffeomorphism $`f:VW\mathrm{I}\mathrm{R}^n`$) such that, writing $`y_i:\mathrm{I}\mathrm{R}^n\mathrm{I}\mathrm{R}`$ for the standard coordinates on $`W`$ and $`𝐞_i`$ for the $`i`$th standard basic vector of $`\mathrm{I}\mathrm{R}^n`$, eq. 3.38 goes into the very simple form
$$\dot{𝐲}=𝐞_n;\mathrm{i}.\mathrm{e}.\dot{y}_n=1,\dot{y}_1=\dot{y}_2=\mathrm{}=\dot{y}_{n1}=0\mathrm{in}W.$$
(3.39)
(In terms of the tangent map (also known as: push-forward) $`f_{}`$ on tangent vectors that is induced by $`f`$: $`f_{}(X)=𝐞_n`$ in $`W`$.) On account of eq. 3.39’s simple form, Arnold suggests the theorem might well be called the ‘rectification theorem’.
We should note two points about the theorem:
(i): The rectifying coordinate system $`f`$ may of course be very hard to find. So the theorem by no means makes all problems “trivially soluble”; cf. again footnote 4.
(ii): The theorem has an immediate corollary about local constants of the motion. Namely: $`n`$ first-order ordinary differential equations have, locally, $`n1`$ functionally independent constants of the motion (also known as: first integrals). They are given, in the above notation, by $`y_1,\mathrm{},y_{n1}`$.
We now apply the rectification theorem, so as to reverse the reasoning in the above example of $`q^n`$ cyclic. That is: assuming $`X`$ is a symmetry, let us rectify it—i.e. let us pass to a coordinate system $`(q)`$ such that eq. 3.35 holds. Then, as above, $`Y^i:=\frac{dX^i}{dt}`$ vanishes; and $`X`$’s being a (variational) symmetry, eq. 3.23, reduces to $`q^n`$ being cyclic; and the momentum conjugate to $`X`$, $`p_X`$ reduces to the familiar conjugate momentum $`p_n=\frac{L}{\dot{q}^n}`$. Thus every symmetry $`X`$ arises locally from a cyclic coordinate $`q^n`$ and the corresponding conserved momentum is $`p_n`$. (But note that this may hold only “very locally”: the domain $`V`$ of the coordinate system $`f`$ in which $`X`$ generates displacements in the direction of the cyclic coordinate $`q^n`$ can be smaller than the set $`U`$ on which $`X`$ is a symmetry.)
In Section 5.3, the fact that every symmetry arises locally from a cyclic coordinate will be important for understanding the Hamiltonian version of Noether’s theorem.
(2): Let us now look at our previous example, the angular momentum of a free particle (eq. 2.6), in the cartesian coordinate system, i.e. a coordinate system without cyclic coordinates. So let $`q_1:=x,q_2:=y,q_3:=z`$. (In this example, subscripts will again be a bit clearer.) Then a small rotation about the $`x`$-axis
$$\delta x=0,\delta y=ϵz,\delta z=ϵy$$
(3.40)
corresponds to a vector field $`X`$ with components
$$X_1=0,X_2=q_3,X_3=q_2$$
(3.41)
so that the $`Y_i`$ are
$$Y_1=0,Y_2=\dot{q}_3,X_3=\dot{q}_2.$$
(3.42)
For the Lagrangian
$$L=\frac{1}{2}m(\dot{q}_1^2+\dot{q}_2^2+\dot{q}_3^2)$$
(3.43)
$`X`$ is a (variational) symmetry since the definition of symmetry eq. 3.23 now reduces to
$$\mathrm{\Sigma }_iX_i\frac{L}{q_i}+\mathrm{\Sigma }_iY_i\frac{L}{\dot{q}_i}=\dot{q}_3\frac{L}{\dot{q}_2}+\dot{q}_2\frac{L}{\dot{q}_3}=0.$$
(3.44)
So Noether’s theorem then tells us that $`X`$’s conjugate momentum is
$$p_X:=\mathrm{\Sigma }_iX_i\frac{L}{\dot{q}_i}=X_2\frac{L}{\dot{q}_2}+X_3\frac{L}{\dot{q}_3}=mz\dot{y}+my\dot{z}$$
(3.45)
which is indeed the $`x`$-component of angular momentum.
#### 3.4.1 A geometrical formulation
We can give a geometric formulation of Noether’s theorem by using the vanishing of the Lie derivative to express constancy along the integral curves of a vector field. There are two vector fields on $`TQ`$ to consider: the dynamical vector field $`D`$ (cf. eq. 2.8), and the lift to $`TQ`$ of the vector field $`X`$ that is the variational symmetry.
I will now write $`\overline{X}`$ for this lift. So given the vector field $`X`$ on $`Q`$
$$X=X^i(q)\frac{}{q^i},$$
(3.46)
the lift $`\overline{X}`$ of $`X`$ to $`TQ`$ is, by eq. 3.21,
$$\overline{X}=X^i(q)\frac{}{q^i}+\frac{X^i(q)}{q^j}\dot{q}^j\frac{}{\dot{q}^i},$$
(3.47)
where the $`q`$ argument of $`X^i`$ emphasises that the $`X^i`$ do not depend on $`\dot{q}`$.
That $`X`$ is a variational symmetry means that in $`TQ`$, the Lie derivative of $`L`$ along the lift $`\overline{X}`$ vanishes: $`_{\overline{X}}L=0`$. On the other hand, we know from eq. 3.33 that the momentum $`p_X`$ conjugate to $`X`$ is the contraction $`<;>`$ of $`\overline{X}`$ with the canonical 1-form $`\theta _L:=\frac{L}{\dot{q}^i}dq^i`$ on $`TQ`$:
$$p_X:=X^i\frac{L}{\dot{q}^i}<\overline{X};\theta _L>.$$
(3.48)
So Noether’s theorem says:
If $`_{\overline{X}}L=0`$, then $`_D<\overline{X};\theta _L>=0`$.
Note finally that eq. 3.48 shows that the theorem has no converse. That is: given that a dynamical variable $`p:TQ\mathrm{I}\mathrm{R}`$ is a constant of the motion, $`_Dp=0`$, there is no single vector field $`\overline{X}`$ on $`TQ`$ such that $`p=<\overline{X};\theta _L>`$. For given such a $`\overline{X}`$, one could get another by adding any field $`\overline{Y}`$ for which $`<\overline{Y};\theta _L>=0`$. However, we will see in Section 5.2 that in Hamiltonian mechanics a constant of the motion does determine a corresponding vector field on the state space.
## 4 Hamiltonian mechanics introduced
### 4.1 Preamble
From now on this paper adopts the Hamiltonian framework. As we shall see, its description of symmetry and conserved quantities is in various ways more straightforward and powerful than that of the Lagrangian framework.
The main idea is to replace the $`\dot{q}`$s by the canonical momenta, the $`p`$s. More generally, the state-space is no longer the tangent bundle $`TQ`$ but a phase space $`\mathrm{\Gamma }`$, which we take to be the cotangent bundle $`T^{}Q`$. (Here, the phrase ‘we take to be’ just signals the fact that eventually, in Section 6.8, we will glimpse a more general kind of Hamiltonian state-space, viz. Poisson manifolds.)
Admittedly, the theory on $`TQ`$ given by Lagrange’s equations eq. 2.1 is equivalent to the Hamiltonian theory on $`T^{}Q`$ given by eq. 4.53 below, once we assume the Hessian condition eq. 2.3.
But of course, theories can be formally equivalent, but different as regards their power for solving problems, their heuristic value and even their interpretation. In our case, two advantages of Hamiltonian mechanics over Lagrangian mechanics are commonly emphasised. (i): The first concerns its greater power or flexibility for describing a given system, that Lagrangian methods can also describe (and so its greater power for solving problems about such a system). (ii): The second concerns the broader idea of describing other systems. In more detail:—
(i): Hamiltonian mechanics replaces the group of point transformations, $`qq^{}`$ on $`Q`$, together with their lifts to $`TQ`$, by a “corresponding larger” group of transformations on $`\mathrm{\Gamma }`$, the group of canonical transformations (also known as, for the standard case where $`\mathrm{\Gamma }=T^{}Q`$: the symplectic group).
This group “corresponds” to the point transformations (and their lifts) in that while for any Lagrangian $`L`$, Lagrange’s equations eq. 2.1 are covariant under all the point transformations, Hamilton’s equations eq. 4.53 below are (for any Hamiltonian $`H`$) covariant under all canonical transformations. And it is a “larger” group because:
(a) any point transformation together with its lift to $`TQ`$ is a canonical transformation: (more precisely: it naturally defines a canonical transformation on $`T^{}Q`$);
(b) not every canonical transformation is thus induced by a point transformation; for a canonical transformation can “mix” the $`q`$s and $`p`$s in a way that point transformations and their lifts cannot.
There is a rich and multi-faceted theory of canonical transformations, to which there are three main approaches—generating functions, integral invariants and symplectic geometry. I will adopt the symplectic approach, but not need many details about it. In particular, we will need only a few details about how the “larger” group of canonical transformations makes for a more powerful version of Noether’s theorem.
(ii): The Hamiltonian framework connects analytical mechanics with other fields of physics, especially statistical mechanics and optics. The first connection goes via canonical transformations, especially using the integral invariants approach. The second connection goes via Hamilton-Jacobi theory; (for a philosopher’s exposition, with an eye on quantum theory, cf. Butterfield (2004b: especially Sections 7-9)).<sup>15</sup><sup>15</sup>15Of course, some aspects of Hamiltonian mechanics illustrate both (i) and (ii). For example, Liouville’s theorem on the preservation of phase space volume illustrates both (i)’s integral invariants approach to canonical transformations and (ii)’s connection to statistical mechanics.
With its theme of symmetry and conservation, this paper will illustrate (i), greater power in describing a given system, rather than (ii), describing other systems. As to (i), we will see two main ways in which the Hamiltonian framework is more powerful than the Lagrangian one. First, cyclic coordinates will “do more work for us” (Section 4.2). Second, the Hamiltonian version of Noether’s theorem is both: more powerful, thanks to the use of the “larger” group of canonical transformations; and more easily proven, thanks to the use of Poisson brackets (Section 5).
So from now on, the broad plan is as follows. After Section 4.2’s deduction of Hamilton’s equations, Section 4.3 introduces symplectic structure, starting from the “naive” form of the symplectic matrix. Section 5 presents Poisson brackets, and the Hamiltonian version of Noether’s theorem. Finally, Section 6 gives a geometric perspective, corresponding to Section 2.2’s geometric perspective on the Lagrangian framework.
### 4.2 Hamilton’s equations
#### 4.2.1 The equations introduced
Recall the vision in (5) of Section 2.2.2: that we seek $`2n`$ new variables, $`\xi ^\alpha `$ say, $`\alpha =1,\mathrm{},2n`$ in which Lagrange’s equations take the simple form
$$\frac{d\xi ^\alpha }{dt}=f_\alpha (\xi ^1,\mathrm{}\xi ^{2n}).$$
(4.49)
We can find the desired variables $`\xi ^\alpha `$ by using the canonical momenta
$$p_i:=\frac{L}{\dot{q}^i}=:L_{\dot{q}^i},$$
(4.50)
to write the $`2n`$ Lagrange equations as
$$\frac{dp_i}{dt}=\frac{L}{dq^i};\frac{dq^i}{dt}=\dot{q}^i.$$
(4.51)
These are of the desired simple form, except that the right hand sides need to be written as functions of $`(q,p,t)`$ rather than $`(q,\dot{q},t)`$. (Here and in the next two paragraphs, we temporarily allow time-dependence, since the deduction is unaffected: the time variable is “carried along unaffected”. In the terms of Section 2.1, this means allowing non-scleronomous constraints and a time-dependent work-function $`U`$.)
For the second group of $`n`$ equations, this is in principle straightforward, given our assumption of a non-zero Hessian, eq. 2.3. This implies that we can invert eq. 4.50 so as to get the $`n`$ $`\dot{q}^i`$ as functions of $`(q,p,t)`$. We can then apply this to the first group of equations; i.e. we substitute $`\dot{q}^i(q,p,t)`$ wherever $`\dot{q}^i`$ appears in any right hand side $`\frac{L}{dq^i}`$.
But we need to be careful: the partial derivative of $`L(q,\dot{q},t)`$ with respect to $`q^i`$ is not the same as the partial derivative of $`\widehat{L}(q,p,t):=L(q,\dot{q}(q,p,t),t)`$ with respect to $`q^i`$, since the first holds fixed the $`\dot{q}`$s, while the second holds fixed the $`p`$s. A comparison of these partial derivatives leads, with algebra, to the result that if we define the Hamiltonian function by
$$H(q,p,t):=p_i\dot{q}^i(q,p,t)\widehat{L}(q,p,t)$$
(4.52)
then the $`2n`$ equations eq. 4.51 go over to Hamilton’s equations
$$\frac{dp_i}{dt}=\frac{H}{q^i};\frac{dq^i}{dt}=\frac{H}{p_i}.$$
(4.53)
So we have cast our $`2n`$ equations in the simple form, $`\frac{d\xi ^\alpha }{dt}=f_\alpha (\xi ^1,\mathrm{}\xi ^{2n})`$, requested in (5) of Section 2.2. More explicitly: defining
$$\xi ^\alpha =q^\alpha ,\alpha =1,\mathrm{},n;\xi ^\alpha =p_{\alpha n},\alpha =n+1,\mathrm{},2n$$
(4.54)
Hamilton’s equations become
$$\dot{\xi }^\alpha =\frac{H}{\xi ^{\alpha +n}},\alpha =1,\mathrm{},n;\dot{\xi }^\alpha =\frac{H}{\xi ^{\alpha n}},\alpha =n+1,\mathrm{},2n.$$
(4.55)
To sum up: a single function $`H`$ determines, through its partial derivatives, the evolution of all the $`q`$s and $`p`$s—and so, the evolution of the state of the system.
#### 4.2.2 Cyclic coordinates in the Hamiltonian framework
Just from the form of Hamilton’s equations, we can immediately see a result that is significant for our theme of how symmetries and conserved quantities reduce the number of variables involved in a problem. In short, we can see that with Hamilton’s equations in hand, cyclic coordinates will “do more work for us” than they do in the Lagrangian framework.
More specifically, recall the basic Lagrangian result from the end of Section 2.1, that the generalized momentum $`p_n:=\frac{L}{\dot{q}^n}`$ is conserved if, indeed iff, its conjugate coordinate $`q^n`$ is cyclic, $`\frac{L}{q^n}=0`$. And recall from Section 3.4 that this result underpinned Noether’s theorem in the precise sense of being “universal” for it. Corresponding results hold in the Hamiltonian framework—but are in certain ways more powerful.
Thus we first observe that the transformation “from the $`\dot{q}`$s to the $`p`$s”, i.e. the transition between Lagrangian and Hamiltonian frameworks, does not involve the dependence on the $`q`$s. More precisely: partially differentiating eq. 4.52 with respect to $`q^n`$, we obtain
$$\frac{H}{q^n}\frac{H}{q^n}_{p;q^i,in}=\frac{L}{q^n}\frac{L}{q^n}_{\dot{q};q^i,in}.$$
(4.56)
(The other two terms are plus and minus $`p_i\frac{\dot{q}^i}{q^n}`$, and so cancel.) So a coordinate $`q^n`$ that is cyclic in the Lagrangian sense is also cyclic in the obvious Hamiltonian sense, viz. that $`\frac{H}{q^n}=0`$. But by Hamilton’s equations, this is equivalent to $`\dot{p}_n=0`$. So we have the result corresponding to the Lagrangian one: $`p_n`$ is conserved iff $`q_n`$ is cyclic (in the Hamiltonian sense).
We will see in Section 5.3 that this result underpins the Hamiltonian version of Noether’s theorem; just as the corresponding Lagrangian result underpinned the Lagrangian version of Noether’s theorem (cf. discussion after eq. 3.37).
But we can already see that this result gives the Hamiltonian formalism an advantage over the Lagrangian. In the latter, the generalized velocity corresponding to a cyclic coordinate, $`q_n`$ will in general still occur in the Lagrangian. The Lagrangian will be $`L(q_1,\mathrm{},q_{n1},\dot{q}_1,\mathrm{},\dot{q}_n,t)`$, so that we still face a problem in $`n`$ variables.
But in the Hamiltonian formalism, $`p_n`$ will be a constant of the motion, $`\alpha `$ say, so that the Hamiltonian will be $`H(q_1,\mathrm{},q_{n1},p_1,\mathrm{},p_{n1},\alpha ,t)`$. So we now face a problem in $`n1`$ variables, $`\alpha `$ being simply determined by the initial conditions. That is: after solving the problem in $`n1`$ variables, $`q_n`$ is determined just by quadrature: i.e. just by integrating (perhaps numerically) the equation
$$\dot{q}_n=\frac{H}{\alpha },$$
(4.57)
where, thanks to having solved the problem in $`n1`$ variables, the right-hand side is now an explicit function of $`t`$.
This result is very simple. But it is an important illustration of the power of the Hamiltonian framework. Indeed, Arnold remarks (1989: 68) that ‘almost all the solved problems in mechanics have been solved by means of’ it!
No doubt his point is, at least in part, that this result underpins the Hamiltonian version of Noether’s theorem. But I should add that the result also motivates the study of various notions related to the idea of cyclic coordinates, such as constants of the motion being in involution (i.e. having zero Poisson bracket with each other), and a system being completely integrable (in the sense of Liouville). These notions have played a large part in the way that Hamiltonian mechanics has developed, especially in its theory of canonical transformations. And they play a large part in the way Hamiltonian mechanics has solved countless problems. But as announced in Section 4.1, this paper will not go into these aspects of Hamiltonian mechanics, since they are not needed for our theme of symmetry and conservation; (for a philosophical discussion of these aspects, cf. Butterfield 2005).
#### 4.2.3 The Legendre transformation and variational principles
To end this Subsection, I note two aspects of this transition from Lagrange’s equations to Hamilton’s. For, although I shall not need details about them, they each lead to a rich theory:
(i): The transformation “from the $`\dot{q}`$s to the $`p`$s” is the Legendre transformation. It has a striking geometric interpretation. In the simplest case, it concerns the fact that one can describe a smooth convex real function $`y=f(x),f^{\prime \prime }(x)>0`$, not by the pairs of its arguments and values $`(x,y)`$, but by the pairs of its gradients at points $`(x,y)`$ and the intercepts of its tangent lines with the $`y`$-axis. Given the non-zero Hessian (eq. 2.3), one readily proves various results: e.g. that the geometric interpretation extends to higher dimensions, and that the transformation is self-inverse, i.e. its square is the identity. For details, cf. e.g.: Arnold (1989: Chapters 3.14, 9.45.C), Courant and Hilbert (1953: Chapter IV.9.3; 1962, Chapter I.6), José and Saletan (1998: 212-217), Lanczos (1986: Chapter VI.1-4). The Legendre transformation is also described using modern geometry’s idea of a fibre derivative; as we will see briefly in Section 6.7.
(ii): The transition to Hamilton’s equations has achieved more than we initially sought with our eq. 4.49. Namely: all the $`f_\alpha `$, all the right hand sides in Hamilton’s equations, are up to a sign, partial derivatives of a single function $`H`$. In the Hamiltonian framework, it is precisely this feature that underpins the possibility of expressing the equations of motion by variational principles; (of course, the Lagrangian framework has a corresponding feature). But as I mentioned, this paper does not discuss variational principles; for details cf. e.g. Lanczos (1986: Chapter VI.4) and Butterfield (2004: especially Section 5.2).
To sum up this introduction to Hamilton’s equations:— Even once we set aside (i) and (ii), these equations mark the beginning of a rich and multi-faceted theory. At the centre lies the $`2n`$-dimensional phase space $`\mathrm{\Gamma }`$ coordinatized by the $`q`$s and $`p`$s: or more precisely, as we shall see later, the cotangent bundle $`T^{}Q`$. The structure of Hamiltonian mechanics is encoded in the structure of $`\mathrm{\Gamma }`$, and thereby in the coordinate transformations on $`\mathrm{\Gamma }`$ that preserve this structure, especially the form of Hamilton’s equations: the canonical transformations. As I mentioned in Section 4.1, these transformations can be studied from three main perspectives: generating functions, integral invariants and symplectic structure—but I shall only need the last.
### 4.3 Symplectic forms on vector spaces
I shall introduce symplectic structure by giving Hamilton’s equations a yet more symmetric appearance. This will lead to some elementary ideas about area in $`\mathrm{I}\mathrm{R}^m`$ and symplectic forms on vector spaces: ideas which will later be “made local” by taking the relevant copy of $`\mathrm{I}\mathrm{R}^m`$ to be the tangent space at a point of a manifold. (As usually formulated, Hamiltonian mechanics is especially concerned with the case $`m=2n`$.)
#### 4.3.1 Time-evolution from the gradient of $`H`$
Writing $`\mathrm{𝟏}`$ and $`\mathrm{𝟎}`$ for the $`n\times n`$ identity and zero matrices respectively, we define the $`2n\times 2n`$ symplectic matrix $`\omega `$ by
$$\omega :=\left(\begin{array}{cc}\mathrm{𝟎}& \mathrm{𝟏}\\ \mathrm{𝟏}& \mathrm{𝟎}\end{array}\right).$$
(4.58)
$`\omega `$ is antisymmetric, and has the properties, writing $`\stackrel{~}{}`$ for the transpose of a matrix, that
$$\stackrel{~}{\omega }=\omega =\omega ^1\text{ so that }\omega ^2=\mathrm{𝟏}\text{ ; }\text{ also }\mathrm{det}\omega =1.$$
(4.59)
Using $`\omega `$, Hamilton’s equations eq. 4.55 get the more symmetric form, in matrix notation
$$\dot{\xi }=\omega \frac{H}{\xi }.$$
(4.60)
In terms of components, writing $`\omega ^{\alpha \beta }`$ for the matrix elements of $`\omega `$, and defining $`_\alpha :=/\xi ^\alpha `$, eq. 4.55 become
$$\dot{\xi }^\alpha =\omega ^{\alpha \beta }_\beta H.$$
(4.61)
Eq. 4.60 and 4.61 show how $`\omega `$ forms, from the naive gradient (column vector) $`H`$ of $`H`$ on the phase space $`\mathrm{\Gamma }`$ of $`q`$s and $`p`$s, the vector field on $`\mathrm{\Gamma }`$ that gives the system’s evolution: the Hamiltonian vector field, often written $`X_H`$. At a point $`z=(q,p)\mathrm{\Gamma }`$, eq. 4.60 can be written
$$X_H(z)=\omega H(z).$$
(4.62)
The vector field $`X_H`$ is also written as $`D`$ (for ‘dynamics’), on analogy with the Lagrangian framework’s vector field $`D`$ of eq. 2.8 in Section 2.2.
In Section 6, we will see how this definition of a vector field from a gradient, i.e. a covector or 1-form field, arises from $`\mathrm{\Gamma }`$’s being a cotangent bundle. More precisely, we will see that any cotangent bundle has an intrinsic symplectic structure that provides, at each point of the base-manifold, a natural i.e. basis-independent isomorphism between the tangent space and the cotangent space. For the moment, we:
(i) note a geometric interpretation of $`\omega `$ in terms of area (Section 4.3.2); and then
(ii) generalize the above discussion of $`\omega `$ into the definition of a symplectic form for a fixed vector space (Section 4.3.3).
#### 4.3.2 Interpretation in terms of areas
Let us begin with the simplest possible case: $`\mathrm{I}\mathrm{R}^2(q,p)`$, representing the phase space of a particle constrained to one spatial dimension. Here, the $`2\times 2`$ matrix
$$\omega :=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)$$
(4.63)
defines the antisymmetric bilinear form on $`\mathrm{I}\mathrm{R}^2`$:
$$A:((q^1,p_1),(q^2,p_2))\mathrm{I}\mathrm{R}^2\times \mathrm{I}\mathrm{R}^2q^1p_2q^2p_1\mathrm{I}\mathrm{R}$$
(4.64)
since
$$q^1p_2q^2p_1=\left(\begin{array}{cc}q^1& p_1\end{array}\right)\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\left(\begin{array}{c}q^2\\ p_2\end{array}\right)=\mathrm{det}\left(\begin{array}{cc}q^1& q^2\\ p_1& p_2\end{array}\right).$$
(4.65)
It is easy to prove that $`A((q^1,p_1),(q^2,p_2))q^1p_2q^2p_1`$ is the signed area of the parallelogram spanned by $`(q^1,p_1),(q^2,p_2)`$, where the sign is positive (negative) if the shortest rotation from $`(q^1,p_1)`$ to $`(q^2,p_2)`$ is anti-clockwise (clockwise).
Similarly in $`\mathrm{I}\mathrm{R}^{2n}`$: the matrix $`\omega `$ of eq. 4.58 defines an antisymmetric bilinear form on $`\mathrm{I}\mathrm{R}^{2n}`$ whose value on a pair $`(q,p)(q^1,\mathrm{}q^n;p_1,\mathrm{},p_n),(q^{},p^{})(q^1,\mathrm{}q^n;p_1^{},\mathrm{},p_n^{})`$ is the sum of the signed areas of the $`n`$ parallelograms formed by the projections of the vectors $`(q,p),(q^{},p^{})`$ onto the $`n`$ pairs of coordinate planes labelled $`1,\mathrm{},n`$. That is to say, the value is:
$$\mathrm{\Sigma }_{i=1}^nq^ip_i^{}q^ip_i.$$
(4.66)
This induction of bilinear forms from antisymmetric matrices can be generalized: there is a one-to-one correspondence between forms and matrices. In more detail: there is a one-to-one correspondence between antisymmetric bilinear forms on $`\mathrm{I}\mathrm{R}^2`$ and antisymmetric $`2\times 2`$ matrices. It is easy to check that any such form, $`\omega `$ say, is given, for any basis $`v,w`$ of $`\mathrm{I}\mathrm{R}^2`$, by the matrix $`\left(\begin{array}{cc}0& \omega (v,w)\\ \omega (v,w)& 0\end{array}\right)`$. Similarly for any integer $`n`$: one easily shows that there is a one-to-one correspondence between antisymmetric bilinear forms on $`\mathrm{I}\mathrm{R}^n`$ and antisymmetric $`n\times n`$ matrices. (In Hamiltonian mechanics as usually formulated, we consider the case where $`n`$ is even and the matrix is non-singular, as in eq. 4.58. But when one generalizes to Poisson manifolds (cf. Section 6.8) one allows $`n`$ to be odd, and the matrix to be singular.)
This geometric interpretation of $`\omega `$ is important for two reasons.
(i): The first reason is that the idea of an antisymmetric bilinear form on a copy of $`\mathrm{I}\mathrm{R}^{2n}`$ is the main part of the definition of a symplectic form, which is the central notion in the usual geometric formulation of Hamiltonian mechanics. More details in Section 4.3.3, for a fixed copy of $`\mathrm{I}\mathrm{R}^{2n}`$; and in Section 6, where the form is defined on many copies of $`\mathrm{I}\mathrm{R}^{2n}`$, each copy being the tangent space at a point in the cotangent bundle $`T^{}Q`$.
(ii): The second reason is that the idea of (signed) area underpins the theory of forms (1-forms, 2-forms etc.): i.e. antisymmetric multilinear functions on products of copies of $`\mathrm{I}\mathrm{R}^n`$. And when these copies of $`\mathrm{I}\mathrm{R}^n`$ are copies of the tangent space at (one and the same) point in a manifold, these forms lead to the whole theory of integration on manifolds. One needs this theory in order to make rigorous sense of any integration on a manifold beyond the most elementary (i.e. line-integrals); so it is crucial for almost any mathematical or physical theory using manifolds. In particular, it is crucial for Hamiltonian mechanics. So no wonder the maestro says that ‘Hamiltonian mechanics cannot be understood without differential forms’ (Arnold 1989, p. 163).
However, it turns out that this paper will not need many details about forms and the theory of integration. This is essentially because we focus only on solving mechanical problems, and simplifying them by appeals to symmetry. This means we will focus on line-integrals: viz. integrating with respect to time the equations of motion; or equivalently, integrating the dynamical vector field on the state space. We have already seen this vector field as $`X_H`$ in eq. 4.62; and we will see it again, for example in terms of Poisson brackets (eq. 5.103), and in geometric terms (Section 6). But throughout, the main idea will be as suggested by eq. 4.62: the vector field is determined by the symplectic matrix, “at” each point in the manifold $`\mathrm{\Gamma }`$, acting on the gradient of the Hamiltonian function $`H`$.
So in short: focussing on line-integrals enables us to side-step most of the theory of forms.<sup>16</sup><sup>16</sup>16But forms are essential for understanding integration over surfaces of dimension two or more: which one needs for the integral invariants approach to Hamiltonian mechanics, and its deep connection with Stokes’ theorem.
#### 4.3.3 Bilinear forms and associated linear maps
We now generalize from the symplectic matrix $`\omega `$ to a symplectic form; in five extended comments.
(1): Preliminaries:—
Let $`V`$ be a (real finite-dimensional) vector space, with basis $`e_1,\mathrm{},e_i,\mathrm{}e_n`$. We write $`V^{}`$ for the dual space, and $`e^1,\mathrm{},e^i,\mathrm{}e^n`$ for the dual basis: $`e^i(e_j):=\delta _j^i`$.
We recall that the isomorphism $`e_ie^i`$ is basis-dependent: for a different basis, the corresponding isomorphism would be a different map. Only with the provision of appropriate extra structure would this isomorphism be basis-independent.
For physicists, the most familiar example of such a structure is the spacetime metric $`𝐠`$ in relativity theory. In terms of components, this basis-independence shows up in the way that $`𝐠`$ and its inverse lower and raise indices. As we will see in a moment, the underlying mathematical point is that because $`𝐠`$ is a bilinear form on a vector space $`V`$, i.e. $`𝐠:V\times V\mathrm{I}\mathrm{R}`$, and is non-degenerate, any $`vV`$ defines, independently of any choice of basis, an element of $`V^{}`$: viz. the map $`uV𝐠(u,v)`$. (In fact, $`V`$ is the tangent space at a spacetime point; but this physical interpretation is irrelevant to the mathematical argument.) We will also see that Hamiltonian mechanics has a non-degenerate bilinear form, viz. a symplectic form, that similarly gives a basis-independent isomorphism between a vector space and its dual. (Roughly speaking, this vector space will be the $`2n`$-dimensional space of the $`q`$s and $`p`$s.)
On the other hand: for any vector space $`V`$, the isomorphism between $`V`$ and $`V^{}`$ given by
$$e_i[e_i]V^{}:e^jV^{}e^j(e_i)=\delta _i^j$$
(4.67)
is basis-independent, and so we identify $`e_i`$ with $`[e_i]`$, and $`V`$ with $`V^{}`$. We will write $`<;>`$ (also written $`<,>`$) for the natural pairing (in either order) of $`V`$ and $`V^{}`$: e.g. $`<e_i;e^j>=<e^j;e_i>=\delta _i^j`$.
A linear map $`A:VW`$ induces (basis-independently) a transpose (aka: dual), written $`\stackrel{~}{A}`$ (or $`A^T`$ or $`A^{}`$), $`\stackrel{~}{A}:W^{}V^{}`$ by
$$\alpha W^{},vV:\stackrel{~}{A}(\alpha )(v)<\stackrel{~}{A}(\alpha );v>:=\alpha (A(v))(\alpha A)(v).$$
(4.68)
If $`A:VW`$ is a linear map between real finite-dimensional vector spaces, its matrix with respect to bases $`e_1,\mathrm{},e_i,\mathrm{}e_n`$ and $`f_1,\mathrm{},f_j,\mathrm{}f_m`$ of $`V`$ and $`W`$ is given by:
$$A(e_i)=A_i^jf_j;\mathrm{i}.\mathrm{e}.\mathrm{with}v=v^ie_i,(A(v))^j=A_i^jv^i.$$
(4.69)
So the upper index labels rows, and the lower index labels columns. Similarly, if $`A:V\times W\mathrm{I}\mathrm{R}`$ is a bilinear form, its matrix for these bases is defined as
$$A_{ij}:=A(e_i,f_j)$$
(4.70)
so that on vectors $`v=v^ie_i,w=w^jf_j`$, we have: $`A(v,w)=v^iA_{ij}w^j.`$
(2): Associated maps and forms:—
Given a bilinear form $`A:V\times W\mathrm{I}\mathrm{R}`$, we define the associated linear map $`A^{\mathrm{}}:VW^{}`$ by
$$A^{\mathrm{}}(v)(w):=A(v,w).$$
(4.71)
Then $`A^{\mathrm{}}(e_i)=A_{ij}f^j`$: for both sides send any $`w=w^jf_j`$ to $`A_{ij}w^j`$. That is: the matrix of $`A^{\mathrm{}}`$ in the bases $`e_i,f^j`$ of $`V`$ and $`W^{}`$ is $`A_{ij}`$:
$$[A^{\mathrm{}}]_{ij}=A_{ij}.$$
(4.72)
On the other hand, we can proceed from linear maps to associated bilinear forms. Given a linear map $`B:VW^{}`$, we define the associated bilinear form $`B^{\mathrm{}}`$ on $`V\times W^{}V\times W`$ by
$$B^{\mathrm{}}(v,w)=<B(v);w>.$$
(4.73)
If we put $`A^{\mathrm{}}`$ for $`B`$ in eq. 4.73, its associated bilinear form, acting on vectors $`v=v^ie_i,w=w^jf_j`$, yields, by eq. 4.71:
$$(A^{\mathrm{}})^{\mathrm{}}(v,w)=<A^{\mathrm{}}(v);w>=A(v,w).$$
(4.74)
One similarly shows that if $`B:VW^{}`$, then $`wW`$:
$$(B^{\mathrm{}})^{\mathrm{}}(v)(w)<(B^{\mathrm{}})^{\mathrm{}}(v);w>=B(v)(w)<B(v);w>\mathrm{so}\mathrm{that}(B^{\mathrm{}})^{\mathrm{}}=B.$$
(4.75)
So the flat and sharp operations, and , are inverses.
(3): Tensor products:—
It will sometimes be helpful to put the above ideas in terms of tensor products. If $`vV,wW`$, we can think of $`v`$ and $`w`$ as elements of $`V^{},W^{}`$ respectively. So we define their tensor product as a bilinear form on $`V^{}\times W^{}`$ by requiring for all $`\alpha V^{},\beta W^{}`$:
$$(vw)(\alpha ,\beta ):=v(\alpha )w(\beta )<v;\alpha ><w;\beta >.$$
(4.76)
Similarly for other choices of vector spaces or their duals. Given $`\alpha V^{},\beta W^{}`$, their tensor product is a bilinear form on $`V\times W`$:
$$(\alpha \beta )(v,w):=\alpha (v)\beta (w)<v;\alpha ><w;\beta >.$$
(4.77)
Similarly, we can think of $`\alpha V^{},wW`$ as elements of $`V^{}`$ and $`W^{}`$ respectively, and so define their tensor product as a bilinear form on $`V\times W^{}`$:
$$(\alpha w)(v,\beta ):=\alpha (v)w(\beta )<v;\alpha ><w;\beta >.$$
(4.78)
In this way we can express the linear map $`A:VW`$ in terms of tensor products. Since
$$A(e_i)=A_i^jf_j\mathrm{iff}<A(e_i);f^j>=A_i^j$$
(4.79)
eq. 4.78 implies that
$$A=A_i^je^if_j.$$
(4.80)
Similarly, a bilinear form $`A:V\times W\mathrm{I}\mathrm{R}`$ with matrix $`A_{ij}:=A(e_i,f_j)`$ (cf. eq. 4.70) is:
$$A=A_{ij}e^if^j$$
(4.81)
The definitions of tensor product eq. 4.76, 4.77 and 4.78 generalize to higher-rank tensors (i.e. multilinear maps whose domains have more than two factors). But we will not need these generalizations.
(4): Antisymmetric and non-degenerate forms:—
We now specialize to the forms and maps of central interest in Hamiltonian mechanics. We take $`W=V`$, dim($`V`$)=$`n`$, and define a bilinear form $`\omega :V\times V\mathrm{I}\mathrm{R}`$ to be:
(i): antisymmetric iff: $`\omega (v,v^{})=\omega (v,v^{})`$;
(ii): non-degenerate iff: if $`\omega (v,v^{})=0v^{}V`$, then $`v=0`$.
The form $`\omega `$ and its associated linear map $`\omega ^{\mathrm{}}:VV^{}`$ now have a square matrix $`\omega _{ij}`$ (cf. eq. 4.72). We define the rank of $`\omega `$ to be the rank of this matrix: equivalently, the dimension of the range $`\omega ^{\mathrm{}}(V)`$.
We will also need the antisymmetrized version of eq. 4.77 that is definable when $`W=V`$. Namely, we define the wedge-product of $`\alpha ,\beta V^{}`$ to be the antisymmetric bilinear form on $`V`$, given by
$$\alpha \beta :(v,w)V\times V(\alpha (v))(\beta (w))(\alpha (w))(\beta (v))\mathrm{I}\mathrm{R}.$$
(4.82)
(The connection with Section 4.3.2, especially eq. 4.66, will become clear in a moment; and will be developed in Section 6.2.A.)
It is easy to show that for any bilinear form $`\omega :V\times V\mathrm{I}\mathrm{R}`$: $`\omega `$ is non-degenerate iff the matrix $`\omega _{ij}`$ is non-singular iff $`\omega ^{\mathrm{}}:VV^{}`$ is an isomorphism.
So a non-degenerate bilinear form establishes a basis-independent isomorphism between $`V`$ and $`V^{}`$; cf. the discussion of the spacetime metric $`𝐠`$ in (1) at the start of this Subsection.
Besides, this isomorphism $`\omega ^{\mathrm{}}`$ has an inverse, suggesting another use of the sharp notation, viz. $`\omega ^{\mathrm{}}`$ is defined to be $`(\omega ^{\mathrm{}})^1:V^{}V`$. The isomorphism $`\omega ^{\mathrm{}}:V^{}V`$ corresponds to $`\omega `$’s role, emphasised in Section 4.3.1, of defining a vector field $`X_H`$ from $`dH`$. (But we will see in a moment that the space $`V`$ implicitly considered in Section 4.3.1 had more structure than being just any finite-dimensional real vector space: viz. it was of the form $`W\times W^{}`$.)
NB: This definition of $`^{\mathrm{}}`$ is of course not equivalent to our previous definition, in eq. 4.73, since:
(i): on our previous definition, carried a linear map to a bilinear form, which reversed the passage by from bilinear form to linear map, in the sense that for a bilinear form $`\omega `$, we had $`(\omega ^{\mathrm{}})^{\mathrm{}}=\omega `$; cf. eq. 4.74;
(ii): on the present definition, carries a bilinear form $`\omega :V\times V\mathrm{I}\mathrm{R}`$ to a linear map $`\omega ^{\mathrm{}}:V^{}V`$, which inverts in the sense (different from (i)) that
$$\omega ^{\mathrm{}}\omega ^{\mathrm{}}=id_V\mathrm{and}\omega ^{\mathrm{}}\omega ^{\mathrm{}}=id_V^{}.$$
(4.83)
So beware: though not equivalent, both definitions are used! But it is a natural ambiguity, in so far as the definitions “mesh”. For example, one easily shows that our second definition, i.e. eq. 4.83, is equivalent to a natural expression:
$$\alpha ,\beta V^{}:<\omega ^{\mathrm{}}(\alpha ),\beta >:=\omega ((\omega ^{\mathrm{}})^1(\alpha ),(\omega ^{\mathrm{}})^1(\beta )).$$
(4.84)
It is also straightforward to show that for any bilinear form $`\omega :V\times V\mathrm{I}\mathrm{R}`$: if $`\omega `$ is antisymmetric of rank $`rn\mathrm{dim}(V)`$, then $`r`$ is even. That is: $`r=2s`$ for some integer $`s`$, and there is a basis $`e_1,\mathrm{},e_i,\mathrm{},e_n`$ of $`V`$ for which $`\omega `$ has a simple expansion as wedge-products
$$\omega =\mathrm{\Sigma }_{i=1}^se^ie^{i+s};$$
(4.85)
equivalently, $`\omega `$ has the $`n\times n`$ matrix
$$\omega =\left(\begin{array}{ccc}\mathrm{𝟎}& \mathrm{𝟏}& \mathrm{𝟎}\\ \mathrm{𝟏}& \mathrm{𝟎}& \mathrm{𝟎}\\ \mathrm{𝟎}& \mathrm{𝟎}& \mathrm{𝟎}\end{array}\right).$$
(4.86)
where $`\mathrm{𝟏}`$ is the $`s\times s`$ identity matrix, and similarly for the zero matrices of various sizes. This normal form of antisymmetric bilinear forms is an analogue of the Gram-Schmidt theorem that an inner product space has an orthonormal basis, and is proved by an analogous argument.
(5): Symplectic forms:—
As usually formulated, Hamiltonian mechanics uses a non-degenerate antisymmetric bilinear form: i.e. $`r=n`$. So eq. 4.86 loses its bottom row and right column consisting of zero matrices, and reduces to the form of Section 4.3.1’s naive symplectic matrix, eq. 4.58. Equivalently: eq. 4.85 reduces to eq. 4.66.
Accordingly, we define: a symplectic form on a (real finite-dimensional) vector space $`Z`$ is a non-degenerate antisymmetric bilinear form $`\omega `$ on $`Z`$: $`\omega :Z\times Z\mathrm{I}\mathrm{R}`$. $`Z`$ is then called a symplectic vector space. It follows that $`Z`$ is of even dimension.
Besides, in Hamiltonian mechanics (as usually formulated) the vector space $`Z`$ is a product $`V\times V^{}`$ of a vector space and its dual. Indeed, this was already suggested by:
(i) the fact in (3) of Section 2.2.2, that the canonical momenta $`p_i:=\frac{L}{\dot{q}^i}`$ transform as a 1-form, and
(ii) Section 4.3.1’s discussion of the one-form field $`H`$ determining a vector field $`X_H`$.
Thus we define the canonical symplectic form $`\omega `$ on $`Z:=V\times V^{}`$ by
$$\omega ((v_1,\alpha _1),(v_2,\alpha _2)):=\alpha _2(v_1)\alpha _1(v_2).$$
(4.87)
So defined, $`\omega `$ is by construction a symplectic form, and so has the normal form given by eq. 4.58.
Given a symplectic vector space $`(Z,\omega )`$, the natural question arises which linear maps $`A:ZZ`$ preserve the normal form given by eq. 4.58. It is straightforward to show that this is equivalent to $`A`$ preserving the form of Hamilton’s equations (for any Hamiltonian); so that these maps $`A`$ are called canonical (or symplectic, or Poisson). But since (as I announced) this paper does not need details about the theory of canonical transformations, I will not go into details about this. Suffice it to say here the following.
$`A:ZZ`$ is symplectic iff, writing $`\stackrel{~}{}`$ for the transpose (eq. 4.68) and using the second definition eq. 4.83 of , the following maps (both from $`Z^{}`$ to $`Z`$) are equal:
$$A\omega ^{\mathrm{}}\stackrel{~}{A}=\omega ^{\mathrm{}};$$
(4.88)
or in matrix notation, with the matrix $`\omega `$ given by eq. 4.58, and again writing $`\stackrel{~}{}`$ for the transpose of a matrix
$$A\omega \stackrel{~}{A}=\omega .$$
(4.89)
(Equivalent formulas are got by taking inverses. We get, respectively: $`\stackrel{~}{A}\omega ^{\mathrm{}}A=\omega ^{\mathrm{}}`$ and $`\stackrel{~}{A}\omega A=\omega `$.)
The set of all such linear symplectic maps $`A:ZZ`$ form a group, the symplectic group, written Sp($`Z,\omega `$).
To sum up this Subsection:— We have, for a vector space $`V`$, dim($`V`$) = $`n`$, and $`Z:=V\times V^{}`$:
(i): the canonical symplectic form $`\omega :Z\times Z\mathrm{I}\mathrm{R}`$; with normal form given by eq. 4.58;
(ii): the associated linear map $`\omega ^{\mathrm{}}:ZZ^{}`$; which is an isomorphism, since $`\omega `$ is non-degenerate;
(iii): the associated linear map $`\omega ^{\mathrm{}}:Z^{}Z`$; which is an isomorphism, since $`\omega `$ is non-degenerate; and is the inverse of $`\omega ^{\mathrm{}}`$; (cf. eq. 4.83).
We will see shortly that Hamiltonian mechanics takes $`V`$ to be the tangent space $`T_q`$ at a point $`qQ`$, so that $`Z`$ is $`T_q\times T_q^{}`$, i.e. the tangent space to the space $`\mathrm{\Gamma }`$ of the $`q`$s and $`p`$s.
## 5 Poisson brackets and Noether’s theorem
We have seen how a single scalar function $`H`$ on phase space $`\mathrm{\Gamma }`$ determines the evolution of the system via a combination of partial differentiation (the gradient of $`H`$) with the symplectic matrix. We now express these ideas in terms of Poisson brackets.
For our purposes, Poisson brackets will have three main advantages; which will be discussed in the following order in the Subsections below. Poisson brackets:
(i) give a neat expression for the rate of change of any dynamical variable;
(ii) give a version of Noether’s theorem which is more simple and powerful (and even easier to prove!) than the Lagrangian version; and
(iii) lead to the generalized Hamiltonian framework mentioned in Section 6.8.
All three advantages arise from the way the Poisson bracket encodes the way that a scalar function determines a (certain kind of) vector field.
### 5.1 Poisson brackets introduced
The rate of change of any dynamical variable $`f`$, taken as a scalar function on phase space $`\mathrm{\Gamma }`$, $`f(q,p)\mathrm{I}\mathrm{R}`$, is given (with summation convention) by
$$\frac{df}{dt}=\dot{q}^i\frac{f}{q^i}+\dot{p}_i\frac{f}{p_i}.$$
(5.90)
(If $`f`$ is time-dependent, $`f:(q,p,t)\mathrm{\Gamma }\times \mathrm{I}\mathrm{R}f(q,p,t)\mathrm{I}\mathrm{R}`$, the right-hand-side includes a term $`\frac{f}{t}`$. But on analogy with how our discussion of Lagrangian mechanics imposed scleronomic constraints, a time-independent work-function etc., we here set aside the time-dependent case.) Applying Hamilton’s equations, this is
$$\frac{df}{dt}=\frac{H}{p_i}\frac{f}{q^i}\frac{H}{q^i}\frac{f}{p_i}.$$
(5.91)
This suggests that we define the Poisson bracket of any two such functions $`f(q,p),g(q,p)`$ by
$$\{f,g\}:=\frac{f}{q^i}\frac{g}{p_i}\frac{f}{p_i}\frac{g}{q^i};$$
(5.92)
so that the rate of change of $`f`$ is given by
$$\frac{df}{dt}=\{f,H\}.$$
(5.93)
In terms of the $`2n`$ coordinates $`\xi ^\alpha `$ (eq. 4.54) and the matrix elements $`\omega ^{\alpha \beta }`$ of $`\omega `$ (eq. 4.61), we can write eq. 5.91 as
$$\frac{df}{dt}=(_\alpha f)\dot{\xi }^\alpha =(_\alpha f)\omega ^{\alpha \beta }(_\beta H);$$
(5.94)
and so we can define the Poisson bracket by
$$\{f,g\}:=(_\alpha f)\omega ^{\alpha \beta }(_\beta g)\frac{f}{\xi ^\alpha }\omega ^{\alpha \beta }\frac{g}{\xi ^\beta }.$$
(5.95)
In matrix notation: writing the naive gradients of $`f`$ and of $`g`$ as column vectors $`f`$ and $`g`$, and writing $`\stackrel{~}{}`$ for transpose, we have at any point $`z=(q,p)\mathrm{\Gamma }`$:
$$\{f,g\}(z)=\stackrel{~}{f}(z).\omega .g(z).$$
(5.96)
With these definitions of the Poisson bracket, we readily infer the following five results. (Later discussion will bring out the significance of some of these; in particular, Section 6.8 will take some of them to jointly define a primitive Poisson bracket for a generalized Hamiltonian mechanics.)
(1): Since the Poisson bracket is antisymmetric, $`H`$ itself is a constant of the motion:
$$\frac{dH}{dt}=\{H,H\}0.$$
(5.97)
(2): The Poisson bracket of a product is given by “Leibniz’s rule”: i.e. for any three functions $`f,g,h`$, we have
$$\{f,hg\}=\{f,h\}g+h\{f,g\}.$$
(5.98)
(3): Taking the Poisson bracket as itself a dynamical variable, its time-derivative is given by a “Leibniz rule”; i.e. the Poisson bracket behaves like a product:
$$\frac{d}{dt}\{f,g\}=\{\frac{df}{dt},g\}+\{f,\frac{dg}{dt}\}.$$
(5.99)
(4): The Jacobi identity (easily deduced from (3)):
$$\{\{f,h\},g\}+\{\{g,f\},h\}+\{\{h,g\},f\}=0.$$
(5.100)
(5): The Poisson brackets for the $`q`$s, $`p`$s and $`\xi `$s are:
$`\{\xi ^\alpha ,\xi ^\beta \}=\omega ^{\alpha \beta };\mathrm{i}.\mathrm{e}.`$ (5.101)
$`\{q^i,p_j\}=\delta _j^i,\{q^i,q^j\}=\{p_i,p_j\}=0.`$ (5.102)
Eq. 5.102 is very important, both for general theory and for problem-solving. The reason is that preservation of these Poisson brackets, by a smooth transformation of the $`2n`$ variables $`(q,p)(Q(q,p),P(q,p))`$, is necessary and sufficient for the transformation being canonical. Besides, in this equivalence ‘canonical’ can be understood both in the usual elementary sense of preserving the form of Hamilton’s equations, for any Hamiltonian function, and in the geometric sense of preserving the symplectic form (explained in (5) of Section 4.3.3, and for manifolds in Section 6).
Note here that, as the phrase ‘for any Hamiltonian function’ brings out, the notion of a canonical transformation is independent of the forces on the system as encoded in the Hamiltonian. That is: the notion is a matter of $`\mathrm{\Gamma }`$’s geometry—as we will emphasise in Section 6.
But (as I announced in Section 4.1) I will not need to go into many details about canonical transformations, essentially because this paper does not aim to survey the whole of Hamiltonian mechanics, or even all that can be said about reducing problems, e.g. by finding simplifying canonical transformations. It aims only to survey the way that symmetries and conserved quantities effect such reductions. In the rest of this Subsection, I begin describing Poisson brackets’ role in this, in particular Noether’s theorem. But the description can only be completed once we have the geometric perspective on Hamiltonian mechanics, i.e. in Section 6.5.
### 5.2 Hamiltonian vector fields
Section 4.3.1 described how the symplectic matrix enabled the scalar function $`H`$ on $`\mathrm{\Gamma }`$ to determine a vector field $`X_H`$. The previous Subsection showed how the Poisson bracket expressed any dynamical variable’s rate of change along $`X_H`$. We now bring these ideas together, and generalize.
Recall that a vector $`X`$ at a point $`x`$ of a manifold $`M`$ can be identified with a directional derivative operator at $`x`$ assigning to each smooth function $`f`$ defined on a neighbourhood of $`x`$ its directional derivative along any curve that has $`X`$ as its tangent vector. Thus recall the Lagrangian definition of the dynamical vector field, eq. 2.8 in Section 2.2. Similarly here: the dynamical vector field $`X_H=:D`$ is a derivative operator on scalar functions, which can be written in terms the Poisson bracket:
$$D:=X_H=\frac{d}{dt}=\dot{q}^i\frac{}{q^i}+\dot{p}_i\frac{}{p_i}=\frac{H}{p_i}\frac{}{q^i}\frac{H}{q^i}\frac{}{p_i}=\{,H\}.$$
(5.103)
But this point applies to any smooth scalar, $`f`$ say, on $`\mathrm{\Gamma }`$. That is: although we think of $`H`$ as the energy that determines the real physical evolution, the mathematics is of course the same for such an $`f`$. So any such function determines a vector field, $`X_f`$ say, on $`\mathrm{\Gamma }`$ that generates what the evolution “would be if $`f`$ was the Hamiltonian”. Thinking of the integral curves as parametrized by $`s`$, we have
$$X_f=\frac{d}{ds}=\{,f\}.$$
(5.104)
$`X_f`$ is called the Hamiltonian vector field of (for) $`f`$; just as, for the physical Hamiltonian, $`fH`$, Section 4.3.1 called $`X_H`$the Hamiltonian vector field’.
The notion of a Hamiltonian vector field will be crucial for what follows, not least for Noether’s theorem in the very next Subsection. For the moment, we just make two remarks which we will need later.
So every scalar $`f`$ determines a Hamiltonian vector field $`X_f`$. But note that the converse is false: not every vector field $`X`$ on $`\mathrm{\Gamma }`$ is the Hamiltonian vector field of some scalar. For a vector field (equations of motion) $`X`$, with components $`X^\alpha `$ in the coordinates $`\xi ^\alpha `$ defined by eq. 4.54
$$\dot{\xi }^\alpha =X^\alpha (\xi ),$$
(5.105)
there need be no scalar $`H:\mathrm{\Gamma }\mathrm{I}\mathrm{R}`$ such that, as required by eq. 4.61,
$$X^\alpha =\omega ^{\alpha \beta }_\beta H.$$
(5.106)
This is the same point as in (ii) of Section 4.2.3: that Hamilton’s equations have the special feature that all the right hand sides are, up to a sign, partial derivatives of a single function $`H`$—a feature that underpins the possibility of expressing the equations of motion by variational principles.
We also need to note under what condition is a vector field $`X`$ Hamiltonian; (this will bear on Noether’s theorem). The answer is: $`X`$ is locally Hamiltonian, i.e. there is locally a scalar $`f`$ such that $`X=X_f`$, iff $`X`$ generates a one-parameter family of canonical transformations. We will give a modern geometric proof of this in Section 6.5. For the moment, we only need to note, as at the end of Section 5.1, that here ‘canonical transformation’ can be understood in the usual elementary sense as a transformation of $`\mathrm{\Gamma }`$ that preserves the form of Hamilton’s equations (for any Hamiltonian); or equivalently, as preserving the Poisson bracket; or equivalently, as preserving the symplectic form (to be defined for manifolds, in Section 6).
### 5.3 Noether’s theorem
#### 5.3.1 An apparent “one-liner”, and three claims
In the Hamiltonian framework, the core of the proof of Noether’s theorem is very simple; as follows. The Poisson bracket is obviously antisymmetric. So for any scalar functions $`f`$ and $`H`$, we have
$$X_f(H)\frac{dH}{ds}\{H,f\}=0\text{ iff }\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0}=\{f,H\}=X_H(f)D(f).$$
(5.107)
In words: $`H`$ is constant under the flow of the vector field $`X_f`$ (i.e. under what the evolution would be if $`f`$ was the Hamiltonian) iff $`f`$ is constant under the dynamical flow $`X_HD`$.
This “one-liner” is the Hamiltonian version of Noether’s theorem! There are three claims here. The first two relate back to the Lagrangian version of the theorem. The third is about the definition of a (continuous) symmetry for a Hamiltonian system, and so about how we should formulate the Hamiltonian version of Noether’s theorem. I will state all three claims, but in this Subsection justify only the first two. For it will be convenient to postpone the third till after we have introduced some modern geometry (Section 6.5).
First, for eq. 5.107 to deserve the name ‘Noether’s theorem’, I need to show that it encompasses Section 3’s Lagrangian version of Noether’s theorem (despite the trivial proof!).
Second, in order to justify my claim that the Hamiltonian version of Noether’s theorem is more powerful than the Lagrangian version, I need to show that eq. 5.107 says more than that version, i.e. that it covers more symmetries.
To state the third claim, note first that we expect a Hamiltonian version of Noether’s theorem to say something like: to every continuous symmetry of a Hamiltonian system, there corresponds a conserved quantity. Here, we expect a ‘continuous symmetry’ to be defined by a vector field on $`\mathrm{\Gamma }`$ (or by its flow). Indeed, a symmetry of a Hamiltonian system is usually defined as a transformation of $`\mathrm{\Gamma }`$ that:
(1) is canonical; (a condition independent of the forces on the system as encoded in the Hamiltonian: a matter of $`\mathrm{\Gamma }`$’s intrinsic geometry); and also
(2) preserves the Hamiltonian function; (a condition obviously dependent on the Hamiltonian).
Accordingly, a continuous symmetry is defined as a vector field on $`\mathrm{\Gamma }`$ that generates a one-parameter family of such transformations; (or as such a field’s flow, i.e. as the family itself).
But with this definition of ‘continuous symmetry’ (of a Hamiltonian system), eq. 5.107 seems to suffer from two lacunae, if taken to express Noether’s theorem, that to every continuous symmetry there corresponds a conserved quantity. Agreed, the rightward implication of eq. 5.107 provides, for a vector field $`X_f`$ with property (2), the conserved quantity $`f`$. But there seem to be two lacunae:
(a): eq. 5.107 is silent about whether $`X_f`$ has property (1), i.e. generates canonical transformations.
(b): eq. 5.107 considers only Hamiltonian vector fields, i.e. vector fields $`X`$ induced by some $`f`$, $`X=X_f`$. But as noted at the end of Section 5.2, there are countless vector fields on $`\mathrm{\Gamma }`$ that are not Hamiltonian. If such a field could be a continuous symmetry, eq. 5.107’s rightward implication would fall short of saying that to every continuous symmetry, there corresponds a conserved quantity.
So the third claim I need is that these lacunae are illusory. In fact, a single result will deal with both (a) and (b). Namely, it will suffice to show that a vector field $`X`$ on $`\mathrm{\Gamma }`$ has property (1), i.e. generates canonical transformations, iff it is Hamiltonian, i.e. induced by some $`f`$, $`X=X_f`$. But I postpone showing this till we have more modern geometry in hand; cf. Section 6.5.
#### 5.3.2 The relation to the Lagrangian version
On the other hand, we can establish the first two claims with the elementary apparatus so far developed. I will concentrate on justifying the first claim; that will also make the second claim clear.
For the first claim, we need to show that:
(i): to any variational symmetry of the Lagrangian $`L`$, i.e. a vector field $`X`$ on $`Q`$ obeying eq. 3.23, there corresponds a vector field $`X_f`$ on $`\mathrm{\Gamma }`$ for which $`X_f(H)=0`$; and
(ii): the correspondence in (i) is such that the scalar $`f`$ can be taken to be (the Hamiltonian version of) the momentum $`p_X`$ conjugate to $`X`$, defined by eq. 3.29 (or geometrically, by 3.48).
It will be clearest to proceed in two stages.
(A): First, I will show (i) and (ii).
(B): Then I will discuss how (A) relates to the usual definition of a symmetry of a Hamiltonian system.
(A): The easiest way to show (i) and (ii) is to use the fact discussed after eq. 3.37, that every variational symmetry $`X`$ arises, around a point where it is non-zero, from a cyclic coordinate in some local system of coordinates. (Recall that this follows from the basic “rectification” theorem securing the local existence and uniqueness of solutions of ordinary differential equations.) That is, there is some coordinate system $`(q)`$ on some open subset of $`X`$’s domain of definition on $`Q`$ such that
(a): $`X`$ being a variational symmetry is equivalent to $`q^n`$ being cyclic, i.e. $`\frac{L}{q^n}=0`$;
(b): the momentum $`p_X`$, which the Lagrangian theorem says is conserved, is the elementary generalized momentum $`p_n:=\frac{L}{\dot{q}^n}`$.
So suppose given a variational symmetry $`X`$, and a coordinate system $`(q)`$ satisfying (a)-(b). Now we recall that the Legendre transformation, i.e. the transition between Lagrangian and Hamiltonian frameworks, does not “involve the dependence on the $`q`$s”. More precisely, we recall eq. 4.56, $`\frac{H}{q^n}=\frac{L}{q^n}`$. Now consider $`p_n:\mathrm{\Gamma }\mathrm{I}\mathrm{R}`$. This $`p_n`$ will do as the function $`f`$ required in (i) and (ii) above, since
$$X_{p_n}(H)\{H,p_n\}=\frac{H}{q^n}=\frac{L}{q^n}=0.$$
(5.108)
Applying eq. 5.107 to eq. 5.108, we deduce that $`p_n`$, i.e. the $`p_X`$ of the Lagrangian theorem, is conserved.
(Hence my remark after eq. 4.56, that the elementary result that $`p_n`$ is conserved iff $`q^n`$ is cyclic, underpins the Hamiltonian version of Noether’s theorem; just as the corresponding Lagrangian result underpins the Lagrangian version of Noether’s theorem: cf. discussion after eq. 3.37.)
(B): I agree that this simple proof seems suspiciously simple. Besides, the suspicion grows when you notice that my argument in (A) has not used a definition of a symmetry, in particular a continuous symmetry, of a Hamiltonian system (contrast Section 3.2). As discussed in Section 5.3.1, we expect a Hamiltonian version of Noether’s theorem to say ‘to every continuous symmetry of a Hamiltonian system there corresponds a conserved quantity’; where a continuous symmetry is a vector field that (1) generates canonical transformations and (2) preserves the Hamiltonian. So the argument in (A) is suspicious since, although eq. 5.108, or the left hand side of eq. 5.107, obviously expresses property (2), i.e. preserving the Hamiltonian, the argument in (A) seems to nowhere use property (1), i.e. the symmetry generating canonical transformations.
But in fact, all is well. The reason why lies in the fact mentioned in (i), (a) of Section 4.1: that every point transformation (together with its lift to $`TQ`$) defines a corresponding canonical transformation on $`T^{}Q`$. That is to say: property (1) is secured by the fact that the Lagrangian Noether’s theorem of Section 3 is restricted to symmetries induced by point transformations.
In other words, in terms of the vector field (variational symmetry) $`X`$ given us by (a) in (A) above: one can check that $`X`$ defines a vector field on $`\mathrm{\Gamma }`$ (equivalently: a one-parameter family of transformations on $`\mathrm{\Gamma }`$) that is canonical, i.e. preserves Hamilton’s equations or equivalently the symplectic form. Indeed, one can easily check that, once we rectify the Lagrangian variational symmetry $`X`$, so that it generates the rectified one-parameter family of point transformations: $`q_i=\mathrm{const},in;q_nq_n+ϵ`$, the vector field that $`X`$ defines on $`\mathrm{\Gamma }`$ is precisely the field $`X_{p_n}`$ chosen above.<sup>17</sup><sup>17</sup>17Details about point transformations on $`Q`$ defining a canonical transformation on $`T^{}Q`$, and lifting the vector field $`X`$ to $`\mathrm{\Gamma }`$, can be found: (i) using traditional terms, in Goldstein et al. (2002: 375-376) and Lanczos (1986: Chapter VII.2); (ii) using modern geometric terms (as developed in Section 6), in Abraham and Marsden (1978: Sections 3.2.10-3.2.12) and Marsden and Ratiu (1999: Sections 6.3-6.4).
Finally, the discussion in (B) also vindicates the second claim in Section 5.3.1: that the Hamiltonian version of Noether’s theorem, eq. 5.107, says more than the Lagrangian version, i.e. covers more symmetries. This follows from the fact (announced in (i) (b) of Section 4.1) that there are canonical transformations not induced by a point transformation (together with its lift).
In elementary discussions, this is often expressed in terms of canonical transformations being allowed to “mix” the $`q`$s and $`p`$s. But a more precise, and geometric, statement is the result announced at the end of Section 5.2 (whose proof is postponed to Section 6.5): that the condition for a vector field on $`\mathrm{\Gamma }`$ to generate a one-parameter family of canonical transformations is merely that it be a Hamiltonian vector field. That is: for any scalar $`f:\mathrm{\Gamma }\mathrm{I}\mathrm{R}`$, the vector field $`X_f`$ generates such a family.
In this sense, canonical transformations are two a penny (also known as: a dime a dozen!). So it is little wonder that most discussions emphasise the other condition, i.e. property (2): that $`X_f`$ preserve the Hamiltonian, $`X_f(H)=0`$. Only very special $`f`$s will satisfy $`X_f(H)=0`$; and if we are given $`H`$ (in certain coordinates $`q,p`$), it can be very hard to find (the coordinate expression of) such an $`f`$.
Indeed, when Jacobi first propounded the theory of canonical transformations, in his Lectures on Dynamics (1842), he was of course aware of this. Accordingly, he pointed out that in theoretical mechanics, it was often more fruitful to first consider an $`f`$ (equivalently: a canonical transformation), and then cast about for a Hamiltonian that it preserved. He wrote: ‘The main difficulty in integrating a given differential equation lies in introducing convenient variables, which there is no rule for finding. Therefore we must travel the reverse path and after finding some notable substitution, look for problems to which it can be successfully applied’; (quoted in Arnold (1989, p. 266)). The fact that Jacobi solved many previously intractable problems bears witness to the power of this strategy, and of his theory of canonical transformations.
We can sum up this Subsection in two comments:—
(1) In Hamiltonian mechanics, Noether’s theorem is a biconditional, an ‘iff’ statement. Not only does a Hamiltonian symmetry—i.e. a vector field $`X`$ on $`\mathrm{\Gamma }`$ that generates canonical transformations (equivalently: preserves the symplectic form, or the Poisson bracket) and preserves the Hamiltonian, $`X(H)=0`$—provide a constant of the motion. Also, given a constant of the motion $`f:\mathrm{\Gamma }\mathrm{I}\mathrm{R}`$, there is a symmetry of the Hamiltonian, viz. the vector field $`X_f`$. (Or if one prefers the integral notion of symmetry: the flow of $`X_f`$). This converse implication, from constant to symmetry, contrasts with the Lagrangian framework; cf. the end of Section 3.4.1.
(2) In elementary Hamiltonian mechanics, Noether’s theorem has a very simple one-line proof, viz. eq. 5.107.
Later, we will return to Noether’s theorem. Section 6.5 will justify the third claim of Section 5.3.1, by showing that a vector field generates a one-parameter family of canonical transformations iff it is a Hamiltonian vector field. Meanwhile, we end Section 5 with a comment about “iterating” Noether’s theorem, and the distinction between such an iteration and the idea of complete integrability.
### 5.4 Glimpsing the “complete solution”
Suppose we “iterate” Noether’s theorem. That is: suppose there are several (continuous) symmetries of the Hamiltonian and so several constants of the motion. Each will confine the system’s time-evolution to a ($`2n1`$)-dimensional hypersurface of $`\mathrm{\Gamma }`$. In general, the intersection of $`k`$ such surfaces will be a hypersurface of dimension $`2nk`$ (i.e. of co-dimension $`k`$); to which the motion is therefore confined. The theory of symplectic reduction (Butterfield 2006) describes how to do a “quotiented dynamics” in this general situation. Here, I just remark on one aspect; which will not be developed in the sequel.
Locally, the rectification theorem secures, for any system, not just several constants of the motion, but “all you could ask for”. Applying the theorem (eq. 3.38 and 3.39) to the Hamiltonian vector field $`X_H`$ on $`\mathrm{\Gamma }`$, we infer that locally there are coordinates $`\xi ^\alpha `$ (maybe very hard to find!) in which $`X_H`$ has $`2n1`$ components that vanish throughout the neighbourhood, while the other component is 1:
$$X_H^\alpha =0\mathrm{for}\alpha =1,2,\mathrm{},2n1;X_H^{2n}=1.$$
(5.109)
So the coordinates $`\xi ^\alpha ,\alpha =1,\mathrm{},2n1,`$ form $`2n1`$ constants of the motion. They are functionally independent, and all other constants of the motion are functions of them; (cf. point (ii) after eq. 3.39). So the motion is confined to the one-dimensional intersection of the $`2n1`$ hypersurfaces, each of co-dimension 1. That is to say, it is confined to the curve given by: $`\xi ^\alpha =\mathrm{const},\alpha =1,\mathrm{},2n1,\xi ^{2n}=t`$.
To this, Noether’s theorem eq. 5.107 adds the physical idea that each such constant of the motion defines a vector field $`X_{\xi ^\alpha }`$ that generates a symmetry of the Hamiltonian:
$$X_{\xi ^\alpha }(H)=0,\mathrm{for}\alpha =1,2,\mathrm{},2n1.$$
(5.110)
In this local sense, the “complete solution” of any Hamiltonian system lies in the local constants of the motion, or equivalently the local symmetries of its Hamiltonian $`H`$.
To sum up: locally, any Hamiltonian system is “completely integrable”. But the scare-quotes here are a reminder that these phrases are usually used with other, stronger, meanings: either that there are $`2n1`$ global constants of the motion or that the system is completely integrable in the sense of Liouville’s theorem.
## 6 A geometrical perspective
In this final Section, we develop the modern geometric description of Hamiltonian mechanics. We will build especially on Sections 4.3; one main aim will of course be to complete the discussion of Noether’s theorem, begun in Section 5.3.
There will be eight Subsections. First, we introduce the cotangent bundle $`T^{}Q`$. Then we collect what we will need about forms. Then we can show that any cotangent bundle is a symplectic manifold. This enables us to formulate Hamilton’s equations geometrically; and to complete the discussion of Noether’s theorem. Then we report Darboux’s theorem, and its relation to reduction of problems. Then we return to the Lagrangian framework, by sketching the geometric formulation of the Legendre transformation. Finally, we “glimpse the landscape ahead” by mentioning the more general framework for Hamiltonian mechanics that uses Poisson manifolds.
### 6.1 Canonical momenta are one-forms: $`\mathrm{\Gamma }`$ as $`T^{}Q`$
So far we have treated the phase space $`\mathrm{\Gamma }`$ informally: saying just that it is a $`2n`$-dimensional space coordinatized by the $`q`$s, a smooth coordinate system on the configuration manifold $`Q`$, and the $`p`$s, which are canonical momenta $`\frac{L}{\dot{q}^i}`$. But we also saw in (3) of Section 2.2.2 that at each point $`qQ`$, the $`p_i`$ transform as a 1-form (eq. 2.12). Accordingly we now take the physical state of the system to be a point in the cotangent bundle $`T^{}Q`$, the $`2n`$-dimensional manifold whose points are pairs $`(q,p)`$ with $`qQ,pT_q^{}`$.
I stress that from now on, the symbol $`p`$ has a (fruitful!) ambiguity, between “dynamics” and “kinematics/geometry”. For $`p`$ represents both:
(A) the conjugate momentum $`\frac{L}{\dot{q}}`$, which of course depends on the choice of $`L`$; and
(B) a point in a fibre $`T_q^{}`$ of the cotangent bundle $`T^{}Q`$ (i.e. a 1-form or covector); or relatedly: the components $`p_i`$ of such a 1-form: notions that are independent of any choice of a Lagrangian or Hamiltonian.
In more detail:—
(A): Recall that in the Lagrangian framework, the basic equations (eq. 2.1, or Newton’s second law!) being second-order in time prompts us to take the initial $`q`$ and $`\dot{q}`$ as chosen independently, with $`L`$ (encoding the forces on the system) then determining the evolution (the Lagrangian dynamical vector field $`D`$)—and so also determining the actual “realized” value of $`\dot{q}`$ at other times as a function of $`q`$, and so ultimately, of $`t`$. Similarly here: Newton’s second law being second-order in time prompts us to take the initial $`q`$ and $`p`$ as independent, with $`H`$ (encoding the forces on the system) then determining the evolution (the Hamiltonian dynamical vector field $`D`$)—and so also determining the actual value of $`p`$ at other times as a function of $`q`$, and so ultimately, of $`t`$. Besides, by passing via the Legendre transformation back to the Lagrangian framework, one can check that the later actual value of $`p`$ is determined to equal $`\frac{L}{\dot{q}}`$.
(B): But $`p`$ also represents any 1-form (so that $`p_i`$ represents the 1-form’s coordinates). Here, we need to recall three points:—
(i): A local coordinate system (a chart) on $`Q`$ defines a basis in the tangent space $`T_q`$ at any point $`q`$ in the chart’s domain. As usual, I write the chart’s coordinate functions as $`q^i`$. So I shall temporarily denote the chart by $`[q]`$, so that there are coordinate functions $`q^i:\mathrm{dom}([q])\mathrm{I}\mathrm{R}`$. I write elements of the coordinate basis as usual, as $`\frac{}{q^i}`$.
(ii): The chart $`[q]`$ thereby also defines a dual basis $`dq^i`$ in the cotangent space $`T_q^{}`$ at any $`q\mathrm{dom}([q])`$.
(Here I recall, en passant, that the isomorphism at each $`q`$ between $`T_q`$ and $`T_q^{}`$, that maps the basis element $`\frac{}{q^i}T_q`$ to the one-form $`dq^i`$ in the dual basis, is basis-dependent. A different basis $`\frac{}{q^i}`$ would give a different isomorphism. Cf. the discussion in (1) of Section 4.3.3.)
(iii): Putting (i) and (ii) together: the chart $`[q]`$ thereby also induces a local coordinate system on a neighbourhood of the cotangent bundle around any point $`(q,p)T^{}Q`$ with $`q\mathrm{dom}([q])`$ and $`pT_q^{}`$.
Putting (i)-(iii) together: the coordinates of any point $`(q,p)`$ in $`T^{}Q`$ in such a coordinate system are usually also written as $`(q,p)`$. That is: $`p`$ is used for the components of any 1-form, in the basis $`dq^i`$ dual to a coordinate basis $`\frac{}{q^i}`$. So, similarly to (i) above: I will write this induced chart on $`T^{}Q`$ as $`[q,p]`$.
(C): Taken together, points (A) and (B) prompt a question:
> Why should an evolution from an arbitrary initial state $`T^{}Q`$ have the property that:—
> if we choose to express
> (i) its configuration, $`q_0`$ say, in terms of an arbitrary initial coordinate system $`[q]`$ on $`Q`$, and
> (ii) its momenta $`\frac{L}{\dot{q}}`$ in terms of the basis $`dq`$ dual to the coordinate basis $`\frac{}{q}`$ at $`q_0`$:—
> then
> the states at a later time $`t`$ have their momenta—which the Lagrangian framework tells us must be $`\frac{L}{\dot{q}}`$ (cf. (A))—equal to their components in the dual basis to the later coordinate basis, i.e. the coordinate basis $`\frac{}{q}`$ at the later configuration $`q_t`$?
>
> In short: why should the state’s components in the dual basis of any coordinate basis continue to be equal, as dynamical evolution goes on, to the values of canonical momenta i.e. $`\frac{L}{\dot{q}}`$?
A good question. The short answer lies in combining Hamilton’s equations for the time-derivative of the $`p_i`$ (eq. 4.53) with Lagrange’s equations, and with the fact that the partial derivatives with respect to $`q^i`$ of the Hamiltonian and Lagrangian, $`H`$ and $`L`$, are negatives of each other (eq. 4.56). Thus we have:
$$\dot{p}_i=\frac{H}{q^i}=\frac{L}{q^i}=\frac{d}{dt}\left(\frac{L}{\dot{q}^i}\right).$$
(6.111)
From this it is clear that for any coordinate system, if at $`t_0`$, $`p_i`$ is chosen to equal $`\frac{L}{\dot{q}^i}`$, then this will be so at later times. For eq. 6.111 forces their time-derivatives to be equal—and so also, their later values must be equal.
So much for the short answer. We will also get more insight into the relations between the Lagrangian and Hamiltonian frameworks in
(i) the fact, expounded in Section 6.3 below, that any cotangent bundle has a natural symplectic structure, independent of the specification of any Lagrangian or Hamiltonian function; and
(ii) some further details about the Legendre transformation, which is further discussed in Section 6.7.
### 6.2 Forms, wedge-products and exterior derivatives
As I said at the end of Section 4.3.2, this paper can largely avoid the theory of forms. For what follows (especially Section 6.5), I need to recall only:
(i) the idea of forms of various degrees, together comprising the exterior algebra, and equipped with operations of wedge-product and contraction (Section 6.2.1);
(ii) the ideas of differential forms, the exterior derivative, and of exact and closed forms (Section 6.2.2).
#### 6.2.1 The exterior algebra; wedge-products and contractions
We begin by recalling some ideas of Sections 4.3.2 and 4.3.3. Let us again begin with the simplest possible case, $`\mathrm{I}\mathrm{R}^2`$, considered as a vector space: not as a manifold with a copy of itself as tangent space at each point.
If $`\alpha ,\beta `$ are covectors, i.e. elements of $`(\mathrm{I}\mathrm{R}^2)^{}`$, we define their wedge-product, an antisymmetric bilinear form on $`\mathrm{I}\mathrm{R}^2`$, by
$$\alpha \beta :(v,w)\mathrm{I}\mathrm{R}^2\times \mathrm{I}\mathrm{R}^2(\alpha (v))(\beta (w))(\alpha (w))(\beta (v))\mathrm{I}\mathrm{R}.$$
(6.112)
Let us write the standard basis elements of $`\mathrm{I}\mathrm{R}^2`$ as $`\frac{}{q}`$ and $`\frac{}{p}`$, with elements of $`\mathrm{I}\mathrm{R}^2`$ having components $`(q,p)`$ in this basis; and let us write the elements of the dual basis as $`dq,dp`$. Recalling the definition of the area form $`A`$, eq. 4.64, we deduce that $`A`$ is $`dqdp`$.
Similarly for $`\mathrm{I}\mathrm{R}^{2n}`$. Recall that the symplectic matrix defines an antisymmetric bilinear form on $`\mathrm{I}\mathrm{R}^{2n}`$ by eq. 4.66. The value on a pair $`(q,p)(q^1,\mathrm{}q^n;p_1,\mathrm{},p_n),(q^{},p^{})(q^1,\mathrm{}q^n;p_1^{},\mathrm{},p_n^{})`$ is the sum of the signed areas of the $`n`$ parallelograms formed by the projections of the vectors $`(q,p),(q^{},p^{})`$ onto the $`n`$ pairs of coordinate planes. This is a sum of $`n`$ wedge-products. That is to say: if we write the standard basis elements as $`\frac{}{q^i}`$ and $`\frac{}{p_i}`$, this form is $`\omega :=\mathrm{\Sigma }_idq^idp_i`$. It has the action on $`\mathrm{I}\mathrm{R}^n\times \mathrm{I}\mathrm{R}^n`$:
$$(q^i\frac{}{q^i}+p_i\frac{}{p_i},q^i\frac{}{q^i}+p_i^{}\frac{}{p_i})\mathrm{\Sigma }_{i=1}^nq^ip_i^{}q^ip_i.$$
(6.113)
In general, if $`V,W`$ are two (real finite-dimensional) vector spaces, we define: $`L(V,W)`$ to be the vector space of linear maps from $`V`$ to $`W`$; $`L^k(V,W)`$ to be the vector space of $`k`$-multilinear maps from $`V\times V\times \mathrm{}.\times V`$ ($`k`$ copies) to $`W`$; and $`L_a^k(V,W)`$ to be the subspace of $`L^k(V,W)`$ consisting of (wholly) antisymmetric maps.
We then define $`\mathrm{\Omega }^k(V):=L_a^k(V,\mathrm{I}\mathrm{R})`$ for $`k=1,2,...,\mathrm{dim}(V)`$, so that $`\mathrm{\Omega }^1(V)=V^{}`$. We also set $`\mathrm{\Omega }^0(V):=\mathrm{I}\mathrm{R}`$. $`\mathrm{\Omega }^k(V)`$ is called the space of (exterior) $`k`$-forms on $`V`$. If dim$`(V)=n`$, then dim$`(\mathrm{\Omega }^k(V))=\left(\begin{array}{c}n\\ k\end{array}\right)`$.
The wedge-product, as defined above, can be extended to be an operation that defines, for $`\alpha \mathrm{\Omega }^k(V),\beta \mathrm{\Omega }^l(V)`$, an element $`\alpha \beta \mathrm{\Omega }^{k+l}(V)`$. We can skip the details: suffice it to say that the idea is to take tensor products as in (3) of Section 4.3.3, and anti-symmetrize.
But to complete our discussion of Noether’s theorem (in Section 6.5), we will need the definition of the contraction, (also known as: interior product), of a $`k`$-form $`\alpha \mathrm{\Omega }^k(V)`$ with a vector $`vV`$. We shall write this as $`𝐢_v\alpha `$. (It is also written with a hook notation.) We define the contraction $`𝐢_v\alpha `$ to be the $`(k1)`$-form given by:
$$𝐢_v\alpha (v_2,\mathrm{},v_k):=\alpha (v,v_2,\mathrm{},v_k).$$
(6.114)
It follows, for example, that contraction distributes over the wedge-product modulo a sign, in the following sense. If $`\alpha `$ is a $`k`$-form, and $`\beta `$ a 1-form, then
$$𝐢_v(\alpha \beta )=(𝐢_v\alpha )\beta +(1)^k\alpha (𝐢_v\beta ).$$
(6.115)
The direct sum of the vector spaces $`\mathrm{\Omega }^k(V),k=0,1,2,\mathrm{},\mathrm{dim}(V)=:n`$, has dimension $`2^n`$. When this direct sum is considered as equipped with the wedge-product $``$ and contraction $`𝐢`$, it is called the exterior algebra of $`V`$, written $`\mathrm{\Omega }(V)`$.
#### 6.2.2 Differential forms; the exterior derivative; the Poincaré Lemma
We extend the discussion given in Section 6.2.1 to a manifold $`M`$ of dimension $`n`$, taking all the tangent spaces $`T_x`$ at $`xM`$ as copies of the vector space $`V`$, and requiring fields of forms to be suitably smooth.
We begin by saying that a (smooth) scalar function $`f:M\mathrm{I}\mathrm{R}`$ is a 0-form field. Its differential or gradient, $`df`$, as defined by its action on all vector fields $`X`$, viz. mapping them to $`f`$’s directional derivative along $`X`$
$$df(X):=X(f)$$
(6.116)
is a 1-form (covector) field, called a differential 1-form.
The set $`(M)`$ of all smooth scalar functions forms an (infinite-dimensional) vector space, indeed a ring, under pointwise operations. We write the set of vector fields on $`M`$ as $`𝒳(M)`$, or as $`𝒯_0^1(M)`$; and the set of covector fields, i.e. differential 1-forms, on $`M`$ as $`𝒳^{}(M)`$, or as $`𝒯_1^0(M)`$. (So superscripts indicate the contravariant order, and subscripts the covariant order.)
Accordingly, we define: $`\mathrm{\Omega }^0(M):=(M)`$; $`\mathrm{\Omega }^1(M)=𝒯_1^0(M)`$; and so on. In short: $`\mathrm{\Omega }^k(M)`$ is the set of smooth fields of exterior $`k`$-forms on the tangent spaces of $`M`$.
The wedge-product, as defined in Section 6.2.1, can be extended to the various $`\mathrm{\Omega }^k(M)`$. We form the direct sum of the (infinite-dimensional) vector spaces $`\mathrm{\Omega }^k(M),k=0,1,2,\mathrm{},\mathrm{dim}(V)=:n`$, and consider it as equipped with this extended wedge-product. We call it the algebra of exterior differential forms on $`M`$, written $`\mathrm{\Omega }(M)`$.
Similarly, contraction, as defined in Section 6.2.1, can be extended to $`\mathrm{\Omega }(M)`$. On analogy with eq. 6.114, we define, for $`\alpha `$ a $`k`$-form field on $`M`$, and $`X`$ a vector field on $`M`$, the contraction $`𝐢_X\alpha `$ to be the $`(k1)`$-form given, at each point $`xM`$, by:
$$𝐢_X\alpha (x):(v_2,\mathrm{},v_k)\alpha (x)(X(x),v_2,\mathrm{},v_k)\mathrm{I}\mathrm{R}.$$
(6.117)
The exterior derivative is a differential operator on $`\mathrm{\Omega }(M)`$ that maps a $`k`$-form field to a $`(k+1)`$-form field. In particular, it maps a scalar $`f`$ to its differential (gradient) $`df`$. Indeed, it is the unique map from the $`k`$-form fields to the $`(k+1)`$-form fields ($`k=1,2,\mathrm{},n`$) that generalizes the elementary notion of gradient $`fdf`$, subject to certain natural conditions.
To be precise: one can show that there is a unique family of maps $`d^k:\mathrm{\Omega }^k(M)\mathrm{\Omega }^{k+1}(M)`$, all of which, for simplicity, we write as $`𝐝`$, such that:
(a): If $`f(M)`$, $`𝐝(f)=df`$.
(b): $`𝐝`$ is $`\mathrm{I}\mathrm{R}`$-linear; and distributes across the wedge-product, modulo a sign. That is: for $`\alpha \mathrm{\Omega }^k(M),\beta \mathrm{\Omega }^l(M)`$, $`𝐝(\alpha \beta )=(𝐝\alpha )\beta +(1)^k\alpha (𝐝\beta ).`$ (Cf. eq. 6.115.)
(c): $`𝐝^2:=𝐝𝐝0`$; i.e. for all $`\alpha \mathrm{\Omega }^k(M)`$ $`d^{k+1}d^k(\alpha )0`$. (This condition looks strong, but is in fact natural. For its motivation, it must here suffice to say that it generalizes the fact in elementary vector calculus, that the curl of any gradient is zero: $`(f)0`$.)
(d): $`𝐝`$ is a local operator; i.e. for any $`xM`$ and any $`k`$-form $`\alpha `$, $`𝐝\alpha (x)`$ depends only on $`\alpha `$’s restriction to any open neighbourhood of $`x`$; more precisely, we define for any open set $`U`$ of $`M`$, the vector space $`\mathrm{\Omega }^k(U)`$ of $`k`$-form fields on $`U`$, and then require that
$$𝐝(\alpha _U)=(𝐝\alpha )_U.$$
(6.118)
To express $`𝐝`$ in terms of coordinates: if $`\alpha \mathrm{\Omega }^k(M)`$, i.e. $`\alpha `$ is a $`k`$-form on $`M`$, given in coordinates by
$$\alpha =\alpha _{i_1\mathrm{}i_k}dx^{i_1}\mathrm{}dx^{i_k}(\mathrm{sum}\mathrm{on}i_1<i_2<\mathrm{}<i_k),$$
(6.119)
then one proves that the exterior derivative is
$$𝐝\alpha =\frac{\alpha _{i_1\mathrm{}i_k}}{x^j}dx^jdx^{i_1}\mathrm{}dx^{i_k}(\mathrm{sum}\mathrm{on}\mathrm{all}j\mathrm{and}i_1,\mathrm{}<i_k),$$
(6.120)
We define $`\alpha \mathrm{\Omega }^k(M)`$ to be:
exact if there is a $`\beta \mathrm{\Omega }^{k1}(M)`$ such that $`\alpha =𝐝\beta `$; (cf. the elementary definition of an exact differential);
closed if $`𝐝\alpha =0`$.
It is immediate from condition (c) above, $`𝐝^2=0`$, that every exact form is closed. The converse is “locally true”. This important result is the Poincaré Lemma; (and we will use it in Section 6.5’s closing discussion of Noether’s theorem).
To be precise: for any open set $`U`$ of $`M`$, we define (as in condition (d) above) the vector space $`\mathrm{\Omega }^k(U)`$ of $`k`$-form fields on $`U`$. Then the Poincaré Lemma states that if $`\alpha \mathrm{\Omega }^k(M)`$ is closed, then at every $`xM`$ there is a neighbourhood $`U`$ such that $`\alpha _U\mathrm{\Omega }^k(U)`$ is exact.
We will also need (again, for Section 6.5’s discussion of Noether’s theorem) a useful formula relating the Lie derivative, contraction and the exterior derivative. Namely: Cartan’s magic formula, which says that if $`X`$ is a vector field and $`\alpha `$ a $`k`$-form on a manifold $`M`$, then the Lie derivative of $`\alpha `$ with respect to $`X`$ (i.e. along the flow of $`X`$) is
$$_X\alpha =\mathrm{𝐝𝐢}_X\alpha +𝐢_X𝐝\alpha .$$
(6.121)
This is proved by straightforward calculation.
### 6.3 Symplectic manifolds; the cotangent bundle as a symplectic manifold
Any cotangent bundle $`T^{}Q`$ has a natural symplectic structure, which is the geometric structure on manifolds corresponding to the symplectic matrix $`\omega `$ introduced by eq. 4.58, and to the symplectic forms on vector spaces defined at the end of Section 4.3.3. (Here ‘natural’ means intrinsic, and in particular, independent of a choice of coordinates or bases.) It is this structure that enables a scalar function to determine a dynamics. That is: the symplectic structure implies that any scalar function $`H:T^{}Q\mathrm{I}\mathrm{R}`$ defines a vector field $`X_H`$ on $`T^{}Q`$.
I first describe this structure (Section 6.3.1), and then show that any cotangent bundle has it (Section 6.3.2). Later subsections will develop the consequences.
#### 6.3.1 Symplectic manifolds
A symplectic structure or symplectic form on a manifold $`M`$ is defined to be a differential 2-form $`\omega `$ on $`M`$ that is closed (i.e. $`𝐝\omega =0`$) and non-degenerate. That is: for any $`xM`$, and any two tangent vectors at $`x`$, $`\sigma ,\tau T_x`$:
$$𝐝\omega =0\text{ and }\tau 0,\sigma :\omega (\tau ,\sigma )0.$$
(6.122)
Such a pair $`(M,\omega )`$ is called a symplectic manifold.
There is a rich theory of symplectic manifolds; but we shall only need a small fragment of it, building on our discussion in Section 4.3.3. (In particular, the fact that we mostly avoid the theory of canonical transformations means we will not need the theory of Lagrangian sub-manifolds.)
First, it follows from the non-degeneracy of $`\omega `$ that $`M`$ is even-dimensional; (cf. eq. 4.86).
It also follows that at any $`xM`$, there is a basis-independent isomorphism $`\omega ^{\mathrm{}}`$ from the tangent space $`T_x`$ to its dual $`T_x^{}`$. We saw this in (2) and (4) of Section 4.3.3, especially eq. 4.71. Namely: for any $`xM`$ and $`\tau T_x`$, the value of the 1-form $`\omega ^{\mathrm{}}(\tau )T_x^{}`$ is defined by
$$\omega ^{\mathrm{}}(\tau )(\sigma ):=\omega (\sigma ,\tau )\sigma T_x.$$
(6.123)
Here we return to the main idea emphasised already in Section 4.3.1: that symplectic structure enables a covector field, i.e. a differential one-form, to determine a vector field. Thus for any function $`H:M\mathrm{I}\mathrm{R}`$, so that $`dH`$ is a differential 1-form on $`M`$, the inverse of $`\omega ^{\mathrm{}}`$ (which we might write as $`\omega ^{\mathrm{}}`$), carries $`dH`$ to a vector field on $`M`$, written $`X_H`$. Cf. eq. 4.62.
So far, we have noted some implications of $`\omega `$ being non-degenerate. The other part of the definition of a symplectic form (for a manifold), viz. $`\omega `$ being closed, $`𝐝\omega =0`$, is also important. We shall see in Section 6.5 that it implies that a vector field $`X`$ on a symplectic manifold $`M`$ preserves the symplectic form $`\omega `$ (i.e. in more physical jargon: generates (a one-parameter family of) canonical transformations) iff $`X`$ is Hamiltonian in the sense of Section 5.2; i.e. there is a scalar function $`f`$ such that $`X=X_f\omega ^{\mathrm{}}(df)`$. Or in terms of the Poisson bracket, with $``$ representing the argument place for a scalar function: $`X()=X_f()\{,f\}`$.
So much by way of introducing symplectic manifolds. I turn to showing that any cotangent bundle $`T^{}Q`$ is such a manifold.
#### 6.3.2 The cotangent bundle
Choose any local coordinates $`q`$ on $`Q`$ (dim($`Q`$)=$`n`$), and the natural local coordinates $`q,p`$ thereby induced on $`T^{}Q`$; (cf. (B) of Section 6.1). We define the 2-form
$$dpdq:=dp_idq^i:=\mathrm{\Sigma }_{i=1}^ndp_idq^i.$$
(6.124)
To show that eq. 6.124 defines the same 2-form, whatever choice we make of the chart $`q`$ on $`Q`$, it suffices to show that $`dpdq`$ is the exterior derivative of a 1-form on $`T^{}Q`$ which is defined naturally (i.e. independently of coordinates or bases) from the derivative (also known as: tangent) map of the projection
$$\pi :(q,p)T^{}QqQ.$$
(6.125)
Thus consider a tangent vector $`\tau `$ (not to $`Q`$, but) to the cotangent bundle $`T^{}Q`$ at a point $`\eta =(q,p)T^{}Q,`$ i.e. $`qQ`$ and $`pT_q^{}`$. Let us write this as: $`\tau T_\eta (T^{}Q)T_{(q,p)}(T^{}Q)`$. The derivative map, $`D\pi `$ say, of the natural projection $`\pi `$ applies to $`\tau `$:
$$D\pi :\tau T_{(q,p)}(T^{}Q)(D\pi (\tau ))T_q.$$
(6.126)
Now define a 1-form $`\theta _H`$ on $`T^{}Q`$ by
$$\theta _H:\tau T_{(q,p)}(T^{}Q)p(D\pi (\tau ))\mathrm{I}\mathrm{R};$$
(6.127)
where in this definition of $`\theta _H`$, $`p`$ is defined to be the second component of $`\tau `$’s base-point $`(q,p)T^{}Q`$; i.e. $`\tau T_{(q,p)}(T^{}Q)`$ and $`pT_q^{}`$.
This 1-form is called the canonical 1-form on $`T^{}Q`$. It is the “Hamiltonian version” of the 1-form $`\theta _L`$ defined by eq. 2.13; and also there called the ‘canonical 1-form’. But Section 6.1’s discussion of the “fruitful ambiguity” of the symbol $`p`$ brings out a contrast. While $`\theta _L`$ as defined by eq. 2.13 clearly depends on $`L`$, the definition of $`\theta _H`$, eq. 6.127, does not depend on any function $`H`$. $`\theta _H`$ is given just by the cotangent bundle structure. Hence the subscript $`H`$ here just indicates “Hamiltonian (as against Lagrangian) version”, not dependence on a function $`H`$.
So much by way of a natural definition of a 1-form. One now checks that in any natural local coordinates $`q,p`$, $`\theta _H`$ is given by
$$\theta _H=p_idq^i.$$
(6.128)
Finally, we define a 2-form by taking the exterior derivative of $`\theta _H`$:
$$𝐝(\theta _H):=𝐝(p_idq^i)dp_idq^i.$$
(6.129)
where the last equation follows immediately from eq. 6.120. One checks that this 2-form is closed (since $`𝐝^2=0`$) and non-degenerate. So $`(T^{}Q,𝐝(\theta _H))`$ is a symplectic manifold.
Referring to eq. 4.66 of Section 4.3, or eq. 4.87 of Section 4.3.3, or eq. 6.113 of Section 6.2, we see that at each point $`(q,p)T^{}Q`$, this symplectic form is, upto a sign, our familiar “sum of signed areas”—first seen as induced by the matrix $`\omega `$ of eq. 4.58.
Accordingly, Section 4.3.3’s definition of a canonical symplectic form is extended to the present case: $`𝐝(\theta _H)`$, or its negative $`𝐝(\theta _H)`$, is called the canonical symplectic form, or canonical 2-form. (The difference from Section 4.3.3’s definition is that on a manifold, the symplectic form is required to be closed.)
(The difference by a sign is of course conventional: it arises from our taking the $`q`$s, not the $`p`$s, as the first $`n`$ out of the $`2n`$ coordinates. For if we had instead taken the $`p`$s, the matrix occurring in eq. 4.60 would have been $`\omega \omega ^1`$: exactly matching the cotangent bundle’s intrinsic 2-form $`𝐝(\theta _H)`$.)
We will see, in Section 6.6, a theorem (Darboux’s theorem) to the effect that locally, any symplectic manifold “looks like” a cotangent bundle: or in other words, a cotangent bundle is locally a “universal” example of symplectic structure. But first we return, in the next two Subsections, to Hamilton’s equations, and Noether’s theorem.
### 6.4 Geometric formulations of Hamilton’s equations
We already emphasised in Sections 4.3 and 5 the main geometric idea behind Hamilton’s equations: that a gradient, i.e. covector, field $`dH`$ determines a vector field $`X_H`$. We first saw this determination via the symplectic matrix, in eq. 4.62 of Section 4.3.1, viz.
$$X_H(z)=\omega H(z);$$
(6.130)
and then via the Poisson bracket, in eq. 5.103 of Section 5.2, viz.
$$D:=X_H=\frac{d}{dt}=\dot{q}^i\frac{}{q^i}+\dot{p}_i\frac{}{p_i}=\frac{H}{p_i}\frac{}{q^i}\frac{H}{q^i}\frac{}{p_i}=\{,H\}.$$
(6.131)
The symplectic structure and Poisson bracket were related by eq. 5.96, viz.
$$\{f,g\}(z)=\stackrel{~}{f}(z).\omega .g(z).$$
(6.132)
And to this earlier discussion, the last Subsection, Section 6.3, added the identification of the canonical symplectic form of a cotangent bundle, eq. 6.129.
Let us sum up these discussions by giving some geometric formulations of Hamilton’s equations at a point $`z=(q,p)`$ in a cotangent bundle $`T^{}Q`$. Let us write $`\omega ^{\mathrm{}}`$ for the (basis-independent) isomorphism from the cotangent space to the tangent space, $`T_z^{}T_z`$, induced by $`\omega :=𝐝(\theta _H)=dq^idp_i`$ (cf. eq. 4.83 and 6.123). Then Hamilton’s equations, eq. 4.62 or 6.130, may be written as:
$$\dot{z}=X_H(z)=\omega ^{\mathrm{}}(𝐝H(z))=\omega ^{\mathrm{}}(dH(z)).$$
(6.133)
Applying $`\omega ^{\mathrm{}}`$, the inverse isomorphism $`T_zT_z^{}`$, to both sides, we get
$$\omega ^{\mathrm{}}X_H(z)=dH(z).$$
(6.134)
In terms of the symplectic form $`\omega `$ at $`z`$, this is (cf. eq. 4.71): for all vectors $`\tau T_z`$
$$\omega (X_H(z),\tau )=dH(z)\tau ;$$
(6.135)
or in terms of the contraction defined by eq. 6.114, with $``$ marking the argument place of $`\tau T_z`$:
$$𝐢_{X_H}\omega :=\omega (X_H(z),)=dH(z)().$$
(6.136)
More briefly, and now for any function $`f`$, it is:
$$𝐢_{X_f}\omega =df.$$
(6.137)
Here is a final example. Recall the relation between the Poisson bracket and the directional derivative (or the Lie derivative $``$) of a function, eq. 5.104 and 6.131: viz.
$$_{X_f}g=dg(X_f)=X_f(g)=\{g,f\}.$$
(6.138)
Combining this with eq. 6.137, we can reformulate the relation between the symplectic form and Poisson bracket, eq. 6.132, in the form:
$$\{g,f\}=dg(X_f)=𝐢_{X_f}dg=𝐢_{X_f}(𝐢_{X_g}\omega )=\omega (X_g,X_f).$$
(6.139)
### 6.5 Noether’s theorem completed
The discussion of Noether’s theorem in Section 5.3 left unfinished business: to prove that a vector field generates a one-parameter family of canonical transformations iff it is a Hamiltonian vector field (and so justify the third claim of Section 5.3.1). Cartan’s magic formula and the Poincaré Lemma, both from Section 6.2, make it easy to prove this, for a vector field on any symplectic manifold $`(M,\omega )`$. ($`(M,\omega )`$ need not be a cotangent bundle.)
We define a vector field $`X`$ on a symplectic manifold $`(M,\omega )`$ to be symplectic (also known as: canonical) iff the Lie-derivative along $`X`$ of the symplectic form vanishes, i.e. $`_X\omega =0`$.<sup>18</sup><sup>18</sup>18As announced in Section 2.2.1, I assume the notion of the Lie-derivative, in particular the Lie-derivative of a 2-form. Suffice it to say, as a sketch, that the flow of $`X`$ defines a map on $`M`$ which induces a map on curves, and so on vectors, and so on co-vectors, and so on 2-forms such as $`\omega `$. Nor will I go into details about the equivalence between this definition of $`X`$’s being symplectic, and $`X`$’s generating (active) canonical transformations, or preserving the Poisson bracket. For as I have emphasised, I will not need to develop the theory of canonical transformations.
Since $`\omega `$ is closed, i.e. $`𝐝\omega =0`$, Cartan’s magic formula, eq. 6.121, applied to $`\omega `$ becomes
$$_X\omega \mathrm{𝐝𝐢}_X\omega +𝐢_X𝐝\omega =\mathrm{𝐝𝐢}_X\omega .$$
(6.140)
So for $`X`$ to be symplectic is for $`𝐢_X\omega `$ to be closed. But by the Poincaré Lemma, if $`𝐢_X\omega `$ is closed, it is locally exact. That is: there locally exists a scalar function $`f:M\mathrm{I}\mathrm{R}`$ such that
$$𝐢_X\omega =df\mathrm{i}.\mathrm{e}.X=X_f.$$
(6.141)
So for $`X`$ to be symplectic is equivalent to $`X`$ being locally Hamiltonian.
So we can sum up Noether’s theorem from a geometric perspective, as follows. We define a Hamilton system to be a triple $`(M,\omega ,H)`$ where $`(M,\omega )`$ is a symplectic manifold and $`H:M\mathrm{I}\mathrm{R}`$, i.e. $`M(M)`$. We define a (continuous) symmetry of a Hamiltonian system to be a vector field $`X`$ on $`M`$ that preserves both the symplectic form, $`_X\omega =0`$, and the Hamiltonian function, $`_XH=0`$. As we have just seen: for any symmetry so defined, there locally exists an $`f`$ such that $`X=X_f`$. So we can apply the “one-liner”, eq. 5.107, i.e. the antisymmetry of the Poisson bracket,
$$X_f(H)\{H,f\}=0\mathrm{iff}X_H(f)\{f,H\}=0,$$
(6.142)
to conclude that $`f`$ is a first integral (constant of the motion). Thus we have
> Noether’s theorem for a Hamilton system If $`X`$ is a symmetry of a Hamiltonian system $`(M,\omega ,H)`$, then locally $`X=X_f`$ and $`f`$ is a constant of the motion. And conversely: if $`f:M\mathrm{I}\mathrm{R}`$ is a constant of the motion, then $`X_f`$ is a symmetry. Besides, this result encompasses the Lagrangian version of the theorem; cf. Sections 3.4 and 5.3.
Example:— For most Hamiltonian systems in euclidean space $`\mathrm{I}\mathrm{R}^3`$, spatial translations and rotations are (continuous) symmetries. For example, consider $`N`$ point-particles interacting by Newtonian gravity. The Hamiltonian is a sum of two terms, which are each individually invariant under these euclidean motions:
(i) a kinetic energy term $`K`$; though I will not go into details, it is in fact defined by the euclidean metric of $`\mathrm{I}\mathrm{R}^3`$ (cf. footnote 4 in Section 2.1), and is thereby invariant; and
(ii) a potential energy term $`V`$; it depends only on the particles’ relative distances, and is thereby invariant.
The corresponding conserved quantities are the total linear and angular momentum.<sup>19</sup><sup>19</sup>19By the way, this Hamiltonian is not invariant under boosts. But as I said in Section 2.2.1 and footnote 8, I restrict myself to time-independent transformations; the treatment of symmetries that “represent the relativity of motion” needs separate discussion.
Finally, an incidental remark which relates to the “rectification theorem”, that on any manifold any vector field $`X`$ can be “straightened out” in a neighbourhood around any point at which $`X`$ is non-zero, so as to have all but one component vanish and the last component equal to 1; cf. eq. 3.39. Using this theorem, it is easy to see that on any even-dimensional manifold any vector field $`X`$ is locally Hamiltonian, with respect to some symplectic form, around a point where $`X`$ is non-zero. (One defines the symplectic form by Lie-dragging from a surface transverse to $`X`$’s integral curves.)
### 6.6 Darboux’s theorem, and its role in reduction
Darboux’s theorem states that cotangent bundles are, locally, a “universal form” of symplectic manifold. That is: Not only is any symplectic manifold ($`M,\omega `$) even-dimensional. Also, it “looks locally like” a cotangent bundle, in that around any $`x`$ in $`M`$, there is a local coordinate system $`(q^1,\mathrm{},q^n;p_1,\mathrm{},p_n)`$—where the use of both upper and lower indices is now just conventional, with no meaning about dual bases!—in which:
(i) $`\omega `$ takes the form $`dq^idp_i`$; and so
(ii) the Poisson brackets of the $`q`$s and $`p`$s take the fundamental form in eq. 5.102.
(The theorem generalizes to the Poisson manifolds mentioned in Section 6.8.)
Besides, the proof of Darboux’s theorem yields further information: information which is important for reducing problems. It arises from the beginning of the proof; and will return us to Section 4.2’s point that the elementary connection between cyclic coordinates and conserved conjugate momenta underpins the role of symmetries and conserved quantities in reductions on symplectic manifolds.
(In fact, Darboux’s theorem also yields two other broad implications about reducing problems; but I will not develop the details here. The second implication concerns the way that a Hamiltonian structure is preserved in the reduced problem. The third implication concerns the requirement that constants of the motion be in involution, i.e. have vanishing Poisson bracket with each other; so it leads to the idea of complete integrability—a topic this paper foreswears.)
Namely, the proof implies that “almost” any scalar function $`f(M)`$ can be taken as the first “momentum” coordinate $`p_1`$; or as the first configurational coordinate $`q^1`$. Here “almost” is not meant in a measure-theoretic sense; it is just that $`f`$ is subject to a mild restriction, that $`df0`$ at the point $`xM`$.
In a bit more detail: The proof of Darboux’s theorem starts by taking any such $`f`$ to be our $`p_1`$, and then constructs the canonically conjugate generalized coordinate $`q^1`$, i.e. the coordinate such that $`\{q^1,p_1\}=1`$: so that $`p_1`$ generates translation in the direction of increasing $`q^1`$. Indeed the construction is geometrically clear. The symplectic structure means that any such $`f`$ defines a Hamiltonian vector field $`X_f`$, and a flow $`\varphi ^f`$. We choose a $`(2n1)`$-dimensional local submanifold $`N`$ passing through the given point $`x`$, and transverse to all the integral curves of $`X_f`$ in a neighbourhood of $`x`$; and we set the parameter $`\lambda `$ of the flow $`\varphi ^f`$ to be zero at all points $`yN`$. Then for any $`z`$ in a suitably small neighbourhood of the given point $`x`$, we define the function $`q^1(z)`$ to be the parameter-value at $`z`$ of the integral curve of $`X_f`$ that passes through $`z`$. So by construction, (i) $`f`$ generates translation in the direction of increasing $`q^1`$, and (ii) defining $`p_1:=f`$, we have $`\{q^1,p_1\}=1`$.
This is just the beginning of the proof. But I will not need details of how it goes on to establish the local existence of canonical coordinates, i.e. coordinates such that analogues of (i) and (ii), also for $`i1`$, hold. In short, the strategy is to use induction on the dimension of the manifold; for details, cf. e.g. Arnold (1989: 230-232).
To see the significance of this for reducing problems, suppose that there is a constant of the motion, and that we take it as our $`f`$, i.e. as the first momentum coordinate $`p_1`$. So the system evolves on a $`(2n1)`$-dimensional manifold given by an equation $`f=`$ constant. So writing $`H`$ in the canonical coordinate system secured by Darboux’s theorem, we conclude that $`0=\dot{f}\frac{H}{q^1}`$. That is, $`q^1`$ is cyclic. So as discussed in Section 4.2, we need only solve the problem in the $`2n2`$ variables $`q^2,\mathrm{},q^n;p_2,\mathrm{},p_n`$. Having done so, we can find $`q^1`$ as a function of time, by solving eq. 4.57 by quadrature.
To put the point in geometric terms:—
(i): The system is confined to a $`(2n1)`$-dimensional manifold $`p_1=\alpha =`$ constant, $`M_\alpha `$ say.
(ii): $`M_\alpha `$ is foliated by a local one-parameter family of $`(2n2)`$-dimensional manifolds labelled by values of $`q^1I\mathrm{I}\mathrm{R}`$, $`M_\alpha =_{q^1I}M_{\alpha ,q^1}`$.
(iii): Of course, the dynamical vector field is transverse to the leaves of this foliation; i.e. $`q^1`$ is not a constant of the motion, $`\dot{q}^10`$. But since $`q^1`$ is ignorable, $`\frac{H}{q^1}=0`$, the problem to be solved is “the same” at points $`x_1,x_2`$ that differ only in their values of $`q^1`$.
### 6.7 Geometric formulation of the Legendre transformation
Let us round off our development of both Lagrangian and Hamiltonian mechanics, by formulating the Legendre transformation as a map from the tangent bundle $`TQ`$ to the cotangent bundle $`T^{}Q`$. In this formulation, the Legendre transformation is often called the fibre derivative.
Again, there is a rich theory to be had here. In part, it relates to the topics mentioned in Section 4.2.3: (i) the description of a function (in the simplest case $`f:\mathrm{I}\mathrm{R}\mathrm{I}\mathrm{R}`$) by its gradients and axis-intercepts, rather than by its arguments and values; (ii) variational principles. But I shall not go into details about this theory: since this paper emphasises the Hamiltonian framework, a mere glimpse of this theory must suffice. (References, additional to those in Section 4.2.3, include: Abraham and Marsden (1978: Sections 3.6-3.8) and Marsden and Ratiu (1999: Sections 7.2-7.5, 8.1-8.3).)
Let us return to the Lagrangian framework. We stressed in Section 2.2 that a scalar on the tangent bundle, the Lagrangian $`L:TQ\mathrm{I}\mathrm{R}`$, “determines everything”: the dynamical vector field $`D=:D_L`$; and so for given initial $`q`$ and $`\dot{q}`$, $`L`$ determines a solution, a trajectory in $`TQ`$, i.e. $`2n`$ functions of time $`q(t),\dot{q}(t)`$ with the first $`n`$ functions determining the latter.
For the Legendre transformation, the fundamental points are that:
(1): $`L`$ also determines at any point $`qQ`$, a preferred map $`FL_q`$ from the tangent space $`T_q`$ to its dual space $`T_q^{}`$. Besides this preferred map:
(2): extends trivially to a preferred map from all of $`TQ`$ to $`T^{}Q`$; this is the Legendre transformation, understood geometrically;
(3): extends, under some technical conditions (about certain kinds of uniqueness, invertibility and smoothness), so as to carry geometric objects of various sorts defined on $`TQ`$ to corresponding objects defined on $`T^{}Q`$, and vice versa.
So under these conditions, the Legendre transformation (together with its inverse) transfers the entire description of the system’s motion between the Lagrangian and Hamiltonian frameworks.
I will explain (1) and (2), but just gesture at (3).
(1): Intuitively, the preferred map $`FL_q`$ from each tangent space $`T_q`$ to its dual space $`T_q^{}`$ is the transition $`\dot{q}p`$. More precisely: since $`L`$ is a scalar on $`TQ`$, any choice of local coordinates $`q`$ on a patch of $`Q`$, together with the induced local coordinates $`q,\dot{q}`$ on a patch of $`TQ`$, defines the partial derivatives $`\frac{L}{\dot{q}}`$. At any point $`q`$ in the domain of the local coordinates, this defines a preferred map $`FL_q`$ from the tangent space $`T_q`$ to the dual space $`T_q^{}`$: $`FL_q:T_qT_q^{}`$. Namely, a vector $`\tau T_q`$ with components $`\dot{q}^i`$ in the coordinate system $`q^i`$ on $`Q`$, i.e. $`\tau =\dot{q}^i\frac{}{q^i}`$ (think of a motion through configuration $`q`$ with generalized velocity $`\tau `$) is mapped to the 1-form whose components in the dual basis $`dq^i`$ are $`\frac{L}{\dot{q}^i}`$. That is
$$FL_q:\tau =\dot{q}^i\frac{}{q^i}T_q\frac{L}{\dot{q}^i}dq^iT_q^{}.$$
(6.143)
One easily checks that because the canonical momenta are a 1-form, this definition is, despite appearances, coordinate-independent.
(2): An equivalent definition, manifestly coordinate-independent and given for all $`qQ`$, is as follows. Given $`L:TQ\mathrm{I}\mathrm{R}`$, define $`FL:TQT^{}Q`$, the fibre derivative, by
$$qQ,\sigma ,\tau T_q:FL(\sigma )\tau =\frac{d}{ds}_{s=0}L(\sigma +s\tau )$$
(6.144)
(We here take $`\sigma ,\tau `$ to encode the identity of the base-point $`q`$, so that we make notation simpler, writing $`FL(\sigma )`$ rather than $`FL((q,\sigma ))`$ etc.) That is: $`FL(\sigma )\tau `$ is the derivative of $`L`$ at $`\sigma `$, along the fibre $`T_q`$ of the fibre bundle $`TQ`$, in the direction $`\tau `$. So $`FL`$ is fibre-preserving: i.e. it maps the fibre $`T_q`$ of $`TQ`$ to the fibre $`T_q^{}`$ of $`T^{}Q`$. In local coordinates $`q,\dot{q}`$ on $`TQ`$, $`FL`$ is given by:
$$FL(q^i,\dot{q}^i)=(q^i,\frac{L}{\dot{q}^i});\mathrm{i}.\mathrm{e}.p_i=\frac{L}{\dot{q}^i}.$$
(6.145)
An important special case involves a free system (i.e. no potential term in the Lagrangian) and a configuration manifold $`Q`$ with a metric $`g=g_{ij}`$ defined by the kinetic energy. (Cf. footnote 4 for the definition of this metric: in short, the constraints being scleronomous (i.e. time-independent, cf. Section 2.1), implies that for any coordinate system on $`Q`$, the kinetic energy is a homogeneous quadratic form in the generalized velocities.) The Lagrangian is then just the kinetic energy of the metric,
$$L(q,\dot{q})L(\dot{q}):=\frac{1}{2}g_{ij}\dot{q}^i\dot{q}^j$$
(6.146)
so that the fibre derivative is given by
$$FL(\sigma )\tau =g(\sigma ,\tau )=g_{ij}\sigma ^i\tau ^j,\mathrm{i}.\mathrm{e}.p_i=g_{ij}\dot{q}^j.$$
(6.147)
(3): We can use $`FL`$ to pull-back to $`TQ`$ the canonical 1-form $`\theta \theta _H`$ and symplectic form $`\omega `$ from $`T^{}Q`$ (eq. 6.127 and 6.128 with $`\omega =𝐝\theta `$, from Section 6.3.B). That is, we can define
$$\theta _L:=(FL)^{}\theta _H\mathrm{and}\omega _L:=(FL)^{}\omega .$$
(6.148)
Since exterior differentiation $`𝐝`$ commutes with pull-backs, $`\omega _L=𝐝\theta _L`$. Furthermore:
(i): As one would hope, $`\theta _L`$, so defined, is Lagrangian mechanics’ canonical 1-form, which we already defined in eq. 2.13 (and which played a central role in the Lagrangian version of Noether’s theorem).
(ii): One can show that $`\omega _L`$ is non-degenerate iff the Hessian condition eq. 2.3 holds. So under this condition, we can analyse Lagrangian mechanics in terms of symplectic structure.
Given $`L`$, we define its energy function $`E:TQ\mathrm{I}\mathrm{R}`$ by
$$v(q,\tau )TQ,E(v):=FL(v)vL(v);$$
(6.149)
or in coordinates
$$E(q^i,\dot{q}^i):=\frac{L}{\dot{q}^i}\dot{q}^iL(q^i,\dot{q}^i)$$
(6.150)
If $`FL`$ is a diffeomorphism, we find that $`E(FL)^1`$ is, as one would hope, the Hamiltonian function $`H:T^{}Q\mathrm{I}\mathrm{R}`$ which we already defined in eq. 4.52.
And accordingly, if $`FL`$ is a diffeomorphism, then the derivative of $`FL`$ carries the dynamical vector field $`\frac{d}{dt}`$ in the Lagrangian description, as defined in eq. 2.8 (Section 2.2, (2)), viz.
$$D_L:=\dot{q}^i\frac{}{q^i}+\ddot{q}^i\frac{}{\dot{q}^i},$$
(6.151)
to the Hamiltonian dynamical vector field, viz.
$$D_H:=\dot{q}^i\frac{}{q^i}+\dot{p}_i\frac{}{p_i}.$$
(6.152)
More generally, one can show if $`FL`$ is a diffeomorphism, there is a bijective correspondence between the various geometric structures used in the Lagrangian and Hamiltonian descriptions. For precise statements of this idea, cf. e.g. Abraham and Marsden (1978: Theorem 3.6.9) and Marsden and Ratiu (1999: Theorem 7.4.3.), and their preceding discussions.
### 6.8 Glimpsing the more general framework of Poisson manifolds
Recall that Section 5.1 listed several properties of the Poisson bracket, as defined by eq. 5.92 or 5.95. We end by briefly describing how the postulation of a bracket that acts on the scalar functions $`F:M\mathrm{I}\mathrm{R}`$ defined on any manifold $`M`$, and possesses four of Section 5.1’s listed properties, provides a sufficient framework for mechanics in Hamiltonian style. The bracket is again called a ‘Poisson bracket’, and the manifold $`M`$ equipped with such a bracket is called a Poisson manifold.
Namely, we require the following four properties. The Poisson bracket is to be bilinear; antisymmetric; and to obey the Jacobi identity (eq. 5.100) for any real functions $`F,G,H`$ on $`M`$, i.e.
$$\{\{F,H\},G\}+\{\{G,F\},H\}+\{\{H,G\},F\}=0;$$
(6.153)
and to obey Leibniz’ rule for products (eq. 5.98), i.e.
$$\{F,HG\}=\{F,H\}G+H\{F,G\}.$$
(6.154)
This generalizes Hamiltonian mechanics: in particular, a Poisson manifold need not be a symplectic manifold. The main idea of the extra generality is that the antisymmetric bilinear map that gives the geometry of the state space (the analogue of Section 4.3’s symplectic form $`\omega `$) can be degenerate. So this map can “have extra zeroes”, as in eq. 4.85 and 4.86. (This map is induced by the generalized Poisson bracket, via an analogue of eq. 5.96.) This means that a Poisson manifold can have odd dimension; while we saw in Section 4.3.3 that any symplectic vector space is even-dimensional—and so, therefore, is any symplectic manifold (Section 6.3.1 and 6.6).
On the other hand, the generalized framework has strong connections with the usual one.<sup>20</sup><sup>20</sup>20Because of these connections, it is natural to still call the more general framework ‘Hamiltonian’; as is usually done. But of course this is just a verbal matter. One main connection is the result that any Poisson manifold $`M`$ is a disjoint union of even-dimensional manifolds, on each of which $`M`$’s degenerate antisymmetric bilinear form (induced by the generalized Poisson bracket) restricts to be non-degenerate; so that there is an orthodox Hamiltonian mechanics on each such ‘symplectic leaf’. Another main connection is that Section 5.3’s “one-liner” version of Noether’s theorem, eq. 5.107, underpins versions of Noether’s theorem for the more general framework.
This generalized framework is important for various reasons; I will just mention two.
(i): For a system whose orthodox Hamiltonian mechanics on a symplectic manifold (dimension $`2n`$, say) depends on $`s`$ real parameters, it is sometimes natural to consider the corresponding $`(2n+s)`$-dimensional space. This is often a Poisson manifold; viz., one foliated into an $`s`$-dimensional family of $`2n`$-dimensional symplectic manifolds. This scenario occurs even for some very familiar systems, such as the pivoted rigid body described by Euler’s equations.
(ii): Poisson manifolds often arise in the theory of symplectic reduction. For when you quotient a symplectic manifold by the action of a group (e.g. a group of symmetries of a Hamiltonian system in the sense of Section 6.5), you often get a Poisson manifold, rather than a symplectic one. Indeed, the pivoted rigid body is itself an example of this.
But this generalized framework is a large topic, which we cannot go into: as mentioned, Butterfield (2006) is a philosopher’s introduction.
For now, we end with a historical point.<sup>21</sup><sup>21</sup>21As mentioned in footnote 10, Olver (2000) gives many details especially about Lie; e.g. Olver (2000: 374-379, 427-428). Cf. also Marsden and Ratiu (1999: 336-338, 430-432), and for a full history, Hawkins (2000). It is humbling, but also I hope inspiring, reflection about one of classical mechanics’ monumental figures. Namely: a considerable part of the modern theory of Poisson manifolds, including their uses for the rigid body and for symplectic reduction, was already contained in Lie (1890)!
Acknowledgements:— I am grateful to the editors, not least for their patience; to audiences in Irvine, Oxford, Princeton and Santa Barbara; and to Katherine Brading, Harvey Brown, Hans Halvorson, David Malament, Wayne Myrvold, David Wallace, and especially Graeme Segal, for conversations, comments—and corrections!
## 7 References
R. Abraham and J. Marsden (1978), Foundations of Mechanics, second edition: Addison-Wesley.
V. Arnold (1973), Ordinary Differential Equations, MIT Press.
V. Arnold (1989), Mathematical Methods of Classical Mechanics, Springer, (second edition).
G. Belot (2003), ‘Notes on symmetries’, in Brading and Castellani (ed.s) (2003), pp. 393-412.
K. Brading and E. Castellani (ed.s) (2003), Symmetry in Physics, Cambridge University Press.
H. Brown and P. Holland (2004), ‘Simple applications of Noether’s first theorem in quantum mechanics and electromagnetism‘, American Journal of Physics 72 p. 34-39. Available at: http://arxiv.org/abs/quant-ph/0302062 and http://philsci-archive.pitt.edu/archive/00000995/
H. Brown and P. Holland (2004a), ‘Dynamical vs. variational symmetries: Understanding Noether’s first theorem’, Molecular Physics, 102, (11-12 Special Issue), pp. 1133-1139.
J. Butterfield (2004), ‘Some Aspects of Modality in Analytical mechanics’, in Formal Teleology and Causality, ed. M. Stöltzner, P. Weingartner, Paderborn: Mentis.
Available at Los Alamos arXive: http://arxiv.org/abs/physics/0210081 or
http://xxx.soton.ac.uk/abs/physics/0210081; and at Pittsburgh archive: http://philsci-archive.pitt.edu/archive/00001192.
J. Butterfield (2004a), ‘Between Laws and Models: Some Philosophical Morals of Lagrangian Mechanics’; available at Los Alamos arXive: http://arxiv.org/abs/physics/0409030 or
http://xxx.soton.ac.uk/abs/physics/0409030; and at Pittsburgh archive: http://philsci-archive.pitt.edu/archive/00001937/.
J. Butterfield (2004b), ‘On Hamilton-Jacobi Theory as a Classical Root of Theory’, in A. Elitzur, S. Dolev and N. Kolenda (eds.), Quo Vadis Quantum Mechanics?, Springer, pp. 239-273; available at Los Alamos arXive: http://arxiv.org/abs/quant-ph/0210140; or at Pittsburgh archive: http://philsci-archive.pitt.edu/archive/00001193/
J. Butterfield (2005), ‘Between Laws and Models: Some Philosophical Morals of Hamiltonian Mechanics’, in preparation.
J. Butterfield (2006), ‘On Symplectic Reduction in Classical Mechanics’, forthcoming in The North Holland Handbook of Philosophy of Physics, ed. J. Earman and J. Butterfield, North Holland.
R. Courant and D. Hilbert (1953), Methods of Mathematical Physics, volume I, Wiley-Interscience (Wiley Classics 1989).
R. Courant and D. Hilbert (1962), Methods of Mathematical Physics, volume II, Wiley-Interscience (Wiley Classics 1989).
E. Desloge (1982), Classical Mechanics, John Wiley.
J. Earman (2003), ‘Tracking down gauge: an ode to the constrained Hamiltonian formalism’, in Brading and Castellani (ed.s) (2003), pp. 140-162.
H. Goldstein et al. (2002), Classical Mechanics, Addison-Wesley, (third edition).
T. Hawkins (2000), Emergence of the Theory of Lie Groups: an essay in the history of mathematics 1869-1926, New York: Springer.
M. Henneaux and C. Teitelboim (1992), Quantization of Gauge Systems, Princeton University Press.
O. Johns (2005), Analytical Mechanics for Relativity and Quantum Mechanics, Oxford University Press, forthcoming.
J. José and E. Saletan (1998), Classical Dynamics: a Contemporary Approach, Cambridge University Press.
H. Kastrup (1987), ‘The contributions of Emmy Noether, Felix Klein and Sophus Lie to the modern concept of symmetries in physical systems’, in Symmetries in Physics (1600-1980), Barcelona: Bellaterra, Universitat Autonoma de Barcelona, p. 113-163.
S. Lie (1890). Theorie der Transformationsgruppen: zweiter abschnitt, Leipzig: B.G.Teubner.
C. Lanczos (1986), The Variational Principles of Mechanics, Dover; (reprint of the 4th edition of 1970).
J. Marsden and T. Ratiu (1999), Introduction to Mechanics and Symmetry, second edition: Springer-Verlag.
G. Morandi et al (1990), ‘The inverse problem of the calculus of variations and the geometry of the tangent bundle’, Physics Reports 188, p. 147-284.
P. Olver (2000), Applications of Lie Groups to Differential Equations, second edition: Springer-Verlag.
D. Wallace (2003), ‘Time-dependent Symmetries: the link between gauge symmetries and indeterminism’, in Brading and Castellani (ed.s) (2003), pp. 163-173.
E. Wigner (1954), ‘Conservation laws in classical and quantum physics’, Progress of Theoretical Physics 11, p. 437-440. |
warning/0507/math0507578.html | ar5iv | text | # The contact process seen from a typical infected site
## 1 Introduction and main results
### 1.1 Introduction
This paper studies contact processes whose underlying lattice is a general countable group. There exists a small body of literature about contact processes on general lattices, but several basic questions have been answered only on specific lattices. In particular, a lot is known about the process on the $`d`$-dimensional integer lattice $`^d`$, and on regular trees. (See \[Lig99\] as a general reference for contact processes on $`^d`$, trees, and other lattices.)
It turns out that the contact process on regular trees behaves quite differently from the contact process on $`^d`$. For the process on $`^d`$, it is known that there is a critical infection rate $`0<\lambda _\mathrm{c}<\mathrm{}`$ such that for $`\lambda \lambda _\mathrm{c}`$, the process dies out, while for $`\lambda >\lambda _\mathrm{c}`$, the process survives with positive probability, and complete convergence holds. On the other hand, on trees, there are two critical values $`0<\lambda _\mathrm{c}<\lambda _\mathrm{c}^{}<\mathrm{}`$ such that in the intermediate regime $`\lambda _\mathrm{c}<\lambda \lambda _\mathrm{c}^{}`$, the process survives, but complete convergence does not hold. The situation is quite similar to the situation for (unoriented) percolation on general transitive lattices, where it is known that one has uniqueness of the infinite cluster whenever the lattice is amenable, while it is conjectured (and proved in several special cases) that on any nonamenable lattice there exists an intermediate parameter regime where there are infinitely many infinite clusters.
While a lot is known nowadays about percolation on general transitive graphs, the same cannot be said for the contact process. In particular, it is not known what is the essential difference between $`^d`$ and trees that causes the observed difference in behavior on these lattices. A natural guess is that the essential feature is amenability ($`^d`$ being amenable, while trees are not). However, as we will see shortly, there are reasons to doubt this.
In the present paper, we study contact processes on general countable groups by means of their exponential growth rate. A simple subadditivity argument shows that the expected number of infected sites of a contact process on a transitive lattice, started with finitely many infected sites, grows at a well-defined exponential rate (independent of the initial state). On $`^d`$, it is known that this exponential growth rate is negative for $`\lambda <\lambda _\mathrm{c}`$ (see \[BG91\] or \[Lig99, Thm I.2.48\]), and zero for any $`\lambda \lambda _\mathrm{c}`$. Indeed, it is easy to see (and prove) that on $`^d`$ there is simply not enough space for a contact process to grow exponentially fast (with positive exponent). On the other hand, one of our main results in this paper is that if a contact process survives on a nonamenable group, then its exponential growth rate must be strictly positive. This result is known for trees; our proof in the case of general nonamenable groups is quite different from the known proof for trees, however. The main idea of our proof is to relate the exponential growth rate of a contact process to the configuration seen by a typical infected site at a typical late time.
Intuition says that a contact process that survives with a positive exponential growth rate behaves very much like a perturbed branching process. On the other hand, contact processes that survive but have a zero exponential growth rate are different. We do not know if (non)amenability is the essential feature here. It is known that there exist exponentially growing groups that are amenable. (A well-known example is the lamplighter group). Although we do not prove it here, it seems plausible that a contact process on such a group, if it survives, must have a positive exponential growth rate. Thus, contact processes on such amenable groups might in some respects show behavior that is more similar to processes on trees than on $`^d`$. It should be noted that (at least) on non-homogeneous lattices, the situation is even more complex. In particular, Pemantle and Stacey \[PS01\] have given examples of non-homogeneous trees of uniformly exponential growth and bounded degree, on which the critical values related to survival and complete convergence of the contact process coincide.
Part of the present work appeared before as Chapter 4 of the author’s habilitation thesis \[Swa07\]. In particular, Proposition 4.3 below is Theorem 4.3 (a) in \[Swa07\].
### 1.2 Set-up
We will study contact processes whose underlying lattice is a general countable group. From the point of view of studying general transitive lattices, this is not quite as general as one might wish; in particular, such lattices are always unimodular. Assuming that the lattice is a group will simplify our proofs, however, so as a first step it seems reasonable.
Our set-up is as follows. We let $`\mathrm{\Lambda }`$ be a finite or countably infinite group, which we refer to as the lattice, with group action $`(i,j)ij`$ and unit element $`0`$, also referred to as the origin. Each site $`i\mathrm{\Lambda }`$ can be in one of two states: healthy or infected. Infected sites become healthy with recovery rate $`\delta 0`$. An infected site $`i`$ infects another site $`j`$ with infection rate $`a(i,j)0`$. We assume that the infection rates are invariant with respect to the left action of the group and summable:
$$\begin{array}{cc}\hfill (\mathrm{i})& a(i,j)=a(ki,kj)(i,j,k\mathrm{\Lambda }),\hfill \\ \hfill (\mathrm{ii})& |a|:=\underset{i}{}a(0,i)<\mathrm{},\hfill \end{array}$$
(1.1)
Here we adopt the convention that sums over $`i,j,k`$ always run over $`\mathrm{\Lambda }`$, unless stated otherwise. Note that we do not assume that $`a(i,j)=a(j,i)`$, i.e., our contact processes are in general asymmetric.
Let $`\eta _t`$ be the set of all infected sites at time $`t0`$. Then $`\eta =(\eta _t)_{t0}`$ is a Markov process in the space $`𝒫(\mathrm{\Lambda }):=\{A:A\mathrm{\Lambda }\}`$ of all subsets of $`\mathrm{\Lambda }`$, called the contact process on $`\mathrm{\Lambda }`$ with infection rates $`a=(a(i,j))_{i,j\mathrm{\Lambda }}`$ and recovery rate $`\delta `$, or shortly the $`(\mathrm{\Lambda },a,\delta )`$-contact process. We equip $`𝒫(\mathrm{\Lambda })\{0,1\}^\mathrm{\Lambda }`$ with the product topology and the associated Borel-$`\sigma `$-field $`(𝒫(\mathrm{\Lambda }))`$, and let $`𝒫_{\mathrm{fin}}(\mathrm{\Lambda }):=\{A\mathrm{\Lambda }:|A|<\mathrm{}\}`$ denote the subspace of finite subsets of $`\mathrm{\Lambda }`$. Under the assumptions (1.1), $`\eta `$ is a well-defined Feller process with cadlag sample paths in the compact state space $`𝒫(\mathrm{\Lambda })`$, and $`\eta _0𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$ implies $`\eta _t𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$ for all $`t0`$ a.s.
Note that we have not assumed any additional structure on $`\mathrm{\Lambda }`$, except for the group structure. In particular, we have not assumed any sort of ‘nearest neighbor’ structure. This may be obtained in the following special case. Assume that $`\mathrm{\Lambda }`$ is finitely generated and that $`\mathrm{\Delta }`$ is a finite, symmetric (with respect to taking inverses), generating set for $`\mathrm{\Lambda }`$. Then the (left) Cayley graph $`𝒢=𝒢(\mathrm{\Lambda },\mathrm{\Delta })`$ associated with $`\mathrm{\Lambda }`$ and $`\mathrm{\Delta }`$ is the graph with vertex set $`𝒱(𝒢):=\mathrm{\Lambda }`$ and edges $`(𝒢):=\{\{i,j\}:i^1j\mathrm{\Delta }\}`$. Examples of Cayley graphs are $`^d`$ and regular trees. (In the case of trees, there are several possible choices for the group structure.) Setting $`a(i,j):=\lambda 1_{\{i^1j\mathrm{\Delta }\}}`$, with $`\lambda >0`$, and choosing $`\delta 0`$, then defines a nearest-neighbor contact process on the Cayley graph $`𝒢(\mathrm{\Lambda },\mathrm{\Delta })`$. In this case, $`\lambda `$ is simply referred to as ‘the’ infection rate. If $`\delta >0`$, then by rescaling time we may set $`\delta =1`$, so it is customary to assume that $`\delta =1`$. If $`\delta =0`$, then $`\eta `$ is a special case of first-passage percolation (see \[Kes86\]).
Returning to our more general set-up, we make the following observation, which is the basis of our analysis. Below, we use the notation $`\eta _t^A`$ to denote the $`(\mathrm{\Lambda },a,\delta )`$-contact process started at time zero in $`\eta _0^A=A`$, evaluated at time $`t0`$.
###### Lemma 1.1
(Exponential growth rate) Let $`\eta `$ be a $`(\mathrm{\Lambda },a,\delta )`$-contact process. Then there exists a constant $`r=r(\mathrm{\Lambda },a,\delta )`$ with $`\delta r|a|\delta `$ such that
$$\underset{t\mathrm{}}{lim}\frac{1}{t}\mathrm{log}𝔼\left[|\eta _t^A|\right]=r(\mathrm{}A𝒫_{\mathrm{fin}}(\mathrm{\Lambda })).$$
(1.2)
We call $`r=r(\mathrm{\Lambda },a,\delta )`$ the exponential growth rate of the $`(\mathrm{\Lambda },a,\delta )`$-contact process. We note that $`r`$ has been defined before in the specific context of nearest-neighbor processes on regular trees. Indeed, $`r=\mathrm{log}\varphi (1)`$, where $`\varphi (\rho )`$ is the function defined in \[Lig99, formula (I.4.23)\].
### 1.3 The exponential growth rate
In this section, we investigate the exponential growth rate $`r(\mathrm{\Lambda },a,\delta )`$ of a contact process defined in Lemma 1.1.
We start by recalling a few basic facts and definitions concerning groups. As before, let $`\mathrm{\Lambda }`$ be a finite or countably infinite group. For $`i\mathrm{\Lambda }`$ and $`A,B\mathrm{\Lambda }`$ we put $`AB:=\{ij:iA,jB\}`$, $`iA:=\{i\}A`$, $`Ai:=A\{i\}`$, $`A^1:=\{i^1:iA\}`$, $`A^0:=\{0\}`$, $`A^n:=AA^{n1}`$ $`(n1)`$, and $`A^n:=(A^1)^n=(A^n)^1`$. We write $`A\mathrm{}B:=(A\backslash B)(B\backslash A)`$ for the symmetric difference of $`A`$ and $`B`$ and let $`|A|`$ denote the cardinality of $`A`$.
By definition, we say that $`\mathrm{\Lambda }`$ is amenable if
$$\begin{array}{c}\text{For every finite nonempty }\mathrm{\Delta }\mathrm{\Lambda }\text{ and }\epsilon >0\text{, there exists a finite}\hfill \\ \text{nonempty }A\mathrm{\Lambda }\text{ such that }|(A\mathrm{\Delta })\mathrm{}A|\epsilon |A|\text{.}\hfill \end{array}$$
(1.3)
If $`\mathrm{\Lambda }`$ is finitely generated, then it suffices to check (1.3) for one finite symmetric generating set $`\mathrm{\Delta }`$. In this case, $`(A\mathrm{\Delta })\mathrm{}A`$ is the set of all $`iA`$ for which there exists a $`jA`$ such that $`i`$ and $`j`$ are connected by an edge in the Cayley graph $`𝒢(\mathrm{\Lambda },\mathrm{\Delta })`$. Thus, we may describe (1.3) by saying that it is possible to find nonempty sets $`A`$ whose surface is small compared to their volume. For example, $`^d`$ is amenable, but regular trees are not.
If $`\mathrm{\Lambda }`$ is a finitely generated group and $`\mathrm{\Delta }`$ is a finite symmetric generating set, then we let $`|i|`$ denote the usual graph distance of $`i`$ to the origin in the Cayley graph $`𝒢(\mathrm{\Lambda },\mathrm{\Delta })`$, i.e., $`|i|:=\mathrm{min}\{n:i\mathrm{\Delta }^n\}`$. The norm $`||`$ depends on the choice of $`\mathrm{\Delta }`$, but any two norms associated with different finite symmetric generating sets are equivalent. It follows from subadditivity that the limit $`lim_n\mathrm{}\frac{1}{n}\mathrm{log}|\{i\mathrm{\Lambda }:|i|n\}|`$ exists; one says that the group $`\mathrm{\Lambda }`$ has exponential (resp. subexponential) growth if this limit is positive (resp. zero). Note that since norms associated with different finite symmetric generating sets are equivalent, having (sub)exponential growth is a property of the group $`\mathrm{\Lambda }`$ only and does not depend on the choice of $`\mathrm{\Delta }`$.
Subexponential growth implies amenabilty, but the converse is not true: as already mentioned, the lamplighter group is an amenable group with exponential growth. See \[MW89, Section 5\] for general facts about amenability and subexponential growth, and \[LPP96\] for a nice exposition of the lamplighter group.
We also need a few definitions concerning contact processes. If $`a=(a(i,j))_{i,j\mathrm{\Lambda }}`$ are infection rates satisfying (1.1), then we define reversed infection rates $`a^{}`$ by $`a^{}(i,j):=a(j,i)`$ $`(i,j\mathrm{\Lambda })`$. We say that the $`(\mathrm{\Lambda },a,\delta )`$-contact process survives if
$$\left[\eta _t^A\mathrm{}t0\right]>0$$
(1.4)
for some, and hence for all $`\mathrm{}A𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$. Using the standard coupling, it is easy to see that if $`\delta <\delta ^{}`$ and the $`(\mathrm{\Lambda },a,\delta ^{})`$-contact process survives, then the $`(\mathrm{\Lambda },a,\delta )`$-contact process survives. We let
$$\delta _\mathrm{c}=\delta _\mathrm{c}(\mathrm{\Lambda },a):=sup\{\delta 0:\text{ the }(\mathrm{\Lambda },a,\delta )\text{-contact process survives}\}$$
(1.5)
denote the critical recovery rate. By comparison with a critical branching process, it is not hard to see that $`\delta _\mathrm{c}|a|`$. Although we do not need this in what follows, we note that if $`\mathrm{\Lambda }`$ is finitely generated and the infection rates $`a`$ are irreducible, then one may use comparison with a one-dimensional nearest-neighbor contact process to show that $`0<\delta _\mathrm{c}`$ (see \[Swa07, Lemma 4.18\]).
Here, we say that infection rates $`a`$ on $`\mathrm{\Lambda }`$ are irreducible if
$$\underset{n0}{}(AA^1)^n=\mathrm{\Lambda }\text{where }A:=\{i\mathrm{\Lambda }:a(0,i)>0\}.$$
(1.6)
At some point, we will need an assumption that is a bit stronger than this. More precisely, we will occasionally use the following assumption (see Lemma 3.7 below):
$$\underset{n0,m0}{}A^nA^m=\mathrm{\Lambda }=\underset{n0,m0}{}A^nA^m,\text{where }A:=\{i\mathrm{\Lambda }:a(0,i)>0\}.$$
(1.7)
Note that this says that for any two sites $`i,j`$ there exists a site $`k`$ from which both $`i`$ and $`j`$ can be infected, and a site $`k^{}`$ that can be infected both from $`i`$ and from $`j`$.
With these definitions, we are ready to formulate our main result.
###### Theorem 1.2
(Properties of the exponential growth rate) Let $`\mathrm{\Lambda }`$ be a finite or countably infinite group, let $`a=(a(i,j))_{i,j\mathrm{\Lambda }}`$ be infection rates satisfying (1.1), and $`\delta 0`$. Let $`r=r(\mathrm{\Lambda },a,\delta )`$ be the exponential growth rate of the $`(\mathrm{\Lambda },a,\delta )`$-contact process, defined in (1.2). Then:
* $`r(\mathrm{\Lambda },a,\delta )=r(\mathrm{\Lambda },a^{},\delta )`$
* The function $`\delta r(\mathrm{\Lambda },a,\delta )`$ is nonincreasing and Lipschitz continuous on $`[0,\mathrm{})`$, with Lipschitz constant 1.
* If the $`(\mathrm{\Lambda },a,\delta )`$-contact process survives, then $`r0`$.
* If $`r>0`$, then the $`(\mathrm{\Lambda },a,\delta )`$-contact process survives.
* If $`\mathrm{\Lambda }`$ is finitely generated and has subexponential growth, and the infection rates satisfy $`_ia(0,i)e^{\epsilon |i|}<\mathrm{}`$ for some $`\epsilon >0`$, then $`r0`$.
* If $`\mathrm{\Lambda }`$ is nonamenable, the $`(\mathrm{\Lambda },a,\delta )`$-contact process survives, and the infection rates satisfy the irreducibility condition (1.7), then $`r>0`$.
Parts (a), (b), and (c) of this theorem are easy. Part (d) follows from a variance calculation, while (e) is proved by some simple large deviation estimates. The proof of part (f) is rather involved. For trees, the statement is a known consequence of \[Lig99, Prop. I.4.27 (b)\]. Our proof for general nonamenable groups is quite different from the methods used there. The basic idea is as follows. If the exponential growth rate of a contact process is zero, then this means that for the process started with one infected site, a ‘typical’ infected site at a ‘typical’ late time produces no net offspring, i.e., the mean number of sites it infects per unit of time is just enough to balance the probability that the site itself recovers. We prove that this implies that the local configuration as seen from this ‘typical’ site is distributed as the upper invariant measure, assuming that the latter is nontrivial. If $`\mathrm{\Lambda }`$ is nonamenable, this leads to a contradiction, since for any finite collection of particles on a nonamenable lattice, a positive fraction of the particles must lie on the ‘outer boundary’ of the collection, hence must see something different from the upper invariant measure.
Parts (b), (d), and (f) of Theorem 1.2 yield the following corollary.
###### Corollary 1.3
(The critical contact process on a nonamenable lattice dies out) If $`\mathrm{\Lambda }`$ is nonamenable and the infection rates satisfy the irreducibility condition (1.7), then $`\delta _\mathrm{c}=\delta _\mathrm{c}(\mathrm{\Lambda },a)>0`$ and the $`(\mathrm{\Lambda },a,\delta _\mathrm{c})`$-contact process dies out.
For nearest-neighbor contact processes on regular trees, this result is known, see \[MSZ94\] or \[Lig99, Proposition I.4.39\]. Like our proof, the proof there is based on showing that a zero exponential growth rate implies extinction (although they use quite different techniques to establish this). The analogue of Corollary 1.3 for (unoriented) percolation says that there are no infinite clusters at criticality on any nonamenable lattice. This has been proved in \[BLPS99\]; it seems that the techniques used there have little in common with the ones used in the present paper. The problem of showing that critical percolation on $`^d`$ has no infinite cluster is still open in dimensions $`3d18`$. Barsky, Grimmett, and Newman \[BGN91\] have shown, however, that at criticality there are no infinite clusters in the half-space $`_+\times ^{d1}`$, and, using similar techniques, Bezuidenhout and Grimmett \[BG90\] have proved that the critical contact process on $`^d`$ dies out.
### 1.4 The process seen from a typical site
There is an intimate relation between the survival probability of a contact process and its upper invariant law. Similarly, there is a relation between the exponential growth rate and certain infinite measures on the space of nonempty subsets of $`\mathrm{\Lambda }`$, which we explain now.
Let $`𝒫_+(\mathrm{\Lambda }):=\{A\mathrm{\Lambda }:A\mathrm{}\}`$ denote the set of all nonempty subsets of $`\mathrm{\Lambda }`$. Note that $`𝒫_+(\mathrm{\Lambda })`$ is a locally compact space in the induced topology from $`𝒫(\mathrm{\Lambda })`$. We say that a measure $`\mu `$ on $`𝒫(\mathrm{\Lambda })`$ or $`𝒫_+(\mathrm{\Lambda })`$ is (spatially) homogeneous if it is invariant under the left action of the group, i.e., if $`\mu (𝒜)=\mu (i𝒜)`$ for each $`i\mathrm{\Lambda }`$ and $`𝒜(𝒫(\mathrm{\Lambda }))`$, where we define $`i𝒜:=\{iA:A𝒜\}`$. We say that a measure $`\mu `$ is an eigenmeasure of the $`(\mathrm{\Lambda },a,\delta )`$-contact process if $`\mu `$ is a nonzero, locally finite measure on $`𝒫_+(\mathrm{\Lambda })`$, and there exists a constant $`\lambda `$ such that
$$\mu (\mathrm{d}A)[\eta _t^A]|_{𝒫_+(\mathrm{\Lambda })}=e^{\lambda t}\mu (t0),$$
(1.8)
where $`|_{𝒫_+(\mathrm{\Lambda })}`$ denotes restriction (of a measure) to $`𝒫_+(\mathrm{\Lambda })`$. We call $`\lambda `$ the associated eigenvalue. As a motivation for this terminology, we observe that if $`G`$ is the generator of the $`(\mathrm{\Lambda },a,\delta )`$-contact process, then formally $`G^{}\mu =\lambda \mu `$. Note that if $`\lambda =0`$ and $`\mu `$ is concentrated on the infinite subsets of $`\mathrm{\Lambda }`$, then the measure on the left-hand side of (1.8) is concentrated on $`𝒫_+(\mathrm{\Lambda })`$, hence in this case (1.8) just says that $`\mu `$ is an invariant measure (though not necessarily a probability measure) for the $`(\mathrm{\Lambda },a,\delta )`$-contact process.
###### Proposition 1.4
(Exponential growth rate and eigenmeasures) For each $`(\mathrm{\Lambda },a,\delta )`$-contact process, the set
$$\begin{array}{cc}\hfill (\mathrm{\Lambda },a,\delta ):=\{\lambda :& \text{there exists a homogeneous eigenmeasure}\hfill \\ & \text{of the }(\mathrm{\Lambda },a,\delta )\text{-contact process with eigenvalue }\lambda \}\hfill \end{array}$$
(1.9)
is a nonempty compact subset of $``$, and $`r(\mathrm{\Lambda },a,\delta )=\mathrm{max}(\mathrm{\Lambda },a,\delta )`$.
In particular, Proposition 1.4 implies that each $`(\mathrm{\Lambda },a,\delta )`$-contact process has a homogeneous eigenmeasure with eigenvalue $`r(\mathrm{\Lambda },a,\delta )`$. It seems natural to conjecture that this eigenmeasure is always unique and the long-time limit of the (suitably rescaled) law of the process started with one infected site, distributed according to counting measure on $`\mathrm{\Lambda }`$. If we condition such an eigenmeasure on the origin being infected, then we can view the resulting probability measure as describing the contact process as seen from a typical infected site, at late times. (Compare Corollary 3.4 and Lemma 4.2 below.)
Recall that the upper invariant measure $`\overline{\nu }`$ of a contact process is the long-time limit law of the process started with all sites infected. It follows from duality (and spatial ergodicity of the graphical representation) that the upper invariant measure of the $`(\mathrm{\Lambda },a,\delta )`$-contact process is nontrivial (i.e., gives zero probability to the empty set) if and only if the $`(\mathrm{\Lambda },a^{},\delta )`$-contact process survives. The next result is an important ingredient in the proof of Theorem 1.2 (f).
###### Theorem 1.5
(Eigenmeasures with eigenvalue zero) Assume that the infection rates satisfy the irreducibility condition (1.7). If the upper invariant measure $`\overline{\nu }`$ of the $`(\mathrm{\Lambda },a,\delta )`$-contact process is nontrivial, then any homogeneous eigenmeasure $`\mu `$ with eigenvalue zero satisfies $`\mu =c\overline{\nu }`$ for some $`c>0`$.
We prove Theorem 1.5 by extending well-known techniques for showing that $`\overline{\nu }`$ is the only nontrivial homogeneous invariant probability measure of a contact process. If $`\overline{\nu }`$ is trivial, then there may exist homogeneous eigenmeasures with eigenvalue zero, which in this case, obviously, are not a multiple of $`\overline{\nu }`$. Indeed, if $`a`$ is symmetric (i.e., $`a=a^{}`$) and the critical process dies out (as we know to be the case on $`^d`$ or on any nonamenable group), then at criticality $`\overline{\nu }`$ is trivial, while by Theorem 1.2 (b), (c), and (d), the exponential growth rate is zero, hence by Proposition 1.4, there exists a homogeneous eigenmeasure with eigenvalue zero.
### 1.5 Discussion, open problems, and outline
The work in this paper started from the question whether it is possible to prove something like ‘uniqueness of the infinite cluster’ in the context of oriented percolation or the (very similar) graphical representation of the contact process. This question is still very much open. See Grimmett and Hiemer \[GH02\] for a weak statement that is proved only on $`^d`$ and Wu and Zhang \[WZ06, Thm 1.4\] or \[Swa07, Lemma 4.5\] for a stronger statement that is proved only in the nearest-neighbor, one-dimensional case.
Whether the methods in the present paper can shed some light on this question I do not know. I have tried to prove the weak statement of Grimmett and Hiemer assuming (only) subexponential growth, but ran into the problem that I would need to replace a size-biased law by a law conditioned on survival, which I do not know how to do (see \[Swa07, Prop 4.4\]).
In fact, although this is not obvious from the presentation above, size-biased laws and Campbell measures, well-known objects from branching theory, are closely related to the eigenmeasures introduced above. (For this connection, see Section 4.3 below.) An interesting feature of the (potentially infinite) eigenmeasures is that they allow one to use some of the simplifications that come from spatial homogeneity while studying processes started in finite initial states.
There are lots of open problems concerning contact processes on general transitive lattices, so we mention just a few.
1. Prove that the $`(\mathrm{\Lambda },a,\delta )`$-contact process has a unique homogeneous eigenmeasure with eigenvalue $`r(\mathrm{\Lambda },a,\delta )`$, which is the long-time limit law of the process started with one infected site distributed according to counting measure on $`\mathrm{\Lambda }`$.
2. Prove that $`\frac{}{\delta }r(\mathrm{\Lambda },a,\delta )<0`$ on $`\{\delta :r(\mathrm{\Lambda },a,\delta )0\}`$. Prove the same statement for all $`\delta `$ if $`\mathrm{\Lambda }`$ is nonamenable. Adapt the known proof for $`^d`$ (see \[BG91\] and \[Lig99, Thm I.2.48\]) that $`r(\mathrm{\Lambda },a,\delta )<0`$ for all $`\delta >\delta _\mathrm{c}`$, to general lattices.
3. Prove that $`\delta _\mathrm{c}>0`$ for some $`(\mathrm{\Lambda },a,\delta )`$-contact process on a group $`\mathrm{\Lambda }`$ that is not finitely generated, e.g. the hierarchical group.
4. Study contact processes on transitive lattices $`\mathrm{\Lambda }`$ that are not groups. In this context, if $`\mathrm{\Lambda }`$ is not unimodular, it is not hard to find examples where a $`(\mathrm{\Lambda },a,\delta )`$-contact process survives but its dual $`(\mathrm{\Lambda },a^{},\delta )`$-contact process dies out. It is an open problem to prove this cannot happen in the unimodular case.
5. Prove some version of uniqueness of the infinite cluster assuming that the exponential growth rate is zero.
6. Prove (or disprove) that $`r(\mathrm{\Lambda },a,\delta )>0`$ whenever the $`(\mathrm{\Lambda },a,\delta )`$-contact process survives and $`r(\mathrm{\Lambda },a,0)>0`$.
The outline of the rest of the paper is as follows. In Section 2, we introduce some basic tools, such as the graphical representation and a martingale problem. In Section 3, we prove Theorem 1.2 (a)–(c), Proposition 1.4, and Theorem 1.5. These results are then used in Section 4 to prove Theorem 1.2 (d)–(f) and Corollary 1.3.
Acknowledgements The author thanks Geoffrey Grimmmett, Olle Häggström, Russel Lyons, Amos Nevo, Yuval Peres, and Roberto Schonmann for useful email conversations about the contact process, oriented percolation, and amenability. In particular, the proof of formula (4.30) is due to Yuval Peres. The author thanks the referee for many useful suggestions.
## 2 Construction and basic properties
### 2.1 Graphical representation
We will, of course, use the graphical representation of the contact process. Let $`\mathrm{\Lambda }\times :=\{(i,t):i\mathrm{\Lambda },t\}`$ and $`\mathrm{\Lambda }\times \mathrm{\Lambda }\times :=\{(i,j,t):i,j\mathrm{\Lambda },t\}`$, where $`t`$ is the time coordinate. Let $`\omega =(\omega ^\mathrm{r},\omega ^\mathrm{i})`$ be a pair of independent, locally finite random subsets of $`\mathrm{\Lambda }\times `$ and $`\mathrm{\Lambda }\times \mathrm{\Lambda }\times `$, respectively, produced by Poisson point processes with intensity $`\delta `$ and $`a(i,j)`$, respectively. We visualize this by plotting $`\mathrm{\Lambda }`$ horizontally and $``$ vertically, marking points $`(i,s)\omega ^\mathrm{r}`$ with a recovery symbol $``$, and drawing an infection arrow from $`(i,t)`$ to $`(j,t)`$ for each $`(i,j,t)\omega ^\mathrm{i}`$. For $`C,D\mathrm{\Lambda }\times `$, say that there is a path from $`C`$ to $`D`$, denoted by $`CD`$, if there exist $`n0`$, $`i_0,\mathrm{},i_n\mathrm{\Lambda }`$, and $`t_0\mathrm{}t_{n+1}`$ with $`(i_0,t_0)C`$ and $`(i_n,t_{n+1})D`$, such that $`(\{i_k\}\times [t_k,t_{k+1}])\omega ^\mathrm{r}=\mathrm{}`$ for all $`k=0,\mathrm{},n`$ and $`(i_{k1},i_k,t_k)\omega ^\mathrm{i}`$ for all $`k=1,\mathrm{},n`$. Thus, a path must walk upwards in time, may follow arrows, and must avoid recoveries. For $`C\mathrm{\Lambda }\times `$, we write $`C\mathrm{}`$ if there is an infinite path with times $`t_k\mathrm{}`$ starting in $`C`$. We define $`\mathrm{}C`$ analogously. Instead of $`\{(i,s)\}`$ and $`\{(j,t)\}`$, simply write $`(i,s)`$ and $`(j,t)`$.
For given $`A𝒫(\mathrm{\Lambda })`$ and $`t_0`$, put
$$\eta _t^{A\times \{t_0\}}:=\{i\mathrm{\Lambda }:A\times \{t_0\}(i,t_0+t)\}(t0).$$
(2.1)
Then $`\eta ^{A\times \{t_0\}}=(\eta _t^{A\times \{t_0\}})_{t0}`$ is a $`(\mathrm{\Lambda },a,\delta )`$-contact process started in $`\eta _0^{A\times \{t_0\}}=A`$.
In analogy with (2.1), put
$$\eta _t^{A\times \{t_0\}}:=\{i\mathrm{\Lambda }:(i,t_0t)A\times \{t_0\}\}(t0).$$
(2.2)
Then $`\eta ^{A\times \{t_0\}}=(\eta _t^{A\times \{t_0\}})_{t0}`$ is a $`(\mathrm{\Lambda },a^{},\delta )`$-contact process started in $`\eta _0^{A\times \{t_0\}}=A`$. Since for any $`st`$ and $`A,B𝒫(\mathrm{\Lambda })`$, the event
$$\left\{\eta _{us}^{A\times \{s\}}\eta _{tu}^{B\times \{t\}}=\mathrm{}\right\}=\left\{A\times \{s\}\rightsquigarrow ̸B\times \{t\}\right\}$$
(2.3)
does not depend on $`u[s,t]`$, it follows (by taking $`s=0`$ and $`u=0,t`$) that the $`(\mathrm{\Lambda },a,\delta )`$-contact process and the $`(\mathrm{\Lambda },a^{},\delta )`$-contact process are dual in the sense that
$$[\eta _t^AB=\mathrm{}]=[A\eta _t^B=\mathrm{}](A,B𝒫(\mathrm{\Lambda }),t0).$$
(2.4)
Here, for brevity, we write $`\eta _t^A:=\eta _t^{A\times \{0\}}`$ and $`\eta _t^B:=\eta _t^{B\times \{0\}}`$.
It is not hard to see that $`|a|:=_ia(0,i)=_ia(i,0)`$ and
$$𝔼\left[|\eta _t^A|\right]|A|e^{|a|t}\text{and}𝔼\left[|\eta _t^A|\right]|A|e^{|a|t}(t0,A\mathrm{\Lambda }).$$
(2.5)
In particular, both the $`(\mathrm{\Lambda },a,\delta )`$-contact process and the $`(\mathrm{\Lambda },a^{},\delta )`$-contact process are well-defined and the processes started from a finite initial state are a.s. finite for all time.
For any $`A\mathrm{\Lambda }`$, we let
$$\rho (A):=[\eta _t^A\mathrm{}t0]=[A\times \{0\}\mathrm{}]$$
(2.6)
denote the survival probability of the $`(\mathrm{\Lambda },a,\delta )`$-contact process started in $`A`$. Similarly, $`\rho ^{}`$ denotes the survival probability of the $`(\mathrm{\Lambda },a^{},\delta )`$-contact process. Setting
$$\overline{\eta }_t:=\{i\mathrm{\Lambda }:\mathrm{}(i,t)\}(t)$$
(2.7)
defines a stationary $`(\mathrm{\Lambda },a,\delta )`$-contact process whose invariant law
$$\overline{\nu }:=[\overline{\eta }_t](t)$$
(2.8)
is uniquely characterized by
$$\left[\overline{\eta }_0A\mathrm{}\right]=\rho ^{}(A)(A𝒫_{\mathrm{fin}}(\mathrm{\Lambda })).$$
(2.9)
(To see this, note that the linear span of the functions $`B1_{\{AB=\mathrm{}\}}`$ with $`A𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$ forms an algebra that separates points, hence by the Stone-Weierstrass theorem is dense in the space of continuous functions on $`𝒫(\mathrm{\Lambda })`$.) It is easy to see that $`\overline{\nu }`$ is nontrivial, i.e., gives zero probability to the empty set, if and only if the $`(\mathrm{\Lambda },a^{},\delta )`$-contact process survives. Moreover, $`\overline{\nu }`$ is the limit law of the process started with all sites infected, i.e., $`\overline{\nu }`$ is the upper invariant law.
### 2.2 Martingale problem
We will need the fact that $`(\mathrm{\Lambda },a,\delta )`$-contact processes started in finite initial states solve a martingale problem. Let
$$𝒮(𝒫_{\mathrm{fin}}(\mathrm{\Lambda })):=\{f:𝒫_{\mathrm{fin}}(\mathrm{\Lambda }):|f(A)|K|A|^k+M\text{ for some }K,M,k0\}.$$
(2.10)
denote the class of real functions on $`𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$ of polynomial growth. Given $`\mathrm{\Lambda },a`$, and $`\delta `$, define a linear operator $`G`$ with domain $`𝒟(G):=𝒮(𝒫_{\mathrm{fin}}(\mathrm{\Lambda }))`$ by
$$\begin{array}{ccc}\hfill Gf(A)& :=& \underset{ij}{}a(i,j)1_{\{iA\}}1_{\{jA\}}\{f(A\{j\})f(A)\}\hfill \\ & & +\delta \underset{i}{}1_{\{iA\}}\{f(A\backslash \{i\})f(A)\}.\hfill \end{array}$$
(2.11)
###### Proposition 2.1
(Martingale problem and moment estimate) The operator $`G`$ maps the space $`𝒮(𝒫_{\mathrm{fin}}(\mathrm{\Lambda }))`$ into itself. For each $`f𝒮(𝒫_{\mathrm{fin}}(\mathrm{\Lambda }))`$ and $`A𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$, the process
$$M_t:=f(\eta _t^A)_0^tGf(\eta _s^A)ds(t0)$$
(2.12)
is a martingale with respect to the filtration generated by $`\eta ^A`$. Moreover, setting $`z^k:=_{i=0}^{k1}(z+i)`$, one has
$$𝔼\left[|\eta _t^A|^k\right]|A|^ke^{k(|a|+(k2)\delta )t}(A𝒫_{\mathrm{fin}}(\mathrm{\Lambda }),k1,t0).$$
(2.13)
Proof Our proof follows the same lines as the proof of \[AS05, Proposition 8\]. It is not hard to see that $`G`$ maps $`𝒮(𝒫_{\mathrm{fin}}(\mathrm{\Lambda }))`$ into itself. Set $`f_k(A):=|A|^k`$. Then, using the fact that $`z^k(z1)^k=kz^{k1}`$, we see that
$$\begin{array}{ccc}\hfill Gf_k(A)& =& \underset{ij}{}a(i,j)1_{\{iA\}}1_{\{jA\}}\{(|A|+1)^k|A|^k\}+\delta \underset{i}{}1_{\{iA\}}\{(|A|1)^k|A|^k\},\hfill \\ & & |a||A|\{(|A|+1)^k|A|^k\}\delta |A|\{|A|^k(|A|1)^k\}\hfill \\ & =& (|a|\delta )k|A|^k+\delta k(k1)|A|^{k1}k(|a|+(k2)\delta )|A|^k=:K_k|A|^k.\hfill \end{array}$$
(2.14)
Define stopping times $`\tau _N:=inf\{t0:|\eta _t^A|N\}`$. The stopped process $`(\eta _{t\tau _N}^A)_{t0}`$ has bounded jump rates, and therefore standard theory tells us that for each $`N1`$ and $`f𝒮(𝒫_{\mathrm{fin}}(\mathrm{\Lambda }))`$, the process
$$M_t^N:=f(\eta _{t\tau _N}^A)_0^{t\tau _N}Gf(\eta _s^A)ds(t0)$$
(2.15)
is a martingale. Moreover, it easily follows from (2.14) that
$$𝔼\left[|\eta _{t\tau _N}^A|^k\right]|A|^ke^{K_kt}(k1,t0),$$
(2.16)
which in turn implies that $`[|\eta _{t\tau _N}^A|N]0`$ as $`N\mathrm{}`$, and hence $`lim_N\mathrm{}\tau _N=\mathrm{}`$. Using the fact that $`G`$ maps $`𝒮(𝒫_{\mathrm{fin}}(\mathrm{\Lambda }))`$ into itself and (2.16) for some sufficiently high $`k`$ (depending on $`f`$), one can show that for fixed $`t0`$, the random variables $`(M_t^N)_{N1}`$ are uniformly integrable. Therefore, letting $`N\mathrm{}`$ in (2.15), one finds that the process in (2.12) is a martingale. Letting $`N\mathrm{}`$ in (2.16) yields (2.13).
### 2.3 Covariance formula
By Proposition 2.1, setting
$$S_tf(A):=𝔼[f(\eta _t^A)](f𝒮(𝒫_{\mathrm{fin}}(\mathrm{\Lambda })),A𝒫_{\mathrm{fin}}(\mathrm{\Lambda }))$$
(2.17)
defines a semigroup $`(S_t)_{t0}`$ of linear operators $`S_t:𝒮(𝒫_{\mathrm{fin}}(\mathrm{\Lambda }))𝒮(𝒫_{\mathrm{fin}}(\mathrm{\Lambda }))`$. Let $``$ be the class of probability measures on $`𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$ such that $`|A|^k\mu (\mathrm{d}A)<\mathrm{}`$ for all $`k1`$. For $`\mu `$ and $`f𝒮(𝒫_{\mathrm{fin}}(\mathrm{\Lambda }))`$, we write $`\mu f:=f(A)\mu (\mathrm{d}A)`$. Note that if $`(\eta _t)_{t0}`$ is a $`(\mathrm{\Lambda },a,\delta )`$-contact process started in an initial law $`[\eta _0]=:\mu `$, then $`[\eta _t]`$ for all $`t0`$ and $`[\eta _t\mathrm{d}A]f(A)=\mu S_tf`$. For this reason, we use the notation $`\mu S_t:=[\eta _t]`$ $`(t0)`$ to denote the law of $`\eta _t`$. For any $`\mu `$ and $`f,g𝒮(𝒫_{\mathrm{fin}}(\mathrm{\Lambda }))`$, we let
$$\mathrm{Cov}_\mu (f,g):=\mu (fg)(\mu f)(\mu g)$$
(2.18)
denote the covariance of $`f`$ and $`g`$ under $`\mu `$, which is always finite.
###### Proposition 2.2
(Covariance formula) For $`f,g𝒮(𝒫_{\mathrm{fin}}(\mathrm{\Lambda }))`$, let
$$\mathrm{\Gamma }(f,g):=\frac{1}{2}\left[G(fg)(Gf)gf(Gg)\right](f,g𝒮(𝒫_{\mathrm{fin}}(\mathrm{\Lambda }))).$$
(2.19)
Then, for any $`\mu `$ and $`f,g𝒮(𝒫_{\mathrm{fin}}(\mathrm{\Lambda }))`$, one has
$$\mathrm{Cov}_{\mu S_t}(f,g)=\mathrm{Cov}_\mu (S_tf,S_tg)+2_0^t\mu S_{ts}\mathrm{\Gamma }(S_sf,S_sg)ds(t0).$$
(2.20)
Proof Set
$$H(s,t,u):=S_s\left((S_tf)(S_ug)\right).$$
(2.21)
We claim that
$$\begin{array}{ccc}\hfill \frac{}{s}H(s,t,u)& =& S_sG\left((S_tf)(S_ug)\right),\hfill \\ \hfill \frac{}{t}H(s,t,u)& =& S_s\left((GS_tf)(S_ug)\right),\hfill \\ \hfill \frac{}{u}H(s,t,u)& =& S_s\left((S_tf)(GS_ug)\right).\hfill \end{array}$$
(2.22)
It follows that
$$\frac{}{t}H(t,Tt,Tt)=2S_t\mathrm{\Gamma }(S_{Tt}f,S_{Tt}g),$$
(2.23)
and therefore
$$\begin{array}{c}\mathrm{Cov}_{\mu S_T}(f,g)\mathrm{Cov}_\mu (S_Tf,S_Tg)\hfill \\ =\left(\mu S_T(fg)(\mu S_Tf)(\mu S_Tg)\right)\left(\mu ((S_Tf)(S_Tg))(\mu S_Tf)(\mu S_Tg)\right)\hfill \\ =\mu \left(S_T(fg)(S_Tf)(S_Tg)\right)\hfill \\ =\mu \left(H(T,0,0)H(0,T,T)\right)=2_0^t\mu S_t\mathrm{\Gamma }(S_{Tt}f,S_{Tt}g)dt.\hfill \end{array}$$
(2.24)
These calculations are standard. However, in order to verify (2.22), we must use some special properties of our model. Let us say that a sequence of functions $`f_n𝒮(𝒫_{\mathrm{fin}}(\mathrm{\Lambda }))`$ converges ‘nicely’ to a limit $`f`$, if $`f_nf`$ pointwise and there exists $`K,M,k0`$ such that $`|f_n(A)|K|A|^k+M`$ for all $`n`$. By (2.13) and dominated convergence, if $`f_nf`$ ‘nicely’, then $`S_tf_nS_tf`$ ‘nicely’, for each $`t0`$. Note also that if $`f_n,f,g𝒮(𝒫_{\mathrm{fin}}(\mathrm{\Lambda }))`$ and $`f_nf`$ ‘nicely’, then $`f_ngfg`$ ‘nicely’. We claim that for each $`f𝒮(𝒫_{\mathrm{fin}}(\mathrm{\Lambda }))`$,
$$\underset{t0}{lim}t^1(S_tff)=Gf,$$
(2.25)
where the convergence happens ‘nicely’. Indeed, by Proposition 2.1,
$$t^1\left(S_tf(A)f(A)\right)=t^1_0^t𝔼\left[(Gf)(\eta _s^A)\right]ds\underset{t0}{}Gf(A),$$
(2.26)
where the ‘niceness’ of the convergence follows from (2.13) and the fact that $`Gf𝒮(𝒫_{\mathrm{fin}}(\mathrm{\Lambda }))`$. It follows from (2.25) that for each $`f𝒮(𝒫_{\mathrm{fin}}(\mathrm{\Lambda }))`$ and $`t0`$,
$$\begin{array}{ccc}\hfill \frac{}{t}S_tf& =& \underset{\epsilon 0}{lim}\epsilon ^1(S_\epsilon 1)S_tf=GS_tf\hfill \\ & =& \underset{\epsilon 0}{lim}\epsilon ^1S_t(S_\epsilon 1)f=S_tGf,\hfill \end{array}$$
(2.27)
where $`1`$ denotes the identity operator. Using (2.27) and the properties of ‘nice’ convergence, (2.22) follows readily.
## 3 The exponential growth rate
### 3.1 Basic facts
Proof of Lemma 1.1 Let us write
$$\pi _t(A):=𝔼\left[|\eta _t^A|\right](A𝒫_{\mathrm{fin}}(\mathrm{\Lambda }),t0).$$
(3.1)
We start by showing that
$$\pi _{s+t}(\{0\})\pi _s(\{0\})\pi _t(\{0\})(s,t0).$$
(3.2)
If $`\eta ^A`$ and $`\eta ^B`$ are defined usng the same graphical representation, then $`\eta _t^A\eta _t^B=\eta _t^{AB}`$. Therefore,
$$𝔼\left[|\eta _t^A|\right]=𝔼\left[\left|\underset{iA}{}\eta _t^{\{i\}}\right|\right]\underset{iA}{}𝔼\left[|\eta _t^{\{i\}}|\right]=|A|𝔼\left[|\eta _t^{\{0\}}|\right],$$
(3.3)
where in the last step we have used shift invariance. As a consequence,
$$\pi _{s+t}(\{0\})=[\eta _s^{\{0\}}\mathrm{d}A]𝔼\left[|\eta _t^A|\right][\eta _s^{\{0\}}\mathrm{d}A]|A|𝔼\left[|\eta _t^{\{0\}}|\right]=\pi _s(\{0\})\pi _t(\{0\}).$$
(3.4)
This proves (3.2). It follows that $`t\mathrm{log}\pi _t(\{0\})`$ is subadditive and therefore, by \[Lig99, Theorem B.22\], the limit
$$\underset{t\mathrm{}}{lim}\frac{1}{t}\mathrm{log}\pi _t(\{0\})=\underset{t>0}{inf}\frac{1}{t}\mathrm{log}\pi _t(\{0\})=:r[\mathrm{},\mathrm{})$$
(3.5)
exists. By monotonicity and (3.3),
$$\pi _t(\{0\})\pi _t(A)|A|\pi _t(\{0\})(A𝒫_{\mathrm{fin}}(\mathrm{\Lambda })).$$
(3.6)
Taking logarithms, dividing by $`t`$, and letting $`t\mathrm{}`$ we arrive at (1.2). Since $`\eta `$ can be bounded from below by a simple death process and from above by a branching process, one has
$$e^{\delta t}𝔼\left[|\eta _t^{\{0\}}|\right]e^{(|a|\delta )t}(t0),$$
(3.7)
which implies that $`\delta r|a|\delta `$.
Proof of Theorem 1.2 (a) By duality (formula (2.4)) and shift invariance,
$$\begin{array}{ccc}\hfill 𝔼\left[|\eta _t^{\{0\}}|\right]& =& \underset{i}{}\left[\eta _t^{\{0\}}\{i\}\mathrm{}\right]=\underset{i}{}\left[\{0\}\eta _t^{\{i\}}\mathrm{}\right]\hfill \\ & =& \underset{i}{}\left[\{i^1\}\eta _t^{\{0\}}\mathrm{}\right]=𝔼\left[|\eta _t^{\{0\}}|\right],\hfill \end{array}$$
(3.8)
which implies that $`r(\mathrm{\Lambda },a,\delta )=r(\mathrm{\Lambda },a^{},\delta )`$.
Proof of Theorem 1.2 (b) Fix a countable group $`\mathrm{\Lambda }`$ and infection rates $`a`$ satisfying (1.1), and for each $`\delta 0`$, write $`\pi (\delta ,t):=𝔼\left[|\eta _t^{\{0\}}|\right]`$, where $`\eta _t^{\{0\}}`$ is the $`(\mathrm{\Lambda },a,\delta )`$-contact process. For $`0\delta <\stackrel{~}{\delta }`$, consider the graphical representations (see Section 2.1) of the $`(\mathrm{\Lambda },a,\delta )`$\- and $`(\mathrm{\Lambda },a,\stackrel{~}{\delta })`$-contact processes, defined by Poisson processes $`(\omega ^\mathrm{r},\omega ^\mathrm{i})`$ and $`(\stackrel{~}{\omega }^\mathrm{r},\stackrel{~}{\omega }^\mathrm{i})`$, respectively. We may couple these graphical representations such that $`\omega ^\mathrm{i}=\stackrel{~}{\omega }^\mathrm{i}`$ and $`\omega ^\mathrm{r}\stackrel{~}{\omega }^\mathrm{r}`$, where $`\stackrel{~}{\omega }^\mathrm{r}\backslash \omega ^\mathrm{r}`$ is an Poisson point process with intensity $`\stackrel{~}{\delta }\delta `$, independent of $`\omega ^\mathrm{i}`$ and $`\omega ^\mathrm{r}`$. Write $``$ and $`\stackrel{~}{}`$ to indicate the existence of an open path in the graphical representations for $`\delta `$ and $`\stackrel{~}{\delta }`$, respectively. Then, if $`\stackrel{~}{\delta }\delta `$ is small, then for each $`t0`$,
$$\begin{array}{c}\pi (\stackrel{~}{\delta },t)=\underset{i}{}[(0,0)\stackrel{~}{}(i,t)]\hfill \\ =\underset{i}{}[(0,0)(i,t)]\hfill \\ (\stackrel{~}{\delta }\delta )_0^t\underset{ij}{}[(0,0)(j,s)(i,t)\text{ and there exists}\hfill \\ \text{no }kj\text{ such that }(0,0)(k,s)(i,t)]\mathrm{d}s+O((\stackrel{~}{\delta }\delta )^2),\hfill \end{array}$$
(3.9)
where the terms order $`(\stackrel{~}{\delta }\delta )^2`$ come from events where two or more recovery symbols in $`\stackrel{~}{\omega }^\mathrm{r}\backslash \omega ^\mathrm{r}`$ are needed to block all paths from $`(0,0)`$ to $`(i,t)`$. Dividing by $`\stackrel{~}{\delta }\delta `$ and letting $`\stackrel{~}{\delta }\delta `$ yields
$$\begin{array}{c}\frac{}{\delta }\pi (\stackrel{~}{\delta },t)\hfill \\ =_0^t\underset{i}{}\left[j\text{ s.t. }(0,0)(j,s)(i,t)\text{ and }kj\text{ s.t. }(0,0)(k,s)(i,t)\right]\mathrm{d}s,\hfill \end{array}$$
(3.10)
which is an analogue of what is known as Russo’s formula in percolation. Since
$$\begin{array}{c}\underset{i}{}\left[j\text{ s.t. }(0,0)(j,s)(i,t)\text{ and }kj\text{ s.t. }(0,0)(k,s)(i,t)\right]\hfill \\ \underset{i}{}[(0,0)(i,t)]=\pi (\delta ,t),\hfill \end{array}$$
(3.11)
it follows that $`0\frac{}{\delta }\pi (\delta ,t)t\pi (\delta ,t)`$ $`(t0)`$, and therefore
$$0\frac{}{\delta }\frac{1}{t}\mathrm{log}\pi (\delta ,t)1.$$
(3.12)
Taking the limit $`t\mathrm{}`$, using (3.5), the claims follow. Note that by Lemma 1.1, $`\delta r|a|\delta `$, so letting $`\delta \mathrm{}`$ we see that the Lipschitz constant $`1`$ is optimal.
Proof of Theorem 1.2 (c) If the $`(\mathrm{\Lambda },a,\delta )`$-contact process survives, then
$$𝔼\left[|\eta _t^{\{0\}}|\right][\eta _t^{\{0\}}\mathrm{}]\underset{t\mathrm{}}{}[\eta _s^{\{0\}}\mathrm{}s0]>0,$$
(3.13)
which implies that $`r0`$.
### 3.2 Eigenmeasures
Recall that a measure $`\mu `$ on a locally compact space is called locally finite if $`\mu (K)<\mathrm{}`$ for all compact sets $`K`$. We need a few basic facts about locally finite measures on $`𝒫_+(\mathrm{\Lambda })`$.
###### Lemma 3.1
(Locally finite measures) Let $`\mu `$ be a measure on $`𝒫_+(\mathrm{\Lambda })`$. Then the following statements are equivalent: 1. $`\mu `$ is locally finite. 2. $`\mu (\mathrm{d}A)1_{\{iA\}}<\mathrm{}`$ for all $`i\mathrm{\Lambda }`$. 3. $`\mu (\mathrm{d}A)1_{\{AB\mathrm{}\}}<\mathrm{}`$ for all $`B𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$.
Proof We will prove that 1$``$3$``$2$``$1. For each $`B𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$, the set $`𝒬(B):=\{A\mathrm{\Lambda }:AB\mathrm{}\}`$ is a compact subset of $`𝒫_+(\mathrm{\Lambda })`$, and $`A1_{\{AB\mathrm{}\}}`$ is a continuous function with compact support $`𝒬(B)`$. It follows that any locally finite measure $`\mu `$ satisfies $`\mu (\mathrm{d}A)1_{\{AB\mathrm{}\}}<\mathrm{}`$ for each $`B𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$. In particular, setting $`B=\{i\}`$ this implies that $`\mu (\mathrm{d}A)1_{\{iA\}}<\mathrm{}`$ for each $`i\mathrm{\Lambda }`$. This proves the implications 1$``$3$``$2. To see that 2 implies 1, let $`\mathrm{\Delta }_n\mathrm{\Lambda }`$ be finite sets increasing to $`\mathrm{\Lambda }`$. Then the $`𝒬(\mathrm{\Delta }_n)`$ increase to $`𝒫_+(\mathrm{\Lambda })`$ and, since the $`𝒬(\mathrm{\Delta }_n)`$ are open sets, each compact subset of $`𝒫_+(\mathrm{\Lambda })`$ is contained in some $`𝒬(\mathrm{\Delta }_n)`$. Therefore, since $`\mu (𝒬(\mathrm{\Delta }_n))=\mu (\mathrm{d}A)1_{\{A\mathrm{\Delta }_n\mathrm{}\}}_{i\mathrm{\Delta }_n}\mu (\mathrm{d}A)1_{\{iA\}}<\mathrm{}`$ for each $`n`$, the measure $`\mu `$ is locally finite.
We equip the space of locally finite measures on $`𝒫_+(\mathrm{\Lambda })`$ with the vague topology, i.e., we say that a sequence of locally finite measures $`\mu _n`$ on $`𝒫_+(\mathrm{\Lambda })`$ converges vaguely to a limit $`\mu `$, denoted as $`\mu _n\mu `$, if $`\mu _n(\mathrm{d}A)f(A)\mu (\mathrm{d}A)f(A)`$ for each continuous compactly supported real function $`f`$ on $`𝒫_+(\mathrm{\Lambda })`$.
###### Lemma 3.2
(Vague convergence) Let $`\mu _n,\mu `$ be locally finite measures on $`𝒫_+(\mathrm{\Lambda })`$. Then the $`\mu _n`$ converge vaguely to $`\mu `$ if and only if $`\mu _n(\mathrm{d}A)1_{\{AB\mathrm{}\}}\mu (\mathrm{d}A)1_{\{AB\mathrm{}\}}`$ for each $`B𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$. The sequence $`\mu _n`$ is relatively compact in the topology of vague convergence if and only if $`sup_n\mu _n(\mathrm{d}A)1_{\{AB\mathrm{}\}}<\mathrm{}`$ for each $`B𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$.
Proof Since for each $`B𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$, the function $`A1_{\{AB\mathrm{}\}}`$ is continuous with compact support, the conditions for convergence and relative compactness given above are clearly necessary. To see that they are also sufficient, let $`\mathrm{\Delta }_m\mathrm{\Lambda }`$ be finite sets increasing to $`\mathrm{\Lambda }`$ and set $`f_m(A):=1_{\{A\mathrm{\Delta }_n\mathrm{}\}}`$. Then the $`f_m`$ are continuous, nonnegative functions with compact supports increasing to $`𝒫_+(\mathrm{\Lambda })`$. It follows that $`\mu _n`$ converges vaguely to $`\mu `$ if and only if for each $`m`$ the weighted measures $`f_m(A)\mu _n(\mathrm{d}A)`$ converge weakly to $`f_m(A)\mu (\mathrm{d}A)`$. Now if $`sup_n\mu _n(\mathrm{d}A)1_{\{AB\mathrm{}\}}<\mathrm{}`$ for each $`B𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$, then by a diagonal argument each subsequence of the $`\mu _n`$ contains a further subsequence such that $`f_m(A)\mu _n(\mathrm{d}A)`$ converges weakly for each $`m`$, hence the $`\mu _n`$ converge vaguely. The linear span of the functions $`B1_{\{AB=\mathrm{}\}}`$ with $`A𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$ forms an algebra that separates points, hence by the Stone-Weierstrass theorem is dense in the space of continuous functions on $`𝒫(\mathrm{\Lambda })`$. It follows that $`\mu _n`$ converges vaguely to $`\mu `$ if and only if $`f_m(A)\mu _n(\mathrm{d}A)1_{\{AB\mathrm{}\}}`$ converges to $`f_m(A)\mu (\mathrm{d}A)1_{\{AB\mathrm{}\}}`$ for each $`m`$ and for each $`B𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$. Since $`f_m(A)1_{\{AB\mathrm{}\}}=1_{\{A\mathrm{\Delta }_n\mathrm{}\}}+1_{\{AB\mathrm{}\}}1_{\{A(B\mathrm{\Delta }_n)\mathrm{}\}}`$, this is in turn implied by the condition that $`\mu _n(\mathrm{d}A)1_{\{AB\mathrm{}\}}\mu (\mathrm{d}A)1_{\{AB\mathrm{}\}}`$ for each $`B𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$.
The next lemma guarantees that expressions as in the left-hand side of (1.8) are well-defined and yield a homogeneous, locally finite measure on $`𝒫_+(\mathrm{\Lambda })`$
###### Lemma 3.3
(Evolution of locally finite measures) If $`\mu `$ is a homogeneous, locally finite measure on $`𝒫_+(\mathrm{\Lambda })`$, then for each $`t0`$, the measure $`\mu (\mathrm{d}A)[\eta _t^A]|_{𝒫_+(\mathrm{\Lambda })}`$ is homogeneous and locally finite on $`𝒫_+(\mathrm{\Lambda })`$. If $`\mu _n`$ are homogeneous, locally finite measures on $`𝒫_+(\mathrm{\Lambda })`$ converging vaguely to a limit $`\mu `$, then
$$\mu _n(\mathrm{d}A)[\eta _t^A]|_{𝒫_+(\mathrm{\Lambda })}\underset{n\mathrm{}}{}\mu (\mathrm{d}A)[\eta _t^A]|_{𝒫_+(\mathrm{\Lambda })}(t0).$$
(3.14)
Proof We start by observing that any homogeneous, locally finite measure $`\mu `$ on $`𝒫_+(\mathrm{\Lambda })`$ satisfies
$$\mu (\mathrm{d}A)1_{\{AB\mathrm{}\}}\underset{iB}{}\mu (\mathrm{d}A)1_{\{iA\}}=|B|\mu (\mathrm{d}A)1_{\{0A\}}<\mathrm{}.$$
(3.15)
Using duality (see (2.4)), it follows that
$$\begin{array}{c}\mu (\mathrm{d}A)[\eta _t^A\mathrm{d}C]1_{\{CB\mathrm{}\}}=\mu (\mathrm{d}A)[\eta _t^AB\mathrm{}]=\mu (\mathrm{d}A)[A\eta _t^B\mathrm{}]\hfill \\ =[\eta _t^B\mathrm{d}C]\mu (\mathrm{d}A)1_{\{AC\mathrm{}\}}[\eta _t^B\mathrm{d}C]|C|\mu (\mathrm{d}A)1_{\{0A\}}\hfill \\ =𝔼\left[|\eta _t^B|\right]\mu (\mathrm{d}A)1_{\{0A\}}<\mathrm{}(B𝒫_{\mathrm{fin}}(\mathrm{\Lambda })).\hfill \end{array}$$
(3.16)
By Lemma 3.1, it follows that the measure $`\mu _n(\mathrm{d}A)[\eta _t^A]|_{𝒫_+(\mathrm{\Lambda })}`$ is locally finite. It is obviously homogeneous. Now if $`\mu _n`$ are homogeneous, locally finite measures on $`𝒫_+(\mathrm{\Lambda })`$ converging vaguely to a limit $`\mu `$, then, by the first three equalities in (3.16),
$$\mu _n(\mathrm{d}A)[\eta _t^A\mathrm{d}C]1_{\{CB\mathrm{}\}}=[\eta _t^B\mathrm{d}C]\mu _n(\mathrm{d}A)1_{\{AC\mathrm{}\}},$$
(3.17)
for each $`B𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$, and this quantity converges to the analogue quantity for $`\mu `$ by dominated convergence, using (2.5), the estimate (3.15), and the fact that the $`\mu _n(\mathrm{d}A)1_{\{0A\}}`$ are uniformly bounded in $`n`$ since they converge. Applying Lemma 3.2, we arrive at (3.14).
Proof of Proposition 1.4 It suffices to prove, for each $`(\mathrm{\Lambda },a,\delta )`$-contact process, the following three claims:
1. There exists a homogeneous eigenmeasure of the $`(\mathrm{\Lambda },a,\delta )`$-contact process, with eigenvalue $`r=r(\mathrm{\Lambda },a,\delta )`$.
2. If $`\lambda `$ is the eigenvalue of a homogeneous eigenmeasure of the $`(\mathrm{\Lambda },a,\delta )`$-contact process, then $`\lambda r`$.
3. The set $`(\mathrm{\Lambda },a,\delta )`$ is closed and bounded from below.
We start with claim 1. Define (by Lemma 3.3 applied to $`\mu =_i\delta _{\{i\}}`$) homogeneous, locally finite measures $`\mu _t`$ on $`𝒫_+(\mathrm{\Lambda })`$ by
$$\mu _t:=\underset{i}{}[\eta _t^{\{i\}}]|_{𝒫_+(\mathrm{\Lambda })}(t0).$$
(3.18)
Let $`\widehat{\mu }_\lambda `$ be the Laplace transform of $`(\mu _t)_{t0}`$, i.e.,
$$\widehat{\mu }_\lambda :=_0^{\mathrm{}}\mu _te^{\lambda t}dt(\lambda >r).$$
(3.19)
We claim that the measures $`\widehat{\mu }_\lambda `$ are locally finite and, properly renormalized, relatively compact in the topology of vague convergence, and that each subsequential limit as $`\lambda r`$ is a homogeneous eigenmeasure of the $`(\mathrm{\Lambda },a,\delta )`$-contact process, with eigenvalue $`r`$.
Note that by duality (see (2.4)),
$$\mu _t(\mathrm{d}A)1_{\{AB\mathrm{}\}}=\underset{i}{}[\eta _t^{\{i\}}B\mathrm{}]=\underset{i}{}[i\eta _t^B]=𝔼\left[|\eta _t^B|\right]=\pi _t^{}(B)$$
(3.20)
$`(t0,B𝒫_{\mathrm{fin}}(\mathrm{\Lambda }))`$, where $`\pi _t^{}(A)`$ is defined in analogy with (3.1) for the $`(\mathrm{\Lambda },a^{},\delta )`$-contact process. It follows that
$$\widehat{\mu }_\lambda (\mathrm{d}A)1_{\{AB\mathrm{}\}}=\widehat{\pi }_\lambda ^{}(B),$$
(3.21)
where
$$\widehat{\pi }_\lambda ^{}(A):=_0^{\mathrm{}}\pi _t^{}(A)e^{\lambda t}dt(\lambda >r,A𝒫_{\mathrm{fin}}(\mathrm{\Lambda })).$$
(3.22)
By (3.5) and Theorem 1.2 (a), which has been proved in Section 3.1, for every $`\epsilon >0`$, there exists a $`T_\epsilon <\mathrm{}`$ such that
$$e^{rt}\pi _t^{}(\{0\})e^{(r+\epsilon )t}(tT_\epsilon ).$$
(3.23)
It follows from the upper bound in (3.23) that $`\widehat{\pi }_\lambda ^{}(\{0\})<\mathrm{}`$ for all $`\lambda >r`$. Hence, by (3.21) and Lemma 3.1, the measures $`\widehat{\mu }_\lambda `$ are locally finite for each $`\lambda >r`$. The lower bound in (3.23) and monotone convergence moreover show that
$$\underset{\lambda r}{lim}\widehat{\pi }_\lambda ^{}(\{0\})=\underset{\lambda r}{lim}_0^{\mathrm{}}\pi _t^{}(\{0\})e^{\lambda t}dt=_0^{\mathrm{}}\pi _t^{}(\{0\})e^{rt}dt=\mathrm{}.$$
(3.24)
Set $`\overline{\mu }_\lambda :=\widehat{\pi }_\lambda ^{}(\{0\})^1\widehat{\mu }_\lambda `$. Then for each $`\lambda >r`$, the measure $`\overline{\mu }_\lambda `$ is homogenous, locally finite, and normalized such that $`\overline{\mu }_\lambda (\mathrm{d}A)1_{\{0A\}}=1`$. Therefore, by Lemma 3.2 and the estimate (3.15), the measures $`\overline{\mu }_\lambda `$ are relatively compact in the topology of vague convergence as $`\lambda r`$. Choose $`\lambda _nr`$ such that $`\overline{\mu }_{\lambda _n}\overline{\mu }_r`$ for some homogenous, locally finite $`\overline{\mu }_r`$. Obviously $`\overline{\mu }_r(\mathrm{d}A)1_{\{0A\}}=1`$ so $`\overline{\mu }_r`$ is nonzero. Then, filling in our definitions, using Lemma 3.3 and the Markov property of the contact process,
$$\begin{array}{c}\overline{\mu }_r(\mathrm{d}A)[\eta _t^A]|_{𝒫_+(\mathrm{\Lambda })}=\underset{n\mathrm{}}{lim}\widehat{\pi }_{\lambda _n}^{}(\{0\})^1_0^{\mathrm{}}e^{\lambda _ns}ds\underset{i}{}[\eta _s^{\{i\}}\mathrm{d}A][\eta _t^A]|_{𝒫_+(\mathrm{\Lambda })}\hfill \\ =e^{rt}\underset{n\mathrm{}}{lim}\widehat{\pi }_{\lambda _n}^{}(\{0\})^1_0^{\mathrm{}}e^{\lambda _n(s+t)}ds\underset{i}{}[\eta _{s+t}^{\{i\}}]|_{𝒫_+(\mathrm{\Lambda })}\hfill \\ =e^{rt}\left(\overline{\mu }_r\underset{n\mathrm{}}{lim}\widehat{\pi }_{\lambda _n}^{}(\{0\})^1_0^te^{\lambda _ns}ds\underset{i}{}[\eta _s^{\{i\}}]|_{𝒫_+(\mathrm{\Lambda })}\right)=e^{rt}\overline{\mu }_r,\hfill \end{array}$$
(3.25)
where in the last step we have used (3.24). This shows that $`\overline{\mu }_r`$ is an eigenmeasure with eigenvalue $`r`$.
To prove claim 2, we observe that if $`\mu `$ is a homogeneous eigenmeasure with eigenvalue $`\lambda `$, then by duality (see (2.4)) and (3.15),
$$\begin{array}{c}e^{\lambda t}\mu (\mathrm{d}A)1_{\{0A\}}=\mu (\mathrm{d}A)[0\eta _t^A]\hfill \\ =\mu (\mathrm{d}A)[\eta _t^{\{0\}}A\mathrm{}]𝔼\left[|\eta _t^{\{0\}}|\right]\mu (\mathrm{d}A)1_{\{0A\}}.\hfill \end{array}$$
(3.26)
By (3.23), for each $`\epsilon >0`$ we can choose $`t`$ large enough such that $`𝔼\left[|\eta _t^{\{0\}}|\right]e^{(r+\epsilon )t}`$. Since $`\mu (\mathrm{d}A)1_{\{0A\}}>0`$, we may divide by it, yielding $`e^{\lambda t}e^{(r+\epsilon )t}`$, which implies $`\lambda r+\epsilon `$. Since $`\epsilon >0`$ is arbitrary, it follows that $`\lambda r`$.
To prove claim 3, finally, we observe that since we may estimate a contact process from below by a simple death process, for any homogeneous eigenmeasure $`\mu `$ with eigenvalue $`\lambda `$, one has
$$e^{\lambda t}\mu (\mathrm{d}A)1_{\{0A\}}=\mu (\mathrm{d}A)[0\eta _t^A]e^{\delta t}\mu (\mathrm{d}A)1_{\{0A\}},$$
(3.27)
which shows that $`(\mathrm{\Lambda },a,\delta )[\delta ,\mathrm{})`$. To show that $`(\mathrm{\Lambda },a,\delta )`$ is closed, assume that $`\lambda _n(\mathrm{\Lambda },a,\delta )`$ and $`\lambda _n\lambda `$. Then we can find homogeneous eigenmeasures $`\mu _n`$ with eigenvalues $`\lambda _n`$. Normalizing such that $`\mu _n(\mathrm{d}A)1_{\{0A\}}=1`$, using Lemma 3.2 and (3.15), we see that the sequence $`\mu _n`$ is relatively compact in the topology of vague convergence, hence has a subsequential limit $`\mu `$, which by Lemma 3.3 is a homogeneous eigenmeasure with eigenvalue $`\lambda `$.
The proof of Proposition 1.4 yields a useful corollary.
###### Corollary 3.4
(Convergence to eigenmeasure) Let $`\mu _t`$ be defined as in (3.18). Then the measures
$$\widehat{\pi }_\lambda ^{}(\{0\})^1_0^{\mathrm{}}\mu _te^{\lambda t}dt$$
(3.28)
are relatively compact as $`\lambda r`$ in the topology of vague convergence, and each subsequential limit as $`\lambda r`$ is a homogeneous eigenmeasure of the $`(\mathrm{\Lambda },a,\delta )`$-contact process, with eigenvalue $`r(\mathrm{\Lambda },a,\delta )`$.
Remark It seems intuitively plausible that the measures $`\mu _t`$, suitably rescaled, converge as $`t\mathrm{}`$ to a vague limit, which by Corollary 3.4 then has to be an eigenmeasure with eigenvalue $`r`$. Indeed, it seems plausible that this convergence is monotone, in a suitable sense, and that these eigenmeasures are the ‘lowest’ possible eigenmeasures, in a suitable stochastic order. Should these conjectures be correct, then these eigenmeasures are quite similar to the ‘second lowest invariant measure’ from \[SS97, SS99\], by which they are inspired. I do not know if these conjectures are correct, or even what kind of stochastic order one should choose here; I spent quite a bit of time in vain trying to prove that the measures $`\mu _t`$ conditioned on the event that the origin is infected, are stochastically increasing in time (in the usual stochastic order).
### 3.3 Proof of Theorem 1.5
We start with a preparatory lemma. We say that a function $`f:𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$ is shift-invariant if $`f(iA)=f(A)`$ for all $`i\mathrm{\Lambda }`$, monotone if $`AB`$ implies $`f(A)f(B)`$, and subadditive if $`f(AB)f(A)+f(B)`$, for all $`A,B𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$. Recall the definition of the generator $`G`$ of the $`(\mathrm{\Lambda },a,\delta )`$-contact process from (2.11). We define $`G^{}`$ analogously, for the $`(\mathrm{\Lambda },a^{},\delta )`$-contact process.
###### Lemma 3.5
(Eigenmeasures and harmonic functions) If $`\mu `$ is a homogeneous eigenmeasure of the $`(\mathrm{\Lambda },a,\delta )`$-contact process with eigenvalue $`\lambda `$, then
$$v(A):=\mu (\mathrm{d}B)1_{\{AB\mathrm{}\}}(A𝒫_{\mathrm{fin}}(\mathrm{\Lambda }))$$
(3.29)
is a shift-invariant, monotone, subadditive function such that $`v(\mathrm{})=0`$, $`v(A)>0`$ for any $`\mathrm{}A𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$, $`v𝒮(𝒫_{\mathrm{fin}}(\mathrm{\Lambda }))`$, and $`G^{}v=\lambda v`$.
Proof The function $`v`$ is obviously shift-invariant, monotone, and satisfies $`v(\mathrm{})=0`$. Since $`\mu `$ is homogeneous and nonzero, $`v(A)>0`$ for any $`\mathrm{}A𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$. The function $`v`$ is subadditive since $`1_{\{(AA^{})B\mathrm{}\}}1_{\{AB\mathrm{}\}}+1_{\{A^{}B\mathrm{}\}}`$. Subadditivity and shift-invariance imply that $`v(A)v(\{0\})|A|`$, so certainly $`v𝒮(𝒫_{\mathrm{fin}}(\mathrm{\Lambda }))`$.
To see that $`G^{}v=\lambda v`$, observe that by duality (see (2.4))
$$\begin{array}{c}𝔼[v(\eta _t^A)]=\mu (\mathrm{d}B)[\eta _t^AB\mathrm{}]=\mu (\mathrm{d}B)[A\eta _t^B\mathrm{}]\hfill \\ =\mu (\mathrm{d}B)[\eta _t^B\mathrm{d}C]1_{\{AC\mathrm{}\}}=e^{\lambda t}\mu (\mathrm{d}C)1_{\{AC\mathrm{}\}}=e^{\lambda t}v(A).\hfill \end{array}$$
(3.30)
Let $`(S_t^{})_{t0}`$ denote the semigroup of the $`(\mathrm{\Lambda },a^{},\delta )`$-contact process. Recall from Section 2.3 that $`S_t^{}`$ maps the space $`𝒮(𝒫_{\mathrm{fin}}(\mathrm{\Lambda }))`$ into itself. Then (3.30) says that $`S_t^{}v=e^{\lambda t}v`$, and hence, by (2.25), $`G^{}v=lim_{\epsilon 0}\epsilon ^1(S_\epsilon ^{}vv)=\lambda v`$.
Recall the definition of the survival probability $`\rho (A)`$ from (2.6). Theorem 1.5 follows from the following, stronger result.
###### Proposition 3.6
(Shift invariant monotone harmonic functions) Assume that the infection rates satisfy the irreducibility condition (1.7) and that the $`(\mathrm{\Lambda },a,\delta )`$-contact process survives. Assume that $`f:𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$ is shift invariant, monotone, $`f(\mathrm{})=0`$, $`f𝒮(𝒫_{\mathrm{fin}}(\mathrm{\Lambda }))`$, and $`Gf=0`$. Then there exists a constant $`c0`$ such that $`f=c\rho `$.
Before we prove this, we first show how this implies Theorem 1.5.
Proof of Theorem 1.5 Let $`\mu `$ be a homogeneous eigenmeasure of the $`(\mathrm{\Lambda },a,\delta )`$-contact process with eigenvalue zero, and let $`v(A):=\mu (\mathrm{d}B)1_{\{AB\mathrm{}\}}`$. By Lemma 3.5, $`v`$ is shift invariant, monotone, $`v(\mathrm{})=0`$, $`v𝒮(𝒫_{\mathrm{fin}}(\mathrm{\Lambda }))`$, and $`G^{}v=0`$. By assumption, the upper invariant measure $`\overline{\nu }`$ of the $`(\mathrm{\Lambda },a,\delta )`$-contact process is nontrivial, hence the $`(\mathrm{\Lambda },a^{},\delta )`$-contact process survives, so by Proposition 3.6, $`v=c\rho ^{}`$ for some $`c0`$, where $`\rho ^{}`$ denotes the survival probability of the $`(\mathrm{\Lambda },a^{},\delta )`$-contact process. By the characterization of the upper invariant measure in (2.9), it follows that $`\mu =c\overline{\nu }`$.
In order to prove Proposition 3.6, we need one more lemma.
###### Lemma 3.7
(Eventual domination of finite configurations) Assume that the infection rates satisfy the irreducibility condition (1.7) and that the $`(\mathrm{\Lambda },a,\delta )`$-contact process survives. Then
$$\underset{t\mathrm{}}{lim}[i\mathrm{\Lambda }\text{ s.t. }\eta _t^AiB|\eta _t^A\mathrm{}]=1(A,B𝒫_{\mathrm{fin}}(\mathrm{\Lambda }),A\mathrm{}).$$
(3.31)
Formula (3.31) says that $`\eta `$ exhibits a form of ‘extinction versus unbounded growth’. More precisely, either $`\eta _t`$ gets extinct or $`\eta _t`$ is eventually larger than a random shift of any finite configuration.
We remark that Lemma 3.7 is no longer true if the infection rates fail to satisfy the first condition in (1.7). Indeed, if $`A`$ is defined as in (1.7) and $`_{n0,m0}A^nA^m\mathrm{\Lambda }`$, then we can find sites $`i,j\mathrm{\Lambda }`$ such that there exists no site $`k`$ from which both $`i`$ and $`j`$ can be infected. In particular, if we set $`B:=\{i,j\}`$, then $`[i^{}\mathrm{\Lambda }\text{ s.t. }\eta _t^{\{0\}}i^{}B]=0`$ for all $`t0`$. For example, this happens if $`\mathrm{\Lambda }`$ is the free group with two generators, say $`g_1`$ and $`g_2`$, $`A=\{g_1,g_2\}`$, and $`B=\{g_1^1,g_2^1\}`$.
Proof of Proposition 3.6 Since the $`(\mathrm{\Lambda },a,\delta )`$-contact process solves the martingale problem for $`G`$, and $`Gf=0`$, the process $`f(\eta _t^A)`$ is a martingale. In particular:
$$f(A)=𝔼[f(\eta _t^A)](A𝒫_{\mathrm{fin}}(\mathrm{\Lambda }),t0).$$
(3.32)
Enumerate the elements of $`\mathrm{\Lambda }`$ an arbitrary way, and for $`A,B𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$, put
$$\widehat{ı}_{A,B}:=\{\begin{array}{cc}\mathrm{min}\{i\mathrm{\Lambda }:AiB\}\hfill & \text{if }\{i\mathrm{\Lambda }:AiB\}\text{ is nonempty,}\hfill \\ 0\hfill & \text{otherwise.}\hfill \end{array}$$
(3.33)
Since $`f`$ is monotone and shift invariant, we have, using Lemma 3.7,
$$\begin{array}{ccc}\hfill f(A)& =& \underset{t\mathrm{}}{lim}𝔼[f(\eta _t^A)]\hfill \\ & & \underset{t\mathrm{}}{lim\; sup}𝔼[1_{\left\{i\mathrm{\Lambda }\text{ s.t. }\eta _t^AiB\right\}}f(\widehat{ı}_{\eta _t^A,B}B)]\hfill \\ & =& f(B)\underset{t\mathrm{}}{lim\; sup}[i\mathrm{\Lambda }\text{ s.t. }\eta _t^AiB]=f(B)\rho (A)(A,B𝒫_{\mathrm{fin}}(\mathrm{\Lambda })).\hfill \end{array}$$
(3.34)
In particular, this shows that
$$f(B)\frac{f(\{0\})}{\rho (\{0\})}<\mathrm{}(B𝒫_{\mathrm{fin}}(\mathrm{\Lambda })),$$
(3.35)
hence $`f`$ is bounded. Now let $`A_n,B_m𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$ be sequences such that $`\rho (A_n)1`$ and $`\rho (B_n)1`$. Then, by (3.34),
$$\underset{n\mathrm{}}{lim\; inf}f(A_n)\underset{n\mathrm{}}{lim\; inf}f(B_m)\rho (A_n)=f(B_m)m,$$
(3.36)
and therefore
$$\underset{n\mathrm{}}{lim\; inf}f(A_n)\underset{m\mathrm{}}{lim\; sup}f(B_m).$$
(3.37)
This proves that the limit
$$\underset{\rho (A_n)1}{lim}f(A_n)=:f(\mathrm{})$$
(3.38)
exists and does not depend on the choice of the sequence $`A_n`$ with $`\rho (A_n)1`$. By the Markov property and continuity of the conditional expectation with respect to increasing limits of $`\sigma `$-fields (see Complement 10(b) from \[Loe63, Section 29\] or \[Loe78, Section 32\]),
$$\rho (\eta _t^A)=[\eta _s^A\mathrm{}s0|\eta _t^A]1_{\left\{\eta _s^A\mathrm{}s0\right\}}\mathrm{a}.\mathrm{s}.\text{as }t\mathrm{}.$$
(3.39)
We conclude that, for all $`A𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$,
$$\begin{array}{c}f(A)=\underset{t\mathrm{}}{lim}𝔼[f(\eta _t^A)]\hfill \\ =\left[\eta _t^A=\mathrm{}\text{ for some }t0\right]f(0)+\left[\underset{t\mathrm{}}{lim}\rho (\eta _t^A)=1\right]f(\mathrm{})=\rho (A)f(\mathrm{}),\hfill \end{array}$$
(3.40)
which shows that $`f`$ is a scalar multiple of $`\rho `$.
The proof of Lemma 3.7 depends on two preparatory lemmas.
###### Lemma 3.8
(Local creation of finite configurations) For each $`B𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$ and $`t>0`$, there exists a finite $`\mathrm{\Delta }\mathrm{\Lambda }`$ and $`j\mathrm{\Lambda }`$ such that
$$\epsilon :=\left[\eta _t^{\{0\}}jB\text{ and }\eta _s^{\{0\}}\mathrm{\Delta }0st\right]>0.$$
(3.41)
Proof It follows from assumption (1.7) that there exists a site $`j^1\mathrm{\Lambda }`$ with $`[\eta _t^{\{j^1\}}B]>0`$, and therefore $`[\eta _t^{\{0\}}jB]>0`$. Since $`_{0st}\eta _s^{\{0\}}`$ is a.s. finite, we can choose a finite but large enough $`\mathrm{\Delta }`$ such that (3.41) holds.
###### Lemma 3.9
(Domination of finite configurations) For each $`B𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$, $`t>0`$, and $`A_n𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$ satisfying $`lim_n\mathrm{}|A_n|=\mathrm{}`$, one has
$$\underset{n\mathrm{}}{lim}[i\mathrm{\Lambda }\text{ s.t. }\eta _t^{A_n}iB]=1.$$
(3.42)
Proof Let $`\mathrm{\Delta }`$, $`j`$, and $`\epsilon `$ be as in Lemma 3.8. We can find $`\stackrel{~}{A}_nA_n`$ such that $`|\stackrel{~}{A}_n|\mathrm{}`$ as $`n\mathrm{}`$, and for fixed $`n`$, the sets $`(k\mathrm{\Delta })_{k\stackrel{~}{A}_n}`$ are disjoint. It follows that
$$\begin{array}{c}[i\mathrm{\Lambda }\text{ s.t. }\eta _t^{A_n}iB]\hfill \\ 1\underset{k\stackrel{~}{A}_n}{}\left(1\left[\eta _t^{\{k\}}kjB\text{ and }\eta _s^{\{k\}}k\mathrm{\Delta }0st\right]\right)\hfill \\ =1(1\epsilon )^{|\stackrel{~}{A}_n|}\underset{n\mathrm{}}{}1,\hfill \end{array}$$
(3.43)
where we have used (3.41) and the fact that events concerning the graphical representation in disjoint parts of space are independent.
Proof of Lemma 3.7 If $`\delta =0`$, then obviously $`lim_t\mathrm{}|\eta _t^A|=\mathrm{}`$ a.s. If $`\delta >0`$, then it is easy to see that $`sup\{\rho (A):|A|M\}<1`$ for all $`M<\mathrm{}`$. Therefore, by (3.39),
$$\eta _t^A=\mathrm{}\text{ for some }t0\text{or}|\eta _t^A|\underset{t\mathrm{}}{}\mathrm{}\mathrm{a}.\mathrm{s}.$$
(3.44)
Fix $`\mathrm{}B𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$ and set $`\psi _t(A):=P[i\mathrm{\Lambda }\text{ s.t. }\eta _t^AiB](A𝒫_{\mathrm{fin}}(\mathrm{\Lambda }),t0)`$. Then, for each $`t>0`$,
$$\underset{T\mathrm{}}{lim}P[i\mathrm{\Lambda }\text{ s.t. }\eta _T^AiB]=\underset{T\mathrm{}}{lim}E[\psi _t(\eta _{Tt}^A)]=\rho (A),$$
(3.45)
where we have used Lemma 3.9 and (3.44).
## 4 Proof of the main results
### 4.1 Exponentially growing processes
In this section, we prove Theorem 1.2 (d). Indeed, we prove the following, more detailed result. Recall that by Theorem 1.2 (a) and Proposition 1.4, there exists a homogeneous eigenmeasure for the $`(\mathrm{\Lambda },a^{},\delta )`$-contact process, with eigenvalue $`r=r(\mathrm{\Lambda },a^{},\delta )=r(\mathrm{\Lambda },a,\delta )`$.
###### Proposition 4.1
(Exponential growth) Let $`\mathrm{\Lambda }`$ be a finite or countably infinite group, let $`a=(a(i,j))_{i,j\mathrm{\Lambda }}`$ be infection rates satisfying (1.1) and let $`\delta 0`$. Let $`\mu `$ be any homogeneous eigenmeasure of the $`(\mathrm{\Lambda },a^{},\delta )`$-contact process, with eigenvalue $`r`$, and let the function $`v`$ be defined in terms of $`\mu `$ as in Lemma 3.5. If $`r>0`$, then, for each $`A𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$, the limit
$$W_A:=\underset{t\mathrm{}}{lim}e^{rt}v(\eta _t^A)$$
(4.1)
exists a.s., and satisfies $`𝔼[W_A]=v(A)`$. If the infection rates satisfy (1.7), then moreover
$$[W_A>0]=[\eta _t^A\mathrm{}t0].$$
(4.2)
Proof Our proof follows a strategy that is familiar from the theory of supercritical branching processes. Using duality (see (2.4)) and the fact that $`\mu `$ is an eigenmeasure, we see that
$$\begin{array}{c}𝔼[v(\eta _t^A)]=\mu (\mathrm{d}B)[\eta _t^AB\mathrm{}]=\mu (\mathrm{d}B)[A\eta _t^B\mathrm{}]\hfill \\ =e^{rt}\mu (\mathrm{d}B)[AB\mathrm{}]=e^{rt}v(A),\hfill \end{array}$$
(4.3)
so, by the Markov property of $`(\eta _t^A)_{t0}`$, the process $`(e^{rt}v(\eta _t^A))_{t0}`$ is a martingale. Every nonnegative martingale converges, so for each $`A𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$, there exists a random variable $`W_A`$ such that (4.1) holds.
To prove that $`𝔼[W_A]=v(A)`$, it suffices to show that the random variables $`\{e^{rt}v(\eta _t^A):t0\}`$ are uniformly integrable. By Proposition 2.2, the variance of $`v(\eta _t^A)`$ is given by
$$\mathrm{Var}(v(\eta _t^A))=2_0^t𝔼\left[\mathrm{\Gamma }(S_sv,S_sv)(\eta _{ts}^A)\right]ds,$$
(4.4)
where $`S_t`$ and $`\mathrm{\Gamma }`$ are defined in (2.17) and (2.19). Formula (4.3) tells us that $`S_tv=e^{rt}v`$, so
$$\mathrm{Var}(v(\eta _t^A))=2_0^t𝔼\left[\mathrm{\Gamma }(e^{rs}v,e^{rs}v)(\eta _{ts}^A)\right]ds=2_0^t𝔼\left[\mathrm{\Gamma }(v,v)(\eta _{ts}^A)\right]e^{2rs}ds.$$
(4.5)
It is not hard to see that for any $`f,g𝒮(𝒫_{\mathrm{fin}}(\mathrm{\Lambda }))`$ and $`A𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$,
$$\begin{array}{ccc}\hfill \mathrm{\Gamma }(f,g)(A)& =& \frac{1}{2}\underset{ij}{}a(i,j)1_{\{iA\}}1_{\{jA\}}\left(f(A\{j\})f(A)\right)\left(g(A\{j\})g(A)\right)\hfill \\ & & +\frac{1}{2}\delta \underset{i}{}1_{\{iA\}}\left(f(A\backslash \{i\})f(A)\right)\left(g(A\backslash \{i\})g(A)\right).\hfill \end{array}$$
(4.6)
Without loss of generality we can normalize $`v`$ such that $`v(\{0\})=1`$. Then, by monotonicity and subadditivity (see Lemma 3.5), $`0v(A\{j\})v(A)1`$ for all $`j,A`$, and therefore
$$\mathrm{\Gamma }(v,v)(A)\frac{1}{2}\underset{ij}{}a(i,j)1_{\{iA\}}1_{\{jA\}}+\frac{1}{2}\delta \underset{i}{}1_{\{iA\}}\frac{1}{2}\left(|a|+\delta \right)|A|.$$
(4.7)
Inserting this into (4.5) yields
$$\begin{array}{ccc}\hfill \mathrm{Var}(e^{rt}v(\eta _t^A))& & \left(|a|+\delta \right)e^{2rt}_0^t𝔼\left[|\eta _{ts}^A|\right]e^{2rs}ds\hfill \\ & & \left(|a|+\delta \right)_0^t𝔼\left[|\eta _{ts}^A|\right]e^{2r(ts)}ds.\hfill \end{array}$$
(4.8)
Since $`r`$ is the exponential growth rate of $`\eta ^A`$ and $`r>0`$, we can find $`K<\mathrm{}`$ such that $`𝔼\left[|\eta _t^A|\right]Ke^{\frac{3}{2}rt}`$ $`(t0)`$. It follows that
$$\mathrm{Var}(e^{rt}v(\eta _t^A))(|a|+\delta )K_0^{\mathrm{}}e^{\frac{1}{2}rs}ds<\mathrm{}(t0),$$
(4.9)
which proves the required uniform integrability.
Set $`f(A):=[W_A>0]`$ and recall that $`\rho (A):=[\eta _t^A\mathrm{}t0]`$. Obviously $`f\rho `$. We have just shown that $`[W_A>0]>0`$ if $`A\mathrm{}`$, so $`\rho (A)f(A)>0`$ for each $`A\mathrm{}`$. Assuming that the infection rates satisfy (1.7), we claim that $`f=\rho `$. We observe that
$$f(\eta _t^A)=\left[\underset{s\mathrm{}}{lim}e^{rs}\eta _s^A>0|\eta _t^A\right].$$
(4.10)
In particular, this shows that $`(f(\eta _t^A))_{t0}`$ is a martingale, hence $`Gf=0`$. It is easy to see that $`f`$ is shift-invariant, monotone, bounded, and satisfies $`f(\mathrm{})=0`$, so applying Proposition 3.6, we see that $`f=c\rho `$ for some $`c0`$. Since $`f\rho `$, we have $`c1`$.
By continuity of the conditional expectation with respect to increasing limits of $`\sigma `$-fields (compare (3.39)), the right-hand side of (4.10) converges a.s. to the indicator function of the event that $`W_A>0`$. Since this event has positive probability, the event $`lim_t\mathrm{}f(\eta _t^A)=1`$ has positive probability. In particular, this shows that for each $`\epsilon >0`$ there exists a finite set $`B`$ with $`f(B)1\epsilon `$. This is possible only if the constant $`c`$ in the equation $`f=c\rho `$ satisfies $`c1`$.
### 4.2 Subexponential lattices
Proof of Theorem 1.2 (e) Consider a branching process on $`\mathrm{\Lambda }`$, started with one particle in the origin, where a particle at $`i`$ produces a new particle at $`j`$ with rate $`a(i,j)`$, and each particle dies with rate $`\delta `$. Let $`B_t(i)`$ denote the number of particles at site $`i\mathrm{\Lambda }`$ and time $`t0`$. It is not hard to see that $`\eta ^{\{0\}}`$ and $`B`$ may be coupled such that
$$1_{\eta _t^{\{0\}}}B_t(t0).$$
(4.11)
Let $`(\xi _t)_{t0}`$ be a random walk on $`\mathrm{\Lambda }`$ that jumps from $`i`$ to $`j`$ with rate $`a(i,j)`$, started in $`\xi _0=0`$. Then it is not hard to see that (compare \[Lig99, Proposition I.1.21\])
$$𝔼[B_t(i)]=[\xi _t=i]e^{(|a|\delta )t}(i\mathrm{\Lambda },t0).$$
(4.12)
Let $`h>0`$ be a constant, to be determined later. It follows from (4.11) and (4.12) that
$$\begin{array}{ccc}\hfill 𝔼\left[|\eta _t^{\{0\}}|\right]& & \underset{i}{}\left(1[\xi _t=i]e^{(|a|\delta )t}\right)\hfill \\ & =& |\{i\mathrm{\Lambda }:|i|ht\}|+[|\xi _t|>ht]e^{(|a|\delta )t}(t0).\hfill \end{array}$$
(4.13)
Let $`(Y_i)_{i1}`$ be i.i.d. $``$-valued random variables with $`[Y_i=k]=\frac{1}{|a|}_{j:|j|=k}a(0,j)`$ $`(k0)`$, let $`N`$ be a Poisson-distributed random variable with mean $`|a|`$, independent of the $`(Y_i)_{i1}`$, and let $`(X_m)_{m1}`$ be i.i.d. random variables with law $`[X_m]=[_{i=1}^NY_i]`$. Since the random walk $`\xi `$ makes jumps whose sizes are distributed in the same way as the $`Y_i`$, and the number of jumps per unit of time is Poisson distributed with mean $`|a|`$, it follows that
$$[|\xi _t|>ht]\left[\frac{1}{t}\underset{m=1}{\overset{t}{}}X_m>h\frac{t}{t}\right](t>0),$$
(4.14)
where $`t`$ denotes $`t`$ rounded up to the next integer. By our assumptions,
$$𝔼\left[\text{e}^{\epsilon X_m}\right]=𝔼\left[\text{e}^{\epsilon _{i=1}^NY_k}\right]=e^{|a|}\underset{n=0}{\overset{\mathrm{}}{}}\frac{|a|^n}{n!}𝔼\left[\text{e}^{\epsilon Y_1}\right]^n=\text{e}^{\left|a\right|\left(1𝔼\left[e^{\epsilon Y_1}\right]\right)}<\mathrm{},$$
(4.15)
for some $`\epsilon >0`$. Therefore, by \[DZ98, Theorem 2.2.3 and Lemma 2.2.20\], for each $`R>0`$ there exists a $`h>0`$ and $`K<\mathrm{}`$ such that
$$\left[\frac{1}{n}\underset{m=1}{\overset{n}{}}X_m>h\right]K\text{e}^{nR}(n1).$$
(4.16)
Choosing $`h`$ such that (4.16) holds for some $`R>|a|\delta `$ yields, by (4.14)
$$\underset{t\mathrm{}}{lim}\left[|\xi _t|>ht\right]e^{(|a|\delta )t}=0.$$
(4.17)
Inserting this into (4.13) we find that the exponential growth rate $`r=r(\mathrm{\Lambda },a,\delta )`$ satisfies
$$r\underset{t\mathrm{}}{lim\; sup}\frac{1}{t}\mathrm{log}|\{i\mathrm{\Lambda }:|i|ht\}|=0,$$
(4.18)
where we have used that the group $`\mathrm{\Lambda }`$ has subexponential growth.
### 4.3 Nonamenable lattices
In this section, we prove Theorem 1.2 (f) and Corollary 1.3. We start by introducing some notation. If $`R`$ is a nonnegative real random variable, defined on some probability space $`(\mathrm{\Omega },,)`$, and $`0<𝔼[R]<\mathrm{}`$, then we define the size-biased law $`\overline{}_R`$ associated with $`R`$ by
$$\overline{}_R(𝒜):=\frac{𝔼[1_𝒜R]}{𝔼[R]}(𝒜).$$
(4.19)
If $`\mathrm{\Delta }`$ is a $`𝒫_{\mathrm{fin}}(\mathrm{\Lambda })`$-valued random variable, defined on some probability space $`(\mathrm{\Omega },,)`$, such that $`0<𝔼[|\mathrm{\Delta }|]<\mathrm{}`$, then we define a probability law $`\widehat{}=\widehat{}_\mathrm{\Delta }`$ on the product space $`\mathrm{\Omega }\times \mathrm{\Lambda }`$ by
$$\widehat{}_\mathrm{\Delta }(𝒜\times \{i\}):=\frac{(\{i\mathrm{\Delta }\}𝒜)}{𝔼[|\mathrm{\Delta }|]}(𝒜,i\mathrm{\Lambda }).$$
(4.20)
We call $`\widehat{}_\mathrm{\Delta }`$ the Campbell law associated with $`\mathrm{\Delta }`$. It is not hard to see that the projection of $`_\mathrm{\Delta }`$ onto $`\mathrm{\Omega }`$ is the size-biased law $`\overline{}_{|\mathrm{\Delta }|}`$. Moreover, if we let $`\iota (\omega ,i):=i`$ denote the projection from $`\mathrm{\Omega }\times \mathrm{\Lambda }`$ to $`\mathrm{\Lambda }`$ and we use the symbol $`\mathrm{\Delta }`$ to denote (also) the random variable on $`\mathrm{\Omega }\times \mathrm{\Lambda }`$ defined by $`\mathrm{\Delta }(\omega ,i):=\mathrm{\Delta }(\omega )`$, then
$$\widehat{}_\mathrm{\Delta }\left[\iota =i|\mathrm{\Delta }\right]=\frac{1}{|\mathrm{\Delta }|}1_\mathrm{\Delta }(i),$$
(4.21)
i.e., conditional on $`\mathrm{\Delta }`$, the site $`\iota `$ is chosen with equal probabilities from all sites in $`\mathrm{\Delta }`$. We may view $`\iota `$ as a ‘typical’ element of $`\mathrm{\Delta }`$. Campbell laws are closely related to the more widely known Palm laws; both play an important role in the theory of branching processes (see, e.g., \[Win99\]).
The next lemma relates Campbell laws to things we have been considering so far. Note that if $`\mu `$ is a locally finite measure on $`𝒫_+(\mathrm{\Lambda })`$ and $`1_{\{0A\}}\mu (\mathrm{d}A)>0`$, then the conditional law
$$\mu (\mathrm{d}A|\mathrm{\hspace{0.17em}0}A):=\frac{1_{\{0A\}}\mu (\mathrm{d}A)}{1_{\{0B\}}\mu (\mathrm{d}B)}$$
(4.22)
is a well-defined probability law.
###### Lemma 4.2
(Campbell law) Let $`\eta `$ be a $`(\mathrm{\Lambda },a,\delta )`$-contact process. For each $`t0`$, let $`\mu _t`$ be defined as in (3.18) and let $`\widehat{}_t:=\widehat{}_{\eta _t^{\{0\}}}`$ be the Campbell law associated with $`\eta _t^{\{0\}}`$. Then
$$\mu _t(\mathrm{d}A|\mathrm{\hspace{0.17em}0}A)=\widehat{}_t\left[\iota ^1\eta _t^{\{0\}}\mathrm{d}A\right]$$
(4.23)
Proof This follows by writing
$$\begin{array}{c}\widehat{}_t\left[\iota ^1\eta _t^{\{0\}}\mathrm{d}A\right]=\underset{i}{}\widehat{}_t\left[i^1\eta _t^{\{0\}}\mathrm{d}A,\iota =i\right]\hfill \\ =\underset{i}{}\frac{\left[i^1\eta _t^{\{0\}}\mathrm{d}A,i\eta _t^{\{0\}}\right]}{𝔼[|\eta _t^{\{0\}}|]}=\frac{_i\left[i^1\eta _t^{\{0\}}\mathrm{d}A,i\eta _t^{\{0\}}\right]}{_i[i\eta _t^{\{0\}}]}\hfill \\ =\frac{_i\left[i^1\eta _t^{\{0\}}\mathrm{d}A,0i^1\eta _t^{\{0\}}\right]}{_i[0i^1\eta _t^{\{0\}}]}=\frac{_i\left[\eta _t^{\{i^1\}}\mathrm{d}A,0\eta _t^{\{i^1\}}\right]}{_i[0\eta _t^{\{i^1\}}]}\hfill \\ =\frac{_j\left[\eta _t^{\{j\}}\mathrm{d}A,0\eta _t^{\{j\}}\right]}{_j[0\eta _t^{\{j\}}]}=\frac{1_{\{0A\}}_j\left[\eta _t^{\{j\}}\mathrm{d}A\right]}{1_{\{0B\}}_j\left[\eta _t^{\{j\}}\mathrm{d}B\right]}\hfill \\ =\frac{1_{\{0A\}}\mu _t(\mathrm{d}A)}{1_{\{0B\}}\mu _t(\mathrm{d}B)}=\mu _t(\mathrm{d}A|\mathrm{\hspace{0.17em}0}A).\hfill \end{array}$$
(4.24)
The next proposition is a direct consequence of Theorem 1.5 and Corollary 3.4. Note that by the remark below Corollary 3.4, we expect the convergence in (4.25) to hold also when the $`\tau _\gamma `$ are replaced by deterministic times tending to infinity.
###### Proposition 4.3
(Convergence of Campbell laws) Assume that the $`(\mathrm{\Lambda },a,\delta )`$-contact process has a nontrivial upper invariant measure $`\overline{\nu }`$, that its exponential growth rate $`r(\mathrm{\Lambda },a,\delta )`$ is zero, and that the infection rates satisfy (1.7). Let $`\eta ^{\{0\}}`$ be the $`(\mathrm{\Lambda },a,\delta )`$-contact process started in $`\{0\}`$ and for $`\gamma 0`$, let $`\tau _\gamma `$ be an exponentially distributed random variable with mean $`\gamma `$, independent of $`\eta ^{\{0\}}`$. For each $`\gamma 0`$, let $`\widehat{}_\gamma =\widehat{}_{\eta _{\tau _\gamma }^{\{0\}}}`$ be the Campbell law associated with $`\eta _{\tau _\gamma }^{\{0\}}`$. Then
$$\widehat{}_\gamma \left[\iota ^1\eta _{\tau _\gamma }^{\{0\}}\mathrm{d}A\right]\underset{\gamma \mathrm{}}{}\overline{\nu }(\mathrm{d}A|\mathrm{\hspace{0.17em}0}A),$$
(4.25)
where $``$ denotes weak convergence of probability measures.
Proof For $`\lambda >0`$, let $`\widehat{\mu }_\lambda `$ denote the Laplace transform of $`\mu _t`$, defined in (3.19). In analogy with Lemma 4.2, it is straightforward to check that
$$\widehat{\mu }_\lambda (\mathrm{d}A|\mathrm{\hspace{0.17em}0}A)=\widehat{}_{1/\lambda }\left[\iota ^1\eta _{\tau _{1/\lambda }}^{\{0\}}\mathrm{d}A\right](\lambda >0).$$
(4.26)
By Theorem 1.5 and Corollary 3.4, the measures $`\widehat{\mu }_\lambda `$, suitably rescaled, converge vaguely to $`\overline{\nu }`$ as $`\lambda 0`$. By (4.26), this implies the weak convergence in (4.25).
The next proposition shows that if the assumptions of Proposition 4.3 are satisfied, then $`\mathrm{\Lambda }`$ must be amenable.
###### Proposition 4.4
(Campbell laws and amenability) Let $`\mathrm{\Lambda }`$ be a countable group, let $`B_n`$ be random nonempty, finite subsets of $`\mathrm{\Lambda }`$, and conditional on $`B_n`$, let $`\iota _n`$ be chosen with equal probabilities from the sites in $`B_n`$. Let $`B`$ be a random subset of $`\mathrm{\Lambda }`$ whose law is nontrivial and homogeneous. Assume that
$$[\iota _n^1B_n]\underset{n\mathrm{}}{}[B|\mathrm{\hspace{0.17em}0}B].$$
(4.27)
Then $`\mathrm{\Lambda }`$ must be amenable.
Proof The idea behind Proposition 4.4 is easy to explain: formula (4.27) says that for large $`n`$, the set $`B_n`$ looks like a random finite piece cut out of the spatially homogeneous configuration $`B`$, such that most points in this piece are far from the boundary. This contradicts nonamenability, since in any finite subset of a nonamenable group, a positive fraction of the points must lie near the boundary.
To make this idea rigorous, we proceed as follows. Assume that $`\mathrm{\Lambda }`$ is nonamenable. Then there exists a finite nonempty $`\mathrm{\Delta }\mathrm{\Lambda }`$ and $`\epsilon >0`$ such that $`|(A\mathrm{\Delta })\mathrm{}A|\epsilon |A|`$ for all finite nonempty $`A\mathrm{\Lambda }`$. Without loss of generality we may assume that $`\mathrm{\Delta }`$ is symmetric. Let $`(\xi _m)_{m0}`$ be a random walk in $`\mathrm{\Lambda }`$, independent of $`B_n`$ and $`B`$, starting in $`\xi _0=0`$, that jumps from a point $`i`$ to a point $`ij`$ with probability $`|\mathrm{\Delta }|^11_{\{j\mathrm{\Delta }\}}`$. Then (4.27) implies that
$$[\iota _n\xi _mB_n]\underset{n\mathrm{}}{}[\xi _mB|\xi _0B](m0).$$
(4.28)
By the stationarity of the process $`(1_{\{\xi _mB\}})_{m0}`$, one has
$$\underset{m\mathrm{}}{lim\; sup}[\xi _mB|\xi _0B]>0.$$
(4.29)
On the other hand, we will show that the nonamenability of $`\mathrm{\Lambda }`$ implies that
$$\underset{m\mathrm{}}{lim}\underset{n0}{sup}[\iota _n\xi _mB_n]=0,$$
(4.30)
which with (4.29) leads to a contradiction in (4.28). Let $`\mathrm{}^2(\mathrm{\Lambda })`$ be the Hilbert space of square summable real functions on $`\mathrm{\Lambda }`$, equipped with the inner product $`x,y:=_ix(i)y(i)`$, let $`P^n(i,j):=[\xi _n=j|\xi _0=i]`$ and $`P^nx(i):=_jP^n(i,j)x(j)`$. Then, by the fact that nearest-neighbor random walk on any nonamenable Cayley graph has a spectral gap (see \[Kes59\] or \[LP08, Thm 6.7\]), there exists a $`0<\theta <1`$ such that
$$|B_n|[\iota _n\xi _mB_n]=\underset{iB_n}{}\underset{jB_n}{}P^m(i,j)=1_{B_n},P^m1_{B_n}\theta ^m1_{B_n},1_{B_n}=\theta ^m|B_n|,$$
(4.31)
which proves (4.30).
Proof of Theorem 1.2 (f) Assume (1.7). Assume that the $`(\mathrm{\Lambda },a,\delta )`$-contact process survives and that its exponential growth rate $`r(\mathrm{\Lambda },a,\delta )`$ is zero. Then the $`(\mathrm{\Lambda },a^{},\delta )`$-contact process has a nontrivial upper invariant law, and, by Theorem 1.2 (a), $`r(\mathrm{\Lambda },a^{},\delta )=0`$. Therefore, by Propositions 4.3 and 4.4, $`\mathrm{\Lambda }`$ must be amenable.
Proof of Corollary 1.3 Let $`𝒮:=\{\delta 0:\text{the }(\mathrm{\Lambda },a,\delta )\text{-contact process survives}\}`$. Note that $`𝒮`$ is nonempty since $`0𝒮`$. By Theorem 1.2 (d) and (f), $`𝒮=\{\delta 0:r(\mathrm{\Lambda },a,\delta )>0\}`$. By Theorem 1.2 (b), the function $`\delta r(\mathrm{\Lambda },a,\delta )`$ is continuous, hence $`𝒮`$ is an open subset of $`[0,\mathrm{})`$. Hence, by monotonicity, $`𝒮=[0,\delta _\mathrm{c})`$, where $`\delta _\mathrm{c}:=sup\{\delta 0:`$ the $`(\mathrm{\Lambda },a,\delta )`$-contact process survives$`\}`$ satisfies $`\delta _\mathrm{c}>0`$. |
warning/0507/hep-th0507076.html | ar5iv | text | # BFT embedding of the Green-Schwarz superstring and the pure spinor formalism
## 1 Introduction
The covariant quantization of GS superstring is an important open problem for string theory and for the general theory of constrained systems. The problem is as old as the initial proposal of the classical action by Green and Schwarz , and many different ways to tackle it with a wide spectrum of techniques was worked along the years. Among them are covariant conversion procedures (with an infinite number of auxiliary fields), BRST program with infinitely reducible constraints , the use of light cone coordinates and the question of conformal invariance , and recently the pure spinor formalism , gauging cosets , and BRST extensions attempting to lift the bosonic spinor constraint of the pure spinor formalism by introducing more ghost variables but a finite number of them .
The problem with the quantization of the GS superstring lies in the fact that we do not know a procedure to separate the first and second class fermionic constraints in a manifestly Lorentz invariant way. The first class sector of these fermionic constraints is responsible for the $`\kappa `$ symmetry of the superstring while the second class constraints appear as in any other fermionic system because the Lagrangian is linear in the time derivatives of the fermionic variables. So we face the problem of trying to quantize the system without splitting the constraints or split them in a non-covariant way and use a somewhat complicated Dirac bracket to perform the quantization. The recent proposed pure spinor formalism is an important step in the direction of a covariant quantization program for the GS superstring. This formulation is covariant but the price to pay is a radical departure from the BRST standard techniques, introducing a constraint in the bosonic ghost sector of the theory, know as the pure spinor constraint. The idea could be related to the use of all the constraints (not only the first class constraints) to construct the BRST operator. The condition $`Q^2=0`$ implies then a constraint in the ghost sector. Nevertheless, this quantization program has been proved to be very useful in many calculations that do not imply the explicit solution of the constrained ghost relation. The other standard formalism to describe the supersymmetric string, the RNS model, has the supersymmetry realized in the worldsheet so the space time supersymmetry is not manifest. As a consequence of this fact the spectrum of the RNS theory is not Lorentz covariant. We need to impose a projection of states in its spectrum, the so called GSO projection to recover a Lorentz invariant spectrum. Without a manifest space time supersymmetric action is very difficult -if not impossible- to describe the superstring in Ramond-Ramond backgrounds.
In what follows we will present a solution to the long standing problem of the implementation of the conversion program of second class constraints into first class constraints by adding to the GS action new fermionic variables with a standard symplectic structure. Our result is completely equivalent to the recent model proposed by Berkovits and Marchioro (BM) to relate the GS superstring with the pure spinor formalism . From our point of view, the key observation given by these authors is that in the extended phase space the Lorentz generators close in the corresponding Lorentz algebra up to and exact BRST term. Moreover, and in spite of the non Lorentz covariant approach, the quantization program can be implemented because the anomaly that comes from the nonlinear products of the new added variables can be canceled using standard BRST techniques. A previous attempt to relate the pure spinor formalism with the GS superstring using standard BRST approach is where the authors were able to modify the GS original action with a non-local term, that when added to the GS action, render it BRST invariant.
In our approach we do not need to fix the gauge (the semi-light cone gauge) nor to add variables other than the ones required by the conversion procedure. Our work is in close relation with a recent proposed action given in where the authors start from a modified GS action doubling the fermion sector and introducing an interaction between the fermionic sectors by hand. These authors claim that theirs modified GS action is equivalent to the BM action by fixing the gauge and performing a Darboux transformation that changes the remaining Dirac bracket to a standard symplectic structure. We will see that this model can be explained and simplified using the conversion approach presented here.
In section 2 we will review the basic ideas of the BFT program. Section 3 will be used to expose our main results and the GS gauged action. In section 4 we present some comments on the relation between our results and the Aisaka-Kazama (AK) action , and in section 5 the conclusion and some notes about possible future work.
## 2 Birds eye to BFT procedure
The aim of the BFT method is to develop a systematic approach to the conversion of a general set of constraints into an algebra of first class constraints by adding to the original phase space an appropriate number of new variables with its own symplectic structure. The method provides us with a procedure for the conversion of second class constraints into first class ones in this extended phase space and a procedure to modify any observable, including any previous first class constraint in such a way that we can construct an effective first class gauge algebra and an effective action consistent with the whole conversion procedure. It is based on homological perturbation theory and is, in this respect, very similar to the iterative construction of the BRST charge given a gauge algebra.
Suppose that we have a set of constraints where some constraints are second class $`\chi _\alpha `$, $`\{\chi _\alpha ,\chi _\beta \}=C_{\alpha \beta }`$ and some are first class $`\varphi _m`$ in a phase space defined by the coordinates $`z^i`$ with the standard symplectic form $`\sigma ^{ij}`$. Now we will add $`\xi _\alpha `$ new variables with symplectic structure $`\omega ^{\alpha \beta }`$ to the original phase space. The idea is to construct a new set of constraints $`\stackrel{~}{\chi }_\alpha `$ satisfying the algebra
$$\{\stackrel{~}{\chi }_\alpha ,\stackrel{~}{\chi }_\beta \}=0.$$
(1)
To solve for $`\stackrel{~}{\chi }_\alpha `$ we propose a solution in power series of the new variables
$$\stackrel{~}{\chi }(z,\xi )=\underset{n}{}X_\alpha ^{(n)},$$
(2)
where the $`n=0`$ term coincides with the original constraint $`\chi _\alpha `$ and $`X_\alpha ^{(n)}`$ is a term proportional to $`\xi ^n`$ in the series expansion. The solution, up to a canonical transformation, in the extended phase space is
$$X_\alpha ^{(0)}=\chi _\alpha ,X_\alpha ^{(1)}=X_{\alpha \gamma }\xi ^\gamma ,X_{\alpha \gamma }\omega ^{\gamma \delta }X_{\beta \delta }=C_{\alpha \beta }(z),$$
(3)
and for the next terms $`n2`$ in the power series
$$X_\alpha ^{(n+1)}=\frac{1}{n+2}\xi ^\beta \omega _{\beta \gamma }X^{\gamma \rho }X_{\rho \alpha }^{(n)},$$
(4)
where
$$X_{\alpha \beta }^{(1)}=\{\chi _{[\alpha },X_{\beta ]}^{(1)}\},X_{\alpha \beta }^{(n)}=\underset{m=0}{\overset{n}{}}\{X_\alpha ^{(nm)},X_\beta ^{(m)}\}+\underset{m=0}{\overset{n2}{}}\{X_\alpha ^{(nm)},X_\beta ^{(m+2)}\}_\xi ,$$
(5)
and the first bracket is evaluated using only the original phase space variables and the second using only the new variables.
The same idea works also to extend any function $`f(z)`$ of the original variables $`z`$ to a new function $`\stackrel{~}{f}(z,\xi )`$ as a solution in power series of the new variables $`\xi `$. This series must satisfy the condition
$$\{\stackrel{~}{\chi }_\alpha ,\stackrel{~}{f}\}=0,\stackrel{~}{f}=\underset{n}{}F^{(n)},$$
(6)
where $`F^{(0)}=\stackrel{~}{f}(z,0)=f(z)`$ and $`F^{(n)}`$ is the term proportional to $`\xi ^n`$. The solution is
$$F^{(n+1)}=\frac{1}{n+1}\xi ^\beta \omega _{\beta \gamma }X^{\gamma \rho }F_\rho ^{(n)},$$
(7)
where
$$F_\alpha ^{(0)}=\{\chi _\alpha ,f(z)\},F_\alpha ^{(1)}=\{X_\alpha ^{(1)},f(z)\}+\{\chi _\alpha ,F^{(1)}\}+\{X_\alpha ^{(2)},F^{(1)}\}_\xi ,$$
(8)
and
$$F_\alpha ^{(n)}=\underset{m=0}{\overset{n}{}}\{X_\alpha ^{(nm)},F^{(m)}\}+\underset{m=0}{\overset{n2}{}}\{X_\alpha ^{(nm)},F^{(m+2)}\}_\xi +\{X_\alpha ^{(n+1)},F^{(1)}\}_\xi .$$
(9)
In particular, we can extend the original first class constraints $`\varphi _m`$ to a new set of constraints $`\stackrel{~}{\varphi }_m`$ in such a way that all the new constrains close in a new gauge algebra in the extended phase space. In what follows we will need only terms up to second order in the new variables.
An interesting corollary of the conversion approach is that the original Dirac bracket can be recovered using
$$\{\stackrel{~}{A},\stackrel{~}{B}\}|_{\xi =0}=\{A,B\}_D,$$
(10)
for any two functions of the original phase space $`A,B`$ that were extended to $`\stackrel{~}{A},\stackrel{~}{B}`$ as can be easily checked.
## 3 BFT embedding of the GS superstring
We start from the GS action that we write in the form
$$S=\frac{1}{2}d^2\zeta \left[\sqrt{g}g^{ij}\mathrm{\Pi }_i^\mu \mathrm{\Pi }_{\mu j}+2\epsilon ^{ij}\mathrm{\Pi }_i^\mu (W_{j\mu }^1W_{j\mu }^2)2\epsilon ^{ij}W_i^{1\mu }W_{j\mu }^2\right],$$
(11)
where
$$W_i^{A\mu }=i\theta ^A\gamma ^\mu _i\theta ^A,\mathrm{\Pi }_i^\mu =_ix^\mu \underset{A}{}W_i^{A\mu }.$$
(12)
As usual the bosonic constraints can be obtained by setting to zero the energy-momentum tensor. Taking the conformal gauge, the first order Lagrangian associated is
$$=\dot{x}^\mu p_\mu +\dot{\theta }_\alpha ^Ap_\alpha ^AH_c\lambda _\alpha ^Ad_\alpha ^A,$$
(13)
where
$$p_\mu =\mathrm{\Pi }_{0\mu }(W_{1\mu }^1W_{1\mu }^2),$$
(14)
and
$$d_\alpha ^1=p_\alpha ^1i(\theta ^1\gamma ^\mu )_\alpha (p_\mu x_\mu ^{}+W_{1\mu }^1),$$
(15)
$$d_\alpha ^2=p_\alpha ^2i(\theta ^2\gamma ^\mu )_\alpha (p_\mu +x_\mu ^{}W_{1\mu }^2),$$
(16)
are the fermionic constraints. The canonical Hamiltonian is
$$H_c=\frac{1}{2}\left[\left(p_\mu +W_{1\mu }^1W_{1\mu }^2\right)^2+\left(x_\mu ^{}\underset{A}{}W_{1\mu }^A\right)^2\right]=\frac{1}{2}\left(\mathrm{\Pi }_0^2+\mathrm{\Pi }_1^2\right).$$
(17)
The bosonic constraints
$$=\frac{1}{2}\left(\mathrm{\Pi }_0^2+\mathrm{\Pi }_1^2\right),_1=\mathrm{\Pi }_0^\mu \mathrm{\Pi }_1^\mu ,$$
(18)
can be written in the form
$$\widehat{T}=+_1=\frac{1}{2}\widehat{\mathrm{\Pi }}^2,T=_1=\frac{1}{2}\mathrm{\Pi }^2,$$
(19)
where
$$\mathrm{\Pi }_\mu =p_\mu x_\mu ^{}+2W_{1\mu }^1,\widehat{\mathrm{\Pi }}_\mu =p_\mu +x_\mu ^{}2W_{1\mu }^2.$$
(20)
The algebra of constraints naturally splits into two sectors
$$\{d_{1\alpha },d_{1\beta }\}=2i\gamma _{\alpha \beta }^\mu \mathrm{\Pi }_\mu ,\{d_{2\alpha },d_{2\beta }\}=2i\gamma _{\alpha \beta }^\mu \widehat{\mathrm{\Pi }}_\mu ,$$
(21)
$$\{d_{1\alpha },d_{2\beta }\}=0,\{d_{1\alpha },\widehat{\mathrm{\Pi }}_\mu \}=0,\{d_{2\alpha },\mathrm{\Pi }_\mu \}=0,\{\mathrm{\Pi }_\mu ,\widehat{\mathrm{\Pi }}_\nu \}=0.$$
(22)
To separate the constraints into first and second class we will write them in the light cone coordinates and divide the $`\alpha ,\beta `$ spinor indices using the spinorial representation of the little group SO(8). The result is that the constraints $`d_a^A=0`$, ($`\mathrm{\Pi }^+0`$) are second class while the rest of the constraints $`d_{\dot{a}}^A=0,T=0,\widehat{T}=0`$ are first class. This fact allow us to count the number of degrees of freedom for the superstring given us the correct result, as expected. In what follows we will not need the details of this first class algebra as our aim is to construct a new effective gauge algebra. To that end we extend the original phase space $`x^\mu ,p_\mu ,\theta _\alpha ^A,p_\alpha ^A`$ by adding the fermionic variables $`S_a`$ with the symplectic structure<sup>3</sup><sup>3</sup>3In what follows we will work on the sector with index 1 and for simplicity we will remove this index from our equations. The other sector can be worked in the same way..
$$\{S_a,S_b\}=i\delta _{ab}.$$
(23)
Searching a solution for the condition (1)
$$\{\stackrel{~}{d}_a,\stackrel{~}{d}_a\}=0,$$
(24)
in power series of $`S`$ give us a very simple result. The solution is linear in $`S`$ and yields
$$\stackrel{~}{d}_a=d_a+i\sqrt{2\mathrm{\Pi }^+}S_a.$$
(25)
The next step consist in the deformation of the other first class constraints $`d_{\dot{a}}=0`$, $`T=0`$ to be consistent with the new $`\stackrel{~}{d}_a`$ constraints. Consider the case of $`d_{\dot{a}}`$. We need to find a solution to the condition
$$\{\stackrel{~}{d}_a,\stackrel{~}{d}_{\dot{a}}\}=0.$$
(26)
The solution has the general form
$$\stackrel{~}{d}_{\dot{a}}=d_{\dot{a}}+A_{\dot{a}b}S_b+B_{\dot{a}[bc]}S_bS_c,$$
(27)
where
$$A_{\dot{a}b}=\frac{2i\gamma _{\dot{a}a}^i\mathrm{\Pi }^i}{\sqrt{2\mathrm{\Pi }^+}},B_{\dot{a}[ac]}=\frac{2\gamma _{\dot{b}[a}^i\gamma _{c]\dot{a}}^i\theta _{\dot{b}}^{}}{\mathrm{\Pi }^+}.$$
(28)
The extended constraint is
$$\stackrel{~}{d}_{\dot{a}}=d_{\dot{a}}+\frac{2i\mathrm{\Pi }^i}{\sqrt{2\mathrm{\Pi }^+}}(\gamma ^iS)_{\dot{a}}+\frac{2(\theta ^{}\gamma ^iS)(\gamma ^iS)_{\dot{a}}}{\mathrm{\Pi }^+},$$
(29)
keeping in mind that the last term has to be antisymmetrized in undoted spinorial indices.
Now to find the extended constraint associated with $`T`$ it is easy to proceed first to the extension of $`\mathrm{\Pi }_\mu `$ defined in eq (20). To do that we need to solve the condition (1) for $`\stackrel{~}{\mathrm{\Pi }}_\mu `$, i.e.,
$$\{\stackrel{~}{d}_a,\stackrel{~}{\mathrm{\Pi }}_\mu \}=0.$$
(30)
A simple check shows that the series in powers of the new variables $`S`$ for $`\stackrel{~}{\mathrm{\Pi }}_\mu `$ stops up to second order terms. The solution is
$$\stackrel{~}{\mathrm{\Pi }}^\mu =\mathrm{\Pi }^\mu +4i\frac{(\theta ^{}\gamma ^\mu S)}{\sqrt{2\mathrm{\Pi }^+}}+i\frac{S\gamma ^\mu S}{\mathrm{\Pi }^+}.$$
(31)
Using this solution we will define the new first class constraint $`\stackrel{~}{T}`$ in such a way that it will close in a Lie algebra with the rest of the new constraints
$$\stackrel{~}{T}=\frac{\stackrel{~}{\mathrm{\Pi }}^2}{4\mathrm{\Pi }^+}=\frac{\mathrm{\Pi }^{}}{4}+\frac{\mathrm{\Pi }^i\mathrm{\Pi }^i}{4\mathrm{\Pi }^+}+2i\frac{\theta _a^{}S_a}{\sqrt{2\mathrm{\Pi }^+}}+i\frac{S_aS_a^{}}{2\mathrm{\Pi }^+}+4i\frac{\mathrm{\Pi }^i(\theta ^{}\gamma ^iS)}{(2\mathrm{\Pi }^+)^{3/2}}2\frac{(\theta ^{}\gamma S)^2}{(\mathrm{\Pi }^+)^2}.$$
(32)
The new effective gauge algebra in the extended space is now
$$\{\stackrel{~}{T},\stackrel{~}{T}\}=0,\{\stackrel{~}{d}_a,\stackrel{~}{d}_b\}=0,\{\stackrel{~}{d}_{\dot{a}},\stackrel{~}{d}_a\}=0,$$
(33)
$$\{\stackrel{~}{d}_a,\stackrel{~}{T}\}=0,\{\stackrel{~}{d}_{\dot{a}},\stackrel{~}{d}_{\dot{b}}\}=8i\stackrel{~}{T}\delta _{\dot{a}\dot{b}}.$$
(34)
The gauged GS first order action is
$$\stackrel{~}{S}=\frac{1}{2}d^2\zeta \left(\dot{x}^\mu p_\mu +\dot{\theta }_\alpha ^Ap_\alpha ^A+\frac{i}{2}\dot{S}_a^AS_a^A\lambda \stackrel{~}{T}\widehat{\lambda }\widehat{\stackrel{~}{T}}\lambda _\alpha ^A\stackrel{~}{d}_\alpha ^A\right),$$
(35)
where we have included the two sectors. Its gauge symmetries are the worldsheet diffeomorphisms that are generated by $`\stackrel{~}{T}`$ and $`\widehat{\stackrel{~}{T}}`$ and a new fermionic gauge symmetry that is generated by $`\stackrel{~}{d}_\alpha ^A`$. Of course the theory is not manifestly Lorentz covariant but the Lorentz invariance is guaranteed up to a BRST trivial transformation .
We have now 17 first class constraints by sector and as expected the model is equivalent to the original GS action and by construction has the same number of degrees of freedom. Two comments are in order. The first is that this embedding of the GS superstring is equivalent to the classical BM action. Its quantization can be performed along the same lines as the quantization of the BM model. Subtitles related to ordering ambiguities must be taken into account for a consistent quantization of this action. Secondly, as the BM model can be related to the pure spinor formalism via similarity transformations between the associated BRST charges, this model can also be related to the pure spinor formalism using the same sequence of similarity transformations between its associated BRST charges. The advantage of our perspective is that we have developed a gauge model in a completely systematic way starting from the plain GS superstring and consequently we have more control over any change in the embedding procedure that can be of help to relate GS and pure spinor formalisms in a more direct way. This procedure can also be of some help to better understand many aspects of pure spinor formalism like its geometrical interpretation, the path integral measure, or the underlying action.
From the other hand our new constraints (25,29,32) are the same as the ones obtained in . It is quite surprising for us that the constraints are exactly the same. The two procedures are very different. In the number of fermions is doubled and an interaction between them was introduced by hand. After fixing the semi-lightcone gauge and making a complicated Darboux transformation simplifying the Dirac bracket, the results of coincides with the BFT embedding presented here. We will try to explain this relation in the next section by extracting more information about how the BFT embedding works.
## 4 Relation with the AK model
That the embedding approach to GS action has something to do with the interacting action proposed in is at first sight very surprising. Here we will try to elaborate on this relation using a slightly modified approach to the conversion procedure. The arguments presented in this section does not apply to the case of a general constrained system but are valid in some special type of systems like the one considered here.
Lets start by noticing that another way to apply the BFT embedding is to seek for new extended coordinates $`\stackrel{~}{x}^\mu ,\stackrel{~}{p}_\mu ,\stackrel{~}{\theta }_\alpha ^A,\stackrel{~}{p}_\alpha ^A`$ such that they satisfy the conditions
$$\{\stackrel{~}{d}_a^A,\stackrel{~}{z}\}=0,$$
(36)
where $`\stackrel{~}{z}(x^\mu ,p_\mu ,\theta _\alpha ,p_\alpha ,S)`$ is any of the phase space extended new coordinates or momenta. If we can solve these conditions then we can use the solutions to extend any observable of the original phase space to the new extended phase space. The procedure is as follows: suppose that we have a function in the original phase space $`A(z)`$. First we write it as $`A(\stackrel{~}{z})`$, then substitute $`\stackrel{~}{z}`$ by the solution to (36) and find $`\stackrel{~}{A}(z,S)`$. For the GS superstring the solution to the conditions (36) are<sup>4</sup><sup>4</sup>4This solutions are for the sector 1. The solution for the other sector has the same form.
$$\stackrel{~}{x}^\mu =x^\mu i\frac{(\theta \gamma ^\mu S)}{\sqrt{2\mathrm{\Pi }^+}},\stackrel{~}{\theta }_a=\theta _a\frac{S_a}{\sqrt{2\mathrm{\Pi }^+}},\stackrel{~}{\theta }_{\dot{a}}=\theta _{\dot{a}},$$
(37)
for configuration space variables. For the bosonic momenta $`p_\mu `$ and the fermionic momenta $`p_\alpha `$ the solutions are
$$\stackrel{~}{p}^\mu =p^\mu +i\left(\frac{\theta \gamma ^\mu S}{\sqrt{2\mathrm{\Pi }^+}}\right)^{},$$
(38)
$$\stackrel{~}{p}_\alpha =p_\alpha i\frac{(\gamma ^\mu S)_\alpha }{\sqrt{2\mathrm{\Pi }^+}}(\mathrm{\Pi }_\mu W_{1\mu }+P_\mu )+i(\gamma ^\mu \theta )_\alpha P_\mu +C_\alpha ,$$
(39)
where
$$P_\mu =2i\frac{\theta ^{}\gamma _\mu S}{\sqrt{2\mathrm{\Pi }^+}}+i\frac{S\gamma _\mu S^{}}{2\mathrm{\Pi }^+}+i\left(\frac{\theta \gamma _\mu S}{\sqrt{2\mathrm{\Pi }^+}}\right)^{},$$
(40)
and
$$C_a=i\sqrt{2\mathrm{\Pi }^+}S_a,C_{\dot{a}}=\frac{2i\mathrm{\Pi }^i}{\sqrt{2\mathrm{\Pi }^+}}(\gamma ^iS)_{\dot{a}}+\frac{2(\theta ^{}\gamma ^iS)(\gamma ^iS)_{\dot{a}}}{\mathrm{\Pi }^+},$$
(41)
where the last term must be antisymmetrized with repect to the spinorial indices without dots. The new momenta $`\stackrel{~}{p}_\mu `$, $`\stackrel{~}{p}_\alpha `$ in (38, 39) are very similar to the Darboux transformation proposed in to simplify the Dirac bracket but there are differences that we will explain below. What is perhaps more interesting is that the transformations (37) in configuration space can be used to obtain, from the GS action (11), the interacting AK action. Indeed, redefine $`S_a/\sqrt{2\mathrm{\Pi }^+}`$ as $`\xi _a`$ and use (37) to get
$$S=\frac{1}{2}d^2\zeta \left[\sqrt{g}g^{ij}\mathrm{\Pi }_i^\mu \mathrm{\Pi }_{\mu j}+2\epsilon ^{ij}\mathrm{\Pi }_i^\mu (W_{j\mu }^1W_{j\mu }^2)2\epsilon ^{ij}W_i^{1\mu }W_{j\mu }^2\right],$$
(42)
where
$$W_i^{A\mu }=i\mathrm{\Theta }^A\gamma ^\mu _i\mathrm{\Theta }^A,\mathrm{\Pi }_i^\mu =_ix^\mu \underset{A}{}W_i^{A\mu }i\underset{A}{}_i(\theta ^A\gamma ^\mu \xi ^A),$$
(43)
with $`\mathrm{\Theta }^A=\theta ^A\xi ^A`$ as in <sup>5</sup><sup>5</sup>5We denote by $`\xi `$ the variable that are denoted as $`\stackrel{~}{\theta }`$ in to avoid some possible confusion with our tilde variables.. Notice that the effect of the substitution of the new configuration variables (37) in terms of the old ones in the original GS action (11) produces a deformation of the symplectic structure and a deformation of the original bosonic constraints. The first order action has now the form
$$S=\frac{1}{2}d^2\zeta \left(\dot{x}^\mu p_\mu +\dot{\theta }_\alpha ^A\mathrm{\Delta }_\alpha ^A+\dot{\xi }_a^A\mathrm{\Xi }_a^A\lambda \tau \widehat{\lambda }\widehat{\tau }\right),$$
(44)
where space-time momenta is
$$p_\mu =\mathrm{\Pi }_{0\mu }(W_{1\mu }^1W_{1\mu }^2),$$
(45)
and the functions in the kinetic term are
$$\mathrm{\Delta }_\alpha ^1=i(\gamma ^\mu \xi ^1)_\alpha p_\mu +i(p^\mu \mathrm{\Pi }_1^\mu W_1^{2\mu })(\mathrm{\Theta }^1\gamma _\mu )_\alpha ,$$
(46)
$$\mathrm{\Delta }_\alpha ^2=i(\gamma ^\mu \xi ^2)_\alpha p_\mu +i(p^\mu +\mathrm{\Pi }_1^\mu +W_1^{1\mu })(\mathrm{\Theta }^2\gamma _\mu )_\alpha ,$$
(47)
and
$$\mathrm{\Xi }_a^1=i(\gamma ^\mu \theta ^1)_ap_\mu +i(p^\mu \mathrm{\Pi }_1^\mu W_1^{2\mu })(\mathrm{\Theta }^1\gamma _\mu )_a,$$
(48)
$$\mathrm{\Xi }_a^2=i(\gamma ^\mu \theta ^2)_ap_\mu +i(p^\mu +\mathrm{\Pi }_1^\mu +W_1^{1\mu })(\mathrm{\Theta }^2\gamma _\mu )_a.$$
(49)
$`\tau ,\widehat{\tau }`$ are the deformed bosonic constraints after the substitution of (37) in the original bosonic constraints $`T`$ and $`\widehat{T}`$. Defining, as usual, the fermionic momenta as the coefficient that multiply $`\dot{\theta }`$ in (42), we find the fermionic constraints
$$D_\alpha ^A=p_\alpha ^A\mathrm{\Delta }_\alpha ^A,$$
(50)
that correspond to the constraints $`d_\alpha ^A`$ of the original GS action. Integration by parts in the kinetic term produces the action
$$S=\frac{1}{2}d^2\zeta \left(\dot{x}^\mu (p_\mu P_\mu )+\dot{\theta }_\alpha ^A(p_\alpha ^AP_\alpha ^A)+i\mathrm{\Pi }^+\dot{\xi }_a^A\xi _a^A\lambda \tau \widehat{\lambda }\widehat{\tau }\lambda _\alpha ^AD_\alpha ^A\right),$$
(51)
where
$$P_\mu =i(\theta \gamma ^\mu \xi )_1^{}+i(\theta \gamma ^\mu \xi )_2^{},$$
(52)
and $`()_1`$ denotes variables of the sector 1 and $`,()_2`$ for sector 2. The fermionic momenta redefinition is
$$P_a=ix_{}^{}{}_{}{}^{+}\xi _a+R_a,P_{\dot{a}}=(\gamma ^i\xi )_{\dot{a}}\left(ix_i^{}+(\theta \gamma _i\xi )_2\right)+R_{\dot{a}},$$
(53)
where
$$R_\alpha =(\theta \gamma ^\mu )_\alpha \left(2(\xi \gamma _\mu \theta ^{})+\xi \gamma _\mu \xi ^{}+(\theta \gamma _\mu \xi )^{}\right)+(\xi \gamma ^\mu )_\alpha \left(3(\theta \gamma _\mu \theta ^{})2(\theta \gamma _\mu \xi ^{})\right).$$
(54)
The redefinition of the space time momenta are the same as the redefinition that we have used in the BFT embedding but the redefinitions associated with the fermionic momenta are slightly different. The reason is that the fermionic constraints $`D_\alpha ^A`$ are not the same as the original fermionic constraints $`d_\alpha `$ after the substitution of the new coordinates (37). Now it is easy to check that the field redefinitions
$$p_\mu p_\mu P_\mu ,p_\alpha ^Ap_\alpha ^AP_\alpha ,\xi _a=S_a/\sqrt{2\mathrm{\Pi }^+},$$
(55)
produces the first order action (35) obtained by the BFT embedding. The efficient way, just presented to analyze the constraint structure of the action (42) is inspired in the Faddeev-Jackiw method that is equivalent to the Dirac method for a wide class of constrained systems . So we have obtained the interacting model from the BFT embedding approach using the configuration variables $`\stackrel{~}{x}^\mu ,\stackrel{~}{\theta }_\alpha ^A`$ obtained by solving the embedding condition (36). Notice that this procedure does not work for a general constrained system. The fact that the GS action is linear in the velocities of the fermionic variables and that the solution to (36) depends only on the configuration variables, after the redefinition that relates $`S`$ with $`\xi `$ are crucial ingredients in the construction of a Lagrangian action compatible with the dynamics of the original Hamiltonian action.
## 5 Conclusions
We have worked out the BFT embedding of the GS formalism using light-cone variables. We have shown that the correction to the second class constraints and the embedding of the first class algebra requires only terms up to second order in the new fermionic variables and its derivatives. The embedding can be performed in a completely systematic way and we do not have to fix the gauge at any stage of our embedding procedure not to add ad-hoc variables other than the ones needed by the BFT embedding approach. The gauged GS action that results from our analysis is equivalent to the BM model proposed recently in . It also explain some aspects of the rationale under the construction of a related model worked in .
The systematics behind our approach can be used to study the relation of the BRST charges and associated cohomologies between this model and the pure spinor formalism. We also have more control on the gauge fixing procedure. It is also a better point of departure to study the supermembrane and other topics not yet well understood in the pure spinor formalism, like the associated action and the measure of the path integral. It could be also of interest to explore the idea of non-Abelian conversion to simplify the relation between the Berkovits pure spinor formalism and the GS embedding. We will return to these aspects elsewhere.
## Acknowledgements
This work was partially supported by Mexico’s National Council of Science and Technology (CONACyT), under grant CONACyT-40745-F, and by DGAPA-UNAM, under grant IN104503-3.
## Appendix
The notation used in the paper is the following: for the GS action (11), $`\epsilon ^{01}=1`$, $`\theta _\alpha ^A`$, is an SO(9,1) spinor with $`A=1,2`$ supersymmetries and $`\alpha =1,2,\mathrm{}16`$ components. They are real Mayorana-Weyl spinors of the same chirality. $`x^\mu `$ $`\mu =0,1,2\mathrm{}9`$ are the space-time configuration variables. The worldsheet coordinates are $`\zeta =(t,\sigma )`$ and the derivatives with respect to time will be denoted by a dot and with respect to sigma by a prime. The $`\gamma ^\mu `$ matrices are $`16\times 16`$ Dirac matrices, real and symmetric. The convention for derivatives are left derivatives and then
$$\{\theta _\alpha ^A,p_\beta ^B\}=\delta _{\alpha \beta }\delta ^{AB},$$
(56)
This convention fixes the order of $`\dot{\theta }`$ and $`p`$ in the kinetic term of the first order action.
Light-cone coordinates: We split the spinorial index $`\alpha `$ according to the the SO(8) chiral and anti-chiral components $`a`$ and $`\dot{a}`$. Space-time indices decompose according to $`\gamma ^\pm =\gamma ^0\pm \gamma ^9`$, $`x^\pm =x^0\pm x^9`$… and $`\gamma ^i`$, $`x^i`$, i=1,2,…8 for the other set of vector components. The Dirac algebra decomposes according to
$$\gamma _{\dot{a}\dot{b}}^+=2\delta _{\dot{a}\dot{b}},\gamma _{ab}^{}=2\delta _{ab},\gamma _{\dot{a}a}^i\gamma _{\dot{a}b}^j+\gamma _{a\dot{a}}^j\gamma _{\dot{a}b}^i=2\delta ^{ij}\delta _{ab},$$
(57)
and
$$\gamma _{a\dot{b}}^i\gamma _{c\dot{d}}^i+\gamma _{a\dot{d}}^i\gamma _{c\dot{b}}^i=2\delta ^{ac}\delta _{\dot{b}\dot{d}},$$
(58)
with $`\gamma _{\dot{a}a}`$ symmetric. We use repeatedly through the text the Fierz identity
$$(\gamma _{(\alpha \beta })^\mu (\gamma _{\gamma )\delta })_\mu =0.$$
(59) |
warning/0507/astro-ph0507273.html | ar5iv | text | # The Large Scale Structure of the Galactic Magnetic Field and High Energy Cosmic Ray Anisotropy
## 1 Introduction
The Galactic Magnetic Field (GMF) is not well known and is hard to study (Han (2002); Vallée (2004); Wielebinski (2005)). The position of the Solar system makes it difficult to measure its global structure and to distinguish local small-scale features from large-scale ones. Faraday rotation measures (RM) of pulsars in our Galaxy and of polarized extragalactic radio sources are one of the best probes of the large scale structure of the GMF in the Galactic disk and the halo. From this measurements it is derived that the GMF has two components: a regular component with strength $``$ few $`\mu `$G, and a turbulent or random component of the same or perhaps even larger strength (Beck (2001)). There seems to be agreement on the spiral structure of the regular field in the Galactic plane between the Galactic arms, although not on the exact shape of the spiral field, axi-symmetric (ASS) or bi-symmetric (BSS) (Han (2003); Vallée (2004)). There is some controversy on the number of field reversals from arm to arm. The controversy extends to the parity (even or odd) of the GMF across the Galactic plane. There is also disagreement on the existence of a possible halo field. An A0 dipole field directed towards the North Galactic Pole (NGP) was suggested as a halo field (Han et al. (1997)).
Cosmic ray propagation in the Galaxy is strongly affected by the GMF. The gyroradius of a proton of energy $`E=10^{18}`$ eV in a $`3\mu `$G field is of the order of 300 pc, the typical thickness of the Galactic disk. For energies $`E<10^{18}`$ eV cosmic rays diffuse in the GMF, they get isotropized, and hence do not reveal the sources where they were accelerated. Such cosmic rays are also fairly insensitive to the large features of the poorly-known GMF. At energies above $`10^{19}`$ eV cosmic rays have long been thought to be of extragalactic origin because there are no astrophysical objects with high magnetic fields on the large scale needed for their acceleration (Cocconi (1956)). Even if their sources were inside the Galaxy, there would exist a clear anisotropy in the arrival direction of cosmic rays in the case of protons that is not supported by data. Heavier nuclei such as iron would not be ruled out as having a Galactic origin, because with their lower rigidity they would become isotropized in the GMF even in this energy range. However, their presence in the cosmic ray spectrum above $`10^{19}`$ eV is less favored by composition measurements (Abbasi (2005)). An extremely interesting energy range is that from $`10^{18}`$ to $`10^{19}`$ eV. In this energy bin cosmic ray propagation through the GMF is thought to change from diffusive to ballistic, cosmic ray composition is thought to change from heavy to light, and cosmic ray origin is thought to change from Galactic to extragalactic (Nagano & Watson (2000)).
The center of our Galaxy, where there is some evidence for the existence of a very massive black hole, provides a natural candidate for acceleration of cosmic rays to very high energies. The high energy astrophysical activity at the Galactic center is supported by the recent observation by the HESS telescope of a TeV gamma-ray source near the location of Saggitarius $`A^{}`$ (Aharonian (2004)).
In this work we demonstrate that protons in the energy range from $`10^{18}`$ to $`10^{19}`$ eV, if accelerated at galactocentric distances typically smaller than the radius of the Solar system orbit around the Galactic center, are sensitive to the large scale structure of the GMF. In particular, if the GMF is of even parity, i.e. does not change sign across the Galactic plane, their distribution in arrival direction reveals the axi-symmetric (ASS) or bi-symmetric (BSS) configuration of the spiral field. In the first case we will show that protons are observed to arrive predominantly from the Northern Galactic hemisphere, while in the second they are seen to come predominantly from the Southern Galactic hemisphere. However, as we demonstrate below, there is no sensitivity to the GMF configuration if the field changes sign across the Galactic plane, i.e. if it is of odd parity.
This letter is structured as follows: In Section 2 we briefly review the current knowledge on the GMF, and describe the GMF configurations used in this work. In Section 3 we briefly describe how we performed our calculations. In section 4 we demonstrate the sensitivity of cosmic ray propagation to the large scale features of the GMF.
## 2 The Galactic Magnetic Field
We briefly summarize here the current knowledge on the GMF. More details can be found in (Han (2003); Vallée (2004); Wielebinski (2005)). At low Galactic latitudes Faraday RMs of pulsars inside the Galaxy reveal that the disk field direction is coherent over a linear scale of at least a few kpc between the Galactic arms. It is not yet clear from an experimental point of view if the field reverses direction from arm to arm (Han (2003); Vallée (2004)). In particular, moving towards the Galactic Center (GC) the field direction – clockwise-going or anticlockwise-going as seen from the North Galactic Pole – between the Perseus and Perseus + I arm at galactocentric distance $`r_{||}12`$ kpc is still controversial. There seems to be agreement on the existence of a clockwise-going field between the Perseus and Carina-Sagittarius arm at $`r_{||}8`$ kpc, close to the position of the Solar System, and with a strength of $`B2\mathrm{}4\mu `$G. The field reverses direction in the next arm towards the Galactic Center (GC) between the Carina-Sagittarius and the Crux-Scutum arms at $`r_{||}6.5`$ kpc. Moving closer to the GC, the field is clockwise-going between the Crux-Scutum and Norma arms at $`r_{||}5`$ kpc. Finally near the Norma arm at $`r_{||}4`$ kpc measurements are again contradictory.
All this evidence gives support for a two-arm logarithmic spiral regular field in the disk. It is not well established whether the disk field is better described by a bi-symmetric spiral (BSS) configuration which allows for multiple field reversals, or an axi-symmetric spiral (ASS) configuration, with the BSS structure slightly favored. These two field configurations are plotted in Fig.1.
Also it has recently been confirmed that the GC contains a highly regular polar field consistent with the presence of a dipole field in the Galaxy and compatible with being generated by an A0 dynamo (Han (2003)). The presence of the dipole field could explain the vertical component of the GMF of strength $`B_z0.3\pm 0.1\mu `$G observed in the vicinity of the Sun.
At high Galactic latitudes the antisymmetric distribution of Faraday RMs is indicative of the odd parity of the field although an even parity is not excluded. At large distances from the disk $`>1`$ kpc or so, the measurements are again very difficult because of contamination by the GMF in the disk.
Given this information, in this paper we do not make strong assumptions about the large scale structure of the spiral GMF and assume all possible combinations, namely ASS without field reversals or BSS, and even or odd parity. The local regular magnetic field in the vicinity of the Solar System is assumed to be $`1.5\mu \mathrm{G}`$ in the direction $`l=90^\mathrm{o}+p`$ where the pitch angle is $`p=10^\mathrm{o}`$ (Han & Qiao (1994)). Some measurements discuss larger total field strengths of up to 6 $`\mu `$G (Beck (2001)), therefore our assumptions should be considered as rather conservative. The field decreases with Galactocentric distance as $`1/r_{||}`$ and it is zero for $`r_{||}>20`$ kpc. In the region around the Galactic center ($`r_{||}<4`$ kpc) the field is highly uncertain, and we assume it is constant and equal to its value at $`r_{||}=4`$ kpc. Following (Stanev (1996)) the spiral field strengths above and below the Galactic plane are taken to decrease exponentially with two scale heights. In this work we also assume a halo field corresponding to an A0 dipole as suggested by (Han (2002)). The dipole field is toroidal and its strength decreases with Galactocentric distance as $`1/r^3`$. The dipole field is very strong in the central region of the Galaxy, but is only 0.3 $`\mu `$G in the vicinity of the Solar system, directed towards the North Galactic Pole. The equations describing the functional form of the field strength for both the spiral and the dipole fields have been published elsewhere (Stanev (1996); Alvarez-Muñiz et al (2002); Prouza & Šmída (2003)).
We also assume a significant turbulent component of the GMF, $`B_{\mathrm{ran}}`$. Its strength is comparable or possibly larger than the regular field (Beck (2001)). To simulate it we add to the spiral and dipole components a random field with a strength of 50% of the local regular field strength with coherence length of 100 pc. We have also performed several runs with a turbulent field twice the local regular field (i.e. four times the standard random field). The possible geometrical offset of the random and regular fields and the possible time dependence of $`B_{\mathrm{ran}}`$ are neglected.
## 3 Calculation technique
We sample protons with energies greater than $`E=10^{18}`$ eV from a $`dN/dEE^{2.7}`$ energy spectrum. We inject 100 protons per source isotropically from sources distributed homogeneously in a ring of radius $`r_{||}=4`$ kpc around the Galactic center in the Galactic plane. We forward (not backward) propagate them from the sources in different models of the GMF by numerically integrating the equations of motion in a magnetic field. There is no energy loss on propagation. We stop the propagation and sample a new proton energy when the proton trajectory intersects our detector – a 1 kpc radius sphere around the Solar system position – when it reaches Galactocentric distances $`r_{||}>20`$ kpc, or when it travels a total pathlength larger than 4 Mpc. The total length of trajectories reaching Earth is always much smaller than this limit. If a proton hits the detector, we keep the proton arrival direction in Galactic coordinates, as well as the position of the source from which it was injected. Our results can be easily re-scaled in rigidity for heavier nuclei. We will show in the next section that our conclusions do not change qualitatively if we reduce the radius of the detector around the Solar System, although the detection efficiency – the ratio of detected to injected protons – decreases as the area of the spherical detector. For this reason we prefer to use a relatively large detector.
Given the still controversial experimental measurements of the GMF, we study the sensitivity of cosmic ray propagation to an ASS and a BSS spiral field, that can be either of even or odd parity across the Galactic plane. In order to better understand the effect of each of the GMF components – spiral, random and dipole – we simulate the propagation of cosmic rays artificially switching on each component, starting with the spiral field, then adding the random component and finally switching all of them on. A total number of 50,000 protons are collected for each configuration of the GMF we have explored.
Numerical Monte Carlo procedures such as the one described above do not usually represent well the power spectrum of the random magnetic field and thus can not be used to predict the energy dependence of the diffusion coefficient. In this study we are concentrating on the propagation of protons in the transition regime between diffusive and ballistic propagation, where the random field properties are less important.
## 4 Results
Tables 1, 2, and 3 summarize the main results of our work. In all of them we give the fraction of cosmic rays in different energy bins coming from the Southern Galactic hemisphere (SGH), i.e. that arrive at the spherical detector from Galactic latitude $`b<0`$. The numbers in parenthesis are the one sigma Poisson limits on the fraction of events, i.e. the probability that the fraction of events is outside the corresponding interval is $`1/e0.37`$. Typically the uncertainty increases with energy as less events are sampled at high energy due to the steep injection spectrum.
Inspection of the tables leads to several important conclusions:
1. If the GMF is of even parity (Table 1 and top half of Table 2), there is a very strong North-South (NS) anisotropy in the arrival direction of protons. In particular, for energies above $`3\times 10^{18}`$ eV, more than $`90\%`$ of the protons that are detected come from the SGH in the BSS model, and from the Northern Galactic hemisphere (NGH) in the ASS model. This is a clear tendency that does not depend very much on the strength of the random or dipole components of the field as can be seen in Table 3. In this table we give the fraction of events coming from the SGH for the BSS even parity model, but using a random field twice the value of the regular field, i.e. four times the standard random field, and also doubling the strength of the dipole field with respect to its nominal value of $`0.3\mu `$G at the position of the Solar system.
This result is very stable, and mostly dependent on the model of the regular field. The insufficient representation of the turbulence of the random field thus does not affect the conclusions of this study.
2. If the field is of odd parity (bottom half of Table 2) there is no sensitivity to the ASS or BSS character of the spiral field. About half of the protons in all energy bins come from the SGH in both the ASS (bottom half of Table 2) and BSS (values not given) odd parity configurations. The non-observation of a large NS anisotropy in the highest energy bins for protons having Galactic longitude in a fairly wide angular bin around $`l=0^{}`$ could be favoring an odd parity GMF.
3. The NS anisotropy seen in Tables 1, 2 and 3 is significantly reduced at energies below $`3\times 10^{18}`$ eV, especially in the ASS model. Clearly in the lower energy bins, proton propagation is more affected by the random and dipole components of the GMF. The dipole field seems to be more important in reducing the NS anisotropy at low energies, possibly due to its large strength near the Galactic center (Yoshiguchi et al. (2004)).
To understand how this strong anisotropy is actually realized, we show in Fig.2 a sample of detected and non detected proton trajectories at $`10^{19}`$ eV in the BSS and ASS (spiral field only) even and odd parity configurations. What is shown is the projection of the proton trajectories through the GMF onto a plane perpendicular to the Galactic disk containing the Solar system position and the Galactic center. In the BSS even parity configuration, the GMF is directed towards $`l270^{}`$ in the first arm that protons encounter on their paths to the Solar system as can be seen in Fig.1, so that their tracks tend to be concave. If a proton is injected from a source towards north, the GMF bends the trajectory so that it escapes from the Galaxy (dashed lines in the top panel of Fig. 2). If the proton is injected towards south, the GMF will bend its track towards north so that it may hit the Solar system and will appear as coming from the Southern Galactic hemisphere (solid lines in the top panel of Fig. 2). The opposite behavior is true for the ASS even parity configuration, i.e. the tracks tend to be convex (middle panel in Fig.2) due to the GMF pointing towards Galactic longitude $`l90^{}`$ in the first magnetic arm between the sources and the Solar system (see right panel of Fig.1). As a consequence only protons injected towards north and bending back towards south can be detected, and will appear to come from the Northern Galactic hemisphere. If the field changes sign across the Galactic plane (odd parity), both concave and convex trajectories are possible and there is no prefered arrival direction. In fact typical proton trajectories arriving at the detector cross the Galactic plane due to the change of sign of the GMF across the Galactic disk as can be seen in the bottom panel of Fig.2.
The use of a 1 kpc radius spherical detector, being the average distance to the sources in the Galactic ring of the order of $`5`$ kpc, induces an unavoidable smearing in the proton arrival angles of $`10^{}`$. However this does not change our conclusions about the strong NS anisotropy we predict. The reason is that, for instance for the BSS even + $`B_{\mathrm{ran}}`$ \+ dipole configuration, $`96\%`$ of the events coming from south in the energy bin $`\mathrm{log}_{10}(E/\mathrm{eV})(18.3,18.4)`$ arrive from latitudes $`b<10^{}`$. As a result the strong NS anisotropy given in Table1 is not expected to change significantly due to the $`10^{}`$ smearing in arrival angle that may possibly shift their arrival directions so that they would actually appear to come from north.
###### Acknowledgements.
We thank R.A. Vázquez for helpful discussions. This research is supported in part by NASA Grant ATP-0000-0080. J. A-M is supported by the Spanish “Ramón y Cajal” program, and acknowledges the Xunta de Galicia (PGIDIT02 PXIC 20611PN), and the MCYT (FPA 2001-3837, FPA 2002-01161 and FPA 2004-01198). We thank the “Centro de Supercomputación de Galicia” (CESGA) for computer resources. |
warning/0507/hep-ph0507241.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The decay constant of a pseudoscalar meson $`D_q`$ consisting of a heavy $`c`$-quark and a light quark $`q=u,d,s`$ is defined through the matrix element of the corresponding pseudoscalar current as follows:
$$<\mathrm{\Omega }|(m_c+m_q)\overline{q}i\gamma _5c(0)|D_q>=f_{D_q}M_{D_q}^2.$$
Recently, two different experimental groups have extracted the decay constant $`f_{D^+}`$ from a direct measurement of the absolute branching fraction for the Cabbing-suppressed leptonic decay mode $`D^+l^+\upsilon `$ . The results are quite different, although there is some overlap of the large experimental errors. Both results are compatible with the upper limit of $`290`$ MeV, established by the MARK-III collaboration at 90% C.L. . For the $`D_s^+`$ meson, there are several measurements of its decay constant $`f_{D^+}`$published in the last decade , with results in the range of $`194430`$ MeV. On the theoretical side, there is a recent estimate based on Borel transformed sum rules $`f_D=195\pm 20`$ MeV , and a preliminary result $`f_D=225+1113\pm 21`$ MeV obtained in three flavor lattice QCD. There appears to be some room for improvement of the accuracy of these results as well as for the systematic study of their uncertainties. On the other hand, the result $`f_{D_s}/f_D=1.18`$ for the ratio of the decay constants appears to be well established in lattice QCD .
In this letter, we estimate the decay constants $`f_D`$ and $`f_{D_s}`$ of the peudoscalar charmed mesons using an alternative method, based on finite energy sum rules, which compare moments of available experimental data with the corresponding QCD theoretical counterpart. In particular, we take linear combinations of positive moments in such a way that the contribution of the data in the region above the resonances turns out to be practically negligible. On the theoretical side we use a large momentum expansion of massive QCD at three loops. This expansion is known up to the seventh power of $`m_c^2/s`$, where $`m_c`$ is the mass of the charm quark and $`s`$ is the square of the CM energy . The expansion makes sense as long as $`s`$ is far enough above the continuum threshold and above the resonances. On the phenomenological side of the sum rule we consider only the lowest lying pseudoscalar $`D_q`$-meson, once the unknown continuum contribution has been removed by a judicious use of quark-hadron duality in our method .
The plan of this note is the following. In section 2 we briefly review the finite energy sum rule method employed. In section 3 we discuss the theoretical and experimental inputs and we present our estimates for the decay constants. Finally, in section 4 we write up the conclusions.
## 2 The method
The two point function associated with the pseudoscalar current is:
$$\mathrm{\Pi }(s=q^2)=i𝑑xe^{iqx}<\mathrm{\Omega }|T(j_5(x)j_5(0))|\mathrm{\Omega }>,$$
(1)
where $`<\mathrm{\Omega }|`$ is the physical vacuum and the current $`j_5(x)`$ is the divergence of the axial-vector current:
$$j_5(x)=(M_Q+m_q):\overline{q}(x)i\gamma _5Q(x):$$
(2)
$`M_Q`$ is the mass of the heavy charm quark $`Q(x)`$ and $`m_q`$ is the mass of the light quarks $`q(x)`$, up, down or strange. The starting point of our sum rules is Cauchy’s theorem applied to the two-point correlation function $`\mathrm{\Pi }(s)`$, weighted with a polynomial $`P(s)`$:
$$\frac{1}{2\pi i}_\mathrm{\Gamma }P(s)\mathrm{\Pi }(s)𝑑s=0.$$
(3)
The integration contour $`\mathrm{\Gamma }`$ extends over a circle of radius $`\left|s\right|=s_0`$, and along both sides of the physical cut $`[s_{phys.},s_0]`$ where $`s_{\text{phys}.}`$ is the physical threshold and $`s_0`$ is a duality parameter to be fixed with some stability criteria. The polynomial $`P(s)`$ does not change the analytical properties of $`\mathrm{\Pi }(s)`$. Therefore we obtain the following sum rule:
$$\frac{1}{\pi }_{s_{\text{phys.}}}^{s_0}P(s)\mathrm{Im}\mathrm{\Pi }(s)𝑑s=\frac{1}{2\pi i}_{\left|s\right|=s_0}P(s)\mathrm{\Pi }(s)𝑑s.$$
(4)
On the left hand side of equation (4) we use the experimental information for $`\mathrm{Im}\mathrm{\Pi }^{\text{DATA}}(s)`$between the physical threshold $`s_{\text{phys}.}`$ and $`s_0`$, whereas on the right hand side we use the asymptotic expansion $`\mathrm{\Pi }^{\mathrm{QCD}}(s)`$ of QCD, including perturbative and non-perturbative terms. The QCD expansion constitutes a good approximation of the two-point correlate on the circle for a large enough integration radius $`s_0`$.
The experimental data are dominated by the first pseudoscalar $`c\overline{q}`$ resonance. In the narrow width approximation, the absorptive part of the two-point function $`\mathrm{Im}\mathrm{\Pi }^{\text{DATA}}(s)`$ can be split in two terms, the contribution of the resonance and the contribution of the hadronic continuum $`\mathrm{Im}\mathrm{\Pi }^{\text{cont}}`$ starting at the physical continuum threshold $`s_{\text{cont}.}>s_{\text{phys.}}=M_{D_q}^2`$, as follows:
$$\frac{1}{\pi }\mathrm{Im}\mathrm{\Pi }^{\text{DATA}}(s)=M_{D_q}^4f_{D_q}^2\delta (sM_{D_q}^2)+\frac{1}{\pi }\mathrm{Im}\mathrm{\Pi }^{\text{cont}}\theta (ss_{\text{cont}.})$$
(5)
where $`M_{D_\mathrm{q}}`$ and $`f_{D_\mathrm{q}}`$ are respectively the mass and the decay constants of the lowest lying pseudoscalar meson $`D_q`$.
For the QCD correlate we write the decomposition,
$$\mathrm{\Pi }^{\mathrm{QCD}}(s)=\mathrm{\Pi }^{\mathrm{pert}.}(s)+\mathrm{\Pi }^{\mathrm{nonpert}.}(s),$$
(6)
We employ the two-point correlation function $`\mathrm{\Pi }^{\text{pert.}}(s)`$ with one massless and one heavy quark given to second order (three loops) in the strong coupling constant $`\alpha _s`$ and expanded in a power series in the pole mass of the heavy quark including terms of order $`(M_c^2/s)^7`$. In the following compact expansion of the two-point function in terms of the pole mass $`M_c`$ can be found
$$\mathrm{\Pi }^{\text{pert}.}(s)=\mathrm{\Pi }^{(0)}(s)+\left(\frac{\alpha _s(M_c)}{\pi }\right)\mathrm{\Pi }^{(1)}(s)+\left(\frac{\alpha _s(M_c)}{\pi }\right)^2\mathrm{\Pi }^{(2)}(s),$$
(7)
where the different terms of the expansion in $`\alpha _s`$ have the form:
$$\mathrm{\Pi }^{(i)}(s)=(M_c+m_q)^2M_c^2\underset{j=1}{\overset{6}{}}\underset{k=0}{\overset{3}{}}A_{j,k}^{(i)}\left(\mathrm{ln}\frac{s}{M_c^2}\right)^k\left(\frac{M_c^2}{s}\right)^j.$$
(8)
In the equations (7,8), $`M_c`$ is the pole mass of the charm quark. The coefficients $`A_{j,k}^{(i)}`$ are explicitly given in . For instance, the one-loop term of the expansion in $`\alpha _s`$ reads :
$`\mathrm{\Pi }^{(0)}(s)`$ $`={\displaystyle \frac{3}{16\pi ^2}}(M_c+m_q)^2s\{32\mathrm{log}\left({\displaystyle \frac{s}{M_c^2}}\right)+4{\displaystyle \frac{M_c^2}{s}}\mathrm{log}\left({\displaystyle \frac{s}{M_c^2}}\right)`$
$`\left[3+2\mathrm{log}\left({\displaystyle \frac{s}{M_c^2}}\right)\right]\left({\displaystyle \frac{M_c^2}{s}}\right)^2+{\displaystyle \frac{2}{3}}\left({\displaystyle \frac{M_c^2}{s}}\right)^3+{\displaystyle \frac{1}{6}}\left({\displaystyle \frac{M_c^2}{s}}\right)^4`$
$`+{\displaystyle \frac{1}{15}}\left({\displaystyle \frac{M_c^2}{s}}\right)^5+{\displaystyle \frac{1}{30}}\left({\displaystyle \frac{M_c^2}{s}}\right)^6+{\displaystyle \frac{2}{105}}\left({\displaystyle \frac{M_c^2}{s}}\right)^7+\mathrm{}..\}`$
The non-perturbative terms in the asymptotic expansion of equation (6) are due to the quark and gluon condensates. We will include terms up to dimension six :
$`\mathrm{\Pi }^{\mathrm{nonpert}.}(s)`$ $`=(M_c+m_q)^2\{M_c\overline{q}q[{\displaystyle \frac{1}{sM_c^2}}(1+2{\displaystyle \frac{\alpha _s}{\pi }})+2{\displaystyle \frac{\alpha _s}{\pi }}\mathrm{ln}{\displaystyle \frac{M_c^2}{s+M_c^2}}]`$
$`{\displaystyle \frac{1}{12}}{\displaystyle \frac{1}{sM_c^2}}{\displaystyle \frac{\alpha _s}{\pi }}G^2{\displaystyle \frac{1}{2}}M_c\left[{\displaystyle \frac{1}{(sM_c^2)^2}}+{\displaystyle \frac{M_c^2}{(sM_c^2)^3}}\right]\overline{q}\sigma Gq`$
$`{\displaystyle \frac{8}{27}}\pi [{\displaystyle \frac{2}{(sM_c^2)^2}}+{\displaystyle \frac{M_c^2}{(sM_c^2)^3}}{\displaystyle \frac{M_c^4}{(sM_c^2)^4}}]\alpha _s\overline{q}q^2\}`$
For the quark condensate we include the $`\alpha _s`$ correction , it turns out to be small but non-negligible.
In order to improve the convergence of the perturbative expansion, we replace the pole mass by the running mass using the $`O(\alpha _s^2)`$ result relating the two . The perturbative piece of order $`(\alpha _s)^i`$ of equation (8) can be rewritten in the form
$$\mathrm{\Pi }^{(i)}(s)=m_c^2(\mu )(m_c(\mu )+m_q(\mu ))^2\underset{j=1}{\overset{6}{}}\underset{k=0}{\overset{3}{}}\stackrel{~}{A}_{j,k}^{(i)}\left(\mathrm{ln}\frac{s}{\mu ^2}\right)^k\left(\frac{m_c^2(\mu )}{s}\right)^j(i=0,1,2).$$
(9)
and similarly for the non-perturbative piece. The coefficients $`\stackrel{~}{A}_{j,k}^{(i)}`$ depend on the mass logarithms $`\mathrm{ln}(m_c^2/\mu ^2)`$ up to the third power. As $`\mathrm{\Pi }(s)^{\mathrm{QCD}}`$ is not known to all orders in $`\alpha _s`$, the results of our analysis will depend to some extend on the choice of the renormalization point $`\mu `$. In the sum rule considered here there are two obvious choices, $`\mu =m_c`$ and $`\mu =\sqrt{s_0}`$. The former choice will sum up the mass logs of the form $`\mathrm{ln}(m_c^2/\mu ^2)`$ and the latter choice the $`\mathrm{ln}(s/\mu ^2)`$ terms. For definiteness, we take $`\mu =m_c`$. With this, the convergence of the perturbative terms is reasonably good. The results differ from taking $`\mu =\sqrt{s_0}`$ by an amount consistent with the general three-loop asymptotic uncertainties, as we will analyze below.
Looking back to equation (4) and taking all the theoretical parameters as well as the mass of the $`D_q`$-meson and the physical continuum threshold as inputs of the calculation, we see that the decay constants can be computed from equation (5) only if we have good control over the hadron continuum contribution of the experimental side.
To cope with this problem we make use of the freedom of choosing the polynomial in equation (4). We take for $`P(s)`$ a polynomial of the form:
$$P_n(s)=a_0+a_1s+a_2s^2+a_3s^3+\mathrm{}+a_ns^n,$$
(10)
where the coefficients are fixed by imposing a normalization condition at threshold
$$P_n\left(s_{\mathrm{phys}.}=M_{D_q}^2\right)=\mathrm{\hspace{0.17em}1},$$
(11)
and requiring that the polynomial $`P_n(s)`$ should minimize the contribution of the continuum in the range $`[s_{\mathrm{cont}.},s_0]`$ in a least square sense, i.e.,
$$_{s_{\mathrm{cont}.}}^{s_0}s^kP_n(s)𝑑s=0\mathrm{for}k=0,\mathrm{}n1,$$
(12)
The polynomials obtained in this way are closely related to the Legendre polynomials. In the appendix the explicit form of the set of polynomials used in this work is given.
This way of introducing the polynomial weight in the sum rule minimizes the continuum contribution $`\frac{1}{\pi }\mathrm{Im}\mathrm{\Pi }^{\mathrm{cont}.}\theta (ss_{\mathrm{cont}.})`$ on the phenomenological side of the sum rule. To the extend that $`\mathrm{Im}\mathrm{\Pi }^{\mathrm{cont}.}`$ can be approximated by an $`n`$-th degree polynomial these conditions lead actually to an exact cancellation of the continuum contribution to the left hand side of equation (4). The role of the $`D_q`$ resonance will be enhanced. We will see in our analysis that this choice of the polynomial has the additional effect of increasing the region of duality characterized by the value of the duality parameter $`s_0`$. In this way the asymptotic expansion of QCD can be used more safely on the circular integration contour. Notice however that increasing the degree of the polynomial $`P_n(s)`$ will require the knowledge of further terms in the mass expansion and in the non-perturbative series. Therefore the polynomial degree cannot be chosen arbitrarily high.
To check the consistency of the method, we have employed polynomials ranging from second degree to fifth degree, verifying that the results are compatible within the range of the errors introduced by the incomplete knowledge of the QCD correlate and other inputs of the calculation. We also have checked explicitly that a smooth continuum contribution had no influence on the result.
Our approach to suppress the continuum has been tested previously in the calculation of the heavy quark masses using analogue sum rules for the vector current correlate where there exists more experimental information on the continuum. In the calculation of the charm quark mass, using the same polynomial method, the continuum, known from the BES II collaboration , was shown to have practically no influence on the predicted mass . Employing the same technique, a very accurate prediction of the bottom quark mass was also obtained using the experimental information of the upsilon system .
After these general considerations we proceed with the analytical calculation. The integrals that we have to evaluate on the right hand side of the sum rule, equation (4), are
$$J(p,k)=\frac{1}{2\pi i}_{\left|s\right|=s_0}s^p\left(\mathrm{ln}\frac{s}{\mu ^2}\right)^k𝑑s,$$
(13)
for $`k=0,1,2,3`$ and $`p=6,5,..,n+1`$. These integrals can be found e.g. in reference . After integration, equation (4) yields the sum rule
$`M_{D_q}^4f_{D_q}^2P(M_{D_q}^2)`$ $`=m_c^2(\mu )(m_c(\mu )+m_q(\mu ))^2`$ (14)
$`\times {\displaystyle \underset{p=0}{\overset{n}{}}}{\displaystyle \underset{i=0}{\overset{2}{}}}{\displaystyle \underset{j=1}{\overset{6}{}}}{\displaystyle \underset{k=0}{\overset{3}{}}}a_p\left({\displaystyle \frac{\alpha _s(\mu )}{\pi }}\right)^i\stackrel{~}{A}_{j,k}^{(i)}m_c^{2j}(\mu )J(pj,k)`$
$`+\text{ non-perturbative terms}`$
where, for brevity, we have not written down the non-perturbative terms explicitly. The contribution of the continuum is neglected as explained above although the continuum threshold is considered in the determination of the coefficients, $`a_p`$ of the polynomials (12).
Plugging the theoretical and experimental inputs (physical threshold, quarks and meson masses, condensates and strong coupling constant) into the sum rule, we obtain the decay constant $`f_{D_q}`$ for various values of the degree $`n`$ of the polynomial and various values of $`s_0`$. Given the correct QCD asymptotic correlate and the correct hadron continuum, the calculation of the decay constant should, of course, not depend either on $`s_0`$ or on the degree $`n`$ of the polynomial in the sum rule (4). Accordingly, for a given $`n`$ we choose the flattest region in the curve $`f_{D_q}(s_0)`$ to extract our prediction for the decay constant. To be specific we choose the point of minimal slope. On the other hand, for different polynomials, the value of $`f_{D_q}`$, extracted in this way, could differ from each other as the cancellation of the continuum may be incomplete or the QCD expansion may not be accurate enough. We find, however, practically the same results for all our polynomials. This additional stability is truly remarkable as the coefficients of the polynomials of order $`n=2,3,4`$ and $`5`$ are completely different and the respective predictions are based on completely different superpositions of finite energy moment sum rules. This extended consistency leads us to attach great confidence in our numbers and associated errors.
## 3 Results
We calculate the decay constants for the $`D`$ and $`D_s`$ heavy mesons. In the first case we take $`m_q=0`$ everywhere. In the second case we retain $`m_q=m_s0`$ in the factor $`(M_c+m_q)^2`$ in front of the correlation function only. Further terms in the power series in $`m_s^2/s`$ in (8) are completely negligible for the integration radii $`s_0`$ we use in the calculation.
The experimental and theoretical inputs are as follows. The physical threshold $`s_{\mathrm{phys}.}`$ is the squared mass of the lowest lying resonance in the $`c\overline{q}`$ channel. For $`q`$ being the light quark $`u`$, we have:
$$s_{\mathrm{phys}.}=M_D^2=3.493\text{ GeV}^2$$
(15)
whereas the continuum threshold $`s_{\mathrm{cont}.}`$ is taken from the next intermediate state $`D\pi \pi `$ in an s-wave $`I=\frac{1}{2}`$, i. e.
$$s_{\mathrm{cont}.}=\left(M_D+2m_\pi \right)^2=4.575\mathrm{GeV}^2.$$
For $`q`$ being the strange quark we take:
$$s_{\mathrm{phys}.}=M_{D_s}^2=3.873\text{ GeV}^2$$
(16)
The continuum threshold starts in this case at the value:
$$s_{\mathrm{cont}.}=\left(M_{D_s}+2m_\pi \right)^2=5.009\mathrm{GeV}^2.$$
On the theoretical side of the sum rule we take the following inputs. The strong coupling constant at the scale of the electroweak $`Z`$ boson mass
$$\alpha _s(M_Z)=0.118\pm 0.003$$
(17)
that is run down to the renormalization scale using the four loop formulas of reference . For the quark and gluon condensates (see for example ) and the mass of the strange quark we take:
$`<\overline{q}q>(2\mathrm{GeV})=(267\pm \mathrm{\hspace{0.17em}17}\mathrm{MeV})^3,`$
$`<{\displaystyle \frac{\alpha _s}{\pi }}GG>=\mathrm{\hspace{0.17em}0.024}\pm \mathrm{\hspace{0.17em}0.012}\mathrm{GeV}^4,`$
$`<\overline{q}\sigma Gq>=m_0^2<\overline{q}q>,\mathrm{with}m_0^2=\mathrm{\hspace{0.17em}0.8}\pm 0.2\mathrm{GeV},`$
$`m_s(2\mathrm{GeV})=\mathrm{\hspace{0.17em}120}\pm 50\mathrm{MeV},`$
$`<\overline{s}s>=(0.8\pm 0.3)<\overline{q}q>.`$ (18)
As discussed above, we fix the renormalization scale to be $`\mu =m_c(m_c)`$. We use a reasonable variation of $`\mu `$ to analyze the corresponding uncertainty in our final result. Finally, for the charm quark, we take the value $`m_c(m_c)=\mathrm{\hspace{0.17em}1.25}\pm 0.10\mathrm{GeV}`$ obtained by similar techniques which is in a generally accepted range.
In order to calculate the decay constant for the pseudoscalar meson $`D`$, we proceed in the way described above. We compute$`f_D`$ as a function of $`s_0`$ with the four different sum rules (4) corresponding to $`n=2,3,4,5`$. The results, plotted in Fig. 1 show remarkable stability properties. We define the optimal value of $`s_0`$ as the center of the stability region (represented by a cross in Fig.1) where the first and/or second derivative of $`f_D(s_0)`$ vanishes. At these values of $`s_0`$ we obtain the following consistent results:
$`f_D`$ $`=176\text{ MeV for }n=2`$ (19)
$`f_D`$ $`=177\text{ MeV for }n=3,4,5`$
Notice from Fig. 1 that for the fifth degree polynomial (n=5) there is a stability region of about $`20\mathrm{G}eV^2`$ around $`s_0`$, where the decay constant changes by less than three percent. From this change we estimate a conservative error inherent to the method of $`\pm 3\mathrm{MeV}`$.
Other sources of errors arising in the calculation of $`f_D`$ are the quark condensates which affect the result by $`\pm 5\mathrm{MeV},`$ and the charm mass which, in the range given above, produces a variation in the decay constant of $`19,+9`$ $`\mathrm{MeV}`$. This is one of the main source of uncertainty in the final result (see table 1). Considering the following results of the perturbative $`\alpha _s`$ expansion:
one-loop calculation $`\text{}f_D^{\left(0\right)}=142\text{ MeV}`$
two-loop calculation $`\text{}f_D^{\left(1\right)}=162\text{ MeV}`$ (20)
three-loop calculation $`\text{}f_D^{\left(2\right)}=177\text{ MeV}`$
we also take as a source of uncertainty the contribution of order $`\alpha _s^2`$, which amounts ten percent of the result. This yields an asymptotic error of $`\pm 15`$ MeV
We point out that the convergence of the asymptotic series in the present calculation of the decay constant of the $`D`$ meson is worse than the one we found for the $`B`$ meson . Finally we have considered the dependence on the renormalization scale in the range $`\mu [2,6]\mathrm{GeV}`$. The error associated to this change in $`\mu `$ is roughly related to the convergence of the asymptotic expansion and therefore it is not considered as an additional one. Other errors due to the QCD side of the sum rule, higher order terms in $`m_c^2/s`$ and the error on $`\alpha _s(m_Z)`$, are negligible in comparison.
From this analysis of errors, we finally quote the following result for the decay constant of $`D`$-meson
$$f_D=\mathrm{\hspace{0.17em}177}\pm \mathrm{\hspace{0.17em}14}(\mathrm{inp}.)\pm 15\left(\text{asymp.}\right)\pm \mathrm{\hspace{0.17em}3}(\mathrm{meth}.)\mathrm{MeV}.$$
(21)
The first error comes from the inputs of the computation, the second to the truncated QCD theory whereas the last one is due to the method itself.
Proceeding in the same fashion, but keeping the mass of the strange quark in the overall factor and the order $`m_s/m_c`$ in the one loop contribution, we find the decay constant for the $`D_s`$ meson,
one-loop calculation $`\text{}f_{D_s}^{\left(0\right)}=163\text{ MeV}`$
two-loop calculation $`\text{}f_{D_s}^{\left(1\right)}=188\text{ MeV}`$ (22)
three-loop calculation $`\text{}f_{D_s}^{\left(2\right)}=205\text{ MeV}`$
and including the analysis of uncertainties we find:
$$f_{D_s}=\mathrm{\hspace{0.17em}205}\pm \mathrm{\hspace{0.17em}14}(\mathrm{inp}.)\pm 17\left(\text{asymp.}\right)\pm \mathrm{\hspace{0.17em}3}(\mathrm{meth}.))\mathrm{MeV}.$$
(see Fig. 2)
In the analysis of theoretical errors the only new ingredient is the uncertainty coming from the strange quark mass which turns out to be negligible.
The ratio of the decay constants $`f_{D_s}`$ and $`f_D`$ (which would be $`1`$ in the chiral limit) is of special interest. We find:
$$\frac{f_{D_s}}{f_D}=\mathrm{\hspace{0.17em}1.16}\pm \mathrm{\hspace{0.17em}0.01}(\mathrm{inp}.)\pm \mathrm{\hspace{0.17em}0.03}(\mathrm{meth}.).$$
(23)
in complete agreement with lattice calculations. The uncertainties due to the theoretical inputs are correlated, so that the final error is very small.
## 4 Conclusions
In this note we have computed the decay constant of $`D_q`$-mesons for $`q`$ being either the strange or the $`u`$ or $`d`$ massless quarks. We have employed a judicious combination of moments in QCD finite energy sum rules in order to minimize the shortcomings of the available experimental data. On the theoretical side of the pseudoscalar two-point function, we have used in the perturbative piece an expansion up to three loops in the strong coupling constant and up to order $`\left(m_c^2/s\right)^7`$ in the mass expansion and in the non-perturbative piece we considered condensates up to dimension six including the $`\alpha _s`$ correction in the leading term. Instead of the commonly adopted pole mass of the bottom quark, we use the running mass to improve convergence of the perturbative series.
In the sum rule, the contour integration of the asymptotic part is performed analytically. This particular fact differs from other computations based on sum rules where the asymptotic QCD is integrated along a cut of the two-point function starting at the pole mass squared. The latter way to proceed is problematic when loop corrections are included and the complete analytical QCD expression along the cut is not known. In this approach QCD has to be extrapolated from low energy to high energy . We also differ from many other sum rule calculation in that we do not require two unrelated sum rules to determine a duality point via an intercept of the curves $`f_D(s_0)`$.
Our results are very sensitive to the value of the running mass, giving most of the theoretical uncertainty. They also turn out to be sensitive to variations of the renormalization scale $`\mu .`$ The uncertainties of other theoretical parameters like quark condensates and coupling constant are less important. Adding quadratically the different estimated errors we have the final results
$`f_D`$ $`=`$ $`\mathrm{\hspace{0.17em}177}\pm \mathrm{\hspace{0.17em}21}\mathrm{MeV}`$
$`f_{D_s}`$ $`=`$ $`\mathrm{\hspace{0.17em}205}\pm \mathrm{\hspace{0.17em}22}\mathrm{MeV}`$ (24)
Comparing (24) with other results in the literature, our results agree within the error bars with the ones obtained using sum rule methods . However, compared with lattice computations they are a bit lower.
## Appendix
For convenience of the reader we list in this appendix the first few polynomials emerging from relations (11,12). From the second condition, namely (12), it is easy to realize that the set of polynomials $`P_n(s)`$ are n-degree orthogonal polynomials in the interval of the variable $`s[s_{\mathrm{cont}.},s_0]`$. With the normalization condition (11) (adopted to emphasize the contribution of the lowest lying resonance in the sum rule) they are related to the Legendre polynomials $`𝒫_n(x)`$ in the interval of the variable $`x[1,1]`$ as follows:
$$P_n(s)=\frac{𝒫_n\left(x(s)\right)}{𝒫_n\left(x(M_{D_q}^2)\right)}$$
(25)
Where the variable $`x(s)`$ is:
$$x(s)=\frac{2s(s_0+s_{\mathrm{cont}.})}{s_0s_{\mathrm{cont}.}}$$
Obviously $`x(s)[1,1]`$ when $`s[s_{\mathrm{cont}.},s_0]`$.
The explicit form of these polynomials is well known and can be found, for instance, in . Nevertheless, for sake of completeness, we quote here the ones we have used in the calculation.
$`𝒫_2\left(x(s)\right)={\displaystyle \frac{1}{2}}(3x^21),`$
$`𝒫_3\left(x(s)\right)={\displaystyle \frac{1}{2}}(5x^33x),`$
$`𝒫_4\left(x(s)\right)={\displaystyle \frac{1}{8}}(35x^430x^2+3),`$
$`𝒫_5\left(x(s)\right)={\displaystyle \frac{1}{8}}(63x^570x^3+15x).`$
Finally, in Fig. 3 and in order to appreciate the suppression of the experimental physical continuum data in the sum rule, we have plotted the form of the polynomials $`P_n(s)`$ for n=2,3,4,5 at the stability values of $`s_0`$ used in the calculation of $`f_D`$. |
warning/0507/nlin0507011.html | ar5iv | text | # Completeness of the cubic and quartic Hénon-Heiles Hamiltonians11footnote 1Corresponding author RC. Preprint S2004/047. nlin.SI/0507011
## 1 Introduction
The considered Hamiltonian originates from celestial mechanics, as a system describing the motion of a star in the axisymmetric potential of the galaxy. Denoting $`q_1`$ the radius and $`q_2`$ the altitude, this “Hénon-Heiles Hamiltonian” (HH) is the sum of a kinetic energy and a potential energy, in which the potential is a cubic polynomial in the position variables $`q_1,q_2`$,
$`H`$ $`=`$ $`{\displaystyle \frac{1}{2}}(p_1^2+p_2^2+q_1^2+q_2^2)+q_1q_2^2{\displaystyle \frac{1}{3}}q_1^3,`$ (1)
it is nonintegrable and displays a strange attractor. However, if one changes the numerical coefficients in the potential, the system may become integrable, and this question (to find all the integrable cases and to integrate them) has attracted a lot of activity in the last three decennia.
A prerequisite is to define the word integrability, and in section 2 we briefly recall its three main acceptations in the context of Hamiltonian systems.
In section 3, we recall all the cases (three “cubic” plus four “quartic”) for which the most general two-degree of freedom classical time-independent Hamiltonian may have a single valued general solution.
Then, discarding the integrated cases (see for a review of the current state of this problem), we focus on the three cases (all “quartic”) for which the general solution is still missing, with the aim of finding this general solution.
In section 4, we build an equivalent fourth order ordinary differential equation (ODE) for $`q_1(t)`$, in the hope of finding it listed in one of the classical tables of explicitly integrated ODEs. This hope is deceived because these tables are not yet finished.
This is why, in the last two sections, we adopt a different strategy. In front of the difficulty to perform the separation of variables in the sense of Arnol’d and Liouville, we establish a birational transformation between the two second order Hamilton equations and a fourth order ODE listed in a classical table established by Cosgrove , whose general solution is single valued.
## 2 Integrability for Hamiltonian systems
Given a Hamiltonian system with a finite number $`N`$ of degrees of freedom, three main definitions of integrability are known,
1. the one in the sense of Liouville, that is the existence of $`N`$ independent invariants $`K_j`$ whose pairwise Poisson brackets vanish, $`\{K_j,K_l\}=0`$,
2. the one in the sense of Arnol’d-Liouville \[2, chap. 9\], which is to find explicitly some canonical variables $`s_j,r_j,j=1,N`$ which “separate” the Hamilton-Jacobi equation for the action $`S`$, which for two degrees of freedom writes as,
$`H(q_1,q_2,p_1,p_2)E=0,p_1={\displaystyle \frac{S}{q_1}},p_2={\displaystyle \frac{S}{q_2}},`$ (2)
3. the one in the sense of Painlevé i.e. the representation of the general solution $`q_j(t)`$ by an explicit, closed form, single valued expression of the time $`t`$.
## 3 The seven Hénon-Heiles Hamiltonians
Given the most general two-degree of freedom classical time-independent Hamiltonian
$`H`$ $`=`$ $`{\displaystyle \frac{1}{2}}(p_1^2+p_2^2)+V(q_1,q_2)=E,`$ (3)
the requirement that the system made of the two Hamilton equations passes the Painlevé test (for at least some integer powers $`q_1^{n_1},q_2^{n_2}`$) selects seven and only seven potentials $`V`$ depending on a finite number of constants, namely
1. three “cubic” potentials (HH3 case) ,
$`H`$ $`=`$ $`{\displaystyle \frac{1}{2}}(p_1^2+p_2^2+\omega _1q_1^2+\omega _2q_2^2)+\alpha q_1q_2^2{\displaystyle \frac{1}{3}}\beta q_1^3+{\displaystyle \frac{1}{2}}\gamma q_2^2,\alpha 0`$ (4)
in which the constants $`\alpha ,\beta ,\omega _1,\omega _2,\gamma `$ can only take three sets of values,
$`\text{(SK)}:`$ $`\beta /\alpha =1,\omega _1=\omega _2,`$ (5)
$`\text{(KdV5)}:`$ $`\beta /\alpha =6,`$ (6)
$`\text{(KK)}:`$ $`\beta /\alpha =16,\omega _1=16\omega _2.`$ (7)
2. four “quartic” potentials (HH4 case) ,
$`H`$ $`=`$ $`{\displaystyle \frac{1}{2}}(P_1^2+P_2^2+\mathrm{\Omega }_1Q_1^2+\mathrm{\Omega }_2Q_2^2)+CQ_1^4+BQ_1^2Q_2^2+AQ_2^4`$ (8)
$`+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\alpha }{Q_1^2}}+{\displaystyle \frac{\beta }{Q_2^2}}\right)\gamma Q_1,B0,`$
in which the constants $`A,B,C,\alpha ,\beta ,\gamma ,\mathrm{\Omega }_1,\mathrm{\Omega }_2`$ can only take the four values (the notation $`A:B:C=p:q:r`$ stands for $`A/p=B/q=C/r=\text{arbitrary}`$),
$`\{\begin{array}{cc}A:B:C=1:2:1,\gamma =0,\hfill & \\ A:B:C=1:6:1,\gamma =0,\mathrm{\Omega }_1=\mathrm{\Omega }_2,\hfill & \\ A:B:C=1:6:8,\alpha =0,\mathrm{\Omega }_1=4\mathrm{\Omega }_2,\hfill & \\ A:B:C=1:12:16,\gamma =0,\mathrm{\Omega }_1=4\mathrm{\Omega }_2.\hfill & \end{array}`$ (13)
All seven cases are integrable in the sense of Liouville, with a second constant of the motion $`K`$ either quadratic or quartic in the momenta $`p_1,p_2`$.
In the sense of Arnol’d-Liouville, the separation of variables has been performed , except in three cases,
1. HH4 1:6:1 $`\alpha \beta `$,
2. HH4 1:6:8 $`\beta \gamma 0`$,
3. HH4 1:12:16 $`\alpha \beta 0`$.
What is remarkable is the fact that, in all cases when the separation of variables is achieved, the equations of Hamilton have the Painlevé property, the general solution being a hyperelliptic function of genus two. The purpose of this work is to prove equally the Painlevé property in the three remaining cases where the separation of variables is not yet performed.
## 4 Equivalent fourth order ODEs
In the cubic case, the two Hamilton equations
$`q_1^{\prime \prime }+\omega _1q_1\beta q_1^2+\alpha q_2^2=0,`$ (14)
$`q_2^{\prime \prime }+\omega _2q_2+2\alpha q_1q_2\gamma q_2^3=0,`$ (15)
together with the Hamiltonian (4), are equivalent to a single fourth order ODE for $`q_1(t)`$,
$`q_1^{\prime \prime \prime \prime }+(8\alpha 2\beta )q_1q_1^{\prime \prime }2(\alpha +\beta )q_1^2{\displaystyle \frac{20}{3}}\alpha \beta q_1^3`$
$`+(\omega _1+4\omega _2)q_1^{\prime \prime }+(6\alpha \omega _14\beta \omega _2)q_1^2+4\omega _1\omega _2q_1+4\alpha E=0,`$ (16)
independent of the coefficient $`\gamma `$ of the nonpolynomial term $`q_2^2`$ and depending on the constant value $`E`$ of the Hamiltonian $`H`$. In the three HH3 cases (5)–(7), this ODE belongs to a list (“classification”) of equations enjoying the Painlevé property, whose general solution is hyperelliptic with genus two.
In the quartic case, the similar fourth order equation is built by eliminating $`Q_2`$ and $`Q_{1}^{\prime \prime \prime }{}_{}{}^{2}`$ between the two Hamilton equations,
$`Q_1^{\prime \prime }+\mathrm{\Omega }_1Q_1+4CQ_1^3+2BQ_1Q_2^2\alpha Q_1^3+\gamma =0,`$ (17)
$`Q_2^{\prime \prime }+\mathrm{\Omega }_2Q_2+4AQ_2^3+2BQ_2Q_1^2\beta Q_2^3=0,`$ (18)
and the Hamiltonian (8), which results in
$`Q_1^{\prime \prime \prime \prime }+2{\displaystyle \frac{Q_1^{}Q_1^{\prime \prime \prime }}{Q_1}}+\left(1+6{\displaystyle \frac{A}{B}}\right){\displaystyle \frac{Q_{1}^{\prime \prime }{}_{}{}^{2}}{Q_1}}2{\displaystyle \frac{Q_{1}^{}{}_{}{}^{2}Q_1^{\prime \prime }}{Q_1^2}}`$
$`+8\left(6{\displaystyle \frac{AC}{B}}BC\right)Q_1^2Q_1^{\prime \prime }+4(B2C)Q_1Q_{1}^{}{}_{}{}^{2}+24C\left(4{\displaystyle \frac{AC}{B}}B\right)Q_1^5`$
$`+\left[12{\displaystyle \frac{A}{B}}\omega _14\omega _2+\left(1+12{\displaystyle \frac{A}{B}}\right){\displaystyle \frac{\gamma }{Q_1}}4\left(1+3{\displaystyle \frac{A}{B}}\right){\displaystyle \frac{\alpha }{Q_1^4}}\right]Q_1^{\prime \prime }`$
$`+6{\displaystyle \frac{A}{B}}{\displaystyle \frac{\alpha ^2}{Q_1^7}}+20{\displaystyle \frac{\alpha }{Q_1^5}}Q_{1}^{}{}_{}{}^{2}12{\displaystyle \frac{A}{B}}{\displaystyle \frac{\gamma \alpha }{Q_1^4}}+4\left(3{\displaystyle \frac{A}{B}}\omega _1\omega _2\right)\left(\gamma {\displaystyle \frac{\alpha }{Q_1^3}}\right)2\gamma {\displaystyle \frac{Q_{1}^{}{}_{}{}^{2}}{Q_1^2}}`$
$`+6\left({\displaystyle \frac{A}{B}}\gamma ^2+2B\alpha 8{\displaystyle \frac{AC}{B}}\alpha \right){\displaystyle \frac{1}{Q_1}}+\left(6{\displaystyle \frac{A}{B}}\omega _1^24\omega _1\omega _28BE\right)Q_1`$
$`+48{\displaystyle \frac{AC}{B}}\gamma Q_1^2+4\left(12{\displaystyle \frac{AC}{B}}B4C\right)\omega _1Q_1^3.`$ (19)
This ODE depends on $`E`$ but not on $`\beta `$ and, as opposed to the cubic case, it does not belong to a classified set of equations, because $`Q_1^{\prime \prime \prime \prime }`$ is not polynomial in $`Q_1`$.
In the three remaining cases, since one is yet unable either to perform the separation of variables or to establish a direct link to a classified ODE, let us build an indirect link to such a classified ODE. This link, which involves soliton equations, is the following.
For each of the seven cases, the two Hamilton equations are equivalent to the traveling wave reduction of a soliton system made either of a single PDE (HH3) or of two coupled PDEs (HH4), most of them appearing in lists established from group theory . Among the various soliton equations which are equivalent to them via a Bäcklund transformation, some of them admit a traveling wave reduction to a classified ODE. This property defines a path which starts from one of the three remaining HH4 cases, goes up to a soliton system of two coupled 1+1-dimensional PDEs admitting a reduction to the considered case, then goes to another 1+1-dim PDE system equivalent under a Bäcklund transformation, finally goes down by reduction to an already integrated ODE or system of ODEs.
## 5 General solution of the quartic 1:6:1 and 1:6:8 cases
Let us denote the two constants of the motion of the 1:6:1 and 1:6:8 cases as,
$`1:6:1\{\begin{array}{cc}H={\displaystyle \frac{1}{2}}(P_1^2+P_2^2)+{\displaystyle \frac{\mathrm{\Omega }}{2}}(Q_1^2+Q_2^2){\displaystyle \frac{1}{32}}(Q_1^4+6Q_1^2Q_2^2+Q_2^4)\hfill & \\ {\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\kappa _1^2}{Q_1^2}}+{\displaystyle \frac{\kappa _2^2}{Q_2^2}}\right)=E,\hfill & \\ K=\left(P_1P_2+Q_1Q_2\left({\displaystyle \frac{Q_1^2+Q_2^2}{8}}+\mathrm{\Omega }\right)\right)^2\hfill & \\ P_2^2{\displaystyle \frac{\kappa _1^2}{Q_1^2}}P_1^2{\displaystyle \frac{\kappa _2^2}{Q_2^2}}+{\displaystyle \frac{1}{4}}\left(\kappa _1^2Q_2^2+\kappa _2^2Q_1^2\right)+{\displaystyle \frac{\kappa _1^2\kappa _2^2}{Q_1^2Q_2^2}},\hfill & \end{array}`$ (24)
and
$`1:6:8\{\begin{array}{cc}H={\displaystyle \frac{1}{2}}(p_1^2+p_2^2)+{\displaystyle \frac{\omega }{2}}(4q_1^2+q_2^2){\displaystyle \frac{1}{16}}(8q_1^4+6q_1^2q_2^2+q_2^4)\hfill & \\ \gamma q_1+{\displaystyle \frac{\beta }{2q_2^2}}=E,\hfill & \\ K=\left(p_2^2{\displaystyle \frac{q_2^2}{16}}(2q_2^2+4q_1^2+\omega )+{\displaystyle \frac{\beta }{q_2^2}}\right)^2{\displaystyle \frac{1}{4}}q_2^2(q_2p_12q_1p_2)^2\hfill & \\ +\gamma \left(2\gamma q_2^24q_2p_1p_2+{\displaystyle \frac{1}{2}}q_1q_2^4+q_1^3q_2^2+4q_1p_2^24\omega q_1q_2^2+4q_1{\displaystyle \frac{\beta }{q_2^2}}\right).\hfill & \end{array}`$ (29)
There is a canonical transformation between the 1:6:1 and 1:6:8 cases, mapping the constants as follows,
$`E_{1:6:8}=E_{1:6:1},K_{1:6:8}=K_{1:6:1},\omega =\mathrm{\Omega },\gamma ={\displaystyle \frac{\kappa _1+\kappa _2}{2}},\beta =(\kappa _1\kappa _2)^2,`$ (30)
therefore one only needs to integrate either case.
The path to an integrated ODE comprises the following three segments.
The coordinate $`q_1(t)`$ of the 1:6:8 case can be identified to the component $`F`$ of the traveling wave reduction $`f(x,\tau )=F(xc\tau ),g(x,\tau )=G(xc\tau )`$ of a soliton system of two coupled KdV-like equations (c-KdV system) denoted c-KdV<sub>1</sub>
$`\{\begin{array}{cc}f_\tau +\left(f_{xx}+{\displaystyle \frac{3}{2}}ff_x{\displaystyle \frac{1}{2}}f^3+3fg\right)_x=0,\hfill & \\ 2g_\tau +g_{xxx}+6gg_x+3fg_{xx}+6gf_{xx}+9f_xg_x3f^2g_x\hfill & \\ +{\displaystyle \frac{3}{2}}f_{xxxx}+{\displaystyle \frac{3}{2}}ff_{xxx}+9f_xf_{xx}3f^2f_{xx}3ff_x^2=0,\hfill & \end{array}`$ (34)
with the identification
$`\{\begin{array}{cc}q_1=F,q_2^2=2\left(F^{}+F^2+2G2\omega \right),\hfill & \\ c=\omega ,K_1=\gamma ,K_2=E,\hfill & \end{array}`$ (37)
in which $`K_1`$ and $`K_2`$ are two constants of integration.
There exists a Bäcklund transformation between this soliton system and another one of the c-KdV type, denoted bi-SH system ,
$`\{\begin{array}{cc}2u_\tau +\left(u_{xx}+u^2+6v\right)_x=0,\hfill & \\ v_\tau +v_{xxx}+uv_x=0.\hfill & \end{array}`$ (40)
This BT is defined by the Miura transformation
$`\{\begin{array}{cc}u={\displaystyle \frac{3}{2}}\left(2gf_xf^2\right),\hfill & \\ v={\displaystyle \frac{3}{4}}\left(2f_{xxx}+4ff_{xx}+8gf_x+4fg_x+3f_x^22f^2f_xf^4+4gf^2\right).\hfill & \end{array}`$ (43)
Finally, the traveling wave reduction $`u(x,\tau )=U(xc\tau ),v(x,\tau )=V(xc\tau )`$ can be identified to the autonomous F-VI equation (a-F-VI) in the classification of Cosgrove ,
$`\text{a-F-VI}:y^{\prime \prime \prime \prime }=18yy^{\prime \prime }+9y_{}^{}{}_{}{}^{2}24y^3+\alpha _{\mathrm{VI}}y^2+{\displaystyle \frac{\alpha _{\mathrm{VI}}^2}{9}}y+\kappa _{\mathrm{VI}}t+\beta _{\mathrm{VI}},\kappa _{\mathrm{VI}}=0,`$ (44)
an ODE whose general solution is meromorphic, expressed with genus two hyperelliptic functions \[11, Eq. (7.26)\]. The identification is
$`\{\begin{array}{cc}U=6\left(y+{\displaystyle \frac{c}{18}}\right),\hfill & \\ V=y^{\prime \prime }6y^2+{\displaystyle \frac{4}{3}}cy+{\displaystyle \frac{16}{27}}c^2{\displaystyle \frac{K_A}{2}},\hfill & \\ \alpha _{\mathrm{VI}}=4c,\beta _{\mathrm{VI}}=K_B2cK_A+{\displaystyle \frac{512}{243}}c^3,\hfill & \end{array}`$ (48)
in which $`K_A,K_B`$ are two constants of integration.
In order to perform the integration of both the 1:6:1 and the 1:6:8 cases, it is sufficient to express $`(F,G)`$ rationally in terms of $`(U,V,U^{},V^{})`$. The result is
$`\{\begin{array}{cc}F={\displaystyle \frac{W^{}}{2W}}+{\displaystyle \frac{K_1}{24W}}[3U_{}^{}{}_{}{}^{2}2(U3c)(12V+(U+3c)^2)\hfill & \\ +36K_B54K_1^2],\hfill & \\ G={\displaystyle \frac{U}{3}}+{\displaystyle \frac{1}{8W}}[(2V+3K_2)(2V^{\prime \prime }+K_1U^{}3K_1^2)\hfill & \\ 2(U3c)(2K_1V^{}+K_1^2(U+3c))],\hfill & \\ W=\left(V+{\displaystyle \frac{3}{2}}K_2\right)^2+{\displaystyle \frac{3}{2}}K_1^2(U3c),\hfill & \\ K_A=K_2.\hfill & \end{array}`$ (55)
Making the product of the successive transformations (37), (55), (48), one obtains a meromorphic general solution for $`Q_1^2,Q_2^2,q_1,q_2^2`$,
$`\{\begin{array}{cc}q_1={\displaystyle \frac{W^{}}{2W}}+{\displaystyle \frac{\gamma }{W}}\left[9j3\left(y+{\displaystyle \frac{4}{9}}\omega \right)(h+E){\displaystyle \frac{9}{4}}\gamma ^2\right],\hfill & \\ q_2^2=16\left(y{\displaystyle \frac{5}{9}}\omega \right)\hfill & \\ +{\displaystyle \frac{1}{W}}[12(y^{}+{\displaystyle \frac{\gamma }{2}})^248y^316\omega y^2+(24E+{\displaystyle \frac{128}{9}}\omega ^2)y+{\displaystyle \frac{1280}{243}}\omega ^3.\hfill & \\ .{\displaystyle \frac{40}{3}}\omega E+{\displaystyle \frac{3}{4}}\beta 24\gamma (y{\displaystyle \frac{5}{9}}\omega )h^{}144\gamma ^2(y{\displaystyle \frac{5}{9}}\omega )^2],\hfill & \\ W=(h+E)^29\gamma ^2\left(y{\displaystyle \frac{5}{9}}\omega \right),\hfill & \\ \alpha _{\mathrm{VI}}=4\omega ,\beta _{\mathrm{VI}}={\displaystyle \frac{3}{4}}\gamma ^2+2\omega E{\displaystyle \frac{3}{16}}\beta {\displaystyle \frac{512}{243}}\omega ^3,\hfill & \\ K_{1,\mathrm{VI}}={\displaystyle \frac{3}{32}}K{\displaystyle \frac{1}{2}}E^2,K_{2,\mathrm{VI}}={\displaystyle \frac{3}{32}}EK{\displaystyle \frac{1}{3}}E^3+{\displaystyle \frac{9}{64}}\beta \gamma ^2,\hfill & \\ K_1=\gamma ,K_2=E,K_A=E,K_B={\displaystyle \frac{3}{16}}\beta +{\displaystyle \frac{3}{4}}\gamma ^2.\hfill & \end{array}`$ (64)
in which $`h`$ and $`j`$ are convenient auxiliary variables \[11, Eqs. (7.4)–(7.5)\],
$`\{\begin{array}{cc}y={\displaystyle \frac{Q(s_1,s_2)+\sqrt{Q(s_1)Q(s_2)}}{2\left(\sqrt{s_1^2C_{\mathrm{VI}}}+\sqrt{s_2^2C_{\mathrm{VI}}}\right)^2}}+{\displaystyle \frac{5}{36}}\alpha _{\mathrm{VI}},\hfill & \\ h={\displaystyle \frac{3}{4}}E_{\mathrm{VI}}{\displaystyle \frac{s_1s_2+C_{\mathrm{VI}}+\sqrt{(s_1^2C_{\mathrm{VI}})(s_2^2C_{\mathrm{VI}})}}{s_1+s_2}}{\displaystyle \frac{F_{\mathrm{VI}}}{2}},\hfill & \\ j={\displaystyle \frac{1}{6}}(2h+F_{\mathrm{VI}})\{y+{\displaystyle \frac{\alpha _{\mathrm{VI}}}{9}}{\displaystyle \frac{E_{\mathrm{VI}}}{4(s_1+s_2)}}.\}\hfill & \end{array}`$ (68)
In the above, the variables $`s_1,s_2`$ are defined by the hyperelliptic system
$`\{\begin{array}{cc}(s_1s_2)s_1^{}=\sqrt{P(s_1)},(s_2s_1)s_2^{}=\sqrt{P(s_2)},\hfill & \\ P(s)=(s^2C_{\mathrm{VI}})Q(s),\hfill & \\ Q(s,t)=(s^2C_{\mathrm{VI}})(t^2C_{\mathrm{VI}}){\displaystyle \frac{\alpha _{\mathrm{VI}}}{2}}(s^2+t^22C_{\mathrm{VI}})+{\displaystyle \frac{E_{\mathrm{VI}}}{2}}(s+t)+F_{\mathrm{VI}},\hfill & \\ Q(s)=Q(s,s).\hfill & \end{array}`$ (73)
The expressions (68) cannot be written as rational functions of $`s_1,s_2,s_1^{},s_2^{}`$ and are nevertheless meromorphic .
The coefficients $`(\alpha ,C,E,F)_{\mathrm{VI}}`$ of the hyperelliptic curve depend algebraically on the parameters of the Hamiltonians $`\beta ,\gamma ,\omega ,E,K`$ \[11, Eqs. (7.9)-(7.12)\]
$`\{\begin{array}{cc}A_{\mathrm{VI}}=4\omega ,\hfill & \\ E_{\mathrm{VI}}^2={\displaystyle \frac{16}{3}}\omega (F_{\mathrm{VI}}2E)\beta +4\gamma ^2,\hfill & \\ C_{\mathrm{VI}}E_{\mathrm{VI}}^2={\displaystyle \frac{4}{3}}(F_{\mathrm{VI}}^24E^2)+K,\hfill & \\ (F_{\mathrm{VI}}2E)^2(F_{\mathrm{VI}}+4E)+{\displaystyle \frac{9K}{4}}(F_{\mathrm{VI}}2E){\displaystyle \frac{27}{4}}\beta \gamma ^2=0,\hfill & \end{array}`$ (78)
and this algebraic dependence could explain the difficulty to separate the variables in the Hamilton-Jacobi equation. Note that, in the particular case $`\beta \gamma =0`$, i.e. $`\kappa _1^2=\kappa _2^2`$, these coefficients become rational, see .
Remark. The F-VI ODE can be written in Hamiltonian form,
$`\{\begin{array}{cc}H=P_2^2+Q_2P_1{\displaystyle \frac{Q_1^4}{6}}+{\displaystyle \frac{3}{2}}Q_1Q_2^2{\displaystyle \frac{13}{1296}}\alpha _{\mathrm{VI}}^3Q_1+{\displaystyle \frac{1}{16}}\alpha _{\mathrm{VI}}^2Q_1^2{\displaystyle \frac{1}{8}}\alpha _{\mathrm{VI}}Q_2^2\hfill & \\ 6\beta _{\mathrm{VI}}Q_16\kappa _{\mathrm{VI}}tQ_1+{\displaystyle \frac{347}{2^93^3}}\alpha _{\mathrm{VI}}^4+{\displaystyle \frac{9}{2}}\alpha _{\mathrm{VI}}\beta _{\mathrm{VI}},\hfill & \\ Q_1=6\left(y{\displaystyle \frac{\alpha _{\mathrm{VI}}}{72}}\right),Q_2=6y^{},P_1=6y^{\prime \prime \prime }108yy^{},P_2=6y^{\prime \prime }.\hfill & \end{array}`$ (82)
In the autonomous case $`\kappa _{\mathrm{VI}}=0`$, the Hamiltonian $`H`$ is a first integral (equal to $`36K_{1,\mathrm{VI}}`$), and the other constant of the motion is cubic in the momenta. However, because of the nonlinear link between $`K_{1,\mathrm{VI}}`$ and the two first integrals of the 1:6:8 case, see (48), there exists no canonical transformation between the variables $`(q_j,p_j)`$ of 1:6:8 and the above canonical variables of a-F-VI.
## 6 General solution of the quartic 1:12:16 case
Let us denote the two constants in involution as,
$`1:12:16\{\begin{array}{cc}H={\displaystyle \frac{1}{2}}(P_1^2+P_2^2)+{\displaystyle \frac{\mathrm{\Omega }}{8}}(4Q_1^2+Q_2^2){\displaystyle \frac{1}{32}}(16Q_1^4+12Q_1^2Q_2^2+Q_2^4)\hfill & \\ {\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\kappa _1^2}{Q_1^2}}+{\displaystyle \frac{4\kappa _2^2}{Q_2^2}}\right)=E,\hfill & \\ K={\displaystyle \frac{1}{16}}\left(8(Q_2P_1Q_1P_2)P_2Q_1Q_2^42Q_1^3Q_2^2+2\mathrm{\Omega }Q_1Q_2^2+32Q_1{\displaystyle \frac{\kappa _2^2}{Q_2^2}}\right)^2\hfill & \\ +\kappa _1^2\left(Q_2^44{\displaystyle \frac{Q_2^2P_2^2}{Q_1^2}}\right).\hfill & \end{array}`$ (87)
Similarly to the 1:6:1-1:6:8 couple, there exists a canonical transformation between the 1:12:16 Hamiltonian and another Hamiltonian , which is however not the sum of a kinetic energy and a potential energy, which we denote similarly as 5:9:4,
$`5:9:4\{\begin{array}{cc}H={\displaystyle \frac{1}{2}}\left(p_1^2+\left(p_2{\displaystyle \frac{3}{2}}q_1q_2\right)^2\right){\displaystyle \frac{1}{8}}(4q_1^4+9q_1^2q_2^2+5q_2^4)+{\displaystyle \frac{\omega }{2}}(q_1^2+q_2^2)\kappa q_1+{\displaystyle \frac{\zeta }{2q_2^2}}=E,\hfill & \\ K={\displaystyle \frac{1}{q_2^2}}\left(2q_2^2p_1+2q_1^2q_2^22q_1q_2p_2q_2^44\kappa q_1\right)^2\hfill & \\ \times \left(2q_2^2p_1+2q_1^2q_2^2+p_2^24q_1q_2p_22q_2^4+\omega q_2^2+4{\displaystyle \frac{\kappa ^2}{q_2^2}}+8\kappa q_14\kappa {\displaystyle \frac{p_2}{q_2}}\right)\hfill & \\ +4(\zeta +4\kappa ^2)((2q_1{\displaystyle \frac{p_2}{q_2}}+4q_1^2+q_2^2+4q_1{\displaystyle \frac{\kappa }{q_2^2}})p_1{\displaystyle \frac{1}{q_2^4}}(q_1^2q_2^2+q_2^4+2\kappa q_1)^2\hfill & \\ +2{\displaystyle \frac{q_1^2}{q_2^2}}(p_2{\displaystyle \frac{3}{2}}q_1q_2)^2+{\displaystyle \frac{(q_1^2+q_2^2)^2}{2}}+q_1^2{\displaystyle \frac{\zeta }{q_2^4}}),\hfill & \\ E_{5:9:4}=E_{1:12:16},K_{5:9:4}=K_{1:12:16},\omega =\mathrm{\Omega },\kappa ={\displaystyle \frac{\kappa _1+\kappa _2}{2}},\zeta =(\kappa _1\kappa _2)^2.\hfill & \end{array}`$ (94)
The path to an integrated ODE is also quite similar and is made of the following three segments .
Firstly, the coordinate $`q_1(t)`$ of (94) is identified to the component $`F`$ of the traveling wave reduction $`f(x,\tau )=F(xc\tau ),g(x,\tau )=G(xc\tau )`$ of a soliton system of two coupled KdV-like equations denoted c-KdVa$`(f,g)`$ ,
$`\{\begin{array}{cc}q_1=F,q_2^2={\displaystyle \frac{2}{5}}\left(F^{}2F^2G+\omega \right),\hfill & \\ c=\omega .\hfill & \end{array}`$ (97)
Secondly, there exists a Bäcklund transformation between this soliton system and another one of the c-KdV type, denoted bi-SK system $`(u,v)`$ , transformation defined by the Miura map
$`\{\begin{array}{cc}u={\displaystyle \frac{3}{10}}\left(3f_xf^2+2g\right),\hfill & \\ v={\displaystyle \frac{9}{10}}\left(f_{xxx}+g_{xx}+f_xgfg_xff_{xx}+g^2\right).\hfill & \end{array}`$ (100)
Finally, the traveling wave reduction $`u(x,\tau )=U(xc\tau ),v(x,\tau )=V(xc\tau )`$ is identified ,
$`\begin{array}{cc}U=3\left(y{\displaystyle \frac{\omega }{30}}\right),V=6y^{\prime \prime }+18y^2{\displaystyle \frac{9}{5}}\omega y+{\displaystyle \frac{1}{10}}\omega ^2{\displaystyle \frac{3}{5}}E,\hfill & \end{array}`$ (102)
to the F-IV equation (or to the F-III as well) in the classification of Cosgrove ,
F-IV $`\{\begin{array}{cc}y^{\prime \prime \prime \prime }=30yy^{\prime \prime }60y^3+\alpha _{\mathrm{IV}}y+\beta _{\mathrm{IV}},\hfill & \\ y={\displaystyle \frac{1}{2}}\left(s_1^{}+s_2^{}+s_1^2+s_1s_2+s_2^2+A\right),\hfill & \\ (s_1s_2)s_1^{}=\sqrt{P(s_1)},(s_2s_1)s_2^{}=\sqrt{P(s_2)},\hfill & \\ P(s)=(s^2+A)^3{\displaystyle \frac{\alpha _{\mathrm{IV}}}{3}}(s^2+A)+Bs+{\displaystyle \frac{\beta _{\mathrm{IV}}}{3}},\hfill & \\ K_{1,\mathrm{IV}}=\left({\displaystyle \frac{3B}{4}}\right)^2,K_{2,\mathrm{IV}}={\displaystyle \frac{9AB^2}{64}},\hfill & \end{array}`$ (108)
in which $`(K_{1,\mathrm{IV}},K_{2,\mathrm{IV}})`$ denote two polynomial first integrals of F-IV. The general solution of this ODE is meromorphic, expressed with genus two hyperelliptic functions .
In order to perform the integration of both Hamiltonians (87) and (94), it is sufficient to express $`(F,G)`$ rationally in terms of $`(U,V,U^{},V^{})`$. The result is
$`\{\begin{array}{cc}F={\displaystyle \frac{W^{}}{2W}}+K_{1,\mathrm{a}}X_2,\hfill & \\ G=F^2X_1X_2+K_{1,\mathrm{a}}{\displaystyle \frac{54U^{}}{X_1}}54K_{1,\mathrm{a}}\left(U+{\displaystyle \frac{3\omega }{20}}\right){\displaystyle \frac{W^{}}{WX_1}}+{\displaystyle \frac{2}{3}}\left(U+{\displaystyle \frac{9\omega }{10}}\right),\hfill & \\ W=X_1^2+108K_{1,\mathrm{a}}^2\left(U+{\displaystyle \frac{3\omega }{20}}\right),\hfill & \\ X_1=V+2U^23\omega U+{\displaystyle \frac{9}{50}}\omega ^2{\displaystyle \frac{27}{5}}E,\hfill & \\ X_2=9(4U_{}^{}{}_{}{}^{2}+{\displaystyle \frac{8}{3}}UV{\displaystyle \frac{8}{25}}\omega U^2+{\displaystyle \frac{2}{5}}\omega V+{\displaystyle \frac{48}{5}}EU\hfill & \\ {\displaystyle \frac{42}{25}}\omega ^2U{\displaystyle \frac{9}{2}}(\kappa _1^2+\kappa _2^2){\displaystyle \frac{9}{2}}K_{1,\mathrm{a}}^2+{\displaystyle \frac{36}{25}}\omega E{\displaystyle \frac{27}{125}}\omega ^3),\hfill & \\ K_{1,\mathrm{a}}=\kappa _1\kappa _2.\hfill & \end{array}`$ (116)
From the point of view of the separation of variables, one should first exhibit a Hamiltonian representation of F-IV. One such structure is that of the cubic SK case. However, since the constant value of the Hamiltonian of the cubic SK case, when expressed only in terms of the parameters $`(E,K,\omega ,\kappa _1,\kappa _2)`$ of the 1:12:16, is not an affine function of $`E`$, there exists no canonical transformation between the cubic SK case and the 1:12:16 case.
## 7 Conclusion, remaining work
The explicit integration of all the seven cases is now achieved in the Painlevé sense (finding a closed form single valued expression for the general solution), and the common features are the following.
1. In all cases, the general solution is hyperelliptic with genus two, and therefore meromorphic.
2. Each case is birationally equivalent to a fourth order ODE which is complete in the Painlevé sense, i.e. which accepts no additional term, under penalty of losing its Painlevé property. Consequently, for each of the seven Hamiltonians, it is impossible to add any term to the Hamiltonian without destroying the Painlevé property, and the seven Hénon-Heiles Hamiltonians are complete.
About the integration in the Arnol’d-Liouville sense (finding the separating variables of the Hamilton-Jacobi equation), two problems remain open.
1. In the 1:6:1-1:6:8 case, the hyperelliptic curve $`y^2=P(s)`$ of F-VI (see (73)) reduces in the separated cases $`\beta \gamma =0`$ to the hyperelliptic curve of the separating variables. Therefore, F-VI is the good ODE to consider, and the only missing item is to find a Hamiltonian structure of F-VI, necessarily distinct from (82), admitting a canonical transformation to 1:6:1-1:6:8.
2. In the 1:12:16-5:9:4 case, the hyperelliptic curve $`y^2=P(s)`$ of F-IV (see (108)) does not reduce in the separated cases $`\kappa _1\kappa _2=0`$ to the hyperelliptic curve of the separating variables, which is
$`\kappa _1\kappa _2=0:P(s)=s^6\omega s^3+2Es^2+{\displaystyle \frac{K}{20}}s+\kappa _1^2+\kappa _2^2=0.`$ (117)
Therefore, F-IV (as well as its birationally equivalent ODE F-III) is not the good ODE to consider, and it should be quite instructive to integrate the fourth order equivalent ODE (19) in that case.
## Acknowledgments
The authors acknowledge the financial support of the Tournesol grant no. T2003.09. CV is a postdoctoral fellow at the FWO-Vlaanderen. |
warning/0507/math0507288.html | ar5iv | text | # References
What is wrong with the Lax-Richtmyer
Fundamental Theorem of Linear
Numerical Analysis ?
Elemer E Rosinger
Department of Mathematics
and Applied Mathematics
University of Pretoria, Pretoria
0002 South Africa
eerosinger@hotmail.com
Abstract
We show that the celebrated 1956 Lax-Richtmyer linear theorem in Numerical Analysis - often called the Fundamental Theorem of Numerical Analysis \- is in fact wrong. Here ”wrong” does not mean that its statement is false mathematically, but that it has a limited practical relevance as it misrepresents what actually goes on in the numerical analysis of partial differential equations. Namely, the assumptions used in that theorem are excessive to the extent of being unrealistic from practical point of view. The two facts which the mentioned theorem gets wrong from practical point of view are :
\- the relationship between the convergence and stability of numerical methods for linear PDEs,
\- the effect of the propagation of round-off errors in such numerical methods.
The mentioned theorem leads to a result for PDEs which is unrealistically better than the well known best possible similar result in the numerical analysis of ODEs. Strangely enough, this fact seems not to be known well enough in the literature. Once one becomes aware of the above, new avenues of both practical and theoretical interest can open up in the numerical analysis of PDEs.
1. Towards a correct relationship between stability and
convergence
It has been shown that in practically relevant situations the converse implication ”convergent $``$ stable” in the Lax-Richtmyer theorem may fail to hold, see Rosinger \[1-8\], Rosinger & van Niekerk , Oberguggenberger & Wang. Thus there need not always be an equivalence between the convergence and stability of a numerical scheme. It may therefore happen that convergence is a weaker property than stability, which means that we may have convergent numerical schemes which nevertheless fail to be stable as well.
In this way, what has become a kind of ”UNIVERSALLY RECITED MANTRA” in the numerical analysis of linear partial differential equations, namely that
( ? ) ”stability and convergence are equivalent”
for linear numerical methods approximating such equations, does in fact lack a valid enough practical reality, and can be replaced with the far more convenient fact, according to which
( ! ) ”convergence need not always imply stability”.
Needless to say, the practical interest in such a possibility is significant, as it can enlarge the class of convergent numerical schemes beyond those which are stable. Examples in this regard are mentioned in Rosinger \[1-8\], Rosinger & van Niekerk , Oberguggenberger & Wang.
By the way, the well known necessary condition for stability, given by von Neumann prior to the Lax-Richtmyer theorem, and which does not require a Banach space setup, can be seen as a further indication of the rather involved relationship between stability and convergence.
A yet more important point about the above mantra is the following. Even if it were true in the linear case - which in fact is the only case addressed by the Lax-Richtmyer theorem - it would still lack relevance in most of the cases when exact solutions of nonlinear PDEs are approximated by respective nonlinear numerical methods. Indeed, as is well known, Kreiss, Stetter, a local linearized stability analysis of nonlinear PDEs and of their nonlinear numerical methods need not in general lead either to necessary, or sufficient convergence conditions.
2. Questions about the implication ”convergent $``$ stable”
Briefly, the Lax-Richtmyer theorem, see below, states the equivalence between the convergence and stability of a linear numerical scheme which is consistent with a well posed linear PDE, see Lax & Richtmyer, Richtmyer, Richtmyer & Morton, as well as a fully detailed analysis and presentation in Rosinger \[2, pp. 1-14\].
The important fact to note is the following.
The proof of the implication ”stable $``$ convergent” is trivial, and certainly, it does in no way require the completeness of space in which it happens.
Therefore, the crux of the Lax-Richtmyer theorem is solely in the proof of the converse implication, namely, ”convergent $``$ stable”.
That converse implication ”convergent $``$ stable”, however, is proved based on the celebrated Principle of Uniform Boundedness of Linear Operators in Banach spaces. And as is well known, see Appendix, that property of uniform boundedness does not necessarily hold in normed spaces which are not complete, thus fail to be Banach spaces.
It is precisely here, with the assumptions which are made in order to secure a Banach space framework, that the Lax-Richtmyer theorem goes twice wrong from practical point of view. Namely, it goes wrong both with respect to the relationship between stability and convergence, as well as regarding the treatment of the essentially nonlinear phenomenon of the propagation of round-off errors.
Furthermore, the proof of the implication ”convergent $``$ stable” is essentially linear, as it makes use of the mentioned linear principle, as well as of a linear concept of stability. This makes the extension of that implication to the fully nonlinear case extremely difficult.
3. Is completeness an appropriate requirement ?
The numerical analysis of a given PDE does typically assume the a priori knowledge of the existence of certain exact solutions of that equation. After all, in case an exact solution does not exist, it is of course nonsensical to try to approximate it numerically. Thus, if and when the existence of an exact solution is known a priori, then the aim of numerical analysis is to construct numerical solutions approximating one or another of such exact solutions. In this way, we are given an exact solution $`U`$ and construct a sequence, say, $`U_{\mathrm{\Delta }t}`$, with $`\mathrm{\Delta }t>0`$, of numerical solutions. Thus the problem is whether or not we have the hoped for convergence property
( * ) $`lim_{\mathrm{\Delta }t0}U_{\mathrm{\Delta }t}=U`$
where the limit holds in some appropriate sense.
Suppose now, as usual, that both the exact solution $`U`$ and the numerical solutions $`U_{\mathrm{\Delta }t}`$ belong to a certain normed space $`(X,||||)`$.
A crucial observation here is the following one. And it is missed by the Lax-Richtmyer theorem.
Clearly, in order to establish whether the above convergence property ( * ) does, or for that matter, does not hold, one does not at all need to assume that the respective normed space $`(X,||||)`$ is complete. Indeed, we have started by assuming that the exact solution $`U`$ exists, thus the hoped for limit value in (\*) exists. Furthermore, the terms $`U_{\mathrm{\Delta }t}`$ of the sequence in (\*) also exist, being the constructed numerical solutions. Finally, the normed space $`X`$ is supposed to be chosen in such a way that both the exact solution $`U`$ and its numerical approximations $`U_{\mathrm{\Delta }t}`$ do belong to it. And then the only problem is whether the constructed numerical solution $`U_{\mathrm{\Delta }t}`$ does indeed happen to converge to the exact solution $`U`$.
Furthermore, often, when for instance the exact solution $`U`$ is a classical solution of the PDE considered, one can choose the normed space $`(X,||||)`$ as constituted from sufficiently smooth functions, since the numerical solutions $`U_{\mathrm{\Delta }t}`$ are typically defined at discrete points, thus they can be extrapolated to functions of required smoothness.
In this way, there does not appear to be any practical reason whatsoever why the normed space $`(X,||||)`$ in which the convergence property (\*) is to be established should be complete.
The alleged reason why nevertheless the completeness of the normed space $`(X,||||)`$ is requested appears to be the claim that it is needed in order to handle the effect of round-off errors as well. Indeed, as it stands, the Lax-Richtmyer theorem is only supposed to deal with the effects of the propagation of truncation errors, since it does not anywhere mention directly round-off errors.
However, as seen in Rosinger , see also Rosinger \[2-4,7\], this claim that the completeness of $`X`$ will give the opportunity to deal as well with the effect of the propagation of round-off errors is simply unrealistic from the point of view of the way round-off errors actually propagate in the computations involved. In particular, this claim leads to the paradox that one obtains a result regarding the effect of round-off errors in the numerical solution of PDEs which is strictly better than the well known best possible corresponding result in the case of the numerical solution of ODEs.
Obviously, in the many usual cases when one approximates classical solutions $`U`$, one can choose the normed space $`(X,||||)`$ made up of sufficiently smooth functions. But then, the completeness requirement in the Lax-Richtmyer theorem obliges one to consider its completion $`(X^\mathrm{\#},||||)`$. And typically, $`X^\mathrm{\#}`$ will be a much larger space, containing a considerable amount of non-smooth functions.
Two important points should be noted here.
First, within this larger and completed space $`X^\mathrm{\#}`$, the original convergence problem ( * ) will remain precisely the same. Indeed, the constructed sequence of numerical solutions $`U_{\mathrm{\Delta }t}X`$ converges to the existing exact solution $`UX`$ in the space $`X`$, if and only if it converges to $`U`$ in the space $`X^\mathrm{\#}`$.
On the other hand, the stability property of the respective numerical method may turn out to lead to a more stringent condition in the larger space $`X^\mathrm{\#}`$, than in the original smaller space $`X`$.
This is indeed of one the issues related to the assumption of completeness, an assumption which is essential in the particular method of proof of the implication ”convergent $``$ stable” in the Lax-Richtmyer theorem.
Otherwise, one simply notes that, in general, the completeness condition does not necessarily belong to the problem of establishing the convergence property (\*).
4. Compactness or Boundedness ?
The above convergence relation ( * ), whenever it holds, clearly implies that the subset
$`\{U_{\mathrm{\Delta }t}|\mathrm{\Delta }t>0\}\{U\}X`$
is compact in $`X`$, regardless of $`X`$ being complete or not. And let us recall that all the elements of this subset are supposed to exist. Indeed, the exact solution $`U`$ of the PDE under consideration exists, otherwise the problem of its numerical approximation would be vacuous. Further, the approximating numerical solutions $`U_{\mathrm{\Delta }t}`$ are effectively constructed by the numerical method employed.
On the other hand, the condition of stability of the numerical methods used in the Lax-Richtmyer theorem, see (5.18) below, is given in terms of boundedness, and as is well known, boundedness does not imply compactness in infinite dimensional normed spaces.
This discrepancy between the association of convergence with compactness, and on the other hand, of stability with boundedness was first pointed out and dealt with in Rosinger , where with an appropriate compactness based definition of stability, a general nonlinear equivalence result was given between convergence and stability.
It should be mentioned here that the above arguments related to stability, convergence, completeness, compactness and boundedness were, back in the early summer of 1979, personally communicated by the author to P D Lax, at a conference at the Tel Aviv University, in Israel.
5. Some details of the Lax-Richtmyer theorem
Let us now, for convenience, recall the Lax-Richtmyer theorem as given in its original formulation, see Lax & Richtmyer, Richtmyer, Richtmyer & Morton, Rosinger \[2, pp. 1-14\]. We consider a linear evolution type PDE
(5.1) $`d/dtU(t)=A(U(t)),t[0,T]`$
with the initial value
(5.2) $`U(0)=u`$
where $`A:DXX`$ is a linear operator defined on the subspace $`D`$ of the Banach space $`X`$, $`uD`$, while $`U:[0,T]D`$ is the sought after solution. Since we deal with an evolution PDE, the operator A is in fact a linear partial differential operator in some space variable $`x^n`$.
Further, one can assume that, when given, linear homogenous boundary conditions have already been incorporated in the definition of $`D`$.
Typically, one can also assume that $`D`$ is dense in $`X`$ and we have satisfied the following exact solution property
(5.3) $`\begin{array}{c}uD:\hfill \\ \\ U:[0,T]X:\hfill \\ \\ )lim_{\mathrm{\Delta }t0}||(U(t+\mathrm{\Delta }t)U(t))/\mathrm{\Delta }tA(U(t))||=0,t[0,T]\hfill \\ \\ )U(0)=u\hfill \end{array}`$
Given now time, respectively space increments $`\mathrm{\Delta }t(0,\mathrm{})`$ and $`\mathrm{\Delta }x(0,\mathrm{})^n`$, we construct a finite difference method
(5.4) $`C_{\mathrm{\Delta }t,\mathrm{\Delta }x}:XX`$
which we assume to be a continuous linear mapping.
The numerical analysis problem we face in the above terms is to characterize the relations
(5.5) $`\mathrm{\Delta }x=\alpha (\mathrm{\Delta }t)`$
where the mapping $`\alpha :(0,\mathrm{})(0,\mathrm{})^n`$ is such that $`lim_{\mathrm{\Delta }t0}\alpha (\mathrm{\Delta }t)=0^n`$, and the convergence property holds
(5.6) $`lim_{\mathrm{\Delta }t0,n\mathrm{},n\mathrm{\Delta }tt}U(t)C_{\mathrm{\Delta }t,\alpha (\mathrm{\Delta }t)}^nu=0`$
uniformly for $`t[0,T]`$, for every $`uD`$, where $`U`$ corresponds to $`u`$ according to (5.3).
As is well known, in general, this is not a trivial problem. In Courant et.al., it was shown for the first time that one cannot in general expect instead of (5.6) the stronger convergence property
(5.7) $`lim_{\mathrm{\Delta }t0,n\mathrm{},n\mathrm{\Delta }tt,\mathrm{\Delta }x0}U(t)C_{\mathrm{\Delta }t,\mathrm{\Delta }x}^nu=0`$
to hold uniformly for $`t[0,T]`$.
Property (5.7), in which $`\mathrm{\Delta }t`$ and $`\mathrm{\Delta }x`$ can simultaneously and independently tend to 0, is called unconditional stability. On the other hand, property (5.6), in which the relation (5.5) ties $`\mathrm{\Delta }x`$ to $`\mathrm{\Delta }t`$ when they both tend to 0, is called conditional stability.
Obviously, in the case of conditional stability, one is interested in numerical methods (5.4) in which $`\alpha `$ tend to 0 as fast as possible, when $`\mathrm{\Delta }t`$ tends to 0. Indeed, in such a situation one can obtain a good space accuracy without increasing too much the computation time.
As a simple and immediate illustration, let us consider the initial value problem for the heat equation
$`U_t=U_{xx},t[0,\mathrm{}),x`$
$`U(0,x)=u(x),x`$
In this case we can take $`(X,||||)=^{\mathrm{}}()`$ and $`A=^2/x^2`$, while
$`D=\{u^{\mathrm{}}()𝒞^2()|Au^{\mathrm{}}()\}`$
and as is well known, John, the exact solution property (5.3) holds.
A simple numerical method for this heat equation is given by
$`(C_{\mathrm{\Delta }t,\mathrm{\Delta }x}u)(x)=u(x)+(u(x+\mathrm{\Delta }x)2u(x)+u(x\mathrm{\Delta }x))\mathrm{\Delta }t/\mathrm{\Delta }x^2`$
for $`uX,x,\mathrm{\Delta }t,\mathrm{\Delta }x>0`$. Also as is well known, Richtmyer, this explicit numerical method will not have the convergence property given by the unconditional stability (5.7), while the weaker convergence property (5.6), called conditional stability, will hold, if and only if
$`2\mathrm{\Delta }t\mathrm{\Delta }x^2`$
which in terms of (5.5) can be written as
$`\alpha (\mathrm{\Delta }t)\sqrt{(}2\mathrm{\Delta }t)`$
This obviously means that $`\alpha 0`$ rather slowly, when $`\mathrm{\Delta }t0`$, which is inconvenient, since a small $`\mathrm{\Delta }x`$ will impose the use of a quadratically smaller $`\mathrm{\Delta }t`$, leading thus to an increased number of time iterations.
Returning to the general case, in view of the exact solution property (5.3), we can define the family of linear mappings
(5.8) $`E_0(t):DX,t[0,T]`$
by
(5.8.1) $`E_0(t)(u)=U(t),t[0,T]`$
The initial value problem (5.1), (5.2) is called properly posed, if and only if the family of linear mappings (5.8) is uniformly bounded, that is, for a certain $`K>0`$, we have
(5.9) $`E_0(t)K,t[0,T]`$
As is well known, since $`D`$ is a dense subspace in $`X`$, one can extend by continuity the family of linear mappings (5.8) to a unique family of linear mappings
(5.10) $`E(t):XX,t[0,T]`$
with the same uniform bound, namely
(5.11) $`E(t)K,t[0,T]`$
In addition, we shall have the semigroup property
(5.12) $`\begin{array}{c}E(t)E(s)=E(t+s),t,s[0,T],t+sT\hfill \\ \\ E(0)=\text{id}_X\hfill \\ \\ lim_{\mathrm{\Delta }t0}E(u)u=0,uX\hfill \end{array}`$
We can note that the semigroup property (5.12) leads to a further extension, this time of the linear mappings (5.10), namely
(5.13) $`E(t):XX,t[0,\mathrm{})`$
where
(5.13.1) $`E(t)=E(t[t/T]T)E(T)^{[t/T]},t[T,\mathrm{})`$
with $`[t/T]`$ denoting the largest integer which is smaller than, or equal to $`t/T`$. In this case, instead of the corresponding above relations, we shall have
(5.14) $`\begin{array}{c}E(t)E(s)=E(t+s),t,s[0,\mathrm{})\hfill \\ \\ E(t)K^{1+[t/T]},t[0,\mathrm{})\hfill \end{array}`$
Returning now to the numerical method (5.4), (5.5), we shall consider as our finite difference scheme the family of continuous linear mappings
(5.15) $`C_{\mathrm{\Delta }t}=C_{\mathrm{\Delta }t,\alpha (\mathrm{\Delta }t)}:XX`$
Here it is important to note that, typically, this family of continuous linear mappings $`C_{\mathrm{\Delta }t}`$, with $`\mathrm{\Delta }t>0`$, is not uniformly bounded for small $`\mathrm{\Delta }t`$.
Now in view of (5.5), (5.6), the finite difference scheme (5.15) is called convergent to the semigroup (5.10) - (5.12) on the time interval $`[0,T]`$, where $`T>0`$ is given, if and only if
(5.16) $`\begin{array}{c}uX,ϵ>0:\hfill \\ \\ \delta >0:\hfill \\ \\ t[0,T],\mathrm{\Delta }t>0,n𝐍,n\mathrm{\Delta }tT:\hfill \\ \\ \mathrm{\Delta }t,|tn\mathrm{\Delta }t|\delta E(t)C_{\mathrm{\Delta }t}^nuϵ\hfill \end{array}`$
Further, the finite difference scheme (5.15) is called consistent with the initial value problem (5.1), (5.2) on the same time interval $`[0,T]`$, if and only if
(5.17) $`\begin{array}{c}uX,ϵ>0:\hfill \\ \\ \theta >0:\hfill \\ \\ t[0,T],\mathrm{\Delta }t>0:\hfill \\ \\ \mathrm{\Delta }t\theta C_{\mathrm{\Delta }t}E(t)E(t+\mathrm{\Delta }t)uϵ\hfill \end{array}`$
Finally, the finite difference scheme (5.15) is called stable on the time interval $`[0,T]`$, if an only if
(5.18) $`\begin{array}{c}L>0:\hfill \\ \\ \mathrm{\Delta }t>0,n𝐍,n\mathrm{\Delta }tT:\hfill \\ \\ C_{\mathrm{\Delta }t}^nL\hfill \end{array}`$
With the above, we have the so called Fundamental Theorem of Linear Numerical Analysis
Theorem ( Lax-Richtmyer, 1956)
Given a properly posed semigroup (5.10) - (5.12) and a finite difference scheme (5.15) which is consistent with it, then the finite difference scheme is convergent to the semigroup, if and only if it is stable.
Remark
The practical interest in the above type of result is in the following. The consistency of a finite difference scheme with a semigroup generated by an initial value problem (5.1), (5.2) is typically easy to establish with the use of a finite Taylor series argument, in case we deal with smooth enough, or classical solutions. Also, what is practically particularly important, the consistency property can be established without the effective knowledge of any specific exact solution of the initial value problem, and only based on the knowledge of the regularity of such solutions, that is, the existence of smooth enough, or classical exact solutions.
The convergence property of such a finite difference scheme is, of course, the main and nontrivial issue, and just like the consistency property, it is a relational property, since it involves the semigroup, or the initial value problems as well. Furthermore, here the fact that, typically, the exact solution is only known to exist, but it is not known effectively - this being the very reason for using numerical analysis - makes it so much more difficult to establish convergence.
On the other hand, the stability property of a finite difference scheme is no longer a relational property, but an intrinsic property which is solely of the finite difference scheme itself, therefore, at least in principle, it can be established alone on the information contained in that finite difference scheme.
In this way, in the study of the convergence of finite difference schemes there is clearly a major interest in establishing a certain connection between the relational property of convergence which is the sought after aim, and on the other hand, the intrinsic property of stability.
The above Lax-Richtmyer theorem does establish such a connection, in fact, an equivalence, between convergence and stability. Unfortunately however, it assumes the completeness of the normed space in which all of this happens, in order to be able to prove the implication ”convergent $``$ stable”.
Appendix
We present a simple counterexample to the celebrated Principle of Uniform Boundedness of Linear Operators in a Banach Space, based on the fact that the respective normed space fails to be complete, that is, Banach. This shows that in this principle, the completeness of the normed space involved is indeed essential.
We take the normed space $`(X,||||)`$ defined as follows
(A.1) $`X=\{x=(x_0,x_1,x_2,...)^{}|\begin{array}{c}m:\hfill \\ n,nm:\hfill \\ x_n=0\hfill \end{array}\}`$
with the norm given by
(A.2) $`x=sup\{|x_n||n\}`$
for $`x=(x_0,x_1,x_2,...)X`$.
Now, for every $`k`$, we define the linear operator $`T_k:XX`$ by
(A.3) $`T_k(x_0,x_1,x_2,...)=(y_0,y_1,y_2,...)`$
where for $`x=(x_0,x_1,x_2,...)X`$, we have
(A.4) $`y_n=\begin{array}{c}kx_k\text{if}n=k\hfill \\ \\ 0\text{if}nk\hfill \end{array}`$
It follows easily that
(A.5) $`T_k=k,k`$
therefore, the family of linear operators $`(T_k|k)`$ is not uniformly bounded.
On the other hand, given any fixed $`x=(x_0,x_1,x_2,...)X`$, there exists $`m`$, such that $`x_n=0`$, for $`n,nm`$. Hence $`T_k(x)=0`$, for $`k,km`$. Consequently
(A.6) $`sup\{T_k(x)|k\}<\mathrm{}`$
In this way, the family of linear operators $`(T_k|k)`$ is bounded at each point $`xX`$, and yet it is not uniformly bounded on $`X`$.
The reason for that is obviously in the fact that the normed space $`(X,||||)`$ is not complete. Indeed, $`(X,||||)`$ is a strict and dense subspace of $`l^{\mathrm{}}`$. |
warning/0507/cond-mat0507054.html | ar5iv | text | # Probing the nanohydrodynamics at liquid-solid interfaces using thermal motion
## Abstract
We report on a new method to characterize nano-hydrodynamic properties at the liquid/solid interface relying solely on the measurement of the thermal motion of confined colloids. Using Fluorescence Correlation Spectroscopy (FCS) to probe the diffusion of the colloidal tracers, this optical technique –equivalent in spirit to the microrheology technique used for bulk properties– is able to achieve nanometric resolution on the slip length measurement. It confirms the no-slip boundary condition on wetting surfaces and shows a partial slip $`b=18\pm 5`$nm on non-wetting ones. Moreover, in the absence of external forcing, we do not find any evidence for large nano-bubble promoted slippage on moderately rough non-wetting surfaces.
Over the recent years the pursuit of scale reduction inherent to nanotechnologies has been extended to the fluidic domain and liquid flow manipulation, with the important development of micro- and nanofluidics Whitesides and Stroock (2001). However, reducing the scale of any system leads invariably to an enhancement of the influence of surface properties with respect to the bulk ones: given the scale reduction, most phenomena take place at the boundaries and a fundamental understanding of how surface properties might affect the overall flow properties has become crucial to design and optimize operational devices Whitesides and Stroock (2001); Stein et al. (2004); Joly et al. (2004).
Classically, one accounts for the influence of these interfaces through effective boundary conditions (BC) in the description of macroscopic hydrodynamics, the most common of those BC being the no-slip assumption (see Ref. Lauga et al. (2005) for an exhaustive review of the litterature). However, the possible deviation from this classical hypothesis, resulting in liquid slippage at the solid surface, has recenlty become a central issue, with immediate perspectives in the micro- and nanofluidics domains Lauga et al. (2005) or in the electrokinetic context Joly et al. (2004). Slippage is usually accounted for by replacing the no-slip BC for the tangential velocity $`v_t`$ by a partial slip BC, in the form $`b_zv_t=v_t`$ (with $`z`$ perpendicular to the planar surface). This generalized BC introduces an extrapolation length $`b`$, usually denoted as the slip length Bocquet and Barrat (1994); Barrat and Bocquet (1999).
Probing the interfacial dynamics of liquids close to solid substrates has accordingly opened new experimental challenges in the recent years: specifically devised experimental tools capable to investigate the nano-hydrodynamics close to the solid substrate have been developed, allowing to address the existence and conditions for slippage. Most recently, two main routes have been followed Lauga et al. (2005) : (i) dissipation methods on one hand on the basis of Surface Force Apparatus -SFA- and Atomic Force Microscopy -AFM- measurements; (ii) flow characterization close to surfaces on the other hand, using optical methods, like FRAP in evanescent waves geometry, or microPIV velocimetry. These various approaches suggest an overall link between the hydrodynamic slippage and the wettability of the solid substrate, in agreement with theoretical and numerical results Barrat and Bocquet (1999); Lauga et al. (2005). They fail however to provide a unified description of the slippage phenomenon, reporting for instance slip lengths $`b`$ on seemingly identical smooth surfaces that span several orders of magnitude, from tens of nanometers to micrometers Lauga et al. (2005). The presence of nanobubbles at the surface, whose existence might be promoted by the driving flow, has been put forward to rationalize these results Vinogradova et al. (1995); Tyrrell and Attard (2001); de Gennes (2002). As demonstrated theoretically Cottin-Bizonne et al. (2004); de Gennes (2002); Lauga and Stone (2003), the existence of such nanobubbles would strongly enhance the measured (apparent) slip length, and could also provide an explanation for shear rate dependent effects reported together with large slip Zhu and Granick (2001). In order to clarify the experimental picture, new non-intrusive approaches to probe the hydrodynamic surface properties are needed.
In this letter, we demonstrate a completely different experimental route, allowing to explore the nano-hydrodynamics of liquids close to surfaces in the absence of any external forcing, thereby avoiding problems inherent, in all previous measurements, to the imposed flow. Rather than measuring the interfacial dissipation or the forced surface flow, we take advantage of the information already included in response to thermal fluctuations Einstein (1905) to extract the interfacial dynamics. This technique, which is analogous –for surfaces– to the passive microrheology technique for bulk characterization Mason and Weitz (1995), proves to be extremely sensitive as it allows us to reach an unprecedented resolution for an optical technique, namely a few nanometers on the slip length measurement. Together with confirming the emerging picture for smooth surfaces of a moderate slippage ($`b=18`$nm) present only in non-wetting situation, it provides a truly non-intrusive technique to explore the possible role of gas pockets (nanobubbles) on slip properties. As a first step, it shows that in the absence of forcing, nano-bubble promoted giant slip is absent on moderately rough non-wetting surfaces.
We first describe the general principle of our approach. The diffusion dynamics of colloidal tracers is measured in a confined geometry between two solid surfaces of interest, using a home-built Fluorescence Correlation Spectroscopy (FCS) device (fig. 1). Tracer dynamics is affected by confinement and this dependence reflects the hydrodynamic boundary conditions which apply on both solid substrates Almeras et al. (2000); Saugey et al. (2005); Lauga and Squires (2005). Results for different nature of the solid substrates (in particular varying wettability and surface roughness) lead to measurable differences in the diffusion coefficient, allowing us to deduce the corresponding surface slippage.
Let us now enter into the details of the experimental setup. Colloidal tracers (polystyrene \[Molecular Probes\] or silica \[Kisker\] fluorescent beads with typical diameters $`2a200`$nm and concentration $`c=1`$bead/$`\mu `$m<sup>3</sup>) in aqueous solution ($`10^5`$M NaOH; $`\mathrm{8\hspace{0.17em}10}^3`$M KCl) are confined between two solid silica surfaces made from a BK7 spherical lens in contact with a Pyrex plane (fig. 1). The thermal diffusion dynamics of the colloids is measured in this confined geometry with a FCS device (fig. 1). With this technique, beads fluorescence is excited with an argon laser (Lasos) focused with a long working distance microscope objective (Leica x40, NA 0.8). Fluorescence is then collected through the same objective and sent to a detector (APD, Perkin-Elmer) connected to a correlator (Correlator.com<sup>©</sup>) via a dichroic mirror and a bandpass filter (Omega). While in bulk measurements, a confocal pinhole is inserted in the detection pathway to get a spatially defined measurement volume $`v`$, here the axial limits are set in pratice by the two confining walls so that $`v=(\pi w^2)e`$, with $`w`$ the beam waist radius and $`e`$ the wall to wall distance. With a bead concentration such that the mean number of tracers in the detection volume is low (typically around 1 here) fluctuations of the collected intensity $`I(t)`$ arise due to the motion of beads entering or leaving the measurement volume. The characteristic time scale for such fluctuations corresponds to the residence time $`\tau _D=w^2/D`$ of a bead within volume $`v`$, where $`D`$ is the bead self-diffusion coefficient. More quantitatively, considering a gaussian radial intensity distribution for the illuminating laser beam, the fluorescence intensity auto-correlation function reads
$$g(\tau )=\frac{I(t)I(t+\tau )}{I(t)^2}=1+\frac{1}{n}\frac{1}{(1+4\tau /\tau _D)},$$
(1)
from which the experimental average number of beads $`n`$ and their residence time can be extracted, as shown in fig. 1. As the bead diffusion coeficient $`D`$ depends on the location $`z`$ within the gap, we eventually obtain an averaged value over the all accessible positions within the thickness $`e`$. To ensure good statistics, values for each confinement $`e/2a`$ were accumulated from over 100000 events (associated with one bead passage) splitted between many different auto-correlation runs and from at least 5 different experimental cells. Runs belonging to different cells may slighlty differ in absolute transit time and average number of beads due to small day to day variations in temperature, concentration or focus location. We accordingly normalized data from each cell by the reference point for (almost) unconfined beads, located at $`e/2a=16.8`$, for which excellent statistics was already achieved in the time frame of a single cell experiment.
For any location of the sphere-plane geometry where we measured $`n`$ and $`\tau _D`$, the wall to wall distance $`e`$ was simultaneously measured using the interference pattern generated by the two confining surfaces (Newton rings). Due to the chosen bead size ($`2a200`$nm), we were able to restrict our data point collection to the dark fringes (setting $`\mathrm{\Delta }e=\lambda /2n_w=183`$nm; with $`\lambda `$ the laser wavelength and $`n_w`$ the optical index of water), therefore optimizing the signal to noise ratio. Owing to the large radii of curvature of the lenses used (from $`250`$ to $`500`$mm), $`e`$ is constant to better than $`0.4`$% over the measurement spot size $`2w1\mu `$m.
The first element we focused on is the evolution of the mean bead number with the confinement defined as $`e/2a`$. This evolution is shown in figure 2 (inset) where we recover a linear behavior for the averaged number of beads in the volume $`v`$ as expected: $`n=c(\pi w^2)e`$. This measurement provides an important check that no depletion or adsorption due to bead-surface interaction is detectable in our system <sup>1</sup><sup>1</sup>1Adsorption was only found when polystyren beads were used together with hydrophobic surface: this couple was accordingly avoided in our measurements.. Note however that the linear evolution does not extrapolate to $`e/2a=1`$ as expected for the sole excluded volume effect but to a slightly higher value ($`1.2`$ in the inset of fig. 2). This is associated with the additionnal “excluded volume” resulting from electrostatic repulsions between the wall and the beads. For silica walls, under the present conditions (electrolytes concentration), the excluded region is $`D_{\mathrm{excl}.}20`$nm thick (for both polystyren and silica beads), while it is found to be thinner $`D_{\mathrm{excl}.}13`$nm for silica beads close to silanized walls. Such a distance agrees with the estimate obtained from the balance between the thermal energy and the electrostatic wall-bead repulsion energy (calculated using the experimental Debye length and typical zeta potentials for silica).
We now come to the measurement of the residence time –or reciprocally to the diffusion coefficient– of the beads as a function of the confinement $`e/2a`$. When the confining walls are hydrophilic, we expect a usual no-slip BC to apply at their surface Barrat and Bocquet (1999); Cottin-Bizonne et al. (2005). In such a situation, the bead mobility should be strongly reduced by the walls proximity Faxen (1924) as was already verified experimentally by a few groups Faucheux and Libchaber (1994); Lin et al. (2000). This wetting configuration was therefore the starting point from which we have explored the influence of surface properties on the hydrodynamics. Our results obtained in such situation (aqueous solution with silica walls) are presented in figure 2 together with theoretical predictions assuming no-slip boundary condition on both walls. For low confinements $`e/2a1`$, a very sound, approximate solution for the mobility can be constructed on the basis of the Faxen approximate solution for one wall Faxen (1924), by adding the independent contribution of each wall to obtain the theoretical bead mobility Saugey et al. (2005). However to avoid theoretical approximations, we conducted numerical resolution of the Stokes equation for a (non-slipping) sphere moving parallel to the two confining solid walls (each characterized by a \[possibly different\] slip length $`b`$). A finite element method was implemented using Femlab<sup>©</sup> (see in Saugey et al. (2005) for details). This theoretical prediction for the mobility was then averaged over accessible $`z`$ (from $`a+D_{\mathrm{excl}.}`$ to $`eaD_{\mathrm{excl}.}`$). As is evidenced in figure 2, the agreement between experiments and theory with no-slip BC on the walls is excellent up to the strongest confinements. It provides another evidence supporting the fact that no-slip applies on our smooth wettable substrates Cottin-Bizonne et al. (2005); Lauga et al. (2005). In addition, the fact that beads with very different chemistry display the very same behavior is another indication that our system is free from specific surface-bead interactions, beside the measured electrostatic repulsion leading to $`D_{\mathrm{excl}.}`$.
We now consider the influence of surface properties on the hydrodynamic BC. For that purpose the silica plane was covalently coated with hydrocarbon chains using octadecyltrichlorosilan (OTS). The resulting residence time of silica beads is plotted for different confinement against the previous results for wetting (plain silica) plane. When beads are confined enough ($`e/2a<5`$), residence times appear to be systematically shorter close to hydrophobic surfaces than to hydrophilic ones. This effect is more easily captured when normalizing this evolution by the theoretical behavior in the absence of slip: $`D_{\mathrm{exp}.}(e/2a)/D_{\mathrm{no}\mathrm{slip}}(e/2a)`$. A failure of the no-slip boundary condition should result in a departure of this ratio from 1. This is indeed what is observed in the inset of figure 3 where the systematic trend described above is best evidenced. Moreover it is possible to render this departure accurately by introducing a finite slip length $`b`$ in the theoretical calculations Saugey et al. (2005); Lauga and Squires (2005). The fitted behavior agrees remarkably well with the experimental data providing a slip length on smooth OTS coated silica planes of $`b=18\pm 5`$nm. This value is in perfect agreement with published results on the exact same surfaces using a dynamic-SFA Cottin-Bizonne et al. (2005). It demonstrates the ability of our aproach to characterize the nanohydrodynamics close to interfaces at zero shear stress together with its unreached sensitivity for an optical method.
This sensitivity was confirmed in a subsequent set of experiments, investigating the role of surface roughness on hydrodynamical properties of interfaces. Our method, relying only on thermal fluctuations, is particularly adapted for such situations as it avoids possible surface modifications (creation and/or coupling with nanobubbles) present in all previously reported methods where external forcing is required. Before the hydrophobic OTS coating, the silica planes were first treated $`30`$min. with a piranha solution (1 vol. H<sub>2</sub>O<sub>2</sub>, 2 vol. H<sub>2</sub>SO<sub>4</sub>) leading to rough topography Eske and Galipeau (1999) with a peak-to-peak height (checked by AFM on $`5\times 5\mu `$m<sup>2</sup> scans) around $`40`$nm, as compared to below $`1`$nm for untreated surfaces. Even on non-wetting surfaces, roughness kills the slippage, leading to reduced beads mobility as compared with smooth non-wetting surfaces. Actually the mobility is found to be below the one for a no-slip BC. This results comes from the fact that the hydrodynamic boundary is shifted toward the top of the roughness while the interference gap measurements provides only the average wall position. Therefore the measured gap $`e`$ (against which data are plotted) is larger than the hydrodynamic gap thus resulting in a translation of the experimental data that might be interpreted as a negative slip length of $`40\pm 20`$nm <sup>2</sup><sup>2</sup>2Roughness inhomogeneity over the scale of the probed region is responsible for the increased uncertainty.. Again this demonstrates the sensitivity of our technique, showing differences for surface properties affecting the nanorheology on scales as low as tens of nanometers. An interesting conclusion deduced from our data is that in the absence of external forcing, only the “negative” effect of roughness on slippage is evidenced. No huge enhancement (reaching micrometric slip length Tretheway and Meinhart (2002)) possibly promoted by trapped gaz is observed.
We have therefore shown that the measure of the thermal motion of colloidal tracers provides an alternative and very sensitive method –achieving nanometric resolution– to address the nano-hydrodynamics of simple liquids close to surfaces. This method is able to give information on interfacial dissipation without any external forcing, by exploiting the intimate fluctuation–dissipation link. Working in the strict zero-shear rate limit, avoids shear induced alteration of the surfaces (such as possible nucleation of nanobubbles). In the present experiments we find no-slip on wetting Pyrex surface, and a finite slip length of $`18`$ nm for water on hydrophobic OTS surfaces. These results are in full agreement with recent findings using a dynamic-SFA Cottin-Bizonne et al. (2005) and with numerical results Barrat and Bocquet (1999). Moreover a nanometric roughness has been shown to kill slippage even on a hydrophobic surface. This absence of large slip efffect at zero shear rate would confirm that the – still debated – presence of nanobubbles might originate in the flow itself, as was already conjectured de Gennes (2002); Lauga and Stone (2003); Zhu and Granick (2002). To answer this question, it would accordingly be interesting to use our equilibrium approach on previously sheared surfaces. Alternatively other types of surfaces might be considered, in particular super-hydrophobic surfaces for which very large slippage is expected Cottin-Bizonne et al. (2004); Ou et al. (2004); Journet et al. (2005). Work along these lines is in progress.
###### Acknowledgements.
We thank E. Charlaix, C. Cottin-Bizonne and J.-L. Barrat for stimulating discussions. |
warning/0507/cond-mat0507245.html | ar5iv | text | # Breakdown of weak-field magnetotransport at a metallic quantum critical point
## Abstract
We show how the collapse of an energy scale in a quantum critical metal can lead to physics beyond the weak-field limit usually used to compute transport quantities. For a density-wave transition we show that the presence of a finite magnetic field at the critical point leads to discontinuities in the transport coefficients as temperature tends to zero. The origin of these discontinuities lies in the breakdown of the weak field Jones-Zener expansion which has previously been used to argue that magneto-transport coefficients are continuous at simple quantum critical points. The presence of potential scattering and magnetic breakdown rounds the discontinuities over a window determined by $`\tau \mathrm{\Delta }<1`$ where $`\mathrm{\Delta }`$ is the order parameter and $`\tau `$ is the quasiparticle elastic lifetime.
Quantum critical points in metallic systems have generated considerable theoretical and experimental interest in recent years coleman\_2005a . They are realized by tuning the finite temperature critical point of a phase transition in a metal to zero temperature. A variety of tuning parameters have been used including hydrostatic pressure julian\_1996a , chemical composition vonlohneysen\_1994a and quantum criticality can occur serendipitously gegenwart\_1999a . In insulating quantum magnets, magnetic fields often provide a ”handle” by which they may be made quantum critical bitko\_1996a and there has recently been much interest in doing the same in the metallic case. Notable examples include the metamagnetic quantum critical endpoint seen in $`\mathrm{Sr}_3\mathrm{Ru}_2\mathrm{O}_7`$ grigera\_2001a and also the antiferromagnetic quantum critical point in $`\mathrm{YbRh}_2\mathrm{Si}_2`$ gegenwart\_2002a . The collapse of the characteristic energy scale for fluctuations can induce deviations from Landau Fermi liquid theory hertz\_1976a ; millis\_1993a seen, for example, in transport quantities. These are usually computed by assuming linear response to driving fields, yet the vanishing energy scale could invalidate that assumption green\_2005b .
In this paper we explicitly demonstrate such a breakdown by considering magneto-transport in the simplest class of quantum critical metal tuned by magnetic field. The presence of the small but finite magnetic field at the transition point leads to modified transport response. This has assumed added significance recently because of the suggestion that the discrepancies between theory and experiment in quantum critical matter coleman\_2005a may originate in part from a transition at the critical point between localized and de-localized spins and with a concomitant change in the Fermi surface volume coleman\_2001a . It is suggested that the Hall coefficient would reveal this volume change and recently evidence paschen\_2004a for such a change has been presented at the field driven quantum critical point in $`\mathrm{YbRh}_2\mathrm{Si}_2`$. Here, we consider the $`T0`$ limit of a spin or charge-density wave transition in finite magnetic field, so the quantum critical fluctuations play a small role in quasi-particle scattering when compared with elastic scattering from impurities. This limit has been considered by Bazaliy et al. norman\_2003a ; bazaliy\_2004a recently for non-field driven quantum critical points under the assumption that a weak-field expansion in magnetic field can be made. They conclude that there are no anomalies in the Hall conductivity and that, in the absence of perfect nesting, changes in other transport quantities are generally linear in the energy gap so change continuously at the critical point.
The essential result of our work is that at a density wave (DW) quantum critical point the weak-field (Jones-Zener) expansion—upon which previous magnetotransport work was based—breaks down because of the vanishing DW gap $`\mathrm{\Delta }`$. This breakdown leads, in the relaxation time approximation, to discontinuities in the magnetoconductivities: jumps in the resistivity and Hall coefficients at the transition, despite the fact that the phase transition itself is continuous. For magnetic fields, $`B>\mathrm{\Delta }/(ev_F^2\tau )`$ we show that the magnetoresistance is non-analytic in the field ($`|B|`$)—a result not obtainable from a weak-field expansion. While our results are valid at any zero temperature DW transition, they assume particular significance at a field driven quantum critical point where $`\mathrm{\Delta }`$ is tuned to zero at finite magnetic field. The interdependence of the gap and magnetic field define a trajectory through Fig. 1(a) \[shown schematically in Figs. 1(b) or (c)\] which inevitably explores the non-analytic region. We also go beyond the relaxation time approximation to consider the effects of disorder and magnetic breakdown very close to the critical point. While these smooth the discontinuity they do not alter the magnitude of the changes in conductivity across the transition at a finite magnetic field.
We consider the following highly simplified model of transport near a density-wave transition (DW). We take the Fermi surface in the paramagnetic state to be circular (and assume two dimensionality for ease of computation). We then consider density wave formation to induce a periodic potential on the otherwise free electrons and assume this potential to be proportional to the mean-field order parameter: $`V_{\pm Q}(B)=\mathrm{\Delta }_0\mathrm{}\sqrt{(B_cB)/B_c}`$. Here, the periodic potential (assumed real so both signs of $`Q`$ are present) has Fourier components at the ordering wave-vector $`Q`$ of the DW and its perturbation on the free electrons describes Bragg scattering off the DW. Degenerate perturbation theory in the presence of this potential will gap the dispersion whenever a resonance condition is met: $`ϵ_\stackrel{}{k}=ϵ_{\stackrel{}{k}\pm n\stackrel{}{Q}}`$. Since transport in the $`T0`$ limit that we will be considering is dominated by the Fermi surface, we need only consider the most important Bragg scattering matrix element. For example, near the point $`ϵ_k=ϵ_{kQ}`$ the dispersion is modified to be
$$E_k=\frac{1}{2}\left(ϵ_k+ϵ_{kQ}\right)\pm \sqrt{\frac{1}{4}\left(ϵ_kϵ_{kQ}\right)+V_Q^2}.$$
(1)
When the chemical potential falls in one of the gaps, as it will if $`2k_F>Q`$, then the Fermi surface is modified by the density wave and it is the change in transport properties induced by this modification that we wish to study.
In Fig. 2 we illustrate the change in Fermi surface topology that we envisage. There are a number of possible scenarios
* $`2k_F<Q`$: no gaps appear at the Fermi surface and transport coefficients will change smoothly through the critical point. We therefore do not consider this case here.
* $`2k_F>Q`$, $`Q`$ incommensurate: gaps appear on the Fermi surface and lead to an open section of Fermi surface \[Fig. 2a\]. In the absence of twinning, transport properties would discriminate between currents parallel and perpendicular to $`\stackrel{}{Q}`$.
* $`2k_F>Q`$, $`Q`$ commensurate with the lattice (e.g. an antiferromagnet). Transport is isotropic through the transition and the Fermi surface remains closed. We treat this case by considering two real DWs with $`\stackrel{}{Q}=(\pm Q_x,\pm Q_y)`$ and $`\stackrel{}{Q}=(\pm Q_x,Q_y)`$ (see Fig. 2b).
* $`2k_F=Q`$ (nested) norman\_2003a . Since the perfect nesting not be met very close to the critical point bazaliy\_2004a this case will ultimately reduce to one of the cases above and is therefore not considered here.
For the initial analysis we use the classical Boltzmann equation in the relaxation time approximation $`\tau `$: we are envisaging the $`T=0`$ limit of the conductivities and ignore inelastic processes. Rather than solve the Boltzmann equation order by order in the magnetic field ziman\_1960a we solve to all orders in the field directly abrikosov\_1988a using the Chamber’s formula chambers\_1969a :
$$\sigma _{ij}=\frac{e^2}{4\pi ^3}\frac{dS}{\mathrm{}|\stackrel{}{v}|}_0^{\mathrm{}}v_i(0)v_j(t)e^{t/\tau }𝑑t.$$
(2)
For each area element of the Fermi surface, $`dS`$, we integrate the velocity, $`\stackrel{}{v}(t)`$ measured along a semiclassical quasiparticle orbit. These orbits are defined by the Lorentz equation of motion
$$\mathrm{}\frac{d\stackrel{}{k}}{dt}=e\stackrel{}{v}\times \stackrel{}{B},$$
(3)
where $`\stackrel{}{v}=\stackrel{}{}_kϵ(\stackrel{}{k})/\mathrm{}`$. In this paper we will always assume the magnetic field is perpendicular to the 2D electron fluid.
Our main numerical results are illustrated in Fig. 3. The circles illustrate the conductivities obtained by integrating the Chambers formula when an DW gap with $`Q`$ parallel to the $`y`$ axis, is suppressed to zero at a critical field corresponding to $`\omega _{}\tau =0.2`$ with $`\mathrm{\Delta }_0/ϵ_F=0.01`$. Here $`\omega _{}=eB/m`$, the precession frequency around the whole circular Fermi surface in the absence of a DW and the transition is occurring in the region where $`\omega _{}\tau 1`$i.e. naïvely the weak field regime. We clearly see discontinuities in the conductivities $`\sigma _{xx}`$ and $`\sigma _{xy}`$ (and a very small one in $`\sigma _{yy}`$) at the critical field. In the case of closed Fermi surfaces (not shown) a clear discontinuity occurs in all three conductivities. Having observed that the conductivities can be discontinuous at a DW transition in a magnetic field we now consider the underlying physics.
Within Boltzmann transport theory the effect of a magnetic field on the conductivity is caused by the precession of the out-of-equilibrium distribution around the Fermi surface according to Eq. 3 between scattering events. Quasi-particles are deflected by the local Hall angle which is related to the local radius of curvature of the Fermi surface ($`R`$) and the Lorentz equation of motion:
$$\delta k=R\theta _\mathrm{H}=ev_FB\tau \theta _\mathrm{H}=\frac{ev_FB}{R}.$$
(4)
This deflects the current, leading to a reduction of the longitudinal conductivity $`\mathrm{\Delta }\sigma _{xx}/\sigma _{xx}(0)\mathrm{cos}\theta _H1\theta _H^2`$ giving the usual magnetoconductance quadratic in applied field. The Hall angle must be small for weak-field response. However at a DW transition a gap opens in the Fermi surface leading to a local radius of curvature of the Fermi surface of order $`\mathrm{\Delta }/v_F`$ which vanishes at the quantum critical point itself \[see Fig. 2(c)\]. The condition for being in the weak-field regime is therefore that $`B_{\mathrm{weak}\mathrm{field}}<\mathrm{\Delta }/(ev_F^2\tau )`$ or equivalently $`\omega _{}\tau <\tau \mathrm{\Delta }/(k_Fl)`$ where $`l`$ is the mean free path. This upper bound vanishes at the critical point so magnetotransport there is never weak-field. Thus the results of Bazaliy et al. are valid only provided the limit $`B/\mathrm{\Delta }0`$ can be taken which is certainly not the case at a field-driven DW quantum critical point.
Instead magneto-transport is dominated by a fraction of the quasiparticles (proportional to $`ev_FB\tau `$) that are deflected around the cusp in the Fermi surface between scattering events. This leads to a magnetoconductance proportional to $`|B|`$i.e. non-analytic in the field and thus beyond any weak-field expansion. This physics is beautifully illustrated by the magnetoresistance of a square Fermi surface calculated by Pippard pippard\_1989a .
We can use Pippard’s method to compute the size of the discontinuities in all three independent components of the conductivity tensor. In the limit that $`\mathrm{\Delta }0^+`$ we can solve the Boltzmann transport equation on each segment of the Fermi surface and then combine the solutions by ensuring that the magnitude of the distribution function is continuous. The result is that
$`\sigma _{xx}`$ $`=`$ $`{\displaystyle \frac{ne^2\tau }{m}}\left[{\displaystyle \frac{1}{1+(\omega _{}\tau )^2}}{\displaystyle \frac{Q^2}{k_F^2}}{\displaystyle \frac{\omega _{}\tau P(\alpha ,\omega _{}\tau )}{\left[1+(\omega _{}\tau )^2\right]^2}}\right],`$ (5)
$`\sigma _{xy}`$ $`=`$ $`\omega _{}\tau \sigma _{xx},`$ (6)
$`\sigma _{yy}`$ $`=`$ $`{\displaystyle \frac{ne^2\tau }{m}}\left[{\displaystyle \frac{1}{1+(\omega _{}\tau )^2}}+{\displaystyle \frac{Q^2}{k_F^2}}{\displaystyle \frac{(\omega _{}\tau )^3P(\alpha ,\omega _{}\tau )}{\left[1+(\omega _{}\tau )^2\right]^2}}\right],`$ (7)
where $`\alpha =\mathrm{cos}^1(Q/2k_F)`$ and
$$P(\alpha ,\omega _{}\tau )=\frac{\mathrm{cosh}\frac{\pi }{2\omega _{}\tau }}{\pi \mathrm{cosh}\frac{\alpha }{\omega _{}\tau }\mathrm{sinh}\frac{\pi /2\alpha }{\omega _{}\tau }}.$$
(8)
On the paramagnetic side ($`\mathrm{\Delta }0^{}`$) the conductivities are as above but without the terms proportional to $`P(\alpha ,\omega _{}\tau )`$. Thus the all three magnetoconductivities, $`\sigma _{xx}`$, $`\sigma _{xy}`$, $`\sigma yy`$, show discontinuities with magnitude of order, $`\omega _{}\tau `$, $`(\omega _{}\tau )^2`$ and $`(\omega _{}\tau )^3`$ respectively so the most dramatic jump is the magnetoresistance footnote. These analytic solutions are shown as the dashed lines in Fig. 3.
We now extend our calculation beyond the relaxation time approximation. If $`\tau \mathrm{\Delta }<1`$ this calculation cannot be valid since the quasiparticles would be unable to notice the Bragg scattering from the DW above the scattering from impurities. One would expect this to washout the discontinuities in magnetoconductivities since quasiparticles will tend to remain on original Fermi surface. Magnetic (Zener) breakdown when $`B>2\mathrm{\Delta }^2/(ev_F^2)`$ abrikosov\_1988a as the same effect. (See Green and Sondhi green\_2005b for $`E`$ field breakdown.) We include this phenomenologically in our calculation in the following fashion. Rather than hybridize the dispersion (as in Eq. 1) we maintain a circular Fermi surface and treat the DW potential in the collision integral as a resonant scatterer between points of the Fermi surface that satisfy the resonance condition. This gives the following transport equation for the out of equilibrium distribution function $`g(ϵ,\theta )`$
$`e\stackrel{}{v}\stackrel{}{E}\tau _ϵf_0+\omega _{}\tau {\displaystyle \frac{g}{\theta }}=g(\tau \mathrm{\Delta })^2[\delta (\mathrm{cos}\theta +\eta )`$
$`+\delta (\mathrm{cos}\theta \eta )\left]\right[g(ϵ,\theta )g(ϵ,\pi \theta )],`$ (9)
where $`f_0`$ is the Fermi function and $`\eta =Q/2k`$. Since the solution is periodic around the Fermi surface it may be solved by Fourier transform and, for a circular Fermi surface gives the following expression for the conductivities
$`\sigma _{xx}`$ $`=`$ $`{\displaystyle \frac{ne^2\tau }{m}}\left[{\displaystyle \frac{1}{1+(\omega _{}\tau )^2}}{\displaystyle \frac{K(\alpha ,\omega _{}\tau )}{(1+(\omega _{}\tau )^2)^2}}\right],`$ (10)
$`\sigma _{xy}`$ $`=`$ $`\omega _{}\tau \sigma _{xx},`$ (11)
$`\sigma _{yy}`$ $`=`$ $`{\displaystyle \frac{ne^2\tau }{m}}\left[{\displaystyle \frac{1}{1+(\omega _{}\tau )^2}}+{\displaystyle \frac{(\omega _{}\tau )^2K(\alpha ,\omega _{}\tau )}{(1+(\omega _{}\tau )^2)^2}}\right],`$ (12)
where
$$K(\alpha ,\omega _{}\tau )=\frac{2A\omega _{}\tau P(\alpha ,\omega _{}\tau )\mathrm{cos}^2\alpha }{\omega _{}\tau P(\alpha ,\omega _{}\tau )+A},$$
(13)
and $`A=\frac{4(\tau \mathrm{\Delta })^2}{\pi \mathrm{sin}\alpha }`$. These conductivities are shown as solid lines on Figs. 3. Note how these expressions interpolate between the paramagnetic solution when $`\tau \mathrm{\Delta }=0`$ (no DW scattering), and the Pippard result in the limit $`\tau \mathrm{\Delta }1`$ where the angle dependent scattering effectively mimics the reconstructed the Fermi surface. Thus we see that disorder washes out the discontinuity over a region in field determined by $`\tau \mathrm{\Delta }<1`$.
We have also considered the case of closed Fermi surfaces. In that case $`\sigma _{xx}=\sigma _{yy}`$ and both show a discontinuity of fractional order $`\omega _{}\tau `$. The other difference from the case of open Fermi surfaces is that, as for all open Fermi surfaces, $`\sigma _{yy}`$ remains finite in the high field limit: $`\omega _{}\tau 1`$. This is the regime where Landau level quantization of the closed Fermi surfaces would also become important and is not considered here.
In summary, we have shown that at a simple density wave (DW) quantum critical point the weak-field regime of magnetotransport collapses to zero field with the size of the gap. At finite field in a clean metal one would expect to see discontinuities in the magnetoresistance of order of the magnitude of the Hall angle: $`\omega _{}\tau `$. This effect will be significant at a field driven quantum critical point where by definition the field is finite at the transition point. The case of $`\mathrm{YbRh}_2\mathrm{Si}_2`$ already shows features in the Hall effect which suggest that it falls outside the class of DW quantum critical points. A prediction from this work is that non-field driven DW quantum critical points should should a low field cross-over to a transverse magnetoresistance that is linear in the applied field. It would be interesting to look for such an effect in the $`Cr_{1x}V_x`$ system under pressure yeh\_2002a where we estimate that $`\tau \mathrm{\Delta }=1`$ at $`(xx_c)/x_c0.1`$ and $`\omega _{}\tau (k_Fl)^1`$ at $`B1`$T. Very recently we have learnt that $`\mathrm{Ca}_3\mathrm{Ru}_2\mathrm{O}_7`$ shows exactly the linear magnetoresistance we predict and is argued to be a small gap density wave system kikugawa\_2005a .
We acknowledge useful discussions with A. Rosch, C. Hooley, Q. Si, Y. Bazaliy and R. Ramazashvili. JF acknowledges the EPSRC for financial support and AJS acknowledges the Royal Society and Leverhulme Trust for financial support and the hospitality of the KITP, Santa Barbara. This research was supported in part by the National Science Foundation under Grant No. PHY99-0794. |
warning/0507/hep-th0507040.html | ar5iv | text | # D-Brane Propagation in Two-Dimensional Black Hole Geometries
## 1 Introduction and Summary
An important yet unsolved problem in string theory is an ab initio formulation of dynamics in time-dependent closed string background. Cosmological or black hole backgrounds are outstanding situations of this sort. In these situations, because there is no globally definable timelike Killing vector, there always arise Bogoliubov excitation of the vacuum, leading to cosmological particle production and Hawking radiation. At present, however, full-fledged string theoretic description of the these processes is unavailable. Open string counterpart turned out more manageable. Decay of unstable D-brane, described by open string tachyon rolling , is amenable to exact conformal field theory approach, thus studied extensively in recent years. Still, some issues are left out unsettled, especially, ambiguity in prescribing timelike conformal field theories and the open string vacuum thereof.
One may hope to find the situation better for the well-known two-dimensional black hole, which admits exact conformal field theory description , and to understand physics inside the horizon and the spacelike singularity. For such purposes, we have learned through a variety of other situations that D-brane serves as a better probe than closed string. Our goal is to investigate dynamics of D0-brane propagating in the two-dimensional black hole geometries.<sup>1</sup><sup>1</sup>1Some years ago, studied extensively closed string propagation in these backgrounds. In this paper, we shall take the first step toward the goal: understanding D0-brane dynamics in the causal region outside the horizon, with particular focus on large curvature regime, where the string worldsheet effects become strong. The D0-brane serves as a local probe of the black hole geometries and, in a certain sense, may be considered as an analog of ZZ-brane in the Liouville theory.
The two-dimensional black hole is often considered as a toy model of more realistic higher-dimensional black holes. This is not necessarily so, since the background is intimately related to the black NS5-branes, whose background is given by
$`\text{d}s^2=\left(1{\displaystyle \frac{r_0^2}{r^2}}\right)\text{d}t^2+\left(1+{\displaystyle \frac{k\alpha ^{}}{r^2}}\right)\left({\displaystyle \frac{\text{d}r^2}{1\frac{r_0^2}{r^2}}}+r^2\text{d}\mathrm{\Omega }_3^2\right)+\text{d}𝐲_^5^2,e^{2\mathrm{\Phi }(r)}=g_s^2\left(1+{\displaystyle \frac{k\alpha ^{}}{r^2}}\right)`$ (1)
along with $`k`$-units of NS-NS $`H_3`$-flux penetrating through $`𝕊^3`$. Thus, $`k`$ refers to the number of NS5-branes at $`r=0`$, $`r=r_0`$ is the location of the event horizon, $`𝐲`$ are the spatial coordinates of the planar NS5-brane worldvolume, and $`g_s`$ is the string coupling constant at infinity. Since the NS5-brane is black, it breaks space-time supersymmetries completely and Hawking radiates.
One hopes to gain intuition by studying classical D0-brane dynamics near the horizon. So, consider taking various near-horizon limits of the black NS5-brane (1). One type of near-horizon limit is $`r_00`$ and $`g_s0`$ independently, leading to the ‘throat geometry’ of extremal NS5-branes :
$`\text{d}s^2=\text{d}t^2+k\alpha ^{}\text{d}\rho ^2+k\alpha ^{}\text{d}\mathrm{\Omega }_3^2+\text{d}𝐲_^5^2,\mathrm{\Phi }=\rho +\text{constant},`$ (2)
where $`r=\sqrt{k\alpha ^{}}\mathrm{exp}\rho `$. This background is describable by the exact conformal field theory involving linear dilaton and $`SU(2)_k`$ super Wess-Zumino-Witten (WZW) model:<sup>2</sup><sup>2</sup>2Here, $`k`$ is the level of total current of super SU(2) WZW models and $`\sqrt{\frac{2}{k}}`$ is the amount of background charge for linear dilaton system.
$`\left[_t\times _{\rho ,\sqrt{\frac{2}{k}}}\times SU(2)_k\right]_{}\times \left[^5\right]_{||}.`$
The first part describes the five-dimensional curved spacetime transverse to the NS5-brane, while the second part describes the flat spatial directions parallel to NS5-brane. The criticality condition is satisfied for any k because
$`\left(1+{\displaystyle \frac{6}{k}}+{\displaystyle \frac{1}{2}}\right)+3\times \left({\displaystyle \frac{k2}{k}}+{\displaystyle \frac{1}{2}}\right)+6\times \left(1+{\displaystyle \frac{1}{2}}\right)=15`$ (3)
D-brane dynamics in this background was studied in via the Dirac-Born-Infeld (DBI) approach, and observed that it strikingly resembles the rolling tachyon of unstable D-brane in ambient flat spacetime . This map, which we refer as ‘radion-tachyon correspondence’, offers a useful guide for understanding D-brane propagating in curved spacetime geometry in terms of known results regarding rolling tachyon dynamics.<sup>3</sup><sup>3</sup>3Subsequent works along the same line include e.g. . Introducing ‘tachyon’ variable $`X=\rho `$, DBI Lagrangian of the D0-brane propagating in the background (2) is recastable to that of rolling tachyon:
$`L_{\mathrm{D0}}=e^\mathrm{\Phi }\sqrt{\left({\displaystyle \frac{\text{d}s}{\text{d}t}}\right)^2}=V(X)\sqrt{1\dot{X}^2}\text{where}V(X)=M_0e^X.`$ (4)
The energy is conserved, so solving $`V(X)/\sqrt{1\dot{X}^2}=1`$, we obtain D0-brane’s geodesic<sup>4</sup><sup>4</sup>4Throughout this work, we refer by ‘geodesic’ of D0-brane the classical trajectory as determined by the DBI action. Recall that equivalence principle does not hold in string theory due to the dilaton coupling — motion of D-brane is different from that of fundamental string, the latter being studied for example in (See also footnote 10.). We would like to thank Itzhak Bars for bringing our attention to these references. as
$`e^X={\displaystyle \frac{e^{\rho _o}}{\mathrm{cosh}(tt_o)}}\text{viz.}e^\rho \mathrm{cosh}(tt_o)=e^{\rho _o}.`$ (5)
This is the Lorentzian counterpart of so-called ‘hairpin’ profile. In the previous works , we constructed exact boundary states of the D-brane and analyzed rolling dynamics of the D-brane in detail.<sup>5</sup><sup>5</sup>5Related analysis via boundary conformal theory was given in . A technical difficulty was that the dilaton blows up at the core, hampering further analysis by the strong coupling singularity.
Another type of near-horizon limit is $`r_00`$ and $`g_s0`$ while keeping the energy density above the extremal configuration $`\mu r_0^2/g_s^2\alpha ^{}`$ fixed. It yields ‘throat geometry’ of the near-extremal NS5-branes (1) :
$`\text{d}s^2=\mathrm{tanh}^2\rho \text{d}t^2+k\alpha ^{}\text{d}\rho ^2+k\alpha ^{}\text{d}\mathrm{\Omega }_3^2+\text{d}𝐲_^5^2,e^{2\mathrm{\Phi }}={\displaystyle \frac{k}{\mu \mathrm{cosh}^2\rho }},`$ (6)
where $`r=r_0\mathrm{cosh}\rho `$. For $`(t,\rho )`$-subspace, the metric and the dilaton coincide with those of the two-dimensional black hole. This Lorentzian black hole is describable by Kazama-Suzuki supercoset conformal field theory $`SL(2;)_k/U(1)`$ (where $`U(1)`$ subgroup is chosen to be the non-compact component (space-like direction)) of central charge $`c=3(1+2/k)`$. Likewise, taking account of the NS-NS $`H_3`$-flux penetrating through $`𝕊^3`$ which is omitted in (6), the angular part $`𝕊^3`$ is describable by the (super) $`SU(2)`$-WZW model. In this way, the string background of the nonextremal NS5-brane is reduced to a solvable superconformal field theory system:<sup>6</sup><sup>6</sup>6Here again, $`k`$ is the level of total current of super WZW models. Namely, $`k+2`$, $`k2`$ are the levels of bosonic $`SL(2)`$ and $`SU(2)`$ currents.
$`\left[{\displaystyle \frac{SL(2;)_k}{U(1)}}\times SU(2)_k\right]_{}\times \left[^5\right]_{||}.`$ (7)
Here, the first part describes the five-dimensional curved spacetime (including the time direction) transverse to the NS5-brane, while the second part describes the flat spatial directions parallel to the NS5-brane. The criticality condition is satisfied for any $`k`$ as in (3). Upon Wick rotation, Euclidean black hole has the ‘cigar geometry’, described by the coset conformal field theory $`SL(2;)_k/U(1)_+^3/`$ (where $`U(1)`$ subgroup is the compact component). Asymptotically, circumference of the cigar geometry is $`2\pi \sqrt{\alpha ^{}k}`$, and is identified with inverse of the Hawking temperature.
Again, by introducing ‘tachyon’ variable $`Y\mathrm{log}\mathrm{sinh}\rho `$, DBI Lagrangian of the D0-brane is castable to that of rolling tachyon:
$`L_{\mathrm{D0}}`$ $`=`$ $`e^\mathrm{\Phi }\sqrt{\left({\displaystyle \frac{\text{d}s}{\text{d}t}}\right)^2}=V(Y)\sqrt{1\dot{Y}^2}\text{where}V(Y)=M_0e^Y.`$ (8)
We now find D0-brane’s geodesic as
$`e^Y={\displaystyle \frac{\mathrm{sinh}\rho _o}{\mathrm{cosh}(tt_o)}}\text{viz.}\mathrm{sinh}\rho \mathrm{cosh}(tt_o)=\mathrm{sinh}\rho _o.`$ (9)
Euclidean counterpart of the geodesic (9) describes D1-brane profile in the Euclidean two-dimensional black hole background. An important point is that, in sharp contrast to the extremal background (2), the dilaton is finite everywhere. Thus, the strong coupling singularity is now capped off by the horizon.
In both cases, it is elementary to understand classical dynamics of the D0-brane: both by gravity and by strong string coupling gradient, D0-brane is pulled in and finds its minimum energy and rest mass at the location of the NS5-brane. This also fits to the observation that the spacetime supersymmetry is completely broken as the NS5-brane and the D0-brane preserve different combinations of the supercharges. Eventually, D0-brane would melt into fluxes of NS5-brane worldvolume gauge field and form a non-threshold bound-state. As the D0-brane is pulled in, acceleration would grow and radiate off the binding energy into closed string modes.
Interestingly, the DBI Lagrangian of D0-brane is identical for both extremal and non-extremal NS5-brane backgrounds, see (4) and (8). Consequently, the geodesics are also identical in $`X`$ or $`Y`$ coordinate. Does this imply that the two-dimensional black hole is featureless and not quite black? As we shall explain in much detail, this coincidence turns out simply an artifact of DBI analysis and disappears in full-fledged boundary conformal field theory description for the D0-brane boundary states. It also hints that the radion-tachyon correspondence (based on DBI analysis alone) would break down in non-extremal NS5-brane background, and we will see such indications from exact boundary conformal field theory analysis.
By construction, D-brane boundary states is obtainable from the coupling of the D-brane to closed string modes. Semiclassically, this can be done by overlapping the geodesic (9) to the mini-superspace wave function of the closed string modes. Thus, in section 2, after recapitulating the mini-superspace wave function in the Euclidean black hole background, we shall construct the boundary state of the Euclidean D1-brane by taking the inner product. Being in Euclidean background, the construction is free from ambiguity.
We are primarily interested in D0-brane’s boundary state in Lorentzian background. Formally, the boundary state is obtainable by Wick rotation of that for the Euclidean D1-brane, as constructed in section 2. However, the Wick rotation is not unique and appropriate analytic continuation has to be specified. Causal region of the Lorentzian black hole background has four boundaries: past and future horizons $`^\pm `$, and past and future asymptotic infinities $`^\pm `$. See Figure 1. The specification amounts to imposing boundary conditions at these four null infinities. In section 3, with careful treatment of the boundary conditions, we construct several boundary states of the D0-brane, corresponding to physically motivated boundary conditions: (1) D0-brane emitted from past horizon, (2) D0-brane absorbed to future horizon, and (3) time-symmetric D0-brane. As a useful mnemonic, via the radion-tachyon correspondence (extended to finite temperature background), (1) and (2) are the counterpart of half S-brane, while (3) is the counterpart of full S-brane.
The D0-brane propagating in curved spacetime emits closed string excitations. One expects that in general radiation rate would be proportional to spacetime curvature. In the eternal black hole background, where the future-directed asymptotic infinities consist of $`^+`$ and $`^+`$, the radiation emitted by the D0-brane propagates out to either of them. In section 4, we estimate radiation distribution for incoming part $`𝒩_{\mathrm{in}}(M)`$ (as measured at $`^+`$) and outgoing part $`𝒩_{\mathrm{out}}(M)`$(as measured at $`^+`$) for a fixed transverse mass $`M`$. Intuitively, the radiation distribution ought to depend sensitively on the background curvature scale,<sup>7</sup><sup>7</sup>7Throughout this work, we shall adopt the convention $`\alpha ^{}=2`$. set by the level $`k`$. Curiously, we find that the radiation rate depends on the Hawking temperature! This has to do with the novel feature that the Hagedorn and the Hawking temperatures, $`T_{\mathrm{Hg}}`$ and $`T_{\mathrm{Hw}}`$, in the two-dimensional black hole background are both set by $`k`$:
$`T_{\mathrm{Hg}}={\displaystyle \frac{1}{\beta _{\mathrm{Hg}}}}={\displaystyle \frac{1}{4\pi \sqrt{1\frac{1}{2k}}}}\text{and}T_{\mathrm{Hw}}={\displaystyle \frac{1}{\beta _{\mathrm{Hw}}}}={\displaystyle \frac{1}{2\pi \sqrt{2k}}}.`$
Especially, the Hawking temperature is independent of the non-extremality $`r_0`$ or $`\mu `$.
Notice that two temperature scales meet at $`k=1`$, viz. the conifold geometry, so we anticipate some sort of cross-over or phase transition. For $`k>1`$, D0-brane energy is emitted almost all to incoming mode toward $`^+`$ and small to outgoing mode toward $`^+`$:
$`{\displaystyle \frac{𝒩_{\mathrm{out}}(M)}{𝒩_{\mathrm{in}}(M)}}{\displaystyle \frac{e^{\frac{1}{2}\beta _{\mathrm{Hw}}M}}{e^{\frac{1}{2}\beta _{\mathrm{Hg}}M}}}1.`$
For $`k<1`$, only small portion of the D0-brane energy is emitted, equally distributed between the incoming and the outgoing parts:
$`{\displaystyle \frac{𝒩_{\mathrm{out}}(M)}{𝒩_{\mathrm{in}}(M)}}1.`$
The distribution does not appear to be the radiation pattern one would expect in black hole background, where the black hole would absorb most of the radiation (as well as D0-brane). This brings out a question: if $`k<1`$ background were not associated with black hole, what would it be?<sup>8</sup><sup>8</sup>8Similar question has been discussed recently in .
In section 5, we argue that $`k=1`$ is a critical point of black hole \- string phase transition, where the branching ratio between incoming and outgoing part of the radiation provides an ‘order parameter’. Our interpretation fits nicely with recent observation of , where the same kind of phase transition was pointed out for $`AdS_3`$ and linear dilaton backgrounds. Given that the NS5-brane background is holographically dual to Little String Theory (LST), it is natural to ask what the holographic dual of the D0-brane falling into the NS5-brane. Adopting the proposal that the dual process is identifiable with decay of a defect or soliton in the LST and taking reasonable assumption concerning $`k`$-dependence of LST’s decay number distribution and density of states, we show that the defect decay rate as computed within LST fits perfectly with the D0-brane radiation rate as computed in the bulk (viz. black hole background). Interestingly, the two-dimensional black hole at $`k=1/2`$ belongs not quite to the black hole phase but to the far extreme of the string phase. Long string condensation in two-dimensional string theory was recently studied via non-singlet matrix model . It would be interesting to study D0-brane dynamics in this matrix model and compare with our results.
In section 6, we extend the boundary state analysis to Ramond-Ramond (R-R) sector by spectral flow, and, in section 7, we also analyze the limit the NS5-brane becomes extremal, thus making contact with our earlier results. In section 8, we return to the issue of boundary conditions and propose yet another physically motivated one: the Hartle-Hawking boundary condition. This boundary condition is particularly compelling because a puzzle concerning origin/fate of the conserved R-R charge of the D0-brane gets around, and also fits well with the radiation rate computed in section 4.
## 2 D1-brane on Euclidean Two-Dimensional Black Hole
In this section, we study boundary state description of the D1-brane profile on the Euclidean two-dimensional black hole geometries. We shall begin with recapitulating aspects of the closed string spectrum relevant for foregoing analysis.
### 2.1 Mini-superspace analysis of closed strings
Consider the Euclidean two-dimensional black hole background, known as ‘cigar geometry’:
$$\text{d}s^2G_{ij}\text{d}x^i\text{d}x^j=2k(\text{d}\rho ^2+\mathrm{tanh}^2\rho \text{d}\theta ^2)\text{and}e^\mathrm{\Phi }=\frac{e^{\mathrm{\Phi }_0}}{\mathrm{cosh}\rho }.$$
(10)
Recall that $`k`$ sets characteristic curvature radius in unit of the string scale and hence string worldsheet effects, while $`e^{\mathrm{\Phi }_0}`$ sets the maximum value of the string coupling at the tip $`\rho =0`$ of the cigar geometry. We shall assume the limit $`k1`$ and $`e^{\mathrm{\Phi }_0}1`$: this limit suppresses both string worldsheet and spacetime quantum effects and facilitates to truncate closed string spectrum to zero-modes, viz. to mini-superspace approximation.
In the mini-superspace approach, difference between bosonic strings (with no worldsheet supersymmetry) and fermionic strings (with $`𝒩=2`$ worldsheet supersymmetry) becomes unimportant. The closed string Hamiltonian $`L_0+\overline{L}_0`$ is reduced in the mini-superspace approximation to the target space Laplacian $`\mathrm{\Delta }_0`$, where:
$`\mathrm{\Delta }_0`$ $`=`$ $`{\displaystyle \frac{1}{e^{2\mathrm{\Phi }}\sqrt{G}}}_i\left(e^{2\mathrm{\Phi }}\sqrt{G}G^{ij}_j\right){\displaystyle \frac{1}{2k}}[_\rho ^2+2\mathrm{coth}2\rho _\rho +\mathrm{coth}^2\rho _\theta ^2].`$ (11)
The Hamiltonian is defined with respect to the volume element:
$`\text{d}\text{Vol}=e^{2\mathrm{\Phi }}\sqrt{G}\text{d}\rho \text{d}\theta :=2k\mathrm{sinh}\rho \mathrm{cosh}\rho \text{d}\rho \text{d}\theta k\mathrm{sinh}2\rho \text{d}\rho \text{d}\theta ,`$ (12)
inherited from the Haar measure on the $`SL(2;)`$ group manifold. In the volume element, the dilaton factor $`e^{2\mathrm{\Phi }}`$ is taken into account, as the inner product for closed string states is defined by the worldsheet two-point correlators on the sphere. The normalized eigenfunctions are obtained straightforwardly . They are:
$$\varphi _n^j(\rho ,\theta )=\frac{\mathrm{\Gamma }^2(j+\frac{|n|}{2})}{\mathrm{\Gamma }(|n|+1)\mathrm{\Gamma }(2j1)}e^{in\theta }\left[\mathrm{sinh}^{|n|}\rho F(j+1+\frac{|n|}{2},j+\frac{|n|}{2};|n|+1;\mathrm{sinh}^2\rho )\right],$$
(13)
where $`F(\alpha ,\beta ;\gamma ;z)`$ is the Gaussian hypergeometric function. These eigenfunctions correspond to the primary state vertex operators of conformal weights
$`h=\stackrel{~}{h}={\displaystyle \frac{j(j+1)}{k2}}+{\displaystyle \frac{n^2}{4k}}\text{or}h=\stackrel{~}{h}={\displaystyle \frac{j(j+1)}{k}}+{\displaystyle \frac{n^2}{4k}}`$ (14)
for bosonic<sup>9</sup><sup>9</sup>9The eigenvalue is actually proportional to $`\frac{j(j+1)}{k}+\frac{n^2}{4k}`$. and fermionic strings, respectively. We shall focus on the continuous series, parametrise the radial quantum number $`j`$ as $`j=\frac{1}{2}+i\frac{p}{2}`$ $`(p)`$, and label the eigenfunctions as $`\varphi _n^p(\rho ,\theta )`$ instead of $`\varphi _n^j(\rho ,\theta )`$. Adopt the convention that, in the asymptotic region $`\rho \mathrm{}`$, the vertex operators with $`p>0`$ corresponds to the incoming waves and those with $`p<0`$ corresponds to the outgoing waves. The eigenfunctions (13) are then normalized as
$`(\varphi _n^p,\varphi _n^{}^p^{})=\delta _{n,n^{}}\left[2\pi \delta (pp^{})+_0(p^{},n)\mathrm{\hspace{0.17em}2}\pi \delta (p+p^{})\right],`$ (15)
where the inner product is defined with respect to the volume element (12). Here, $`_0(p,n)`$ refers to the reflection amplitude of the mini-superspace analysis:
$`_0(p,n)={\displaystyle \frac{\mathrm{\Gamma }(+ip)\mathrm{\Gamma }^2(\frac{1}{2}\frac{ip}{2}+\frac{n}{2})}{\mathrm{\Gamma }(ip)\mathrm{\Gamma }^2(\frac{1}{2}+\frac{ip}{2}+\frac{n}{2})}}.`$ (16)
That is, from the definition (13), the reflection amplitude is seen to obey the mini-superspace reflection relation:
$`\varphi _n^p(\rho ,\theta )=_0(p,|n|)\varphi _n^{+p}(\rho ,\theta ).`$ (17)
We shall refer $`_0(p,n)`$ as ‘mini-superspace’ reflection amplitude, valid strictly within mini-superspace approximation at $`k\mathrm{}`$, and anticipate string worldsheet effects at finite $`k`$. Notice that no winding states wrapping around $`\theta `$-direction are present since by definition the mini-superspace approximation retains states with zero winding only.
Utilizing the analytic continuation formula of the hypergeometric functions:
$`F(\alpha ,\beta ;\gamma ;z)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(\gamma )\mathrm{\Gamma }(\beta \alpha )}{\mathrm{\Gamma }(\beta )\mathrm{\Gamma }(\gamma \alpha )}}(z)^\alpha F(\alpha ,\alpha +1\gamma ;\alpha +1\beta ;1/z)`$ (18)
$`+`$ $`{\displaystyle \frac{\mathrm{\Gamma }(\gamma )\mathrm{\Gamma }(\alpha \beta )}{\mathrm{\Gamma }(\alpha )\mathrm{\Gamma }(\gamma \beta )}}(z)^\beta F(\beta ,\beta +1\gamma ;\beta +1\alpha ;1/z),`$ (19)
the eigenfunction (13) is decomposable into
$`\varphi _n^p(\rho ,\theta )=\varphi _{L,n}^p(\rho ,\theta )+_0(p,|n|)\varphi _{R,n}^p(\rho ,\theta ),`$ (20)
where
$`\varphi _{L,n}^p(\rho ,\theta )`$ $``$ $`e^{in\theta }(\mathrm{sinh}\rho )^{1ip}F({\displaystyle \frac{1}{2}}+{\displaystyle \frac{ip+n}{2}},{\displaystyle \frac{1}{2}}+{\displaystyle \frac{ipn}{2}};1+ip;{\displaystyle \frac{1}{\mathrm{sinh}^2\rho }}),`$ (21)
$``$ $`e^\rho e^{ip\rho +in\theta }\text{at}\rho +\mathrm{}`$
and
$`\varphi _{R,n}^p(\rho ,\theta )`$ $``$ $`e^{in\theta }(\mathrm{sinh}\rho )^{1+ip}F({\displaystyle \frac{1}{2}}{\displaystyle \frac{ip+n}{2}},{\displaystyle \frac{1}{2}}{\displaystyle \frac{ipn}{2}};1ip;{\displaystyle \frac{1}{\mathrm{sinh}^2\rho }})`$ (22)
$``$ $`e^\rho e^{ip\rho +in\theta }\text{at}\rho +\mathrm{}`$
refer to the left- and the right-movers, respectively, at $`\rho +\mathrm{}`$, and $`_0(p,|n|)`$ is defined in (16). Obviously, they are related to each other under the reflection of radial momentum: $`\varphi _{R,n}^{+p}=\varphi _{L,n}^p`$, which is also evident from (20) and (16). These mini-superspace wave functions (20) constitute the starting point of constructing boundary states of D-brane in the Euclidean two-dimensional black hole background.
We close the mini-superspace analysis with remarks concerning Wick rotation of the results to the Lorentzian background and string worldsheet effects present at finite $`k`$.
1. The decomposition of $`\varphi _n^p`$ into $`\varphi _{L,n}^p`$ and $`\varphi _{R,n}^p`$ is not globally definable over the entire cigar geometry. They are ill-defined around the tip $`\rho =0`$, and the reflection relation (17) implies that $`\varphi _n^p`$ is not independent of $`\varphi _n^{+p}`$. Therefore, of the continuous series, only the eigenfunctions $`\varphi _n^p`$ with $`p>0,n\text{Z}`$ span the physical Hilbert space of the closed strings on the Euclidean two-dimensional black hole. On the other hand, the situation will become further complicated once Wick rotated to the Lorentzian two-dimensional black hole.
2. Notice that $`\varphi _n^p`$ is not analytic with respect to the angular quantum number $`n`$ as it depends on its absolute value, $`|n|`$. This leads to the ambiguity for Wick rotation from Euclidean to Lorentzian background, under which roughly speaking $`in`$ is replaced by energy $`\omega `$. As for the mini-superspace reflection amplitude $`_0(p,n)`$, since $`_0(p,n)=_0(p,n)`$ holds for all $`n\text{Z}`$, it is unnecessary to take absolute value $`|n|`$ in (17), (20). When taking Wick rotation, we will start from the expression $`_0(p,|n|)`$. In other words, we analytically continue $`_0(p,n)`$ if $`n>0`$ and $`_0(p,n)`$ if $`n<0`$.
3. It is evident that $`|_0(p,n)|=1`$, viz, the mini-superspace reflection amplitude is purely a phase shift in the Euclidean black hole background. It is of utmost importance that, in the Lorentzian black hole background, $`n`$ is analytically continued to pure imaginary value, and the modulus of the reflection amplitude becomes less than unity.
4. For the fermionic Euclidean $`SL(2;)/U(1)`$ conformal field theory, exact result for the reflection amplitude (i.e. taking account of all string worldsheet effects) is known . In our notations, it is
$`(j,m,\stackrel{~}{m})=\nu (k)^{2j1}{\displaystyle \frac{\mathrm{\Gamma }(1+\frac{2j+1}{k})}{\mathrm{\Gamma }(1\frac{2j+1}{k})}}{\displaystyle \frac{\mathrm{\Gamma }(2j+1)\mathrm{\Gamma }(j+m)\mathrm{\Gamma }(j\stackrel{~}{m})}{\mathrm{\Gamma }(2j1)\mathrm{\Gamma }(j+1+m)\mathrm{\Gamma }(j+1\stackrel{~}{m})}},`$ (23)
where
$`\nu (k){\displaystyle \frac{1}{\pi }}{\displaystyle \frac{\mathrm{\Gamma }(1\frac{1}{k})}{\mathrm{\Gamma }(1+\frac{1}{k})}},m={\displaystyle \frac{kw+n}{2}},\stackrel{~}{m}={\displaystyle \frac{kwn}{2}}.`$
Denoting by $`\mathrm{\Phi }_{j;m,\stackrel{~}{m}}`$ the vertex operator with conformal weights $`h=\frac{m^2j(j+1)}{k}`$, $`\stackrel{~}{h}=\frac{\stackrel{~}{m}^2j(j+1)}{k}`$, the exact reflection relation reads
$`\mathrm{\Phi }_{(j+1);m,\stackrel{~}{m}}=((j+1),m,\stackrel{~}{m})\mathrm{\Phi }_{j;m,\stackrel{~}{m}},`$ (25)
The mini-superspace reflection amplitude $`_0(p,n)`$ is then related to the exact one $`(j,m,\stackrel{~}{m})`$ by taking the $`k\mathrm{}`$ limit as mentioned above (up to overall constant):
$`_0(p,n)=\underset{k+\mathrm{}}{lim}(j={\displaystyle \frac{1}{2}}+{\displaystyle \frac{ip}{2}},m={\displaystyle \frac{n}{2}},\stackrel{~}{m}={\displaystyle \frac{n}{2}}).`$ (26)
### 2.2 Boundary state of Euclidean D1-brane
We shall now study D1-brane in the Euclidean two-dimensional black hole background. Classically, profile of the D1-brane follows the geodesic curve
$`\mathrm{cos}(\theta \theta _0)\mathrm{sinh}\rho =\mathrm{sinh}\rho _0,`$ (27)
and is known as the ‘hairpin brane’. Here, $`\theta _0`$, $`\rho _0`$ are free parameters characterizing the geodesic curves. The ‘hairpin brane’ is obtainable as a descendant of the Euclidean $`AdS_2`$-brane in the Euclidean $`AdS_3`$ space, described by $`SL(2;)`$ ($`_+^3`$) Wess-Zumino-Witten model. Correspondingly, exact boundary state of the D1-brane was constructed in from the boundary conformal field theory analysis of the $`SL(2;)`$ Wess-Zumino-Witten model . See also the closely related works e.g. , and for a review. For the case of the Euclidean $`SL(2;)/U(1)`$ supercoset conformal field theory, the relevant boundary state of the NS-NS sector is given by
$`{}_{\mathrm{D1}}{}^{}B;\rho _0,\theta _0|`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\text{d}p}{2\pi }}{\displaystyle \underset{n}{}}\mathrm{\Psi }_{D1}(\rho _0,\theta _0;p,n)p,n|,`$
$`\mathrm{\Psi }_{\mathrm{D1}}(\rho _0,\theta _0;p,n)`$ $`=`$ $`𝒩(k){\displaystyle \frac{\mathrm{\Gamma }(ip)\mathrm{\Gamma }\left(1+\frac{ip}{k}\right)}{\mathrm{\Gamma }\left(\frac{1}{2}+\frac{ip+n}{2}\right)\mathrm{\Gamma }\left(\frac{1}{2}+\frac{ipn}{2}\right)}}e^{in\theta _0}\left[e^{ip\rho _0}+(1)^ne^{ip\rho _0}\right].`$ (28)
Here, $`p,n|`$ refers to the Ishibashi state constructed over the primary state whose mini-superspace wave function is given by $`\varphi _n^p(\rho ,\theta )`$. Also, $`𝒩(k)`$ is a normalization factor. Since it would not affect foregoing analysis, we will set it to $`2\pi `$ for simplicity. One can readily check that the boundary wave function (28) is consistent with the exact reflection amplitude (23).
The result (28) can be understood intuitively as follows. A D-brane boundary wave function is the weighted sum of the wave function of closed string states restricted to the location of the D-brane. In the mini-superspace approximation, as is implicit in , the weighted sum equals to the overlap between the mini-superspace wave function and the delta function constraint enforcing $`(\rho ,\theta )`$ coordinates over the hairpin trajectory (27) (with respect to the volume element (12)). The result is
$`{\displaystyle _0^{\mathrm{}}}\mathrm{sinh}\rho \text{d}\mathrm{sinh}\rho {\displaystyle _{\frac{\pi }{2}+\theta _0}^{\frac{\pi }{2}+\theta _0}}\text{d}\theta \delta \left(\mathrm{cos}(\theta \theta _0)\mathrm{sinh}\rho \mathrm{sinh}\rho _0\right)\varphi _n^p(\rho ,\theta )={\displaystyle _{\frac{\pi }{2}}^{\frac{\pi }{2}}}\text{d}\theta ^{}{\displaystyle \frac{\mathrm{sinh}\rho _0}{\mathrm{cos}^2\theta ^{}}}\varphi _n^p(\widehat{\rho }(\rho _0,\theta ^{}),\theta ^{})e^{in\theta _0},`$
where $`\theta ^{}=(\theta \theta _0)`$ and $`\widehat{\rho }(\rho _0,\theta ^{})`$ refers to the solution of $`\mathrm{cos}\theta ^{}\mathrm{sinh}\rho =\mathrm{sinh}\rho _0`$. Using the decomposition (20), we are then to evaluate integrals:
$`{\displaystyle _{\frac{\pi }{2}}^{\frac{\pi }{2}}}\text{d}\theta {\displaystyle \frac{\mathrm{sinh}\rho _0}{\mathrm{cos}^2\theta }}\varphi _{L,n}^p(\widehat{\rho }(\rho _0,\theta ),\theta )={\displaystyle \frac{2\pi \mathrm{\Gamma }(ip)}{\mathrm{\Gamma }\left(\frac{1}{2}+\frac{ip+n}{2}\right)\mathrm{\Gamma }\left(\frac{1}{2}+\frac{ipn}{2}\right)}}e^{ip\rho _0},`$
$`{\displaystyle _{\frac{\pi }{2}}^{\frac{\pi }{2}}}\text{d}\theta {\displaystyle \frac{\mathrm{sinh}\rho _0}{\mathrm{cos}^2\theta }}\varphi _{R,n}^p(\widehat{\rho }(\rho _0,\theta ),\theta )={\displaystyle \frac{2\pi \mathrm{\Gamma }(ip)}{\mathrm{\Gamma }\left(\frac{1}{2}\frac{ip+n}{2}\right)\mathrm{\Gamma }\left(\frac{1}{2}\frac{ipn}{2}\right)}}e^{+ip\rho _0}.`$ (29)
Details of the computation are relegated in Appendix B. Using the mini-superspace reflection amplitude (16), we then obtain
$`\mathrm{\Psi }_{\mathrm{D1}}^{(0)}(\rho _0,\theta _0;p,n)={\displaystyle \frac{2\pi \mathrm{\Gamma }(ip)}{\mathrm{\Gamma }\left(\frac{1}{2}+\frac{ip+n}{2}\right)\mathrm{\Gamma }\left(\frac{1}{2}+\frac{ipn}{2}\right)}}e^{in\theta _0}\left(e^{ip\rho _0}+(1)^ne^{+ip\rho _0}\right).`$ (30)
We see the result (30) reproduces the exact result (28) modulo the factor $`\mathrm{\Gamma }\left(1+i\frac{p}{k}\right)`$. Importantly, this missing factor depends on $`k`$ (measured in string unit) and hence corresponds precisely to the corrections due to string worldsheet effects. The mini-superspace approximation sets $`k\mathrm{}`$, so this factor is consistently dropped out. Equivalently, this missing factor can be reinstated to the D-brane boundary wave function by demanding consistency of the wave function in the mini-superspace approximation with the exact reflection amplitude (23).
## 3 D0-Brane in Lorentzian Two-Dimensional Black Hole
### 3.1 Analytic continuation of boundary states
In this section, we shall construct the exact boundary state describing the D0-brane moving in the Lorentzian two-dimensional black hole background. Recall that the Lorentzian two-dimensional black hole (‘Lorentzian cigar’) background is obtainable by the Wick rotation $`\theta =it`$ of the Euclidean one (10)
$$\text{d}s^2=2k(\text{d}\rho ^2\mathrm{tanh}^2\rho \text{d}t^2)\text{and}e^\mathrm{\Phi }=\frac{e^{\mathrm{\Phi }_0}}{\mathrm{cosh}\rho }.$$
(31)
Wick-rotating the geodesic of the Euclidean D1-brane, we found the geodesic of the Lorentzian D0-brane in (9) as<sup>10</sup><sup>10</sup>10Another familiar parametrization of the two-dimensional black hole is the analogue of the Kruskal coordinates
$$u=\mathrm{sinh}\rho e^t,v=\mathrm{sinh}\rho e^t,\text{d}s^2=2k\frac{\text{d}u\text{d}v}{1uv},$$
and the geodesic (32) is just a straight line in these coordinates. This is also pointed out in .
$`\mathrm{cosh}(tt_0)\mathrm{sinh}\rho =\mathrm{sinh}\rho _0,`$ (32)
where $`t_0`$, $`\rho _0`$ are free parameters. Notice that the D0-brane reaches the horizon $`\rho =0`$ at $`t\pm \mathrm{}`$ irrespective of the values of $`\rho _0`$ and $`t_0`$. Thus, formally, the Lorentzian D0-brane boundary state is obtainable by Wick rotation of the Euclidean D1-brane boundary state (27).<sup>11</sup><sup>11</sup>11Some classical analysis of D-brane dynamics was attempted in within the Dirac-Born-Infeld approach.
Reconstructing boundary states of the Lorentzian D-brane from those of the Euclidean D-brane is generically not unique. Rather, the following potential subtleties need to be faced:
* The Euclidean momentum $`n`$ along the asymptotic circle of cigar is quantized, while the corresponding quantum number in the Lorentzian theory (i.e. the energy) takes a continuous value.
* The Wick rotations of primary states are not necessarily unique. Often, appropriate boundary conditions should be specified.
As for the first point, which has to do with Matsubara formulation, we can formally avoid the difficulty of quantized momentum by the following heuristic consideration. Suppose the boundary wave function $`\widehat{f}(n,\alpha )`$ ($`n\text{Z}`$ is the quantized Euclidean energy, and $`\alpha `$ denotes the remaining quantum numbers not touched here) is given by the Fourier transform of a periodic function $`f(x+2\pi ,\alpha )=f(x,\alpha )`$. We then obtain
$`B|={\displaystyle \underset{\alpha }{}}{\displaystyle \underset{n}{}}\stackrel{~}{f}(n,\alpha )n,\alpha |`$ $`=`$ $`{\displaystyle \underset{\alpha }{}}{\displaystyle \underset{n}{}}{\displaystyle \frac{1}{2\pi }}{\displaystyle _\pi ^\pi }\text{d}xf(x,\alpha )e^{inx}n,\alpha |`$ (33)
$`=`$ $`{\displaystyle \underset{\alpha }{}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{\text{d}q}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\text{d}xf(x,\alpha )e^{iqx}q,\alpha |,`$
where we used the identity $`_n\delta (qn)=_me^{2\pi imq}`$ in obtaining the last expression. Assuming that $`f(x,\alpha )`$ is analytic along the entire real $`x`$ axis, the Wick rotation can be performed. Often, $`f(x,\alpha )`$ is non-analytic over the real $`x`$ axis, and the integral in the last expression is ill-defined. This turns out to be the case for the boundary wave function of the Euclidean D1-brane (28): in the coordinate space, the wave function has branch cuts and singularities along the real $`x`$-axis. In such cases, the best we can do is to adopt the slightly deformed integration contour $`𝒞`$ in $`x`$-space<sup>12</sup><sup>12</sup>12To be more precise, we should allow to use some decomposition
$$f(x,\alpha )=f_1(x,\alpha )+f_2(x,\alpha )+\mathrm{},$$
and to take the different contours for each piece $`f_i(x,\alpha )`$. to render the Fourier integral well-defined:
$`B^{}||_{\mathrm{Euclidean}}:={\displaystyle \underset{\alpha }{}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{\text{d}q}{2\pi }}{\displaystyle _𝒞}\text{d}xf(x,\alpha )e^{iqx}q,\alpha |.`$ (34)
Likewise, disk one-point function of vertex operator $`\mathrm{\Phi }_{q,\alpha }^{\mathrm{Euclidean}}`$ (associated with the Ishibashi state $`q,\alpha |`$) is evaluated as the deformed contour integral:
$`<\mathrm{\Phi }_{q,\alpha }^{\mathrm{Euclidean}}>_{\text{disk}}={}_{\mathrm{E}}{}^{}B^{}|q,\alpha ={\displaystyle }_𝒞\text{d}xf(x,\alpha )e^{iqx}.`$ (35)
Assuming sufficient analyticity, one then defines Wick rotation of the states (34) by the contour deformation of $`𝒞`$ accompanied by the continuation $`qi\omega ,xit`$;
$`B^{}||_{\mathrm{Lorentzian}}:={\displaystyle \underset{\alpha }{}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{i\text{d}\omega }{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}i\text{d}tf(it,\alpha )e^{i\omega t}i\omega ,\alpha |.`$ (36)
This is essentially the procedure . Of course, we potentially have an ambiguity in the choice of the contour $`𝒞`$, and the correct choice should be determined by the physics under study.
In the present case $`B|`$ corresponds to (28) and $`B^{}|`$ is given by
$`{}_{\mathrm{D1}}{}^{}B^{};\rho _0,\theta _0|={\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\text{d}p}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{\text{d}q}{2\pi }}\mathrm{\Psi }_{\mathrm{D1}}^{}(\rho _0,\theta _0;p,q)p,q|,`$ (37)
where
$`\mathrm{\Psi }_{\mathrm{D1}}^{}(\rho _0,\theta _0;p,q)`$ $`=`$ $`{\displaystyle \frac{\mathrm{sinh}(\pi p)}{\left|\mathrm{cosh}\left(\pi \frac{p+iq}{2}\right)\right|^2}}{\displaystyle \frac{\pi \mathrm{\Gamma }(ip)\mathrm{\Gamma }\left(1+\frac{ip}{k}\right)}{\mathrm{\Gamma }\left(\frac{1}{2}+\frac{ip+q}{2}\right)\mathrm{\Gamma }\left(\frac{1}{2}+\frac{ipq}{2}\right)}}e^{iq\theta _0}\left[e^{ip\rho _0}+{\displaystyle \frac{\mathrm{cosh}\left(\pi \frac{pi|q|}{2}\right)}{\mathrm{cosh}\left(\pi \frac{p+i|q|}{2}\right)}}e^{ip\rho _0}\right]`$
$``$ $`B({\displaystyle \frac{1}{2}}{\displaystyle \frac{ipq}{2}},{\displaystyle \frac{1}{2}}{\displaystyle \frac{ip+q}{2}})\mathrm{\Gamma }\left(1+{\displaystyle \frac{ip}{k}}\right)e^{iq\theta _0}\left[e^{ip\rho _0}+{\displaystyle \frac{\mathrm{cosh}\left(\pi \frac{pi|q|}{2}\right)}{\mathrm{cosh}\left(\pi \frac{p+i|q|}{2}\right)}}e^{ip\rho _0}\right].`$ (38)
Here $`B(p,q)\mathrm{\Gamma }(p)\mathrm{\Gamma }(q)/\mathrm{\Gamma }(p+q)`$ denotes Euler’s beta function. The integration contour $`𝒞`$ we choose is shown in Figure 2 . As in (29), we separately evaluated the integrals of $`\varphi _{L,q}^p`$ and $`\varphi _{R,q}^p`$ based on the decomposition (20). For the convergence of integrals, we choose the contour $`𝒞^+`$ for $`\varphi _{L,q}^p`$ ($`p>0`$ sector) and $`𝒞^{}`$ for $`\varphi _{R,q}^p`$ ($`p<0`$ sector). Such choice of integration contours rendered an extra damping factor $`\mathrm{sinh}(\pi p)/|\mathrm{cosh}\left(\pi \frac{p+iq}{2}\right)|^2`$, which improves the ultraviolet behavior of the wave function and makes it possible to take the Wick rotation sensibly. The non-trivial phase factor $`\mathrm{cosh}\left(\pi \frac{pi|q|}{2}\right)/\mathrm{cosh}\left(\pi \frac{p+i|q|}{2}\right)`$ in the second term originates from the reflection amplitude, and it reduces to $`(1)^n`$ when $`q=n\text{Z}`$.
The second subtlety implies that $`i\omega ,\alpha |`$ is not uniquely defined in (36). This is the issue that arises in a background with horizon, equivalently, non-existence of globally definable timelike Killing vector. As such, this subtlety did not arise for the extremal NS5-brane geometry (described asymptotically by free linear dilaton theory ) considered in . In the next section, within the mini-superspace analysis for the Lorentzian two-dimensional black hole, we shall clarify this subtlety.
An alternative, sensible prescription of the analytic continuation is to define the disk one-point correlator directly via the Lorentzian Fourier transform:
$`<\mathrm{\Phi }_{\omega ,\alpha }^{\text{Lorentzian}}>_{\text{disk}}={\displaystyle _{\mathrm{}}^{\mathrm{}}}\text{d}tf(it,\alpha )e^{i\omega t}.`$ (39)
This is not always equivalent to the the former method elaborated above. In fact, the latter method does not necessarily assert that the boundary state constructed so is expandable in terms of the Lorentzian Ishibashi states that are analytically continued from the Euclidean ones.
### 3.2 Lorentzian mini-superspace wave functions
The Wick rotation of the mini-superspace eigenfunctions in the Euclidean cigar geometry (13) is not so trivial. Fortuitously, the Lorentzian eigenfunctions are already classified thoroughly in . The complete basis for waves outside the black hole horizon are spanned by the following four types of eigenfunctions<sup>13</sup><sup>13</sup>13Here we adopt slightly different normalization from . of the Lorentzian Klein-Gordon operator. For those with the eigenvalue $`\frac{p^2}{4k}\frac{\omega ^2}{4k}+\frac{1}{4k}`$ of the Klein-Gordon operator, the four eigenfunctions are
$`U_\omega ^p(\rho ,t)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }^2(\nu _+)}{\mathrm{\Gamma }(1i\omega )\mathrm{\Gamma }(ip)}}e^{i\omega t}(\mathrm{sinh}\rho )^{i\omega }F(\nu _+,\nu _{}^{};1i\omega ;\mathrm{sinh}^2\rho )`$ (40)
$``$ $`e^{i\omega ti\omega \mathrm{ln}\rho }\text{as}\rho \mathrm{\hspace{0.17em}0},`$
$`V_\omega ^p(\rho ,t)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }^2(\nu _+^{})}{\mathrm{\Gamma }(1+i\omega )\mathrm{\Gamma }(ip)}}e^{i\omega t}(\mathrm{sinh}\rho )^{i\omega }F(\nu _+^{},\nu _{};1+i\omega ;\mathrm{sinh}^2\rho )`$ (41)
$``$ $`e^{i\omega t+i\omega \mathrm{ln}\rho }\text{as}\rho \mathrm{\hspace{0.17em}0},`$
$`L_\omega ^p(\rho ,t)`$ $`=`$ $`e^{i\omega t}(\mathrm{sinh}\rho )^{1ip}F(\nu _+^{},\nu _{}^{};1+ip;{\displaystyle \frac{1}{\mathrm{sinh}^2\rho }})`$ (42)
$``$ $`e^\rho e^{ip\rho i\omega t}\text{as}\rho \mathrm{},`$
$`R_\omega ^p(\rho ,t)`$ $`=`$ $`e^{i\omega t}(\mathrm{sinh}\rho )^{1+ip}F(\nu _+,\nu _{};1ip;{\displaystyle \frac{1}{\mathrm{sinh}^2\rho }})`$ (43)
$``$ $`e^\rho e^{+ip\rho i\omega t}\text{as}\rho \mathrm{}`$
with the notations
$`\nu _\pm ={\displaystyle \frac{1}{2}}i\left({\displaystyle \frac{p}{2}}\pm {\displaystyle \frac{\omega }{2}}\right).`$
These eigenfunctions are defined by the following analytic continuations of the mini-superspace Euclidean eigenfunctions:
$`U_\omega ^p(\rho ,t)=\{\begin{array}{cc}\varphi _{n=+i\omega }^p(\rho ,\theta =+it)\hfill & (\omega >0,n<0)\hfill \\ \varphi _{n=i\omega }^p(\rho ,\theta =it)\hfill & (\omega <0,n>0)\hfill \end{array}`$ (46)
$`V_\omega ^p(\rho ,t)=\{\begin{array}{cc}\varphi _{n=i\omega }^p(\rho ,\theta =it)\hfill & (\omega >0,n<0)\hfill \\ \varphi _{n=+i\omega }^p(\rho ,\theta =+it)\hfill & (\omega <0,n>0)\hfill \end{array}`$ (49)
$`L_\omega ^p(\rho ,t)=\varphi _{L,n=i\omega }^p(\rho ,\theta =+it)`$
$`R_\omega ^p(\rho ,t)=\varphi _{R,n=i\omega }^p(\rho ,\theta =+it),`$ (50)
where the $`n<0`$ and $`n>0`$ ranges are mapped to $`\omega >0`$ and $`\omega <0`$, respectively.
As discussed in , only two out of the four eigenfunctions are linearly independent. In particular,
$`V_\omega ^p(\rho ,t)=U_\omega ^p(\rho ,t)\text{and}R_\omega ^p(\rho ,t)=L_\omega ^p(\rho ,t).`$
The reason why we introduce the above four eigenfunctions is because they encode four possible boundary conditions (We here assume $`p>0`$) in the Lorentzian black hole background. Recall that, for the region outside the horizon of the eternal black hole, the boundaries consist of four segments: ‘future (past) horizon’ $`t=+\mathrm{},\rho =0`$ ($`t=\mathrm{},\rho =0`$) by $`^+`$ ($`^{}`$), and the ‘future (past) infinity’ $`t=+\mathrm{},\rho =+\mathrm{}`$ ($`t=\mathrm{},\rho =+\mathrm{}`$) by $`^+`$ ($`^{}`$). The four eigenfunctions $`U,V,L,R`$ are the ones obeying boundary conditions:
$`U_\omega ^p=0\text{at}^{},V_\omega ^p=0\text{at}^+,L_\omega ^p=0(R_\omega ^p=0)\text{at}^+,R_\omega ^p=0(L_\omega ^p=0)\text{at}^{}`$
for $`\omega >0`$ ($`\omega <0`$). See Figure 3.
By Wick rotating the mini-superspace reflection relations (17), we obtain linear relations among the Lorentzian eigenfunctions:
$`U_\omega ^p=L_\omega ^p+_0(p,\omega )R_\omega ^p\text{and}V_\omega ^p=R_\omega ^p+_0^{}(p,\omega )L_\omega ^p.`$ (51)
Equivalently,
$`L_\omega ^p={\displaystyle \frac{1}{1|_0(p,\omega )|^2}}\left\{U_\omega ^p_0(p,\omega )V_\omega ^p\right\}\text{and}R_\omega ^p={\displaystyle \frac{1}{1|_0(p,\omega )|^2}}\left\{V_\omega ^p_0^{}(p,\omega )U_\omega ^p\right\}.`$ (52)
Here, the mini-superspace reflection amplitude $`_0(p,\omega )`$ in Lorentzian theory is given by
$`_0(p,\omega )={\displaystyle \frac{\mathrm{\Gamma }(+ip)\mathrm{\Gamma }^2(\nu _+)}{\mathrm{\Gamma }(ip)\mathrm{\Gamma }^2(\nu _{}^{})}}{\displaystyle \frac{B(\nu _+,\nu _{})}{B(\nu _+^{},\nu _{}^{})}}{\displaystyle \frac{\mathrm{cosh}\pi \left(\frac{p\omega }{2}\right)}{\mathrm{cosh}\pi \left(\frac{p+\omega }{2}\right)}}.`$ (53)
Notice that, in sharp contrast to the Euclidean black hole, the reflection amplitude is less than unity due to the second factor:
$`|_0(p,\omega )|^2={\displaystyle \frac{\mathrm{cosh}^2\pi \left(\frac{p\omega }{2}\right)}{\mathrm{cosh}^2\pi \left(\frac{p+\omega }{2}\right)}}1.`$ (54)
The inequality is saturated at $`p=\omega =0`$. The inequality (54) shall play a prominent role for understanding string dynamics in the Lorentzian black hole background. The mini-superspace reflection relations for $`U_\omega ^p`$, $`V_\omega ^p`$ are also expressible in a form similar to the Euclidean ones. Recalling that $`_0(p,\omega )_0(+p,\omega )=1`$,
$`U_\omega ^p(\rho ,t)=_0(p,\omega )U_\omega ^p(\rho ,t)\text{and}V_\omega ^p(\rho ,t)=_0^{}(p,\omega )V_\omega ^p(\rho ,t),`$ (55)
while $`L_\omega ^p`$ and $`R_\omega ^p`$ are simply related by reflection:
$`L_\omega ^p(\rho ,t)=R_\omega ^{+p}(\rho ,t).`$ (56)
Moreover, $`U_\omega ^p`$ and $`V_\omega ^p`$ are linearly independent except for the special kinematic regime, $`\omega =0`$. Notice also, in the relation (54), the reflection amplitude involves the mini-superspace contribution only, not the full-fledged stringy one.
Before proceeding further, we shall here collect explicitly relations among inner products of Lorentzian primary fields, where the inner product is defined with respect to the Lorentzian measure $`\text{d}v_L=k\mathrm{sinh}2\rho \text{d}\rho \text{d}t`$. Taking quantum numbers $`p`$, $`\omega `$ fixed and dropping off delta function factors $`2\pi \delta (pp^{})`$, $`2\pi \delta (\omega \omega ^{})`$ for notational simplicity, we have
$`(U_\omega ^p,U_\omega ^p)=(V_\omega ^p,V_\omega ^p)=N_0(p,\omega ),N_0(p,\omega ){\displaystyle \frac{1+|_0(p,\omega )|^2}{2}}`$
$`(U_\omega ^p,V_\omega ^p)=_0^{}(p,\omega ),`$
$`(L_\omega ^p,L_\omega ^p)=(R_\omega ^p,R_\omega ^p)={\displaystyle \frac{1}{2}},(L_\omega ^p,R_\omega ^p)=0,`$
$`(U_\omega ^p,L_\omega ^p)=(V_\omega ^p,R_\omega ^p)={\displaystyle \frac{1}{2}},(R_\omega ^p,U_\omega ^p)=(V_\omega ^p,L_\omega ^p)={\displaystyle \frac{_0(p,\omega )}{2}}.`$ (57)
The inner products involving $`L_\omega ^p`$ and $`R_\omega ^p`$ are readily evaluated since dominant contributions are supported in the asymptotic region $`\rho 0`$, yielding the volume factor $`2\pi \delta (0)`$. The remaining inner products are then extractable from the linear relations (51), (52).<sup>14</sup><sup>14</sup>14We checked these inner products numerically using MATHEMATICA. We also fixed the overall normalization factors from consistency with the Euclidean inner product (15) under the $`\omega \mathrm{\hspace{0.17em}0}`$ limit. Notice also that
$`N_0(p,\omega )=\left|_0(p,\omega )\right|^2N_0(+p,\omega ),`$ (58)
as is consistent with the mini-superspace reflection relation (55).
It is easy to construct the exact string vertex operators or primary states corresponding to the mini-superspace eigenfunctions $`U`$, $`V`$, $`L`$, $`R`$. To be specific, we shall consider primarily the fermionic $`SL_k(2,)/U(1)`$ supercoset conformal field theory.<sup>15</sup><sup>15</sup>15For the bosonic $`SL(2;)_\kappa /U(1)`$ coset conformal field theory, we instead have $`h=\stackrel{~}{h}=\frac{p^2}{4(\kappa 2)}\frac{\omega ^2}{4\kappa }+\frac{1}{4(\kappa 2)}`$, and $`(p,\omega )_0(p,\omega )\frac{\mathrm{\Gamma }\left(1+\frac{ip}{\kappa 2}\right)}{\mathrm{\Gamma }\left(1\frac{ip}{\kappa 2}\right)}`$. The primary states $`|U_\omega ^p`$, $`|V_\omega ^p`$ are the ones of conformal weights $`h=\stackrel{~}{h}=\frac{p^2}{4k}\frac{\omega ^2}{4k}+\frac{1}{4k}`$ and obey the exact reflection relations
$`|U_\omega ^p=(p,\omega )|U_\omega ^p,|V_\omega ^p=^{}(p,\omega )|V_\omega ^p,`$ (59)
and the exact reflection amplitude is given by
$`(p,\omega )_0(p,\omega ){\displaystyle \frac{\mathrm{\Gamma }\left(1+\frac{ip}{k}\right)}{\mathrm{\Gamma }\left(1\frac{ip}{k}\right)}}.`$ (60)
Notice that the string worldsheet effect entering through the $`1/k`$-correction is a pure phase. Thus, the exact reflection probability $`|(p,\omega )|^2`$ remains unmodified from the mini-superspace approximation result $`|_0(p,\omega )|^2`$ given in (54). We shall normalize the primary states $`|U_\omega ^p`$, $`|V_\omega ^p`$ ($`p>0`$) as
$`U_\omega ^p|U_\omega ^{}^p^{}=V_\omega ^p|V_\omega ^{}^p^{}=N(p,\omega )\mathrm{\hspace{0.17em}2}\pi \delta (pp^{})2\pi \delta (\omega \omega ^{}),`$
$`V_\omega ^p|U_\omega ^{}^p^{}=^{}(p,\omega )\mathrm{\hspace{0.17em}2}\pi \delta (pp^{})2\pi \delta (\omega \omega ^{}),`$ (61)
where the new normalization factor $`N(p,\omega )`$ is simply defined by replacing $`_0`$ with $``$ in $`N_0(p,\omega )`$. The primary states $`|L_\omega ^p`$, $`|R_\omega ^p`$ are also definable by using the linear relations (51) or (52) but now with $`_0`$ replaced by $``$. Notice that $`|U_\omega ^p`$, $`|V_\omega ^p`$ are the ones analytically continuable to the Euclidean primary states $`|\varphi _n^{\pm p}`$, so often referred as the ‘Hartle-Hawking vacua’. On the other hand, the states $`|L_\omega ^p`$, $`|R_\omega ^p`$ does not have Euclidean counterparts. Recall that, over the Euclidean black hole background, $`\varphi _{L,n}^p`$, $`\varphi _{R,n}^p`$ behave badly in the vicinity of $`\rho =0`$ and hence ill-defined.
We also find it useful to introduce the dual basis $`\widehat{U_\omega ^p|}`$, $`\widehat{V_\omega ^p|}`$ ($`p,p^{}>0`$) with inner products
$`\widehat{U_\omega ^p|}U_\omega ^{}^p^{}=\widehat{V_\omega ^p|}V_\omega ^{}^p^{}=2\pi \delta (pp^{})2\pi \delta (\omega \omega ^{}),\widehat{U_\omega ^p|}V_\omega ^{}^p^{}=\widehat{V_\omega ^p|}U_\omega ^{}^p^{}=0.`$ (62)
Explicitly, they are given by
$`\widehat{U_\omega ^p|}={\displaystyle \frac{2}{1\left|(p,\omega )\right|^2}}\left\{L_\omega ^p|^{}(p,\omega )R_\omega ^p|\right\},\widehat{V_\omega ^p|}={\displaystyle \frac{2}{1\left|(p,\omega )\right|^2}}\left\{R_\omega ^p|(p,\omega )L_\omega ^p|\right\}.`$ (63)
As such, these dual basis obey the following exact reflection relations:
$`\widehat{U_\omega ^p|}=(p,\omega )\widehat{U_\omega ^p|}\text{and}\widehat{V_\omega ^p|}=(p,\omega )^{}\widehat{V_\omega ^p|}.`$ (64)
A remark is in order. The dual basis $`\widehat{U_\omega ^p|}`$, $`\widehat{V_\omega ^p|}`$ are not Wick rotatable to the Euclidean dual basis $`\varphi _n^{+p}|`$, $`\varphi _n^p|`$, since $`|(p,\omega )|=1`$ for $`\omega i`$. The correct procedure would be that we first define Wick rotations for the ‘ket’ states, and then define their dual states within the Lorentzian Hilbert space. Nevertheless, one-point correlators in the Lorentzian theory, from which a set of physical observables can be computed, ought to be always analytically continuable to the one-point correlators in the Euclidean theory. Roughly speaking, ambiguities inherent to the Wick rotation of dual states drop out upon taking inner product.
Having obtained the Lorentzian primary states, we shall now construct several interesting class of boundary states for a D0-brane propagating in the black hole background. We have seen that the D0-brane propagates along the trajectory (32). The two-dimensional black hole is eternal, so, in addition to the past and the future asymptotic infinities, the causal propagation region has the past horizon $`^{}`$ surrounding the white hole singularity and the future horizon $`^+`$ surrounding the black hole singularity. As such, by taking variety of possible boundary conditions, we can construct interesting class of boundary states.
### 3.3 Boundary state of D0-brane absorbed to future horizon
Consider first the boundary state obeying the boundary condition $`\psi (\rho ,t)\mathrm{\hspace{0.17em}0}`$ at the past horizon $`^{}`$, viz. the primary states $`|U_\omega ^p`$. This boundary condition is relevant for scattering of a D0-brane off the black hole, since the condition represents absorption only and no emission of the D0-brane by the black hole. D0-brane boundary state obeying such absorbing boundary condition is then expanded solely by the Ishibashi states $`{}_{}{}^{\widehat{U}}p,\omega |`$, $`|p,\omega ^U`$ that are associated with the primary states $`\widehat{U_\omega ^p|}`$, $`|U_\omega ^p`$:
$`{}_{\text{absorb}}{}^{}B;\rho _0,t_0|={\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\text{d}p}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{\text{d}\omega }{2\pi }}\mathrm{\Psi }_{\text{absorb}}(\rho _0,t_0;p,\omega ){}_{}{}^{\widehat{U}}p,\omega |,`$
$`|B;\rho _0,t_0_{\text{absorb}}={\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\text{d}p}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{\text{d}\omega }{2\pi }}\mathrm{\Psi }_{\text{absorb}}^{}(\rho _0,t_0;p,\omega )|p,\omega ^U.`$ (65)
The boundary wave function $`\mathrm{\Psi }_{\text{absorb}}(\rho _0,t_0;p,\omega )`$ is then interpreted as the disk one-point correlators:
$`\mathrm{\Psi }_{\text{absorb}}(\rho _0,t_0;p,\omega )`$ $`=`$ $`U_\omega ^p_{\text{disk}}{}_{\text{absorb}}{}^{}B;\rho _0,t_0|U_\omega ^p,`$ (66)
The boundary wave function (66) is then obtained by taking the Wick rotation $`qi\omega `$ ($`qi\omega `$) for $`q<0`$ ($`q>0`$) in (38) (recall (50)):<sup>16</sup><sup>16</sup>16In reality, there is a further overall factor $`i`$, but, for notational simplicity, we will absorb it to the definition of the Ishibashi states.
$`\mathrm{\Psi }_{\text{absorb}}(\rho _0,t_0;p,\omega )=B(\nu _+,\nu _{})\mathrm{\Gamma }\left(1+{\displaystyle \frac{ip}{k}}\right)e^{i\omega t_0}\left[e^{ip\rho _0}{\displaystyle \frac{\mathrm{cosh}\left(\pi \frac{p\omega }{2}\right)}{\mathrm{cosh}\left(\pi \frac{p+\omega }{2}\right)}}e^{ip\rho _0}\right],`$ (67)
The relative minus sign in the second term of $`\mathrm{\Psi }_{\text{absorb}}(\rho _0,t_0;p,\omega )`$ originates from the fact that the contour rotation defining the Wick rotation has opposite directions for $`𝒞^+`$ (suitable for $`p>0`$) and $`𝒞^{}`$ (suitable for $`p<0`$). See figure 2. This boundary wave function (67) satisfies the exact reflection relation
$`\mathrm{\Psi }_{\text{absorb}}(\rho _0,t_0;p,\omega )=(p,\omega )\mathrm{\Psi }_{\text{absorb}}(\rho _0,t_0;p,\omega ).`$ (68)
With such boundary condition, the boundary wave function (67) would have no overlap with D0-brane’s trajectory (32) in the far past region $`tt_0`$. In fact, the trajectory (32) starts from the past horizon $`^{}`$ at $`t=\mathrm{}`$, reaches the time-symmetric point $`\rho =\rho _0`$ at $`t=t_0`$, and then falls back the future horizon $`^+`$ at $`t=+\mathrm{}`$, while the wave function $`U_\omega ^p`$ does not have any component outgoing from $`^{}`$. We thus interpret that the boundary state (67) describes the future half of the classical trajectory (32). We shall hence call it the ‘absorbed D-brane’.
By utilizing the radion-tachyon correspondence, the rolling radion (as described by the boundary state (67)) is also interpretable as the rolling tachyon. In the latter interpretation, the D0-brane absorbed to the future horizon is the counterpart of the future-half S-brane , in which the tachyon rolls down the potential hill at asymptotic future $`t+\mathrm{}`$ and emits radiation.
### 3.4 Boundary state of D0-brane emitted from past horizon
Consider next the boundary condition: $`\psi (\rho ,t)\mathrm{\hspace{0.17em}0}`$ at $`^+`$, viz. use the basis $`|p,\omega ^V`$, $`{}_{}{}^{\widehat{V}}p,\omega |`$ instead of $`|p,\omega ^U`$, $`{}_{}{}^{\widehat{U}}p,\omega |`$. Utilizing the reflection relation, we can first rewrite (38) as the form which only includes the $`p<0`$ Ishibashi states by means of the reflection relation. Then, we can analytically continue the states $`|\varphi _q^p`$ ($`p>0`$) into $`|V_\omega ^p`$. The resultant boundary state is obtained by simply replacing $`pp`$, $`\omega \omega `$ in (67);
$`{}_{\text{emitted}}{}^{}B;\rho _0,t_0|={\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\text{d}p}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{\text{d}\omega }{2\pi }}\mathrm{\Psi }_{\text{emitted}}(\rho _0,t_0;p,\omega ){}_{}{}^{\widehat{V}}p,\omega |.`$
$`|B;\rho _0,t_0_{\text{emitted}}={\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\text{d}p}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{\text{d}\omega }{2\pi }}\mathrm{\Psi }_{\text{emitted}}^{}(\rho _0,t_0;p,\omega )|p,\omega ^V.`$ (69)
where
$`\mathrm{\Psi }_{\text{emitted}}(\rho _0,t_0;p,\omega )=B(\nu _+^{},\nu _{}^{})\mathrm{\Gamma }\left(1{\displaystyle \frac{ip}{k}}\right)e^{i\omega t_0}\left[e^{ip\rho _0}{\displaystyle \frac{\mathrm{cosh}\left(\pi \frac{p\omega }{2}\right)}{\mathrm{cosh}\left(\pi \frac{p+\omega }{2}\right)}}e^{ip\rho _0}\right].`$
Obviously, the emitted D0-brane wave function is the time-reversal of the absorbed D0-brane wave function (67):
$`\mathrm{\Psi }_{\text{emitted}}(\rho _0,t_0;p,\omega )=\mathrm{\Psi }_{\text{absorb}}^{}(\rho _0,t_0;p,\omega ).`$
Namely, it describes the D0-brane emitted from the past horizon at asymptotic past $`t=\mathrm{}`$. By the choice of the boundary condition, this boundary state (69) describes only the past half of the classical D0-brane trajectory (32).
The exact reflection relation has the form
$`\mathrm{\Psi }_{\text{emitted}}(\rho _0,t_0;p,\omega )=^{}(p,\omega )\mathrm{\Psi }_{\text{emitted}}(\rho _0,t_0;p,\omega ).`$ (70)
Again, in light of the radion-tachyon correspondence, the D0-brane emitted from the past horizon is the counterpart of the past-half S-brane in tachyon rolling. The radiation creeps up the tachyon potential hill from past infinity and forms an unstable D-brane.
### 3.5 Boundary state of time-symmetric D0-brane
The third possible boundary state is obtainable by directly taking the analytic continuation in the disk one-point amplitudes, as we already mentioned. Recalling (50), we shall analytically continue the disk amplitudes as (assume $`p>0`$)
$`\varphi _q^{+p}_{\text{disk}}U_\omega ^p_{\text{disk}}\text{and}\varphi _q^p_{\text{disk}}V_\omega ^p_{\text{disk}}.`$ (71)
The Euclidean one-point amplitudes $`\varphi _q^{\pm p}_{\text{disk}}`$ are given in (38), and can be expressed in contour integrals as in (35). Recall that $`\varphi _{L,q}^p_{\text{disk}}`$, $`\varphi _{R,q}^p_{\text{disk}}`$ are prescribed by the contour integrals over $`𝒞^+`$, $`𝒞^{}`$ in figure 2. We shall thus analytically continue them to the real time axis (imaginary $`x`$-axis). In this way, we extract the Lorentzian disk one-point amplitudes as
$`U_\omega ^p_{\text{disk}}=U_\omega ^p_{\text{disk}}^{(\text{absorb})}\text{and}V_\omega ^p_{\text{disk}}=V_\omega ^p_{\text{disc}}^{(\text{emitted})},`$ (72)
where the right-hand sides are simply the amplitudes associated with the ‘absorbed’ and ‘emitted’ D0-branes considered in the previous subsections and explicitly given in (67) and (69). Since $`U_\omega ^p`$ and $`V_\omega ^p`$ constitute the complete set of basis for Lorentzian primary fields, the amplitudes (72) would yield yet another Lorentzian D0-brane boundary states. As is obvious from the above construction, this state keeps the time-reversal symmetry manifest and reproduces the entire classical trajectory (32), that is, it describes a D0-brane emitted from the past horizon and reabsorbed to the future horizon. From the viewpoint of the boundary conformal theory, this would be considered the most natural one since it captures the entire classical trajectory of the D0-brane. In the radion-tachyon correspondence, this state is the counterpart of the full S-brane .
Explicitly, the time-symmetric boundary states are given by
$`{}_{\text{symm}}{}^{}B;\rho _0,t_0|={}_{\text{absorb}}{}^{}B;\rho _0,t_0|+{}_{\text{emitted}}{}^{}B;\rho _0,t_0|`$
$`={\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\text{d}p}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{\text{d}\omega }{2\pi }}[2\mathrm{\Psi }_{\text{symm}}(\rho _0,t_0;p,\omega ){}_{}{}^{L}p,\omega |+2\mathrm{\Psi }_{\text{symm}}^{}(\rho _0,t_0;p,\omega ){}_{}{}^{R}p,\omega \left|\right]`$
$`|B;\rho _0,t_0_{\text{symm}}=|B;\rho _0,t_0_{\text{absorb}}+|B;\rho _0,t_0_{\text{emitted}}`$
$`={\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\text{d}p}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{\text{d}\omega }{2\pi }}\left[2\mathrm{\Psi }_{\text{symm}}^{}(\rho _0,t_0;p,\omega )\right|p,\omega ^L+2\mathrm{\Psi }_{\text{symm}}(\rho _0,t_0;p,\omega )|p,\omega ^R],`$ (73)
where
$`\mathrm{\Psi }_{\text{symm}}(\rho _0,t_0;p,\omega )=B(\nu _+,\nu _{})\mathrm{\Gamma }\left(1+{\displaystyle \frac{ip}{k}}\right)e^{ip\rho _0i\omega t_0}`$
and $`{}_{}{}^{L}p,\omega |`$, $`|p,\omega ^L`$, $`{}_{}{}^{R}p,\omega |`$, $`|p,\omega ^R`$ are the Ishibashi states constructed over the primary states $`L_\omega ^p|`$, $`|L_\omega ^p`$, $`R_\omega ^p|`$, $`|R_\omega ^p`$,<sup>17</sup><sup>17</sup>17The extra factor of ‘2’ was introduced for convenience. Recall (57). respectively. One can readily check that the second lines in (73) are indeed correct by evaluating the disk one-point amplitudes from them. For instance, using (57), we obtain
$`U_\omega ^p_{\text{disk}}^{(\text{symm})}`$ $`=`$ $`{}_{\text{symm}}{}^{}B;\rho _0,t_0|U_\omega ^p`$
$`=`$ $`\mathrm{\Psi }_{\text{symm}}(\rho _0,t_0;p,\omega )+(p,\omega )\mathrm{\Psi }_{\text{symm}}^{}(\rho _0,t_0;p,\omega )`$
$`=`$ $`B(\nu _+,\nu _{})\mathrm{\Gamma }\left(1+{\displaystyle \frac{ip}{k}}\right)e^{i\omega t_0}\left[e^{ip\rho _0}{\displaystyle \frac{\mathrm{cosh}\left(\pi \frac{p\omega }{2}\right)}{\mathrm{cosh}\left(\pi \frac{p+\omega }{2}\right)}}e^{ip\rho _0}\right]`$
$`=`$ $`U_\omega ^p_{\text{disk}}^{(\text{absorb})}{}_{\text{absorb}}{}^{}B;\rho _0,t_0|U_\omega ^p.`$
Other one-point amplitudes can be checked analogously.
Two remarks are in order. First, notice that, though the disk one-point amplitudes are, the symmetric boundary states (73) by themselves are not analytically continuable to the Euclidean boundary state (38). This should not be surprising as the Lorentzian Hilbert space is generated by twice as many generators as the Euclidean theory. In other words, the Lorentzian bases $`|U_\omega ^p`$, $`|V_\omega ^p`$ correspond to $`|\varphi _n^p`$, $`|\varphi _n^p`$ in the Euclidean theory, which were however linearly dependent due to the reflection relation. Nevertheless, the boundary state (73) is a consistent one and yields disk one-point amplitudes that can be correctly continued to the Euclidean ones. Second, the full Lorentzian Hilbert space is decomposed as
$`=^U^V\text{and}\widehat{}=\widehat{^U}\widehat{^V},`$ (74)
where $`^U`$ ($`^V`$) is spanned by $`|U_\omega ^p`$ , ($`|V_\omega ^p`$) and their descendants. The dual space $`\widehat{^U}`$ ($`\widehat{^V}`$) is similarly spanned by $`\widehat{U_\omega ^p|}`$, ($`\widehat{V_\omega ^p|}`$). Here, the Hilbert subspaces $`^U`$, $`\widehat{^U}`$ ($`^V`$, $`\widehat{^V}`$) correspond to the boundary condition $`\psi (\rho ,t)\mathrm{\hspace{0.17em}0}`$ at $`^{}`$ ($`^+`$). The ‘absorbed’ and ‘emitted’ D0-brane boundary states (67), (69) are consistent only in the subspaces $`^U`$, $`^V`$ ($`\widehat{^U}`$, $`\widehat{^V}`$), while the ‘symmetric’ D0-brane boundary state (73) is well-defined in the entire Hilbert space $``$ ($`\widehat{}`$). We thus have simple relations
$`|B;\rho _0,t_0_{\text{absorb}}=P_U|B;\rho _0,t_0_{\text{symm}}`$ and $`{}_{\text{absorb}}{}^{}B;\rho _0,t_0|={}_{\text{symm}}{}^{}B;\rho _0,t_0|\widehat{P_U},`$
$`|B;\rho _0,t_0_{\text{emitted}}=P_V|B;\rho _0,t_0_{\text{symm}}`$ and $`{}_{\text{emitted}}{}^{}B;\rho _0,t_0|={}_{\text{symm}}{}^{}B;\rho _0,t_0|\widehat{P_V},`$ (75)
where $`P_{U,V}`$ ($`\widehat{P_{U,V}}`$) denotes projection of the Hilbert space $``$ to $`^{U,V}`$ ($`\widehat{^{U,V}}`$).
## 4 Radiation out of D0-Brane Rolling in the Black Hole Background
In the background of the black hole, the D0-brane moves along the geodesic and we have constructed a variety of boundary states describing the geodesic motion, specified by appropriate boundary conditions.
Both by gravity and by strong string coupling gradient, the D$`p`$-brane is pulled in and finds its minimum energy and mass at the location of the NS5-brane. The D$`p`$-brane is supersymmetric in flat spacetime, but preserves no supersymmetry in black NS5-brane background. Even in extremal NS5-brane background, until the D$`p`$-brane dissociates into the NS5-brane and form a non-threshold bound-state, the spacetime supersymmetry is completely broken. In these respects,the D$`p`$-brane propagating in the NS5-brane background is much like excited D$`p`$-brane (many excited open strings attached on it) in flat spacetime. Decay of the latter via closed string emission was studied extensively for $`p=1`$ : the decay spectrum was found to match exactly with the Hawking radiation of the non-extremal black hole made out of these excited D-branes, and the effective temperature of excited open string modes agrees exactly with the Hawking temperature. In this section, we shall find certain analogous results for the closed string radiation off the rolling D0-brane, though special features also arise.
As the D0-brane is pulled in, acceleration would grow and radiate off the binding energy into closed string modes. Details of the radiation spectra would differ for different choice of the boundary conditions, viz. for different boundary states of the D0-brane. In this section, as a probe of the black hole geometry and D-brane dynamics therein, we shall analyze spectral distribution of the closed string radiation off the rolling D0-particle.
By applying the optical theorem, the radiation rate during the radion-rolling process is obtainable as the imaginary part of the annulus amplitude in the closed string channel.<sup>18</sup><sup>18</sup>18For the tachyon rolling process in flat spacetime background, the amplitude was evaluated first in . Denote the differential number density $`\text{d}𝒩(p,M)`$ of the radiation at a fixed value of the radial momentum $`p`$ and the mass-level $`M`$. By the definition of the D-brane boundary state, the radiation number density $`\text{d}𝒩`$ is then given in terms of the boundary wave functions:
$`\text{d}𝒩(p,M)`$ $`:=`$ $`{\displaystyle \frac{\text{d}p}{2\pi }}{\displaystyle \frac{\text{d}M}{(2\pi )^d}}{\displaystyle \text{d}\omega <\mathrm{\Psi }(\omega ,p,M)\left|\delta (L_0+\overline{L}_0)\right|\mathrm{\Psi }(\omega ,p,M)>}`$ (76)
$`=`$ $`{\displaystyle \frac{\text{d}p}{2\pi }}{\displaystyle \frac{\text{d}M}{(2\pi )^d}}{\displaystyle \frac{1}{2\omega (p,M)}}\left|\mathrm{\Psi }(p,\omega (p,M))\right|^2.`$
Here, $`\omega ,p`$ are the energy and the radial momentum in two-dimensional Lorentzian background, $`M`$ is the total mass (conformal weight) of the remaining subspaces of dimension $`d`$ (including mass gap), $`\mathrm{\Psi }(\omega ,p,M)`$ is the boundary wave function (including that of the remaining subspace), and $`\omega (p,M)(>0)`$ is the on-shell energy of the radiated closed string state determined by the on-shell condition $`L_0+\overline{L}_0=0`$ including the ghost contribution. From the kinematical consideration, it is obvious that the differential number density (76) is nonzero only when the D-brane is rolling. Of particular physical interest is the spectral distribution in the phase-space, as measured by the independent moments, e.g.
$`<\omega ^mM^n>`$ $`=`$ $`{\displaystyle \frac{\text{d}p}{2\pi }\frac{\text{d}M}{(2\pi )^d}\omega ^m(p,M)M^n\frac{1}{2\omega (p,M)}\left|\mathrm{\Psi }(p,\omega (p,M))\right|^2}`$
for $`m,n=0,1,2,\mathrm{}`$. We shall evaluate these spectral observables by first evaluating the integral over the radial momentum $`p`$ by saddle-point approximation. In doing so, we pay particular attention to the asymptotic behavior as the mass-level $`M`$ becomes asymptotically large. We shall then evaluate the integral over the mass-level (conformal weight) $`M`$, and extract the spectral observables.
Consider the boundary state (67) describing a D0-brane absorbed by the future horizon. The radiation emitted by the D0-brane is decomposable into ‘incoming’ (toward the horizon) and ‘outgoing’ (toward the null infinity) components in the far future. The positive energy sector is expanded by the wave function $`U_\omega ^p`$, and has the following asymptotic behavior at $`t+\mathrm{}`$:
$`U_\omega ^p(\rho ,t)e^{i\omega \mathrm{ln}\rho i\omega t}+d(p,\omega )e^\rho e^{+ip\rho i\omega t}\text{where}|d(p,\omega )|e^{\pi p}.`$ (77)
Here, we assumed $`\omega M0`$. The first and the second terms correspond to the incoming wave supported around $`\rho =0`$ and the outgoing wave supported in the region $`\rho +\mathrm{}`$, respectively. The damping factor $`d(p)`$ originates from the exact reflection amplitude $`(p,\omega )`$. (See (51), (52).) To obtain the radiation number density, we need to evaluate $`\left|\mathrm{\Psi }(p,\omega )\right|^2\times |U_\omega ^p(\rho ,t)|^2`$. At far future infinity, the interference term in $`|U_\omega ^p|^2`$ drops off upon taking the $`p`$-integral. Therefore, after integrating over the radial momentum $`p`$, the partial radiation distribution is seen to consist of the ‘incoming’ and ‘outgoing’ parts:
$`𝒩(M)_{\text{in}}`$ $``$ $`{\displaystyle _0^M}\text{d}M{\displaystyle \frac{\text{d}𝒩_{\text{in}}}{\text{d}M}}={\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\text{d}p}{2\pi }}{\displaystyle \frac{1}{2\omega (p,M)}}\left|\mathrm{\Psi }(p,\omega (p,M))\right|^2`$
$`𝒩(M)_{\text{out}}`$ $``$ $`{\displaystyle _0^M}\text{d}M{\displaystyle \frac{\text{d}𝒩_{\text{out}}}{\text{d}M}}={\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\text{d}p}{2\pi }}{\displaystyle \frac{1}{2\omega (p,M)}}\left|d(p)\right|^2\left|\mathrm{\Psi }(p,\omega (p,M))\right|^2.`$ (78)
We shall now evaluate the branching ratio between the two radiation rates (78) with emphasis on possible string worldsheet effects. To this end, consider the conformal field theory defined by $`SL(2;)/U(1)\times `$, where $`SL(2;)/U(1)`$ denotes the (super)coset model and $``$ denotes a unitary (super)conformal field theory of central charge $`c_{}`$. Such (super)conformal field theory covers a variety of interesting string theory backgrounds. For the fermionic string, superconformal invariance asserts that the central charge ought to be critical:
$`3\left(1+{\displaystyle \frac{2}{k}}\right)+c_{}=15,`$
where $`k`$ denotes the level of the super $`SL(2;)`$ current algebra. If the background describes a stack of black NS5-branes, $`=SU(2)_k\times ^5`$ where $`k`$ equals to the NS5-brane charge. Likewise, for the bosonic string case, conformal invariance asserts that the central charge should take the critical value:
$`2+{\displaystyle \frac{6}{\kappa 2}}+c_{}=26,`$ (79)
where now $`\kappa `$ refers to the level of the bosonic $`SL(2;)`$ current algebra. For the background describing the black hole in two-dimensional string theory, $``$ is empty and $`\kappa `$ should be set to $`9/4`$.
It would be illuminating to analyze the branching ratio for the ‘rolling closed string’, viz. a closed string state of fixed transverse mass $`M`$ and radial momentum $`p`$ propagating in black hole geometry. The branching ratio is simply given by the reflection amplitude (see (60)):
$`{\displaystyle \frac{𝒩_{\mathrm{out}}(p,\omega )}{𝒩_{\mathrm{in}}(p,\omega )}}|_{\mathrm{closed}\mathrm{string}}=|(p,\omega )|^2={\displaystyle \frac{\mathrm{cosh}^2\pi \left(\frac{\omega p}{2}\right)}{\mathrm{cosh}^2\pi \left(\frac{\omega +p}{2}\right)}}.`$ (80)
As emphasized below (60), string worldsheet effects are present for the reflection amplitude $``$ itself but, being an overall phase, it drops out of (80). The $`k`$-dependence enters in the branching ratio (80) only through the on-shell dispersion relation $`\omega =\sqrt{p^2+2kM^2}`$. For two-dimensional case, first studied in and , $`k=1/2`$, $`M=0`$ and $`\omega =p`$, so the scattering probability is exponentially suppressed as the energy increases.
For a fixed transverse mass $`M`$ and the forward radial momentum $`p`$, the reflection probability of the infalling D0-brane is given precisely by the same result as (80):
$`{\displaystyle \frac{𝒩_{\mathrm{out}}(p,\omega )}{𝒩_{\mathrm{in}}(p,\omega )}}|_{\mathrm{D0}\mathrm{brane}}=|(p,\omega )|^2={\displaystyle \frac{\mathrm{cosh}^2\pi \left(\frac{\omega p}{2}\right)}{\mathrm{cosh}^2\pi \left(\frac{\omega +p}{2}\right)}}.`$ (81)
This is simply because back-scattering of the boundary wave function originates from that of the closed string wave function: roughly speaking, the boundary wave function is defined by overlap of the closed string wave function with the classical trajectory of the D0-brane.
Radiation out of the falling D0-brane is coherent, so we integrate over the radial momentum $`p`$ as in (78) in extracting the branching ratio. We shall first analyze the partial radiation distribution at large mass-level, $`M\mathrm{}`$. More precisely, we shall examine asymptotic behavior of $`𝒩(M)`$ multiplied by the phase-space ‘degeneracy factor’ $`\rho (M)e^{\frac{1}{2}M\beta _{\mathrm{Hg}}}`$, where $`\beta _{\mathrm{Hg}}`$ denotes inverse of the Hagedorn temperature. The closed string states that couple to the boundary states are left-right symmetric, so we need to take the square root of the usual degeneracy factor in the closed string sector. Here, inverse of the Hagedorn temperature is given by
$`\beta _{\mathrm{Hg}}=4\pi \sqrt{1{\displaystyle \frac{1}{2k}}},`$ (82)
for the superstring theory, and
$`\beta _{\mathrm{Hg}}=4\pi \sqrt{2{\displaystyle \frac{1}{2(\kappa 2)}}},`$ (83)
for the bosonic string theory, where the $`1/k`$ $`(1/\kappa )`$-correction is interpreted as the string worldsheet effects of the two-dimensional background. These results are derivable from the Cardy formula with the ‘effective central charge’ $`c_{\text{eff}}=c24h_{\text{min}}`$ , where $`h_{\text{min}}`$ refers to the lowest conformal weight of normalizable primary states.
### 4.1 Radiation distribution in superstring theory
Begin with the spectral distribution in superstring theories. We shall focus exclusively on the NS-NS sector of the radiation and defer the analysis of the R-R sector to section 6. The on-shell condition of closed string state in NS-NS sector is given by
$`{\displaystyle \frac{\omega ^2}{4k}}+{\displaystyle \frac{p^2}{4k}}+{\displaystyle \frac{1}{4k}}+\mathrm{\Delta }_{}={\displaystyle \frac{1}{2}},`$ (84)
where $`\mathrm{\Delta }_{}`$ denotes the conformal weight of the $``$-part. The on-shell energy is given by
$`\omega \omega (p,M)=\sqrt{p^2+2kM^2}\text{where}M^22\left(\mathrm{\Delta }_{}+{\displaystyle \frac{1}{4k}}{\displaystyle \frac{1}{2}}\right).`$
Consider now a D0-brane propagating outside the black hole and absorbed into the future horizon. The relevant boundary wave function was constructed in (67) and, from them, the differential radiation number distributions (78) can be computed. At large $`\omega `$ and $`p`$, using Stirling’s approximation, we find that
$`𝒩(M)_{\text{in}}`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\text{d}p}{2\pi }}{\displaystyle \frac{1}{2\omega (p,M)}}\left|\mathrm{\Psi }_{\mathrm{absorb}}(\rho _0,t_0;p,\omega (p,M))\right|^2`$ (85)
$``$ $`{\displaystyle \frac{1}{M}}{\displaystyle _0^{\mathrm{}}}\text{d}pe^{+\pi \left(1\frac{1}{k}\right)p\pi \sqrt{p^2+2kM^2}}`$
$`𝒩(M)_{\text{out}}`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\text{d}p}{2\pi }}\left|d(p,\omega (p,M))\right|^2{\displaystyle \frac{1}{2\omega (p,M)}}\left|\mathrm{\Psi }_{\mathrm{absorb}}(\rho _0,t_0;p,\omega (p,M))\right|^2`$ (86)
$``$ $`{\displaystyle \frac{1}{M}}{\displaystyle _0^{\mathrm{}}}\text{d}pe^{\pi \left(1+\frac{1}{k}\right)p\pi \sqrt{p^2+2kM^2}}.`$
In the second lines, we have taken $`M`$ large, viz. $`\omega p1`$, and keep the leading terms only. Thus, for each fixed but large $`M`$, the partial number distributions take the forms:
$`𝒩(M)_{\text{in}}{\displaystyle _0^{\mathrm{}}}\text{d}p\sigma _{\text{in}}(p)e^{\frac{1}{2}\beta _{\mathrm{Hw}}M}\text{and}N(M)_{\text{out}}{\displaystyle _0^{\mathrm{}}}\text{d}p\sigma _{\text{out}}(p)e^{\frac{1}{2}\beta _{\mathrm{Hw}}M},`$ (87)
where
$`\beta _{\mathrm{Hw}}=2\pi \sqrt{2k},`$ (88)
is the inverse Hawking temperature of the fermionic two-dimensional black hole. As discussed above, the radiation off the D-brane in NS5-brane background is analogous to the decay of excited D-brane in flat ambient spacetime. Indeed, asymptotic expression (87) suggests that open string excitations of energy $`M`$ on the rolling D0-brane are populated as the distribution function $`\mathrm{exp}(\frac{1}{2}\beta _{\mathrm{Hw}}M)`$ and decay into closed string radiation. In this interpretation, the distribution function encodes change of available states for open string excitations on the D0-brane after emitting radiations of energy $`M`$. Curiously, ‘effective temperature’ of the excited closed strings is set by the Hawking temperature of the nonextremal NS5-brane, not that of a black hole that would have been made of the D0-brane. It is tempting to interpret this as indicating that the D0-brane represents a class of possible excitation modes of the black NS5-brane. The closed string states of energy $`M`$ emitted by the D0-brane are certainly coherent, but according to this interpretation, they still can be recasted in effective thermal distribution set by the Hawking temperature of the two-dimensional black hole. In the next subsection, we shall account for the origin of such effective thermal behavior of the rolling D0-brane from the viewpoints of Euclidean cylinder amplitude between D1-brane, extending the argument of for the Hawking radiation of closed strings in the black hole background.
The functions $`\sigma _{\mathrm{in}}`$ and $`\sigma _{\mathrm{out}}`$ are interpretable as the black hole ‘greybody’ factors for incoming and outgoing parts of the radiation. The factor 1/2 in the exponent of the Boltzmann distribution function reflects the fact that only left-right symmetric closed string states can appear in the boundary states and the radiated closed string modes.
The ‘greybody factors’ $`\sigma _{}(p)`$ depend on the radial momentum $`p`$ exponentially, so the radiation distribution would be modified once the radial momentum $`p`$ is integrated out. Below, we shall show this explicitly. We are primarily interested in keeping track of string worldsheet effects set by the value of the level $`k`$. We shall consider different ranges of the level $`k`$ separately, and focus on the asymptotic behaviors at large $`M`$ via the saddle point methods.
This is the case for the black NS5-brane background. Consider first the incoming part. Since $`1\frac{1}{k}>0`$, the dominant contribution in the $`p`$-integral arises from the saddle point:
$`pp_{}={\displaystyle \frac{k1}{\sqrt{1\frac{1}{2k}}}}M.`$
Substituting this to (85), we obtain
$`𝒩(M)_{\text{in}}e^{2\pi M\sqrt{1\frac{1}{2k}}}=e^{\frac{1}{2}M\beta _{\mathrm{Hg}}},`$ (89)
up to pre-exponential powers of $`M`$. Taking account of the density of states $`\rho (M)e^{\frac{1}{2}M\beta _{\mathrm{Hg}}}`$, we find that $`\rho (M)𝒩(M)_{\text{in}}`$ scales with powers of $`M`$, and is independent of $`k`$. More explicitly, for the black NS5-brane $`=SU(2)_k\times ^5`$, the incoming radiation distribution of the D$`p`$-brane parallel to the NS5-brane yields
$`𝒩(M)_{\text{in}}`$ $``$ $`{\displaystyle \frac{1}{M}}{\displaystyle \frac{\text{d}^{5p}𝐤_{}}{(2\pi )^{5p}}_0^{\mathrm{}}\text{d}pe^{\pi (1\frac{1}{k})p\pi \sqrt{p^2+2k(M^2+𝐤_{}^2)}}}`$
$``$ $`M^{2\frac{p}{2}}e^{2\pi M\sqrt{1\frac{1}{2k}}}.`$
Taking account of the density of states $`\rho (M)M^3e^{2\pi M\sqrt{1\frac{1}{2k}}}`$, the average radiation number distribution is given by
$`{\displaystyle \frac{\overline{𝒩}_{\text{in}}}{V_p}}{\displaystyle ^{M_\mathrm{D}}}{\displaystyle \frac{\text{d}M}{M}}M^{\frac{p}{2}}\text{where}M_\mathrm{D}𝒪({\displaystyle \frac{1}{g_{\mathrm{st}}}}).`$ (90)
This result coincides with the computations of , and corroborates with the radion-tachyon correspondence. Interestingly, the incoming part of the radiation number distribution in the the nonextremal NS5-brane background is exactly the same as the distribution in the extremal NS5-brane background. Later, we shall examine carefully taking the extremal limit and its consequence in section 7. As in the extremal case, (90) implies that nearly all the D0-brane potential energy is released into closed string radiations before it falls into the black hole.
On the other hand, for the outgoing radiation, the far infrared $`p0`$ dominates the momentum integral. We thus obtain
$`𝒩(M)_{\text{out}}e^{2\pi M\sqrt{\frac{k}{2}}}=e^{\frac{1}{2}M\beta _{\mathrm{Hw}}},`$
displaying effective thermal distribution set by the Hawking temperature. Taking account of the density of states,
$`\rho (M)𝒩(M)_{\text{out}}e^{\frac{1}{2}M\left(\beta _{\mathrm{Hg}}\beta _{\mathrm{Hw}}\right)}=e^{2\pi M\left(\sqrt{1\frac{1}{2k}}\sqrt{\frac{k}{2}}\right)}.`$
This is ultraviolet finite for any $`k`$ since
$`\left(1{\displaystyle \frac{1}{2k}}\right){\displaystyle \frac{k}{2}}={\displaystyle \frac{1}{2k}}\left(k1\right)^2<0.`$ (91)
We thus conclude that the radiation number distribution is mostly in the incoming part:
$`{\displaystyle \frac{𝒩_{\mathrm{out}}(M)\rho (M)}{𝒩_{\mathrm{in}}(M)\rho (M)}}|_{\mathrm{falling}\mathrm{D0}}{\displaystyle \frac{e^{\frac{1}{2}\beta _{\mathrm{Hw}}M}}{e^{\frac{1}{2}\beta _{\mathrm{Hg}}M}}}=e^{2\pi M\left(\sqrt{1\frac{1}{2k}}\sqrt{\frac{k}{2}}\right)}1.`$
Intuitively, this may be understood as follows: for the absorbed boundary state, the boundary condition is such that the D0-brane flux is directed from past null infinity to the future horizon. This also corroborates the observation that $`T_{t\rho }`$-component of D0-brane’s energy-momentum tensor is nonzero and increases monotonically as the D0-brane approaches the future horizon. The outgoing part of the distribution is exponentially small compared to the incoming part and exhibits effective thermal distribution at the Hawking temperature. Notice that, despite being so, this outgoing part has nothing to do with the Hawking radiation of the black hole. The latter is the feature of the background by itself. A priori, the outgoing radiation could be in a distribution characterized by a temperature different from the Hawking temperature. As mentioned above, it is tempting to interpret coincidence of the two temperatures as a consequence of maintaining equilibrium between the black NS5-brane and the D0-brane.
This is the regime which includes the conifold geometry at $`k=1`$. Since $`1\frac{1}{k}0`$, the dominant contribution to the momentum integral is from $`p0`$, not only for the outgoing radiation but also for the incoming one. We thus obtain
$`𝒩(M)_{\text{in}}𝒩(M)_{\text{out}}e^{2\pi M\sqrt{\frac{k}{2}}}e^{\frac{1}{2}M\beta _{\mathrm{Hw}}},`$ (92)
viz. both are in effective thermal distribution set by the Hawking temperature. All spectral moments are manifestly ultraviolet finite since, at large $`M`$, exponential growth of the density of the final closed string states is insufficient to overcome the suppression by the distribution. Thus,
$`{\displaystyle \frac{𝒩_{\mathrm{out}}(M)\rho (M)}{𝒩_{\mathrm{in}}(M)\rho (M)}}|_{\mathrm{falling}\mathrm{D0}}1.`$
We interpret this as indicating that the D0-brane does not radiate off most of its energy before falling into the horizon.
This special case corresponds to empty $``$. The two-dimensional background permits no transverse degrees of freedom of the string. The physical spectrum includes massless tachyon only, with $`M=0`$ and $`\rho (M)=1`$. We now have a crucial difference from the previous cases for the on-shell configurations. The radial momentum $`p`$ is fixed by the on-shell condition as $`\omega =\pm p`$, so it should not be integrated over for the final states. Consequently, we cannot decompose the radiation distribution into incoming and outgoing radiations, and only the total distribution is physically relevant.
We thus obtain the following large $`\omega `$ behavior of the radiation distribution:
$`𝒩(\omega )e^{2\pi \omega }e^{\omega \beta _{\mathrm{Hw}}}.`$ (93)
Again, we have found effective thermal distribution at the Hawking temperature! Notice the absence of extra 1/2-factor in contrast to the previous regimes. This is not a contradiction. In the present case, the transverse oscillators are absent and the string behaves as a point particle. Again, the D0-brane does not radiate off most of its energy before falling across the black hole horizon.
### 4.2 Radiation distribution in bosonic string theory
The analysis for the bosonic string case proceeds quite the same route. The boundary state for the infalling D0-brane includes the string worldsheet correction factor $`\mathrm{\Gamma }\left(1+i\frac{p}{\kappa 2}\right)`$, where again $`\kappa `$ refers to the level of bosonic $`SL(2;)/U(1)`$ coset model. The on-shell condition now reads
$`{\displaystyle \frac{\omega ^2}{4\kappa }}+{\displaystyle \frac{p^2}{4(\kappa 2)}}+{\displaystyle \frac{1}{4(\kappa 2)}}+\mathrm{\Delta }_{}=1,`$ (94)
where $`\mathrm{\Delta }_{}`$ denotes the conformal weight in the $``$-sector. This is solved by
$`\omega \omega (p,M)=\sqrt{{\displaystyle \frac{\kappa }{\kappa 2}}p^2+2\kappa M^2}\text{where}M^22\left(\mathrm{\Delta }_{}+{\displaystyle \frac{1}{4(\kappa 2)}}1\right).`$ (95)
The partial radiation number distribution at large $`M`$ limit is given by:
$`𝒩(M)_{\text{in}}{\displaystyle \frac{1}{M}}{\displaystyle _0^{\mathrm{}}}\text{d}pe^{+\pi \left(1\frac{1}{\kappa 2}\right)p\pi \sqrt{\frac{\kappa }{\kappa 2}p^2+2\kappa M^2}}.`$ (96)
$`𝒩(M)_{\text{out}}{\displaystyle \frac{1}{M}}{\displaystyle _0^{\mathrm{}}}\text{d}pe^{\pi \left(1+\frac{1}{\kappa 2}\right)p\pi \sqrt{\frac{\kappa }{\kappa 2}p^2+2\kappa M^2}}.`$ (97)
Thus, as in the superstring case, there can arise several distinct behaviors depending on how stringy the background is.
Consider first the incoming radiation part. Since $`1\frac{1}{\kappa 2}>0`$, the dominant contribution to the momentum integral in (96) is from the saddle point
$`pp_{}={\displaystyle \frac{\kappa 3}{\sqrt{2\frac{1}{2(\kappa 2)}}}}M.`$
We thus obtain, up to pre-exponential powers of $`M`$,
$`𝒩(M)_{\text{in}}e^{2\pi M\sqrt{2\frac{1}{2(\kappa 2)}}}=e^{\frac{1}{2}M\beta _{\mathrm{Hg}}},`$
where $`\beta _{\mathrm{Hg}}`$ denotes the Hagedorn temperature of the bosonic string theory (83). In this way, we again find the power-law behavior of $`\rho (M)𝒩(M)_{\text{in}}`$ at large $`M`$, independent of the level $`\kappa `$.
For the outgoing radiation part, again the $`p0`$ dominates the momentum integral in (97). The result is
$`𝒩(M)_{\text{out}}e^{2\pi M\sqrt{\frac{\kappa }{2}}}=e^{\frac{1}{2}M\beta _{\mathrm{Hw}}}.`$
Here,
$`\beta _{\mathrm{Hw}}2\pi \sqrt{2\kappa }`$
is the Hawking temperature of the bosonic two-dimensional black hole. We then obtain
$`\rho (M)𝒩(M)_{\text{out}}e^{\frac{1}{2}\left(\beta _{\mathrm{Hg}}\beta _{\mathrm{Hw}}\right)M}=e^{2\pi M\left[\sqrt{2\frac{1}{2(\kappa 2)}}\sqrt{\frac{\kappa }{2}}\right]}.`$
As in the superstring case, the exponent is always negative definite:
$`\left(2{\displaystyle \frac{1}{2(\kappa 2)}}\right){\displaystyle \frac{\kappa }{2}}={\displaystyle \frac{\left(\kappa 3\right)^2}{2(\kappa 2)}}0.`$
so the outgoing radiation distribution (as well as spectral moments) is manifestly ultraviolet finite.
Physical interpretation of the above results is the same as the superstring case: The D0-brane falling into the black hole has nonzero component $`T_{t\rho }`$ of the energy-momentum tensor, and entails that dominant part of the closed string radiation is incoming toward the future horizon. The outgoing part of the radiation is exponentially suppressed, and is in effective thermal distribution set by the Hawking temperature. Again, this distribution is distinct from the Hawking radiation of the two-dimensional black hole. As for the fermionic string, the branching ratio is exponentially suppressed.
In this regime, $`1\frac{1}{\kappa 2}<0`$ and the momentum integrals for both incoming and outgoing radiation distributions are dominated by $`p0`$:
$`𝒩(M)_{\text{in}}𝒩(M)_{\text{out}}e^{2\pi M\sqrt{\frac{\kappa }{2}}}e^{\frac{1}{2}M\beta _{\mathrm{Hw}}}.`$
Both are in effective thermal distribution at the Hawking temperature, and all spectral moments are manifestly ultraviolet finite since, at large $`M`$, the growth of the density of state does not overcome the suppression by the distribution. The branching ratio remains order unity.
This is the most familiar situation: black hole in two-dimensional bosonic string theory, originally studied in . The physical spectrum of closed string consists only of the massless tachyon, so we again need to set $`M=0`$ and $`\rho (M)=1`$. The calculation is slightly more complicated than the supersymmetric case: The canonically normalized energy is
$`E={\displaystyle \frac{\sqrt{2}}{3}}\omega =\sqrt{2}p,`$
so we obtain
$`𝒩(E)e^{\pi \left(1\frac{1}{1/4}\right)p\pi \sqrt{\frac{9/4}{1/4}}p}=e^{3\sqrt{2}\pi E}e^{E\beta _{\mathrm{Hw}}}.`$
It again shows effective thermal distribution of the radiated closed string modes at the Hawking temperature: $`\beta _{\mathrm{Hw}}=2\pi \sqrt{2\kappa }=3\pi \sqrt{2}`$.
### 4.3 Radiation distribution for emitted or time-symmetric boundary states
The closed string radiations for the other types boundary states, viz. the ‘emitted’ (69) or the ‘symmetric’ (73) D0-branes, can be studied analogously.
For the emitted D0-brane boundary state (69), by the time-reversal, we should observe the radiation distribution at the far past: $`t\mathrm{}`$. The relevant decomposition corresponding to (77) is given by (assuming $`\omega >0`$, $`p>0`$)
$`V_\omega ^p(\rho ,t)e^{i\omega \mathrm{ln}\rho i\omega t}+d^{}(p,\omega )e^\rho e^{ip\rho i\omega t},`$ (98)
where the first term is supported near the past horizon and the second term corresponds to the incoming wave from the null infinity. Obviously we find precisely the same behavior of the radiation distribution as the absorbed D0-brane once the role of ‘in’ and ‘out’ states are reversed. So, for $`k>1`$, $`𝒩(M)_{\mathrm{in}}\mathrm{exp}(\frac{1}{2}\beta _{\mathrm{Hw}}M)`$ while $`𝒩(M)_{\mathrm{out}}\mathrm{exp}(\frac{1}{2}\beta _{\mathrm{Hg}}M)`$ and, for $`1k>1/2`$, $`𝒩(M)_{\mathrm{in}}`$, $`𝒩(M)_{\mathrm{out}}\mathrm{exp}(\frac{1}{2}\beta _{\mathrm{Hw}}M)`$.
Consider next the boundary state describing D0-brane in symmetric boundary condition (73). Recalling the relations (72), one finds that the radiation rates are simply obtained by adding contributions from ‘absorbed’ and ‘emitted’ D0-brane boundary states. So, the radiation distributions behave as $`𝒩(M)_{\mathrm{in}}`$, $`𝒩(M)_{\mathrm{out}}\mathrm{exp}(\frac{1}{2}\beta _{\mathrm{Hg}}M)`$ for $`k>1`$ and the dependence on Hawking temperature disappeared.<sup>19</sup><sup>19</sup>19Dependence on the Hawking temperature exponentially suppressed, so completely negligible compared to other power-suppressed subleading terms. We then find that the ‘detailed balance’ $`𝒩(M)_{\text{in}}=𝒩(M)_{\text{out}}`$ is obeyed. This is as expected since the boundary state (73) is defined so that it keeps the time-reversal symmetry and the one-particle state unitarity manifest.
### 4.4 Revisit to the radiation distribution: thermal string propagator
We shall revisit the radiation distribution and discuss salient features of the distribution from another different angle. Argument we shall present here would be somewhat heuristic, but we feel it quite helpful for grasping physical intuition and for understanding how the effective thermal behavior of the radiation comes about. This argument is similar to the one given in , where thermal distribution of the Hawking radiation in the two-dimensional black hole background was observed via the closed string thermal propagator. Our foregoing discussion is an extension of theirs to the open string sector.
Consider the thermal cylinder amplitude for the D1-brane on the Euclidean cigar (28).<sup>20</sup><sup>20</sup>20To be more precise, we consider the fermionic black-hole of level $`k`$ and focus on the space-time bosons. If the space-time fermions are considered, the thermal Kaluza-Klein momenta should be half integer-valued $`n1/2+\text{Z}`$ instead of being integer-valued $`n\text{Z}`$ as for the bosons. This change leads to the Fermi-Dirac distribution $`(e^{\beta _{\text{Hw}}\omega _{p,M}}+1)^1`$ instead of the Bose-Einstein $`(e^{\beta _{\text{Hw}}\omega _{p,M}}1)^1`$ in the following argument. Schematically, the amplitude is evaluated as (we omit the parameters $`\rho _0`$, $`\theta _0`$ for simplicity)
$`𝒜_{\text{cylinder}}^{(E)}`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}\text{d}T{}_{D1}{}^{}<B\left|e^{\pi TH^{(c)}}\right|B>_{D1}{\displaystyle \underset{M}{}}{\displaystyle \underset{n}{}}{\displaystyle }{\displaystyle \frac{\text{d}p}{2\pi }}{\displaystyle \frac{1}{p^2+\left(\frac{2\pi n}{\beta _{\text{Hw}}}\right)^2+M^2}}\sqrt{\rho _c(M)}|\mathrm{\Psi }_{D1}(p,n)|^2`$ (99)
$`=`$ $`\beta _{\text{Hw}}{\displaystyle \underset{M}{}}{\displaystyle \frac{\text{d}p}{2\pi }\frac{\text{d}q}{2\pi }\sqrt{\rho _c(M)}\frac{\left|\mathrm{\Psi }_{D1}(p,\frac{q}{2\pi }\beta _{\text{Hw}})\right|^2}{p^2+q^2+M^2}\left(1+\underset{m_{>0}}{}e^{i\beta _{\text{Hw}}mq}+\underset{m_{>0}}{}e^{i\beta _{\text{Hw}}mq}\right)}.`$
Here, $`p`$ is the radial momentum, $`n`$ is the Kaluza-Klein momentum around the asymptotic circle (thermal circle) of the cigar geometry, $`M`$ is again the transverse mass in the $``$-sector, $`\rho _c(M)`$ is the density of the closed string states, and $`\beta _{\text{Hw}}2\pi \sqrt{2k}`$ is the inverse Hawking temperature. We now Wick rotate the cylinder amplitude by contour deformation of the $`q`$-integration in the manner similar to . By formal manipulation $`qi\omega `$ <sup>21</sup><sup>21</sup>21Here, to make comparison with easier, we normalized $`\omega `$, $`p`$ as $`L_0=\frac{1}{2}\omega ^2+\frac{1}{2}p^2+\mathrm{}`$, rather than $`L_0=\frac{1}{4k}\omega ^2+\frac{1}{4k}p^2+\mathrm{}`$. and $`𝒜_{\text{cylinder}}^{(E)}i𝒜_{\text{cylinder}}^{(L)}`$, we obtain
$`𝒜_{\text{cylinder}}^{(L)}`$ $``$ $`\beta _{\text{Hw}}{\displaystyle \underset{M}{}}{\displaystyle \frac{\text{d}p}{2\pi }\sqrt{\rho _c(M)}\left[\frac{\text{d}\omega }{2\pi }\frac{\left|\mathrm{\Psi }_{\mathrm{D1}}(p,\frac{i\omega }{2\pi }\beta _{\text{Hw}})\right|^2}{p^2+M^2\omega ^2+iϵ}\frac{2\pi i}{\omega _{p,M}}\frac{\left|\mathrm{\Psi }_{D1}(p,\frac{i\omega _{p,M}}{2\pi }\beta _{\text{Hw}})\right|^2}{e^{\beta _{\text{Hw}}\omega _{p,M}}1}\right]},`$ (100)
where we denoted $`\omega _{p,M}\sqrt{p^2+M^2}`$ for the on-shell energy and used the identity
$$\left|\mathrm{\Psi }_{D1}(p,\frac{i\omega _{p,M}}{2\pi }\beta _{\text{Hw}})\right|^2=\left|\mathrm{\Psi }_{D1}(p,\frac{i\omega _{p,M}}{2\pi }\beta _{\text{Hw}})\right|^2.$$
Since $`\left|\mathrm{\Psi }_{D1}(p,\frac{i\omega }{2\pi }\beta _{\text{Hw}})\right|^2`$ is proportional to $`\mathrm{exp}(\frac{1}{2}\beta _{\text{Hw}}|\omega |)`$, the first term (including the Feynmann propagator) gives rise to an ultraviolet divergent contribution. This is not surprising and reveals the reason why the naive Wick-rotation of (28) is not viable. The second term ($`𝒜_{\text{thermal}}^{(L)}`$) exhibits an effective thermal distribution. More pertinently, this term contributes to the imaginary part of the thermal cylinder amplitude we are interested in. Indeed, it yields the anticipated behavior:
$`\text{Im}𝒜_{\text{thermal}}^{(L)}`$ $``$ $`{\displaystyle \frac{1}{\omega _{p,M}}}{\displaystyle \frac{1}{e^{\beta _{Hw}\omega _{p,M}}1}}\sqrt{\rho _c(M)}\left|\mathrm{\Psi }_{D1}(p,{\displaystyle \frac{i\omega }{2\pi }}\beta _{\text{Hw}})\right|^2`$ (101)
$``$ $`{\displaystyle \frac{\sqrt{\rho _c(M)}}{\omega _{p,M}}}\sigma (p)e^{\frac{1}{2}\beta _{\text{Hw}}\omega _{p,M}},`$
and reproduces the previous results (87) including the correct greybody factor $`\sigma (p)`$ and the density of the radiated closed string states $`\sqrt{\rho _c(M)}\rho (M)`$. Recall that, in our construction of the Lorentzian boundary states, the damping factor was crucial, which reads in the present conventions as $`\mathrm{sinh}(\pi \sqrt{2k}p)\mathrm{cosh}^1(\pi \sqrt{\frac{k}{2}}(p+\omega ))\mathrm{cosh}^1(\pi \sqrt{\frac{k}{2}}(p\omega ))`$. At large $`\omega `$ or large $`p`$, this damping factor shows the same asymptotic behavior as the Boltzmann distribution function $`(e^{\beta _{\text{Hw}}\omega }\pm 1)^1`$. In this sense, our prescription of Wick-rotating the Euclidean boundary states would be roughly identified with the prescription of keeping only the finite second term in (100). This then explains origin of the effective thermal distribution as derived from the Lorentzian boundary states.
As yet another viewpoint, consider the thermal cylinder amplitude (99) in the open string channel. For simplicity, concentrate on the asymptotic region $`\rho 0`$. The hairpin D1-brane (28) appears just as two halves of the $`D1`$-$`\overline{D}1`$ system, which put open strings around the thermal circle to obey Dirichlet boundary condition (so, identified as the ‘s$`D`$-s$`\overline{D}`$ system’ ), as pointed out in . In this set up, for simple kinematical reasons, we find on-shell closed string states in the cylinder amplitude, while only off-shell states in the open string channel. As discussed e.g. in , using the modular transformation, it can be shown that the thermal distribution of physical closed string states emitted/absorbed by the s$`D`$-s$`\overline{D}`$ system is captured by the unphysical open string winding modes along the thermal circle.<sup>22</sup><sup>22</sup>22This is a simple extension of the standard argument concerning the thermal toroidal partition functions . For instance, the Hagedorn behavior is interpretable as the tachyonic instability due to the unphysical winding modes around the thermal circle. Especially, unit of the winding energy should determine temperature of the thermal distribution of closed string states coupled with the s$`D`$-s$`\overline{D}`$ system. In the present case, it is identified with the interval of the hairpin $`(=\frac{1}{2}\beta _{\text{Hw}})`$, which is just associated to the open string stretched between $`D1`$-$`\overline{D}1`$. (Notice that, taking suitably the GSO projection into account, we can check that the zero winding modes, i.e. the $`D1`$-$`D1`$ or $`\overline{D}1`$-$`\overline{D}1`$ strings, are canceled out. See .) This would be the simplest explanation for the reason we get the effective thermal distribution $`\mathrm{exp}(\frac{1}{2}\beta _{\text{Hw}}\omega _{p,M})`$ from the cylinder amplitude (99).
As already noted in footnote 9, all the regular geodesics of the $`D0`$-brane motion are just straight lines in the Kruskal coordinates. Once Wick-rotated back to the hairpin profiles of Euclidean D1-brane, this means that they all have the same interval $`\frac{1}{2}\beta _{\text{Hw}}`$ around the thermal circle. This observation leads us again to the same effective thermal behavior (87) characterized by the Hawking temperature (before integrating $`p`$ out),<sup>23</sup><sup>23</sup>23One might ask why the D0-brane motion with different ‘temperature’ is not considered. Such case corresponds to singular hairpin profiles and hence to singular Lorentzian trajectories of the D0-brane. They cannot be solutions of the D0-brane’s DBI action because of discontinuity of the velocity at the singular points. Quite interestingly, this feature is strikingly similar to the original Hawking’s prescription for black hole temperature: demanding the Euclidean geometry smooth, we can fix asymptotic periodicity of the Euclidean time and read off the temperature characterizing the radiation from the black-hole. as is already pointed out.
## 5 ‘String - Black Hole’ Transition
It has been a recurrent theme that an elementary particle or a string is a black hole: a configuration consisting of (multiple) strings with high enough total mass is equivalent to a black hole of the same mass and other conserved charges. This brings a question whether a given configuration is most effectively described in terms of strings or black holes. By the string - black hole transition, we will refer to such change of the effective description for a configuration involving massive string excitations. Roughly speaking, the string is dual to the black hole and vice versa.
An immediate, interesting question is whether the two-dimensional black hole geometries is also subject to the string - black hole transition and if so what precisely the dual of the geometries would be. In this section, we shall investigate this transition by studying rolling dynamics of a D0-brane placed on the background. If the background undergoes the transition between the black hole and the string configurations, propagation of a probe D0-brane would be affected accordingly. The transition is triggered by $`k`$ or $`\kappa `$, which measures characteristic curvature scale of the background measured in sting unit and hence string worldsheet effects. We shall explore a signal of the transition by examining spectral distribution of the closed string radiation out of the rolling D0-brane. Other physical observables associated with D0-brane would certainly be equally viable probes. Though straightforward to analyze, in this work, we shall not consider them.
### 5.1 Probing ‘string - black hole’ transition via D-brane
In the previous section, we observed that $`𝒩(M)_{\text{in}}𝒩(M)_{\text{out}}`$ for both the supersymmetric and bosonic string theories in case the string worldsheet effects are weak enough, viz. $`k>1`$ and $`\kappa >3`$, respectively. Obviously, such behavior is interpretable as indicating that the background on which the radiative process takes place is indeed a black-hole: D0-brane falls into the horizon and subsequent radiation is mostly absorbed by the black hole. On the other hand, the behavior that $`𝒩(M)_{\text{in}}𝒩(M)_{\text{out}}\rho (M)^1`$ for $`k<1`$ or $`\kappa <3`$ does not seem to bear features present in the black hole background: while D0-brane falls inward, subsequent radiation is not mostly absorbed by the black hole but disperse away. Since this is the regime where the string worldsheet effects are significant, the background may be described most effectively in terms of strings. We are thus led to conclude that the background, whose stringy effects are controlled by the parameter $`k`$ or $`\kappa `$, would make a phase-transition between the black hole and the string across $`k=1`$ or $`\kappa =3`$. In a different physical context, this so-called ‘black hole-string transition’ was studied recently . What distinguishes our consideration and result from is that we are probing possible phase-transition of the (closed string) background by introducing a D0-brane in it and studying open string dynamics.
Possible existence of such a phase transition was first hinted in in the closed string sector, where they observed that the $`𝒩=2`$ Liouville superpotential becomes normalizable once $`k>1`$ and it violates the Seiberg bound. Recall that the marginal interaction term is
$`S^\pm =\psi ^{}e^{\frac{1}{𝒬}(\varphi \pm iY)},(𝒬=\sqrt{2/k})`$ (102)
for the $`𝒩=2`$ Liouville theory, and
$`S^\pm =e^{\frac{1}{𝒬}(\varphi \pm \sqrt{1+𝒬^2}iY)}e^{\sqrt{\frac{\kappa 2}{2}}\varphi \sqrt{\frac{\kappa }{2}}iY},(𝒬=\sqrt{2/(\kappa 2)}),`$ (103)
for the bosonic sine-Liouville theory, respectively. Both interactions are normalizable (exponentially falling off in the asymptotic far region) if the curvature is sufficiently small that $`k>1`$ or $`\kappa >3`$ is satisfied. As is well-known, $`𝒩=2`$ Liouville or sine-Liouville theory is T-dual to the $`SL(2;)/U(1)`$ coset theory , so the condition on the level $`k`$ or $`\kappa `$ ought naturally to descend to the two-dimensional black hole description. Indeed, such aspect was discussed in purely in the language of the $`SL(2;)/U(1)`$ coset theory (see also ). Their reasoning is closely related to the non-formation of the black hole in two-dimensional string theory (see also for the discussion concerning this issue from the matrix model viewpoint).<sup>24</sup><sup>24</sup>24Another interesting observation related to the $`k=1`$ transition is the following. If we consider a two-dimensional $`U(1)`$ gauge theory in the ultraviolet that flows to $`SL(2;)/U(1)`$ coset theory in the infrared (as was introduced in to prove the mirror duality to the $`𝒩=2`$ Liouville theory), the central charge of the $`U(1)`$ gauge theory is given by $`9`$. Since the IR $`SL(2;)/U(1)`$ coset theory has a central charge $`c=3(1+\frac{2}{k})`$, there is an apparent contradiction to Zamolodchikov’s $`c`$-theorem if the level $`k<1`$ is considered. However, we should note that $`SL(2;)/U(1)`$ coset theory is dilatonic so that the effective central charge is always given by $`3`$. In the strong curvature regime, $`k<1`$, the background is described more effectively in terms of the $`𝒩=2`$ Liouville theory as it is weakly coupled. Evidently, the black hole interpretation of the $`SL(2;)/U(1)`$ theory is less clear in this region, because the classical $`𝒩=2`$ Liouville theory does not admit an interpretation in terms of black hole geometry in any obvious way.
We emphasize that such string - black hole transition is not likely to arise perturbatively and could arise only from nonperturbative string worldsheet effects. For instance, tree-level closed string amplitudes are manifestly analytic with respect to the level $`k`$. These amplitudes exhibit a finite absorption rate (thus displaying the non-unitarity of the reflection amplitudes) regardless of the value of $`k`$. In fact, finite-$`k`$ correction to the amplitudes yield an irrelevant phase-factor .
However, as was first observed in , situation changes drastically if we consider the closed string radiation from the rolling D-brane in such a background. In , it was shown that the distribution of radiation off D0-brane in extremal NS5-brane background becomes ultraviolet finite for $`k<1`$. In the previous section, extending the analysis of , we have shown that the $`k=1`$ transition shows up manifestly in the open string sector in the sense that branching ratio between the incoming and the outgoing radiation distribution (as well as spectral moments) behaves very differently across $`k=1`$. Remarkably, retaining finite $`1/k`$-correction, which originated from consistency with the exact reflection relations, was crucial in obtaining physically sensible results even for $`k1`$. Cancellation between the radiation distribution and the exponential growth of the density of states at large $`M`$ is quite nontrivial, and relied crucially on precise functional dependence on $`k`$.
An ‘order-parameter’ of the transition is thus provided by the radiation distribution of rolling D-brane. The phase transition across $`k=1`$ is that while the radiation distribution from the falling D-brane exhibits powerlike ultraviolet divergence for $`k>1`$, it becomes finite for $`k<1`$. Thus, the rolling D-brane in the $`k<1`$ regime does not yield a large back-reaction unlike the $`k>1`$ case. This is also consistent with the assertion that black hole cannot be formed in the two-dimensional string theory: It seems difficult to construct two-dimensional black hole by injecting D-branes to the linear dilaton (or usual Liouville) theory.<sup>25</sup><sup>25</sup>25Such a possibility was proposed in .
It is also worth mentioning that the radion-tachyon correspondence is likely to fail in the two-dimensional string theory ($`k=1/2`$). In fact, had we have such a correspondence, the rolling radion of the D0-brane could be identified with the rolling tachyon of the ZZ-brane in the Liouville theory. On the other hand, it is known that the radiation distribution of the-ZZ brane exhibits a powerlike ultraviolet divergence at leading order in string perturbation theory, while that of the falling D0-brane does not.
### 5.2 Holographic Viewpoint
The string - black hole transition across $`k=1`$ also has a natural interpretation in terms of the holographic principle, as recently discussed in . Adding $`Q_1`$ fundamental strings to $`k`$ NS5-branes, one obtains the familiar bulk geometry of the $`AdS_3/CFT_2`$-duality. In this context, the density of states of the dual conformal field theory is given by the naive Cardy formula $`S=2\pi \sqrt{\frac{cL_0}{6}}+2\pi \sqrt{\frac{\overline{c}\overline{L}_0}{6}}`$ with $`c=6kQ_1`$ for $`k>1`$, but not for $`k<1`$. Rather, the central charge that should be used in the Cardy formula is replaced by an effective one $`c_{\mathrm{eff}}=6Q_1(2\frac{1}{k})`$ . The similar effects also showed up in the double scaling limit of the ‘little string theory’(LST) .<sup>26</sup><sup>26</sup>26Even though the original ‘little string theory’ is the theory of NS5-brane, so $`k`$ should be positive integer-valued, one can also consider models with fractional value of the level $`k`$, which is less than 1 generically. This is achieved by considering the wrapped NS5-brane backgrounds, or compactifications on a Calabi-Yau threefold having rational singularity . From the regularized torus partition function, one can prove that there is no normalizable massless states (corresponding to the ‘Lehmann-Symanzik-Zimmerman-poles’ ) in such string vacua if $`k<1`$, as was discussed in e.g. . We shall now show that such change of the central charge is also imperative for reproducing the closed string radiation distribution correctly from the dual holographic picture.
It is an interesting attempt to reproduce the phase transition in the radiation distribution of rolling D-brane across $`k=1`$ from the holographic viewpoint. In , it was proposed that the rolling D-brane should correspond to the decay of a certain defect in the dual LST. We shall now extend that analysis to the $`k<1`$ case and explore the phase-transition. The relevant holographic description is based on the following two assumptions.
1. fixed radiation number distribution: The radiation distribution for a fixed mass $`M`$ is determined by large $`k`$ behavior of the pressure in the far future (past). This is equivalent to the statement that the decay of the radion is described by a ‘holographic tachyon condensation’. We assume that there is no phase transition at $`k=1`$ for a fixed mass $`M`$.<sup>27</sup><sup>27</sup>27Theoretically, there is no reason to exclude a finite $`1/k`$ correction here. We only need this assumption phenomenologically in order to reproduce the ten-dimensional calculation even for $`k>1`$. A priori, the tachyon condensation (in the critical bosonic string) itself may receive large string worldsheet corrections. In the Dirac-Born-Infeld action analysis, such potential corrections were completely dropped out. In our convention, the distribution is given by
$$𝒩(M)_{\mathrm{LST}}e^{2\pi M\sqrt{\frac{k}{2}}}.$$
(104)
2. change of density of states: The final density of closed little string states in the ‘holographic tachyon condensation’ is given by the square root of the full nonperturbative density of states in LST. As is discussed in , the full nonperturbative density of states of the LST is believed to exhibit a phase transition at $`k=1`$: for $`k>1`$, the density of states is related to the Hawking temperature as
$$n(M)_{\mathrm{LST}}e^{4\pi M\sqrt{\frac{k}{2}}}.$$
(105)
In other words, the Hagedorn temperature in LST should be equated with the Hawking temperature (see also, e.g. ).
On the other hand, for $`k<1`$, because of the non-normalizability of the black hole excitation, the nonperturbative density of states of the LST is equivalent to the density of states of the (dual) perturbative string theory :
$$n(M)_{\mathrm{LST}}e^{4\pi M\sqrt{1\frac{1}{2k}}}.$$
(106)
With these assumptions, we can estimate the average radiation number of the ‘holographic tachyon condensation’ to be
$`\overline{𝒩}_{\mathrm{LST}}={\displaystyle _0^{\mathrm{}}}\text{d}M𝒩(M)\sqrt{n(M)_{\mathrm{LST}}}.`$
Note that, in contrast to the bulk string theory calculation, we have no integration over the radial momentum. Substituting (104) and (105) or (106) according to the value of $`k`$, we obtain
$`\overline{𝒩}_{\mathrm{LST}}{\displaystyle ^{\mathrm{}}}\text{d}Me^{2\pi M\sqrt{\frac{k}{2}}+2\pi M\sqrt{\frac{k}{2}}}`$
for $`k>1`$, showing powerlike ultraviolet divergent behavior because of the complete cancellation in the exponent, and
$`\overline{𝒩}_{\mathrm{LST}}{\displaystyle ^{\mathrm{}}}\text{d}Me^{2\pi M\sqrt{\frac{k}{2}}+2\pi M\sqrt{1\frac{1}{2k}}},`$
for $`k<1`$, showing exponential suppression in the ultraviolet. It is easy to see that this holographic dual computation reproduces the bulk computation presented in section 4.1 up to a subleading power dependence (89), (91).<sup>28</sup><sup>28</sup>28The exact determination of the pre-exponential power part is beyond the scope of the rough estimate presented here. It requires the full computational ability in the LST.
It should be noted, however, that the cancellation between the radiation distribution and the density of states has a different origin in the dual holographic description as compared to the bulk side. In the holographic description, the origin of the phase transition is the nonperturbative density of the states in LST while the radiation distribution at a fixed mass-level $`M`$ keeps its functional form unchanged. On the other hand, in the bulk theory, origin of the cancellation was that the radiation distribution changes at $`k=1`$ due to the disappearance of the non-trivial saddle point in the integration of the radial momentum $`p`$, while the density of states is always given by the same formula. Thus the agreement between the two descriptions is quite non-trivial and we believe that our results provide yet another evidence of the holographic duality for the NS5-brane and blackhole physics.
Though we presented the dual description based on some assumptions, we can turn the logic around and regard our results as a support for such assumptions. In particular the quantum gravity phase transition at $`k=1`$ in the dual theory proposed in is crucial for understanding the radiation distribution out of a defect decay in the dual LST. We thus propose our discussion in this section as a strong support for string - black hole transition.
## 6 Boundary States and Radiation in the Ramond-Ramond Sector
In the case of fermionic black hole background, the rolling D0-brane would also radiate off closed string states in the Ramond-Ramond (R-R) sector. In this section, we shall construct R-R boundary state of the D0-brane and compute radiation rates. Since the worldsheet theory corresponds to $`𝒩=2`$ superconformal field theory, correlation functions of the R-R sector and boundary states are readily obtainable by performing the standard $`𝒩=2`$ spectral flow.
We shall begin with discussion regarding properties of reflection amplitudes for the R-R sector (see in the context of 2D black hole). Recall that the reflection relation was given in the NS-NS sector as
$`U_\omega ^p(\rho ,t)^{\mathrm{NS}}=^{\mathrm{NS}}(p,\omega )U_\omega ^p(\rho ,t)^{\mathrm{NS}}\text{and}V_\omega ^p(\rho ,t)^{\mathrm{NS}}=^{\mathrm{NS}}(p,\omega )V_\omega ^p(\rho ,t)^{\mathrm{NS}},`$
where the exact reflection amplitude $`^{\mathrm{NS}}(p,\omega )`$ was defined by
$`^{\mathrm{NS}}(p,\omega )={\displaystyle \frac{\mathrm{\Gamma }(1+\frac{ip}{k})\mathrm{\Gamma }(+ip)\mathrm{\Gamma }^2(\frac{1}{2}i\frac{p+\omega }{2})}{\mathrm{\Gamma }(1\frac{ip}{k})\mathrm{\Gamma }(ip)\mathrm{\Gamma }^2(\frac{1}{2}+i\frac{p\omega }{2})}}.`$
To obtain the reflection relation of the R-R sector, we shall perform the spectral flow by half unit of the $`𝒩=2`$ $`U(1)`$ current.
In sharp contrast to the $`𝒩=2`$ Liouville theory, the reflection amplitude now depends on the spin structure of the R-R sector.<sup>29</sup><sup>29</sup>29This is because, in the $`𝒩=2`$ Liouville theory, the reflection amplitudes for the momentum modes have a symmetry under $`\omega \omega `$. Explicitly, the spectral flow is defined as $`\omega \omega \pm i`$, where the $`+`$ sign corresponds to spin ($`+,`$) states and $``$ sign corresponds to spin ($`,+`$) states (in the $`(\frac{1}{2},\frac{1}{2})`$ picture): in the $`\rho \mathrm{}`$ limit, they are described by $`S^\pm e^\rho e^{ip\rho i\omega t}`$ and the conformal weight is given by $`h=\frac{p^2\omega ^2+1}{4k}+\frac{1}{8}`$.
Therefore, for the R-R states with spin ($`+,`$), the exact reflection amplitudes become
$$^{\mathrm{R}+}(p,\omega )=\frac{\mathrm{\Gamma }(1+\frac{ip}{k})\mathrm{\Gamma }(+ip)\mathrm{\Gamma }^2(1i\frac{p+\omega }{2})}{\mathrm{\Gamma }(1\frac{ip}{k})\mathrm{\Gamma }(ip)\mathrm{\Gamma }^2(1+i\frac{p\omega }{2})}.$$
(107)
Equivalently, if we take spin ($`,+`$) R-R states, the exact reflection amplitudes become
$$^\mathrm{R}(p,\omega )=\frac{\mathrm{\Gamma }(1+\frac{ip}{k})\mathrm{\Gamma }(+ip)\mathrm{\Gamma }^2(i\frac{p+\omega }{2})}{\mathrm{\Gamma }(1\frac{ip}{k})\mathrm{\Gamma }(ip)\mathrm{\Gamma }^2(+i\frac{p\omega }{2})}.$$
(108)
It is important to notice that the latter amplitudes have a second order zero in the light-cone direction $`p=\omega >0`$ (recall that $`p>0`$ in our convention). Similarly, we could derive the reflection relation for $`(\pm ,\pm )`$ spin structure, but the resultant amplitudes are compatible only with the analytic continuation to the ‘winding time’ (in the interior of the singularity), so we would not delve into details anymore.
Consider next the boundary wave function of the R-R sector. For definiteness, we shall take the absorbed D0-brane (67) (We focus on the $`t_0=0`$ case for simplicity.)
$`{}_{\text{absorb}}{}^{}B,\mathrm{NS};\rho _0|={\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\text{d}p}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{\text{d}\omega }{2\pi }}\mathrm{\Psi }_{\text{absorb:NS}}(\rho _0;p,\omega ){}_{}{}^{\widehat{U}}p,\omega |,`$
where
$`\mathrm{\Psi }_{\text{absorb:NS}}(\rho _0;p,\omega )={\displaystyle \frac{\mathrm{\Gamma }(\frac{1}{2}i\frac{p+\omega }{2})\mathrm{\Gamma }(\frac{1}{2}i\frac{p\omega }{2})}{\mathrm{\Gamma }(1ip)}}\mathrm{\Gamma }\left(1+{\displaystyle \frac{ip}{k}}\right)\left[e^{ip\rho _0}{\displaystyle \frac{\mathrm{cosh}\left(\pi \frac{p\omega }{2}\right)}{\mathrm{cosh}\left(\pi \frac{p+\omega }{2}\right)}}e^{+ip\rho _0}\right].`$
The boundary wave functions of the R-R sector are then derived by applying the $`𝒩=2`$ spectral flow $`\omega \omega \pm i`$:
$`\mathrm{\Psi }_{\text{absorb:R}+}(\rho _0;p,\omega ){\displaystyle \frac{\mathrm{\Gamma }(i\frac{p+\omega }{2})\mathrm{\Gamma }(1i\frac{p\omega }{2})}{\mathrm{\Gamma }(1ip)}}\mathrm{\Gamma }\left(1+{\displaystyle \frac{ip}{k}}\right)\left[e^{ip\rho _0}+{\displaystyle \frac{\mathrm{sinh}\left(\pi \frac{p\omega }{2}\right)}{\mathrm{sinh}\left(\pi \frac{p+\omega }{2}\right)}}e^{+ip\rho _0}\right],`$
and
$`\mathrm{\Psi }_{\text{absorb:R}}(\rho _0;p,\omega )={\displaystyle \frac{\mathrm{\Gamma }(1i\frac{p+\omega }{2})\mathrm{\Gamma }(i\frac{p\omega }{2})}{\mathrm{\Gamma }(1ip)}}\mathrm{\Gamma }\left(1+{\displaystyle \frac{ip}{k}}\right)\left[e^{ip\rho _0}+{\displaystyle \frac{\mathrm{sinh}\left(\pi \frac{p\omega }{2}\right)}{\mathrm{sinh}\left(\pi \frac{p+\omega }{2}\right)}}e^{+ip\rho _0}\right],`$
for the two opposite spin structures. These boundary wave functions are of course consistent with the exact reflection amplitudes (107),(108).
From these boundary wave functions, we can deduce some physical properties of the boundary states in the R-R sector:
* For $`k>\frac{1}{2}`$, in the saddle point approximation of the radial momentum integral, radiation distribution of the R-R sector behaves the same as that of the NS-NS sector. In particular, the absolute value of the reflection amplitudes behave in the similar manner. Thus, the radiation distribution of the R-R sector is the same as that of the NS-NS sector.
* For $`k=\frac{1}{2}`$, viz. the two-dimensional black hole, considerable differences arise. Both boundary wave function and reflection amplitudes show singularity (or zero) when we take particular spin structure. It is not clear what the origin of these singularities of lightlike on-shell states $`p=\omega `$ would be. We note that some related discussions were given in .
* In the mini-superspace limit $`k\mathrm{}`$, the mass gap in the R-R sector vanishes. Therefore, it is well-posed to question radiation of the massless R-R states off the R-R charge. ¿From the boundary states given above, we observe that, assuming $`p,\omega >0`$, there is no lightlike pole in $`R+`$ state while there is a pole at $`p=+\omega `$ in the $`R`$ state. It is also interesting to note that, in the subleading contribution proportional to $`e^{+ip\rho _0}`$, the pole from the gamma function is cancelled by the zero in the $`\mathrm{sinh}(\pi \frac{p\omega }{2})`$ factor.
A possible interpretation is that, roughly speaking, R-R charge is localized on the incoming light-cone $`p=\omega `$.<sup>30</sup><sup>30</sup>30This is true only in the asymptotic region $`\rho \mathrm{}`$ since the distribution near $`\rho =0`$ is further related to the basis of Ishibashi states used in the expansion. In the case of ‘absorbed’ basis, there is no contribution from the past horizon. In addition, because the reflection amplitude vanishes in the $`R`$ sector, an observer at $`\rho \mathrm{}`$ do not detect any outgoing wave.
## 7 Back to Extremal NS5-Brane Background
By tuning off $`\mu 0`$, we are back to the extremal NS5-brane background. Roughly speaking, the extremal background is described by the free linear dilaton theory, but crucial differences from the non-extremal counterpart studied in this work are the followings:
* We have no reflection relation, and the $`p>0`$ and $`p<0`$ states should be treated as independent states.<sup>31</sup><sup>31</sup>31In this sense, the arguments given in are not completely precise, although the main part of physical results, say, the closed string radiation rates, are not altered.
* The conformal field theory description is not effective in the entire space-time: the string coupling diverges at the location of the NS5-brane. We cannot completely trace the classical trajectory of the D0-brane (32) without facing strong coupling problem.
We thus have to keep it in mind that the validity of the conformal field theory description of extremal NS5-brane is limited to the sufficiently weak string coupling region.
For the extremal NS5-brane, since the relevant conformal field theory involves linear dilaton and hence is a free theory, we can introduce the basis of the Ishibashi states as $`|p,\omega `$, $`(p,\omega )`$ associated with the wave function $`\psi _\omega ^p(\rho ,t)e^\rho e^{ip\rho i\omega t}`$. Another non-trivial difference from the non-extremal case is the volume form of the space-time. Since we have the linear dilaton $`\mathrm{\Phi }=\text{const}\rho `$ and a flat metric $`G_{ij}=\eta _{ij}`$, the relevant volume form becomes
$`\text{d}\text{Vol}=e^{2\mathrm{\Phi }}\sqrt{G}\text{d}\rho \text{d}t=e^{2\rho }\text{d}\rho \text{d}t.`$ (109)
Now, the classical trajectory of D0-brane in the extremal NS5-brane is given by :
$`2\mathrm{cosh}(tt_0)e^\rho =e^{\rho _0}.`$ (110)
The boundary state describing the D0-brane moving along (110) ought to have the following form:
$`B;\rho _0,t_0|={\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{\text{d}p}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{\text{d}\omega }{2\pi }}\mathrm{\Psi }(\rho _0,t_0;p,\omega )p,\omega |.`$ (111)
The boundary wave function is evaluated as
$`\mathrm{\Psi }(\rho _0,t_0;p,\omega )`$ $``$ $`{\displaystyle \text{d}v\delta \left(2\mathrm{cosh}(tt_0)e^\rho e^{\rho _0}\right)e^{\rho ip\rho i\omega t}}`$
$`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}\text{d}te^{ip\rho _0}e^{i\omega t}\left[2\mathrm{cosh}(tt_0)\right]^{ip1}`$
$`=`$ $`{\displaystyle \frac{1}{2}}B({\displaystyle \frac{1}{2}}i{\displaystyle \frac{p+\omega }{2}},{\displaystyle \frac{1}{2}}i{\displaystyle \frac{p\omega }{2}})e^{ip\rho _0i\omega t_0}.`$
In the last expression, we used the formula (A.2). This is essentially the calculation given in . Finally, by restoring the important ‘worldsheet correction factor’ $`\mathrm{\Gamma }\left(1+i\frac{p}{k}\right)`$,<sup>32</sup><sup>32</sup>32Since in this case we do not have the reflection relation, the inclusion of the factor $`\mathrm{\Gamma }\left(1+i\frac{p}{k}\right)`$ may sound less affirmative than the nonextremal NS5-brane background. We argue that the procedure is actually justified by considering the limit from the non-extremal case. we obtain the boundary wave function
$`\mathrm{\Psi }(\rho _0,t_0;p,\omega )={\displaystyle \frac{1}{2}}B(\nu _+,\nu _{})\mathrm{\Gamma }(1+i{\displaystyle \frac{p}{k}})e^{ip\rho _0i\omega t_0}.\text{where}\nu _\pm {\displaystyle \frac{1}{2}}i{\displaystyle \frac{p\pm \omega }{2}},`$
This is the extremal counterpart of the ‘symmetric D0-brane’ in the non-extremal NS5-brane background (73).
We can also consider the ‘half S-brane’ counterpart by taking the Hartle-Hawking contours depicted in the Figures 4 and 5. Namely, for the ‘absorbed brane’, we obtain
$`{}_{\text{absorb}}{}^{}B;\rho _0,t_0|=[{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\text{d}p}{2\pi }}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\text{d}\omega }{2\pi }}+{\displaystyle _{\mathrm{}}^0}{\displaystyle \frac{\text{d}p}{2\pi }}{\displaystyle _{\mathrm{}}^0}{\displaystyle \frac{\text{d}\omega }{2\pi }}]\mathrm{\Psi }(\rho _0,t_0;p,\omega )p,\omega |,`$ (113)
and for the ‘emitted brane’,
$`{}_{\text{emitted}}{}^{}B;\rho _0,t_0|=[{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\text{d}p}{2\pi }}{\displaystyle _{\mathrm{}}^0}{\displaystyle \frac{\text{d}\omega }{2\pi }}+{\displaystyle _{\mathrm{}}^0}{\displaystyle \frac{\text{d}p}{2\pi }}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\text{d}\omega }{2\pi }}]\mathrm{\Psi }(\rho _0,t_0;p,\omega )p,\omega |.`$ (114)
They are regarded as the counterparts of (118) and (119).
The radiation rates were already evaluated in .<sup>33</sup><sup>33</sup>33 In this paper, we scaled energy and momentum differently from . In light of normalization as in (84), $`\omega ,p`$ in this work should be read as $`2\sqrt{k}`$ times $`\omega ,p`$ in . Crucial differences from the non-extremal case are the followings: We have the ‘forward radiations’ (e.g., the incoming radiation for the absorbed D-brane (113)) only and no ‘backward radiations’ (e.g., the outgoing radiation for the absorbed D-brane). This is because there is no reflection relation in the extremal case. The forward radiations behave in the completely same way as the non-extremal case (that is, in a fermionic two-dimensional black hole with $`k>1`$), giving rise to the Hagedorn-like ultraviolet divergence again. At fixed but large $`M`$ before integrating over $`p`$, the partial radiation number distribution takes again exactly the same asymptotic form as in (87) except that now the coefficient $`2\pi \sqrt{2k}`$ is not interpretable as the inverse Hawking temperature of the black hole.<sup>34</sup><sup>34</sup>34An obvious alternative interpretation could be that, even for extremal background, the falling D0-brane excites the NS5-brane above the extremality. Again, this has to do with the peculiarity that the Hawking temperature of the two-dimensional black hole is set by the level $`k`$, not by the nonextremality $`\mu `$. On the other hand, the absence of the backward radiation matches with the extremality of the background; there is no Hawking radiation.
## 8 More on Physical Interpretations : Hartle-Hawking States
We shall now revisit the boundary states we constructed in this work and elaborate further on their physical interpretations with particular emphasis on analogy with the rolling tachyon problem via the radion-tachyon correspondence. We also elaborate on the fate of R-R charge carried by the D0-brane. To be concrete, we shall focus on the cases $`k2`$ admitting interpretation in terms of near horizon geometry of black NS5 branes.
The boundary state (67) describes the late-time rolling ($`tt_0`$) of the D0-brane rolling into the black NS5 branes. The relevant D0-brane has the initial condition $`\rho =\rho _0`$, $`\frac{d\rho }{dt}=0`$ at $`t=t_0`$ and starts to roll down toward the black hole. After sufficiently long coordinate time elapsed, the D0-brane gets close to the future horizon ($`^+`$). As examined in section 4, almost all energy of the D0-brane is absorbed by the black hole in the form of incoming radiation. The incoming radiation is dominated by very massive, and hence highly non-relativistic closed string excitations. Via the radion-tachyon correspondence, these states are identifiable with the ‘tachyon matter’ in the rolling tachyon problem in flat spacetime. On the other hand, we have seen that a small part of energy escapes to the spatial infinity ($`^+`$) as the outgoing radiation. We have seen that the spectral distribution is characterized by the Hawking temperature, and is necessarily dominated by light modes. This interpretation is quite natural from the viewpoint of the radion-tachyon correspondence for the extremal NS5-brane background . Since we are now working with the non-extremal NS5-brane background, our analysis may be considered as an evidence that the correspondence is valid even at finite temperature.
What about evolution in the far past $`t<t_0`$? Here, we face a subtlety. Recall that the boundary condition defining (67) does not allow contributions from the past horizon ($`^{}`$), namely, the basis of Ishibashi states $`|p,\omega ^U`$ does not reproduce the past half of the classical trajectory (32). Rather, the NS-NS sector of the D0-brane boundary wave function appears widely distributed in the space-time in the far past. This may be interpreted as radiations imploding to $`\rho =\rho _0`$ from spatial infinity, but then it is subtle to trace the R-R charge carried by the D0-brane, created out of the imploding radiation. Classically, the D0-brane charge density ought to be localized along the classical trajectory (32) and hence emanates from the past horizon. Once stringy effects are taken into account, the charge appears to originate from asymptotic infinity along the light-cone coordinate. Complete understanding of this curious feature is highly desirable but we shall relegate it to future study. Here, instead, we present a simple prescription of avoiding this subtlety: a version of ‘Hartle-Hawking’ boundary condition.
We shall first focus on the absorbed D0-brane boundary state (67). Formally, by construction, we can regard the boundary wave function specified by the time-integration over the ‘real contour’ $`𝒞=`$ as in (36). Now, let us discuss what happens if we choose the ‘Hartle-Hawking’ type contour instead of the real contour, which connect the Euclidean time with the future or past half of real time axis at $`t=t_0`$:
$`𝒞_{\text{future}}^\pm =\left(t_0+i_{}\right)\left(t_0+_+\right),𝒞_{\text{past}}^\pm =\left(t_0+i_{}\right)\left(t_0+_{}\right).`$ (115)
More precisely, we should avoid suitably the branch cuts on $`t_0+i`$ to render the integral convergent. See Figures 4 and 5 for details. The superscript $`+`$ $`()`$ is associated with the positive (negative) energy sector. Note that the phase-factor $`e^{i\omega t}`$ behaves well on the lower (upper) half of complex $`t`$-plane if $`\omega `$ is positive (negative). Let us pick up $`𝒞_{\text{future}}`$. Following the traditional interpretation of the Hartle-Hawking type wave function, we may suppose that both the D0-brane and black NS5-brane are created from ‘nothing’ at $`t=t_0`$, and then the D0-brane starts to fall down toward the future horizon along the classical trajectory (32). In this prescription, the subtlety we mentioned above is completely circumvented.
One may paraphrase the prescription as follows: choosing the Hartle-Hawking contour $`𝒞_{\text{future}}`$, we explicitly obtain
$`{}_{HH+,\text{absorb}}{}^{}B;\rho _0,t_0|={\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\text{d}p}{2\pi }}[{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\text{d}\omega }{2\pi }}\mathrm{\Psi }_{\text{symm}}(\rho _0,t_0;p,\omega )+{\displaystyle _{\mathrm{}}^0}{\displaystyle \frac{\text{d}\omega }{2\pi }}(p,\omega )\mathrm{\Psi }_{\text{symm}}^{}(\rho _0,t_0;p,\omega )]{}_{}{}^{\widehat{U}}p,\omega |,`$
(116)
where $`\mathrm{\Psi }_{\text{symm}}(\rho _0,t_0;p,\omega )`$ is defined in (73). In fact, by taking $`𝒞_{\text{future}}`$, we are only left with the $`L_\omega ^p`$ ($`R_\omega ^p`$)-part of the one-point function for the $`\omega >0`$ ($`\omega <0`$) sector. See the figure 4. This boundary wave function is formally regarded as the limit of (67) under $`t_0\mathrm{}`$, $`\rho _0+\mathrm{}`$ while keeping $`|\rho _0|/|t_0|`$ finite. Note that the second (first) term $`e^{ip\rho _0i\omega t_0}`$ ($`e^{ip\rho _0i\omega t_0}`$) in (67) oscillates very rapidly in this limit for $`\omega >0`$ ($`\omega <0`$) and hence drops off.<sup>35</sup><sup>35</sup>35More precise argument would be as follows: The disk amplitude for a wave packet e.g. $`\frac{\text{d}p}{2\pi }\frac{\text{d}\omega }{2\pi }f(p,\omega )|L_\omega ^p`$ is evaluated as $`lim_{\rho _0+\mathrm{},t_0\mathrm{}}\frac{\text{d}p}{2\pi }\frac{\text{d}\omega }{2\pi }f(p,\omega )\mathrm{\Psi }(\rho _0,t_0;p,\omega )`$. Then, the rapidly oscillating term in the boundary wave function $`\mathrm{\Psi }(\rho _0,t_0;p,\omega )`$ cannot contribute for any $`L^2`$-normalizable wave packet $`f(p,\omega )`$. The limit just means that the D0-brane moving along the trajectory (32) is coming from the past infinity $`(^{})`$, and falling into the future horizon ($`^+`$). Everything is supposed to be localized over the classical trajectory in this case.
Adopting the past Hartle-Hawking contour $`𝒞_{\text{past}}`$ for the boundary state of emitted D0-brane (69) is completely parallel. We take the time-reversal of the above:
$`{}_{HH,\text{emit}}{}^{}B;\rho _0,t_0|={\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\text{d}p}{2\pi }}[{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\text{d}\omega }{2\pi }}\mathrm{\Psi }_{\text{symm}}^{}(\rho _0,t_0;p,\omega )+{\displaystyle _{\mathrm{}}^0}{\displaystyle \frac{\text{d}\omega }{2\pi }}^{}(p,\omega )\mathrm{\Psi }_{\text{symm}}(\rho _0,t_0;p,\omega )]{}_{}{}^{\widehat{V}}p,\omega |,`$
(117)
which is regarded as the $`t_0+\mathrm{}`$, $`\rho _0+\mathrm{}`$ limit of (69). It describes the trajectory of D0-brane emitted from the past horizon $`^{}`$ and escaping to the future infinity $`^+`$.
Let us turn to the ‘symmetric’ D0-brane (73). Naively, it appears that the prescription is that
$`{}_{HH+,\text{symm}}{}^{}B;\rho _0,t_0|^{}={\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\text{d}p}{2\pi }}[{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\text{d}\omega }{2\pi }}\mathrm{\hspace{0.17em}2}\mathrm{\Psi }_{\text{symm}}(\rho _0,t_0;p,\omega ){}_{}{}^{L}p,\omega |+{\displaystyle _{\mathrm{}}^0}{\displaystyle \frac{d\omega }{2\pi }}\mathrm{\hspace{0.17em}2}\mathrm{\Psi }_{\text{symm}}^{}(\rho _0,t_0;p,\omega ){}_{}{}^{R}p,\omega \left|\right]`$
(118)
for the future Hartle-Hawking contour $`𝒞_{\text{future}}`$, and
$`{}_{HH,\text{symm}}{}^{}B;\rho _0,t_0|^{}={\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\text{d}p}{2\pi }}[{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\text{d}\omega }{2\pi }}\mathrm{\hspace{0.17em}2}\mathrm{\Psi }_{\text{symm}}^{}(\rho _0,t_0;p,\omega ){}_{}{}^{R}p,\omega |+{\displaystyle _{\mathrm{}}^0}{\displaystyle \frac{d\omega }{2\pi }}2\mathrm{\Psi }_{\text{symm}}(\rho _0,t_0;p,\omega ){}_{}{}^{L}p,\omega \left|\right]`$
(119)
for the past Hartle-Hawking contour $`𝒞_{\text{past}}`$. However, this cannot be the whole story. The existence of Euclidean part of the Hartle-Hawking path-integral enforces the boundary states to be expanded by the basis smoothly connected to the Euclidean ones, while $`|L_\omega ^p`$, $`|R_\omega ^p`$ do not possess such a property. Consequently, to achieve the correct Hartle-Hawking states, we ought to make further the projection to $`^U`$, ($`\widehat{^U}`$) for the contour $`𝒞_{\text{future}}`$, and to $`^V`$, ($`\widehat{^V}`$) for $`𝒞_{\text{past}}`$. We thus obtain as the correct Hartle-Hawking states:
$`{}_{HH+,\text{symm}}{}^{}B;\rho _0,t_0|={}_{HH+,\text{symm}}{}^{}B;\rho _0,t_0|^{}\widehat{P_U}{}_{HH+,\text{absorb}}{}^{}B;\rho _0,t_0|,`$
$`{}_{HH,\text{symm}}{}^{}B;\rho _0,t_0|={}_{HH,\text{symm}}{}^{}B;\rho _0,t_0|^{}\widehat{P_V}{}_{HH,\text{emitted}}{}^{}B;\rho _0,t_0|,`$ (120)
where the right-hand sides are already given in (116), (117).
Remarkably, this feature resembles much that of the S-branes discussed in . Namely, it was shown there that
$$\text{half S-brane}\text{full S-brane with the Hartle-Hawking contour}.$$
(121)
In our case, (73) corresponds to the full S-brane, while the Hartle-Hawking state (116) ((117)) is identifiable as the analogue of the half S-brane describing unstable D-brane decay (creation) process. The equalities (120) suggest that we have roughly identical relation to (121).
Notice that the parameters $`\rho _0`$, $`t_0`$ appear just as phase factors of boundary wave functions in (116), (117) contrary to (67), (69). Namely, the choice of parameters $`\rho _0`$, $`t_0`$ does not cause any physical difference for the Hartle-Hawking type states : They all can be regarded as describing the D0-brane moving from $`^{}`$ to $`^+`$ (from $`^{}`$ to $`^+`$) for (116) (for (117)) irrespective of $`\rho _0`$, $`t_0`$. These two parameters merely parameterize displacing the trajectory in two-dimensional black hole background. Similar feature comes about for the full S-brane with Hartle-Hawking contour as well: It is equivalent to the half S-brane not depending on any shift of the origin (the point connecting the real and imaginary times).
Finally, we remark a comment from the viewpoints of boundary conformal field theory: in contrast to the original ones (67), (69) and (73), the Hartle-Hawking boundary states (120) (or equivalently (116), (117)) are not compatible with the reflection relations. One may regard the boundary states (67) and (69) as the ‘completions’ of the Hartle-Hawking states (120) so that they satisfy the reflection relations.
Note Added : After this work was published, the work was brought to our attention. However, the work claimed results discrepant with ours concerning treatment of spectral amplitudes and the closed string radiations. We made in another publication further investigation from the open string channel. The result of completely supports the result given in this work, and clarifies where the errors originate in the work .
###### Acknowledgments.
We thank Changrim Ahn, Dongsu Bak, Tohru Eguchi, Yasuaki Hikida, Jaemo Park, Jongwon Park, Steve Shenker, Hiromitsu Takayanagi, Tadashi Takayanagi and Jung-Tay Yee for useful discussions. We are especially grateful to H. Takayanagi for his collaboration in the early stage of this work. SJR is supported in part by the KRF BK-21 Physics Divison, KRF Leading Scientist Grant, KOSEF SRC Program “Center for Quantum Spacetime” (R11-2005-021), and by F.W. Bessel Research Award from Alexander von Humboldt Foundation. YN was supported in part by JSPS Research Fellowships for Young Scientists. YS was supported by the Ministry of Education, Culture, Sports, Science and Technology of Japan.
## Appendix A Useful Formulae
$`{\displaystyle _{\frac{\pi }{2}}^{\frac{\pi }{2}}}(2\mathrm{cos}\theta )^{a1}e^{ib\theta }\text{d}\theta =\pi {\displaystyle \frac{\mathrm{\Gamma }(a)}{\mathrm{\Gamma }\left(\frac{1}{2}+\frac{a+b}{2}\right)\mathrm{\Gamma }\left(\frac{1}{2}+\frac{ab}{2}\right)}},(\text{Re}a>0,\left|\text{Re}b\right|<\text{Re}a+1),`$ (A.1)
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}(2\mathrm{cosh}t)^{a1}e^{ibt}\text{d}t={\displaystyle \frac{1}{2}}B({\displaystyle \frac{1}{2}}{\displaystyle \frac{a+ib}{2}},{\displaystyle \frac{1}{2}}{\displaystyle \frac{aib}{2}}){\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{1}{2}\frac{a+ib}{2}\right)\mathrm{\Gamma }\left(\frac{1}{2}\frac{aib}{2}\right)}{\mathrm{\Gamma }(1a)}},`$
$`(\text{Re}a<1,\left|\text{Im}b\right|<1\text{Re}a).`$ (A.2)
The integral (A.2) follows from the more general formula:
$`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\mathrm{cosh}(2at)}{\mathrm{cosh}^{2\beta }(pt)}}\text{d}t=4^{\beta 1}p^1B(\beta +{\displaystyle \frac{a}{p}},\beta {\displaystyle \frac{a}{p}}),(p>0,\text{Re}\left(\beta \pm {\displaystyle \frac{a}{p}}\right)>0),`$ (A.3)
given in . It is also derivable from (A.1) by contour deformation, as was shown in .
## Appendix B Proof of (29)<sup>36</sup><sup>36</sup>36The results of this section is due to H. Takayanagi.
Here we would like to evaluate explicitly the integral (29) for any $`\rho _0`$ (strictly speaking, we need to assume $`\mathrm{sinh}\rho _0>1`$). We begin with series expansion of the hypergeometric function in $`\varphi _n^p(\rho ,\theta )`$:
$`F({\displaystyle \frac{1}{2}}+{\displaystyle \frac{ip+n}{2}},{\displaystyle \frac{1}{2}}+{\displaystyle \frac{ipn}{2}};ip+1;{\displaystyle \frac{\mathrm{cos}^2\theta }{\mathrm{sinh}^2\rho _0}})`$
$`={\displaystyle \underset{\mathrm{}=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{\Gamma }(ip+1)}{\mathrm{\Gamma }(\frac{1}{2}+\frac{ip+n}{2})\mathrm{\Gamma }(\frac{1}{2}+\frac{ipn}{2})}}{\displaystyle \frac{\mathrm{\Gamma }(\frac{1}{2}+\frac{ip+n}{2}+\mathrm{})\mathrm{\Gamma }(\frac{1}{2}+\frac{ipn}{2}+\mathrm{})}{\mathrm{\Gamma }(ip+1+\mathrm{})}}{\displaystyle \frac{(1)^{\mathrm{}}}{\mathrm{}!}}\left({\displaystyle \frac{\mathrm{cos}\theta }{\mathrm{sinh}\rho _0}}\right)^2\mathrm{}.`$ (B.1)
Using the formula (A.1), we can perform, in the $`\mathrm{}`$-th sector, the integral (29) as
$`\mathrm{\Psi }_{\mathrm{}}=g(\mathrm{}){\displaystyle _{\frac{\pi }{2}}^{\frac{\pi }{2}}}\text{d}\theta e^{in\theta }(\mathrm{cos}\theta )^{ip1+2\mathrm{}}={\displaystyle \frac{g(\mathrm{})}{2^{ip1+2\mathrm{}}}}{\displaystyle \frac{\pi \mathrm{\Gamma }(ip1+2\mathrm{}+1)}{\mathrm{\Gamma }(\frac{1}{2}+\mathrm{}+\frac{ip+n}{2})\mathrm{\Gamma }(\frac{1}{2}+\mathrm{}+\frac{ipn}{2})}},`$
where $`g(\mathrm{})`$ refers to
$$g(\mathrm{})=\frac{(1)^{\mathrm{}}}{\mathrm{}!}(\mathrm{sinh}\rho _0)^{ip2\mathrm{}}\frac{\mathrm{\Gamma }(ip+1)}{\mathrm{\Gamma }(\frac{1}{2}+\frac{ip+n}{2})\mathrm{\Gamma }(\frac{1}{2}+\frac{ipn}{2})}\frac{\mathrm{\Gamma }(\frac{1}{2}+\frac{ip+n}{2}+\mathrm{})\mathrm{\Gamma }(\frac{1}{2}+\frac{ipn}{2}+\mathrm{})}{\mathrm{\Gamma }(ip+1+\mathrm{})}.$$
(B.2)
Then the total integral (29) is
$$\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}\mathrm{\Psi }_{\mathrm{}}=\frac{\pi }{2^{ip1}(\mathrm{sinh}\rho _0)^{ip}}\frac{\mathrm{\Gamma }(ip+1)}{\mathrm{\Gamma }(\frac{1}{2}+\frac{ip+n}{2})\mathrm{\Gamma }(\frac{1}{2}+\frac{ipn}{2})}\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}\frac{(1)^{\mathrm{}}}{\mathrm{}!}\frac{1}{2^2\mathrm{}(\mathrm{sinh}\rho _0)^2\mathrm{}}\frac{\mathrm{\Gamma }(ip+2\mathrm{})}{\mathrm{\Gamma }(ip+1+\mathrm{})}.$$
(B.3)
We can rewrite the summation into a hypergeometric function by using
$$\mathrm{\Gamma }(ip+2\mathrm{})=\frac{2^{ip1+2\mathrm{}}}{\sqrt{\pi }}\mathrm{\Gamma }\left(\frac{ip}{2}+\mathrm{}\right)\mathrm{\Gamma }\left(\frac{1}{2}+\frac{ip}{2}+\mathrm{}\right).$$
(B.4)
and then obtain
$$\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}\mathrm{\Psi }_{\mathrm{}}=\frac{\sqrt{\pi }}{(\mathrm{sinh}\rho _0)^{ip}}\frac{\mathrm{\Gamma }(\frac{ip}{2})\mathrm{\Gamma }(\frac{1}{2}+\frac{ip}{2})}{\mathrm{\Gamma }(\frac{1}{2}+\frac{ip+n}{2})\mathrm{\Gamma }(\frac{1}{2}+\frac{ipn}{2})}F(\frac{ip}{2},\frac{1}{2}+\frac{ip}{2};ip+1;\frac{1}{\mathrm{sinh}^2\rho _0}).$$
(B.5)
Making use of the formula
$`F(2\alpha ,2\beta ;\alpha +\beta +{\displaystyle \frac{1}{2}};z)=F(\alpha ,\beta ;\alpha +\beta +{\displaystyle \frac{1}{2}};4z(1z)).`$ (B.6)
$`|z|<{\displaystyle \frac{1}{2}},|z(1z)|<{\displaystyle \frac{1}{4}}`$
we find that
$$\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}\mathrm{\Psi }_{\mathrm{}}=\frac{\sqrt{\pi }}{(\mathrm{sinh}\rho _0)^{ip}}\frac{\mathrm{\Gamma }(\frac{ip}{2})\mathrm{\Gamma }(\frac{1}{2}+\frac{ip}{2})}{\mathrm{\Gamma }(\frac{1}{2}+\frac{ip+n}{2})\mathrm{\Gamma }(\frac{1}{2}+\frac{ipn}{2})}F(ip,ip+1;ip+1;\frac{1}{2}\frac{\mathrm{cosh}\rho _0}{2\mathrm{sinh}\rho _0}).$$
(B.7)
Note that the second and third arguments of the hypergeometric function are the same. The function is thus simplified as
$$F(ip,ip+1;ip+1;\frac{1}{2}\frac{\mathrm{cosh}\rho _0}{2\mathrm{sinh}\rho _0})=\left(\frac{\mathrm{sinh}\rho _0+\mathrm{cosh}\rho _0}{2\mathrm{sinh}\rho _0}\right)^{ip}=\left(2e^{\rho _0}\mathrm{sinh}\rho _0\right)^{ip},$$
(B.8)
because of the relation
$$(1z)^\alpha =F(\alpha ,\beta ;\beta ;z).$$
(B.9)
In this way, we finally obtain
$`{\displaystyle \underset{\mathrm{}=0}{\overset{\mathrm{}}{}}}\mathrm{\Psi }_{\mathrm{}}=\sqrt{\pi }e^{ip\rho _0}2^{ip}{\displaystyle \frac{\mathrm{\Gamma }(\frac{ip}{2})\mathrm{\Gamma }(\frac{1}{2}+\frac{ip}{2})}{\mathrm{\Gamma }(\frac{1}{2}+\frac{ip+n}{2})\mathrm{\Gamma }(\frac{1}{2}+\frac{ipn}{2})}}={\displaystyle \frac{2\pi \mathrm{\Gamma }(ip)}{\mathrm{\Gamma }(\frac{1}{2}+\frac{ip+n}{2})\mathrm{\Gamma }(\frac{1}{2}+\frac{ipn}{2})}}e^{ip\rho _0},`$ (B.10)
and this is the desired formula. |
warning/0507/cond-mat0507511.html | ar5iv | text | # Size–sensitive melting characteristics of gallium clusters: Comparison of Experiment and Theory for Ga₁₇⁺ and Ga₂₀⁺
## I Introduction
It is now well established that the melting points of particles with thousands of atoms decrease smoothly with decreasing particle size, due to the increase in the surface to volume ratio. Pawlow ; Nature-1977 However, unlike particles with thousands of atoms or the bulk material, probing the finite–temperature properties of small clusters (with $`<`$ 500 atoms) is non–trivial and remains a challenging task. Experimental studies of the melting transitions of clusters in small size regime have only recently become possible. Haberland-1997 ; Haberland-Nature ; Haberland-PRL-2003 ; Haberland-PRL-2005 ; Jarold-Tin ; Jarrold-Gallium1 ; Jarrold-Jacs ; Jarrold-Al ; Jarrold-Tinfrag Several interesting phenomena have been observed, including melting temperatures that rise above the bulk value Jarold-Tin ; Jarrold-Gallium1 and strong size–dependent variations in the melting temperatures. Jarrold-Jacs ; Jarrold-Al . These experimental findings have motivated many theoretical investigations on the finite–temperature behavior of clusters. Ga-prl ; Our-PRBsn10 ; Our-PRBsn20 ; James-Tin ; Eur-Phys.J ; Na-PRB ; Na-JCP Simulations based on first principles have been particularly successful in quantitatively explaining the factors behind the size–dependent variations in the melting behavior of clusters. Na-PRB Thus, a confluence of recent advances in experimental methods and theoretical studies using first principles methods have set the stage for a major increase in our understanding of phase transitions in these small systems.
Coming to the present work on gallium clusters, it is by now well known that gallium clusters not only melt at substantially higher temperatures than the bulk ($`T_{m[\mathrm{bulk}]}=`$ 303 K), Jarrold-Gallium1 but they also exhibit wide variations in the temperature dependences of their specific–heats, with some clusters showing strong peaks (due to the latent heat), while others (apparent ”non-melters”) showing no peak. Jarrold-Jacs . These features show a strong dependence on cluster size, where the addition of a single atom can change a cluster with no peak in the specific–heat into a “magic melter” with a very distinct peak. This behavior has been observed for gallium clusters, Ga<sub>n</sub>, with $`n=`$ 30–55.
In the present work, we probe the melting behavior of small gallium cluster ions where we show that the “non-melting” and “melting” features in the specific–heats are observed in clusters as small as Ga$`{}_{17}{}^{}^+`$ and Ga$`{}_{20}{}^{}^+`$, respectively. Prior experimental results for Ga$`{}_{}{}^{+}{}_{17}{}^{}`$ over a limited temperature range showed no evidence for a melting transition. Jarrold-Gallium1 The experimental results in this case were specific–heat measurements performed using multi–collision induced dissociation, where a peak in the specific–heat due to the latent heat was the signature of melting. On the other hand, recent simulations for Ga<sub>17</sub> show a broad peak in the specific–heat centered around 600 K. The previous specific–heat measurements for Ga$`{}_{}{}^{+}{}_{17}{}^{}`$ extended only up to 700 K, so one possible explanation for this apparent discrepancy is that the melting transition occurred at a slightly higher temperature than examined in the experiments. Here, we report specific–heat measurements for Ga$`{}_{}{}^{+}{}_{17}{}^{}`$ over a more extended temperature range, along with specific–heat measurements for Ga$`{}_{}{}^{+}{}_{20}{}^{}`$. While no peak is observed in the heat–capacities for Ga$`{}_{}{}^{+}{}_{17}{}^{}`$, a peak is observed for Ga$`{}_{}{}^{+}{}_{20}{}^{}`$.
To further probe the melting transitions in these clusters, ion mobility measurements were performed for Ga$`{}_{}{}^{+}{}_{17}{}^{}`$ and Ga$`{}_{}{}^{+}{}_{20}{}^{}`$ as a function of temperature. The ion mobility measurements provide average collision cross sections which can reveal information about the shape and volume changes that occur on melting. For example, a cluster with a non–spherical geometry might be expected to adopt a spherical shape (a liquid droplet) on melting. If there is not a significant shape change, there may still be a volume change on melting. Most bulk materials expand when they melt (the liquid is less dense than the solid). Even in the absence of a significant shape or volume change, the cross sections might show an inflection at the melting transition due to the thermal coefficient of expansion of the liquid cluster being larger than for the solid (in the macroscopic regime most liquids have larger coefficients of expansion than the corresponding solids). An inflection is observed in the cross sections for both Ga$`{}_{}{}^{+}{}_{17}{}^{}`$ and Ga$`{}_{}{}^{+}{}_{20}{}^{}`$. Thus, the ion mobility measurements suggest that Ga$`{}_{}{}^{+}{}_{17}{}^{}`$ as well as Ga$`{}_{}{}^{+}{}_{20}{}^{}`$ are in a liquidlike state above 800 K.
To explore the reasons behind the behavior outlined above (i.e., to determine why Ga$`{}_{}{}^{+}{}_{17}{}^{}`$ apparently melts without a peak in its specific–heat, while a peak is observed for Ga$`{}_{}{}^{+}{}_{20}{}^{}`$), we have carried out first principles Density–Functional (DF) Molecular–Dynamics (MD) calculations on both clusters. The ground–state structure and the bonding within the clusters is analyzed. The ionic specific–heat is computed using multiple histogram method. MH ; amv-review The calculated specific–heats for Ga$`{}_{}{}^{+}{}_{17}{}^{}`$ show three broad low intensity maxima that extend from 300 to 1400 K. This resembles the experiment results where the measured specific–heats are relatively featureless. In contrast, the calculated specific–heats for Ga$`{}_{}{}^{+}{}_{20}{}^{}`$ show a clear peak around 750 K. This is in excellent agreement with the peak obtained from experimental measurements (which occurs at around 700 K). Finally, our theoretical results show that the features in the specific–heat curves are influenced by the ground–state geometry, the bonding of the atoms within the ground–state structure, and the isomer distribution that becomes accessible as the temperature is raised.
In Sec. II, we present the experimental methods and the computational details. In Sec. III we discuss the experimental and theoretical results on both clusters. We conclude our results in Sec. IV.
## II Methodology
Specific-heats were measured using the recently developed multi–collision induced dissociation approach. The cluster ions are generated by laser vaporization of a liquid gallium target in a continuous flow of helium buffer gas. After exiting the laser vaporization region of the source the clusters travel through a 10 cm long temperature variable extension where their temperature is set. Cluster ions that exit the extension are focused into a quadrupole mass spectrometer where a particular cluster size is selected. The size selected clusters are then focused into a collision cell containing 1 Torr of helium. As the clusters enter the collision cell they undergo numerous collisions with the helium, each one converting a small fraction of the ions translational energy into internal energy. If the initial kinetic energy is high enough some of the cluster ions may be heated to the point where they dissociate. The dissociated and undissociated cluster ions are swept across the collision cell by a small electric field and some of them exit through a small aperture. The ions that exit are analyzed in a second quadrupole mass spectrometer and then detected by an off–axis collision dynode and dual microchannel plates. The fraction of the ions that dissociate is determined from the mass spectrum. Measurements are performed as a function of the ions initial kinetic energy, and the initial kinetic energy required for 50% dissociation (IKE50%D) is determined from a linear regression. IKE50%D is measured as a function of the temperature of the temperature–variable extension on the source. IKE50%D decreases as the temperature is raised because hotter clusters have more internal energy, and hence less energy needs to be added in order to cause dissociation. At the melting transition a sharp decrease in IKE50%D is expected due to the latent heat. The derivative of IKE50%D with respect to temperature is approximately proportional to the specific–heat. The proportionality constant is the fraction of the clusters initial kinetic energy that is converted into internal energy, which is estimated from an impulsive collision model. A drop in the IKE50%D values due to the latent heat of a melting transition leads to a peak in the specific–heat.
Ion mobility measurements can provide information on the shape and volume changes that occur when clusters melt. For the ion mobility measurements, the collision cell is replaced by a 7.6 cm long drift tube. 50 $`\mu `$s pulses of cluster ions are injected into the drift tube and the drift time distribution is obtained by recording the ions arrival times at the detector with a multichannel scalar. Average collision cross sections are obtained from the drift time distributions using standard methods. masonmcdaniel
All the simulations are performed using Born–Oppenheimer molecular–dynamics based on Kohn–Sham formulation of Density–Functional Theory (DFT). KS We have used Vanderbilts’ ultrasoft pseudopotentials uspp-vanderbilt within the GGA approximation, as implemented in the vasp package vasp for both clusters. For all calculations, we use only $`4s^2`$ and $`4p^1`$–electrons as valence, taking the 3$`d`$–electrons d-electron as a part of the ionic core. An energy cutoff of about $`10`$ Ry is used for the plane–wave expansion of the wavefunction, with a convergence in the total energy of the order of 0.0001 eV. Cubic supercells of lengths 20 and 25 Å are used for Ga$`{}_{}{}^{+}{}_{17}{}^{}`$ and Ga$`{}_{}{}^{+}{}_{20}{}^{}`$, respectively. For examining the finite–temperature behavior, the ionic phase space of the clusters is sampled by isokinetic MD where kinetic energy is held constant via a velocity scaling method. For both the clusters, we split the total temperature range from 100–1400 K into 15 different temperatures. We maintain the cluster at each temperature for a period of at least 90 ps, leading to total simulation times of the order of 1 ns. The resulting trajectory data were used to compute standard thermodynamic indicators as well as the ionic specific–heat, via a multihistogram technique. Details can be found in Ref. amv-review, ; abhijat-mh, .
## III Results and Discussion
Specific–heats measured for Ga$`{}_{}{}^{+}{}_{17}{}^{}`$ and Ga$`{}_{}{}^{+}{}_{20}{}^{}`$ as a function of temperature are shown in the lower half of Fig. 1. The points are the experimental values, while the dashed line is the prediction of a modified Debye model. In the case of Ga$`{}_{}{}^{+}{}_{17}{}^{}`$, the specific–heats shown in Fig. 1 appear to gradually increase up to around 900 K. The sharp decrease in the specific–heats above 900 K is an artifact due to evaporative cooling, the spontaneous unimolecular dissociation of the cluster ions as they travel between the source extension and the collision cell. For Ga$`{}_{}{}^{+}{}_{20}{}^{}`$, the specific–heats show a broad maximum, around 400 K wide, centered at around 725 K. The peak for Ga$`{}_{}{}^{+}{}_{20}{}^{}`$ is significantly broader than observed for larger clusters (like Ga$`{}_{}{}^{+}{}_{39}{}^{}`$ and Ga$`{}_{}{}^{+}{}_{40}{}^{}`$) where the peak was attributed to a melting transition. However, it is well known that the melting transition, and the corresponding peak in the specific–heats, becomes broader with decreasing cluster size. Thus, even though the peak in the specific-heats for Ga$`{}_{}{}^{+}{}_{20}{}^{}`$ is around 400 K wide, it is appropriate to assign it to a finite–size analog of a bulk melting transition. The center of the peak is at around 725 K, this is well above the bulk melting point (303 K). This continues a trend reported for larger cluster sizes ($`n=`$ 30–55) where the melting temperatures are also significantly above the bulk value. The unfilled red circles in Fig. 1 show the average collision cross sections determined for Ga$`{}_{}{}^{+}{}_{17}{}^{}`$ and Ga$`{}_{}{}^{+}{}_{20}{}^{}`$ as a function of the temperature. The cross sections are expected to systematically decrease with increasing temperature because the long range attractive interactions between the cluster ion and the buffer gas atoms becomes less important, and the collisions become harder as the temperature is raised. The thick dashed red line in the figures show the expected exponential decrease in the cross sections with increasing temperature. There is an inflection in the cross sections for Ga$`{}_{}{}^{+}{}_{20}{}^{}`$ that appears to slightly precede the peak in the specific–heat for this cluster. The inflection in consistent with a melting transition where the liquid cluster has a larger coefficient of thermal expansion than the solid. There is also an inflection in the cross sections for Ga$`{}_{}{}^{+}{}_{17}{}^{}`$. This suggests that a solid–liquid transition also occurs for Ga$`{}_{}{}^{+}{}_{17}{}^{}`$, but without a significant peak in the specific–heat.
To understand the reason behind the different behavior observed for Ga$`{}_{17}{}^{}^+`$ and Ga$`{}_{20}{}^{}^+`$, we have carried out a detailed analysis of structure and bonding in both clusters. As will be become apparent from the following discussion, that the ground–state geometry and the nature of bonding plays a crucial role in determining the finite–temperature behavior of the cluster. We begin with a discussion of the ground–state geometries of cationic Ga<sub>17</sub> and Ga<sub>20</sub> clusters. We have obtained more than 50 distinct equilibrium geometries by quenching more than 200 structures, selected from a few high temperature MD runs, for both sizes. In Fig. 2, we show the lowest energy structure along with some low lying excited state geometries of both clusters. The lowest energy geometry of the Ga$`{}_{17}{}^{}^+`$ cluster (see Fig. 2–a(1)) is similar to that of Ga<sub>17</sub> reported in our earlier work. Ga-prl It has a distorted decahedral structure, which suggests the possibility of further cluster growth to a 19–atom double decahedron. In contrast, the ground–state geometry of Ga$`{}_{20}{}^{}^+`$, shown in Fig. 2–b(1), is more symmetric. It can be described as a double decahedral structure of 19 atoms, with the bottom capped atom merging into the pentagonal plane to form a hexagonal ring. In addition, an atom from the top pentagon and the upper capped atom rearrange to accommodate the 20<sup>th</sup> atom, leading to a dome–shaped hexagonal ring.
We now analyze the structural properties in detail to get an insight into the features that influence the melting characteristics. An analysis of the bond–length distribution shows that there are 12 bonds, for each cluster, having distances less than 2.55 Å. metallic-bond Interestingly, for Ga$`{}_{17}{}^{}^+`$, these short bonds are spread all over the cluster, whereas for Ga$`{}_{20}{}^{}^+`$, they form the upper and the lower hexagonal rings. The distribution of coordination numbers coordination-number indicate that for Ga$`{}_{20}{}^{}^+`$, almost all the atoms in the rings (about 16), have a coordination number of 4. The Ga$`{}_{17}{}^{}^+`$ cluster, however, does not have such a uniform distribution of coordination numbers. Thus, the ground–state geometry of Ga$`{}_{17}{}^{}^+`$ might be considered to be “disordered”, while that of Ga$`{}_{20}{}^{}^+`$ exhibits a more–ordered structure.
Striking differences are also observed in the low energy isomers and their distribution on the potential–energy surface. As mentioned above, we have obtained more than 50 distinct isomers spanning an energy range of about 1.0 eV above the ground–states for each cluster. In Fig. 3, we plot the energies of these isomers relative to the ground–state, arranged in an ascending order. The isomers for the Ga$`{}_{17}{}^{}^+`$ cluster appear to exhibit an almost continuous energy distribution. While a few of these isomers are severe distortions of the ground–state geometry, the rest do not show any resemblance (see Fig. 2a). It appears that for the Ga$`{}_{17}{}^{}^+`$ isomers in this low energy regime, small rearrangements of the atoms, costing just a small amount of energy, lead to several close-lying isomers, so that the isomer distribution is almost continuous. In contrast, the isomers of Ga$`{}_{20}{}^{}^+`$ cluster are distributed in three groups, separated by an energy gap of about 0.2 eV (Fig. 3). The first group of isomers have slightly different orientations of atoms in the hexagonal rings and are nearly degenerate with the ground–state. The second group consists of structures having only the lower hexagonal ring while the third group has no rings. This indicates that the hexagonal units of Ga$`{}_{20}{}^{}^+`$ cluster are stable and difficult to break. The stability of the ring–pattern of Ga$`{}_{20}{}^{}^+`$ and the isomer distribution for both clusters should have a substantial effect on the melting characteristics. Indeed, as we shall see further below, these features play a crucial role in the finite–temperature characteristics. It should be mentioned that although these observations are based on rather limited search, we believe that the general features described here are essentially correct.
The most important difference between the two clusters is the nature of the bonding. We use the concept of an electron localization function (ELF), elf-silvi to describe the nature of bonding. This function is normalized to a value between zero and unity; a value of 1 represents a perfect localization of the valence charge while the value for the uniform electron gas is 1/2. The locations of maxima of this function are called *attractors*, since other points in space can be connected to them by paths of maximum gradient. The set of all such points in space that are attracted by a maximum is defined to be the *basin* of that attractor. Basin formations are usually observed as the value of the ELF is lowered from its maximum, at which there are as many basins as the number of atoms in the system. Typically, the existence of an isosurface or a basin in the bonding region between two atoms at a high ELF value, say $`0.70`$, signifies a localized bond in that region.
We have analyzed the electron localization functions for Ga$`{}_{17}{}^{}^+`$ and Ga$`{}_{20}{}^{}^+`$ clusters for values $`0.85`$. In Table 1, we give the number of basins containing two or more atoms, for selected ELF values. The table clearly shows a fragmented growth pattern of the basins for Ga$`{}_{17}{}^{}^+`$, each containing very few atoms as compared to that of Ga$`{}_{20}{}^{}^+`$. For instance, at an isovalue of 0.75, while Ga$`{}_{17}{}^{}^+`$ has three basins each having 2 atoms, Ga$`{}_{20}{}^{}^+`$ has just two basins each containing 5 and 7 atoms that corresponds to the two hexagonal rings. The ELF contours for the isovalue of 0.75 are shown in Fig. 4. The merged basins structures are shown by the black lines. It may be inferred that the bonds between atoms in the hexagonal rings of Ga$`{}_{20}{}^{}^+`$ are strong and covalent in nature with similar strengths, while the fragmented basin growth pattern in Ga$`{}_{17}{}^{}^+`$ indicates inhomogeneity of the bond strengths.
The calculated, normalized, canonical specific–heats are shown in Fig. 5 plotted against temperature. The plot for Ga$`{}_{17}{}^{}^+`$ exhibits a broad feature (apparently consisting of three components) which extends from 300 K to 1400 K. For Ga$`{}_{20}{}^{}^+`$, the calculated specific–heat remains nearly flat up to about 600 K, it then increases sharply and peaks at about 800 K, in excellent agreement with the experimental results described above. Thus, interesting size–sensitive features seen in the experimental heat–capacities are reproduced in our simulations. This behavior can be understood from our earlier discussion of the bond–length distributions, coordination numbers, isomer–distributions, and the nature of bonding in these clusters. While the Ga$`{}_{17}{}^{}^+`$ cluster shows no real evidence for ordered behavior, the Ga$`{}_{20}{}^{}^+`$ cluster has well–ordered ring–patterns. Thus, when Ga$`{}_{17}{}^{}^+`$ is heated, the bonds soften gradually, and the cluster hops through all its isomers continuously. This is clearly demonstrated by the ionic motion as a function of temperature, which shows that this cluster evolves through all isomers smoothly from 300 K to 1400 K. On the other hand, the ionic motion for Ga$`{}_{20}{}^{}^+`$ shows only minor rearrangements of the atoms until 600 K, and then the cluster visits all the isomers corresponding to the first group of isomers described above. At about 700 K, the upper hexagonal ring breaks, while at about 800 K, the lower ring breaks. Thus, melting of Ga$`{}_{20}{}^{}^+`$ cluster is associated with the breaking of the well–ordered covalently bonded hexagonal units.
We have also analyzed the melting characteristics via traditional parameters such as, the root–mean–squared bond–length–fluctuations ($`\delta _{\mathrm{rms}}`$) and the mean–squared ionic displacements (MSD). In Fig. 6, we show the $`\delta _{\mathrm{rms}}`$ for Ga$`{}_{17}{}^{}^+`$ and Ga$`{}_{20}{}^{}^+`$ clusters. This plot correlates well with the specific–heat curve shown in Fig. 5. The $`\delta _{\mathrm{rms}}`$ for Ga$`{}_{17}{}^{}^+`$ rises gradually from 300 K, while for Ga$`{}_{20}{}^{}^+`$, it rises sharply at about 700 K, and finally saturates to the same value for both clusters. It may be inferred from this observation that the behavior of both clusters at temperatures say, $`T800`$ K, are similar and that both clusters can be considered to be in *liquidlike* states. This conclusion is further substantiated by the MSD plots (figures not shown), which saturate at $``$ 21$`\AA ^2`$, at about 1200 K for both clusters.
## IV Summary and conclusion
It is evident from the present study that the nature of the ground–state geometry and bonding strongly influences the finite–temperature characteristics of Ga$`{}_{17}{}^{}^+`$ and Ga$`{}_{20}{}^{}^+`$. At high temperatures, $`T800`$ K, both Ga$`{}_{17}{}^{}^+`$ and Ga$`{}_{20}{}^{}^+`$ have similar root–mean–squared bond–length–fluctuations and the mean–squared ionic displacements so that both of them can be considered to be in liquidlike states. The experimental results show that while Ga$`{}_{17}{}^{}^+`$ apparently undergoes a solid–liquid transition without a significant peak in the specific–heat, Ga$`{}_{20}{}^{}^+`$ melts with a relatively sharp peak. The simulations show that if the cluster is “ordered” (i.e. a large fraction of the constituent atoms show similar bonding, coordination numbers, and bond energies) then it is likely to show a sharp melting transition with a significant peak in the specific–heat. On the other hand, if the cluster is “disordered” (i.e. the constituent atoms occur in a wide distribution of bonding environments) it will probably undergo a solid–liquid transition without a significant peak in the specific–heat. In the latter case, the number of isomers or conformations sampled by the cluster increases steadily as the temperature is raised, instead of the abrupt increase that occurs when a cluster undergoes a sharp melting transition.
These observations have interesting consequences for the finite–temperature behavior of small clusters as a function of cluster size. It is likely that as clusters grow in size their structures evolve from one well–ordered structure to another, passing on the way through some cluster sizes that have “disordered” structures. For instance, the 13–atom gallium cluster is predicted to have a highly symmetric decahedron structure with a bonding pattern that is similar to that found here for Ga$`{}_{20}{}^{}^+`$Ga-prl So in the present case, cluster growth from Ga<sub>13</sub> (a decahedron) to (Ga$`{}_{20}{}^{}^+`$, a distorted double–decahedron) proceeds via a disordered Ga$`{}_{17}{}^{}^+`$ structure. Such behavior is also observed for sodium clusters in 40 to 55 atom size range; the ground–state geometries of Na<sub>40</sub> and Na<sub>55</sub> are either icosahedron or close to icosahedron while that of Na<sub>50</sub> has no particular symmetry. Na-PRB ; Na-JCP In such cases, we expect that the specific–heats should change from showing a well–defined peak to a rather broad one, and back again to well–defined. We believe this behavior to be generic as it has not only been observed in case of gallium clusters Jarrold-Gallium1 ; Jarrold-Jacs but also in case of aluminum clusters Jarrold-Al , experimentally, and for sodium clusters in the simulations mentioned above.
## V Acknowledgment
The financial support of the Indo–French Center for Promotion of Advanced Research is greatfully acknowledged. C–DAC (Pune) is acknowledged for providing us with the supercomputing facilities. We gratefully acknowledge partial support of this work by the US National Science Foundation. |
warning/0507/hep-ph0507137.html | ar5iv | text | # 2005 International Linear Collider Workshop - Stanford, U.S.A. Split Supersymmetry at Colliders
## I INTRODUCTION
Split supersymmetry SpS is a possibility to evade many of the phenomenological constraints that plague generic supersymmetric extensions of the Standard Model. By splitting the supersymmetry-breaking scale between the scalar and the gaugino sector, the squarks and sleptons are rendered heavy (somewhere between several $`\mathrm{TeV}`$ and the GUT scale), while charginos and neutralinos may still be at the $`\mathrm{TeV}`$ scale or below. This setup eliminates dangerous flavor-changing neutral current transitions, electric dipole moments, and spurious proton-decay operators without the need for mass degeneracy between the sfermion generations. The benefits of the supersymmetry paradigm, in particular the unification of gauge groups at a high scale and the successful dark-matter prediction, are retained.
In the Higgs sector, the split-supersymmetry scenario requires a fine-tuning that pulls the Higgs vacuum expectation value down to the observed electroweak scale. The extra Higgses of a supersymmetric model are located at the sfermion mass scale. This fine-tuning is obviously unnatural, but it may (or may not) find a convincing explanation in ideas beyond the realm of particle physics landscape ; drees .
In this talk (for more details, see pheno ), we investigate the split-supersymmetry scenario from a purely phenomenological point of view. We ask ourselves the question whether and how (i) the particles present in the low-energy spectrum can be detected, and (ii) the underlying supersymmetric nature of the model can be verified. The two tasks require combining LHC and ILC data.
The low-energy effective theory is particularly simple. In addition to the Standard Model spectrum including the Higgs boson, the only extra particles are the four neutralinos, two charginos, and a gluino. Since all squarks are very heavy, the gluino is long-lived. The gluino can be produced at the LHC only, while the charginos and neutralinos can be accessible both at the LHC and at the ILC.
Renormalization group running without sfermions and heavy Higgses lifts the light Higgs mass considerably above the LEP limit, solving another problem of the MSSM. Still, the Higgs boson is expected to be lighter than about $`200\mathrm{GeV}`$. Apart from this Higgs mass bound, the only trace of supersymmetry would be the mutual interactions of Higgses, gauginos, and Higgsinos, i.e., the chargino and neutralino Yukawa couplings. These couplings are determined by the gauge couplings at the matching scale $`\stackrel{~}{m}`$, where the scalars are integrated out.
## II RENORMALIZATION GROUP EVOLUTION
Although, in the absence of sfermions, the overall phenomenology of split-supersymmetry models does not depend very much on the particular spectrum, for a quantitative analysis we have to select a specific scenario. To this end, we note that the popular assumption of radiative symmetry breaking (i.e., a common scalar mass parameter for sfermions and Higgses) has to be dropped, and the Higgsino mass parameter $`\mu `$ is an independent quantity. Assuming gauge coupling unification and gaugino mass unification, we start from the following model parameters at the grand unification scale $`M_{\mathrm{GUT}}=6\times 10^{16}\mathrm{GeV}`$:
$`M_1(M_{\mathrm{GUT}})=M_2(M_{\mathrm{GUT}})`$ $`=M_3(M_{\mathrm{GUT}})=120\mathrm{GeV},`$ $`\mu (M_{\mathrm{GUT}})`$ $`=90\mathrm{GeV},`$ $`\mathrm{tan}\beta `$ $`=4.`$ (1)
For the SUSY-breaking scale we choose $`\stackrel{~}{m}=10^9\mathrm{GeV}`$. Figure 1 displays the solutions of the renormalization group equations with the input parameters set in eq.(1pheno . At the low scale $`Q=m_Z`$, we extract the mass parameters:
$`M_1(Q=m_Z)`$ $`=74.8\mathrm{GeV}`$ $`M_3^{\overline{\text{DR}}}(Q=1\mathrm{TeV})`$ $`=690.1\mathrm{GeV}`$
$`M_2(Q=m_Z)`$ $`=178.1\mathrm{GeV}`$ $`\mu (Q=m_Z)`$ $`=120.1\mathrm{GeV}`$ (2)
The resulting physical gaugino and Higgsino masses are:
$`m_{\stackrel{~}{\chi }_1^0}`$ $`=71.1\mathrm{GeV},`$ $`m_{\stackrel{~}{\chi }_2^0}`$ $`=109.9\mathrm{GeV},`$ $`m_{\stackrel{~}{\chi }_3^0}`$ $`=141.7\mathrm{GeV},`$ $`m_{\stackrel{~}{\chi }_4^0}`$ $`=213.7\mathrm{GeV},`$
$`m_{\stackrel{~}{\chi }_1^+}`$ $`=114.7\mathrm{GeV},`$ $`m_{\stackrel{~}{\chi }_2^+}`$ $`=215.7\mathrm{GeV},`$ $`m_{\stackrel{~}{g}}`$ $`=807\mathrm{GeV}`$ (3)
These mass values satisfy the LEP constraints. Virtual effects of a split supersymmetry spectrum on Standard Model observables have recently been discussed in Martin:2004id .
The neutralinos $`\stackrel{~}{\chi }_{1,2,3,4}^0`$ are predominantly bino, Higgsino, Higgsino, and wino, respectively. The Higgsino content of the lightest neutralino is $`h_f=0.2`$, so the dark-matter condition dm is satisfied. To our given order the Higgs mass is $`m_H=150\mathrm{GeV}`$.
Because we integrate out the heavy scalars, the neutralino and chargino Yukawa couplings deviate from their usual MSSM prediction, parameterized by four anomalous Yukawa couplings $`\kappa `$. We can extract their weak-scale values from Fig. 1:
$`{\displaystyle \frac{\stackrel{~}{g}_u}{g\mathrm{sin}\beta }}`$ $`1+\kappa _u=1+0.018`$ $`{\displaystyle \frac{\stackrel{~}{g}_d}{g\mathrm{cos}\beta }}`$ $`1+\kappa _d=1+0.081`$
$`{\displaystyle \frac{\stackrel{~}{g}_u^{}}{g^{}\mathrm{sin}\beta }}`$ $`1+\kappa _u^{}=10.075`$ $`{\displaystyle \frac{\stackrel{~}{g}_d^{}}{g^{}\mathrm{cos}\beta }}`$ $`1+\kappa _d^{}=10.17`$ (4)
## III LONG-LIVED GLUINOS
Since the standard cascade decays of initial squarks and gluinos are absent in this model, there are only two sources of new particles left. The gluino is produced in pairs from gluon-gluon and quark-antiquark annihiliation. Pairs of charginos and neutralinos are produced through a Drell–Yan $`s`$-channel $`Z`$ boson, photon, or $`W`$ boson. These cross sections are known to next-to-leading order precision LHC-NLO .
Unless we have a-priori knowledge about the sfermion scale $`\stackrel{~}{m}`$, the gluino lifetime is undetermined. Figure 2 compares this scale with other relevant scales of particle physics. Once $`\stackrel{~}{m}10^3\mathrm{GeV}`$, the gluino hadronizes before decaying. The resulting states that consist of either of a gluino and pairs or triplets of quarks, or of a gluino bound to a gluon, are called $`R`$-hadrons Rhadrons . For gluinos produced near threshold, the formation of gluino-pair bound states (gluinonium) is also possible and leads to characteristic signals gluinonium .
For $`\stackrel{~}{m}>10^6\mathrm{GeV}`$, the gluino travels a macroscopic distance. If $`\stackrel{~}{m}>10^7\mathrm{GeV}`$, strange $`R`$-hadrons can also decay weakly, and gluinos typically leave the detector undecayed or are stopped in the material. For even higher scales, $`\stackrel{~}{m}>10^9\mathrm{GeV}`$, $`R`$-hadrons could become cosmologically relevant, since they affect nucleosynthesis if their abundance in the early Universe is sufficiently high SpS .
If gluino decays can be observed, their analysis yields information about physics at the scale $`\stackrel{~}{m}`$ and thus allows us to draw conclusions about the mechanism of supersymmetry breaking gldecay .
Without relying on gluino decays, there are two strategies for detecting the corresponding $`R`$-hadrons pheno ; Hewett:2004nw . (i) The production of a stable, charged, $`R`$-hadron will give a signal much like the production of a stable charged weakly-interacting particle. This signal consists of an object that looks like a muon but arrives at the muon chambers significantly later than a muon owing to its large mass. (ii) While for stable neutral $`R`$-hadrons there will be some energy loss in the detector, there will be a missing transverse energy signal due to the escape of the $`R`$-hadrons. As leptons are unlikely to be produced in this process, the signal will be the classic SUSY jets with missing transverse energy signature.
In Fig. 3, we show the expected discovery reach for both channels, based on models for the $`R`$-hadron spectrum and the $`R`$-hadron interaction in the detector that we have implemented in HERWIG HERWIG .
## IV CHARGINO AND NEUTRALINO YUKAWA COUPLINGS
If split supersymmetry should be realized in nature, the observation of the gluino, charginos and neutralinos will only be the first task. Once these states are discovered, we will have to show that at the scale $`\stackrel{~}{m}`$ they constitute a supersymmetric Lagrangian.
A quantitative trace of this is given by the off-diagonal elements in the mass matrices that derive from the gaugino-higgsino-Higgs couplings (4). They determine the mixing of gauginos and higgsinos into charginos and neutralinos as mass eigenstates. Simultaneously, they also constitute the neutralino and chargino Yukawa couplings. In split supersymmetry, the renormalization flow below the sfermion scale $`\stackrel{~}{m}`$ induces non-zero values of order $`\kappa _i^{()}=0.2\mathrm{}0.2`$. If we are able to detect deviations of this size at a collider, we can both establish the supersymmetric nature of the model and verify the matching condition to the MSSM at $`\stackrel{~}{m}`$.
To measure the neutralino and chargino mixing matrices, a precise mass measurement is necessary. This is possible (for mass differences, at least) at the LHC and, to a better accuracy, at the ILC. Without gaugino–Higgsino mixing the mass matrices would be determined by the MSSM parameters $`M_1,M_2`$ and $`\mu `$. The gaugino–Higgsino mixing adds terms of the order of $`M_Z`$ and introduces the additional parameter $`\mathrm{tan}\beta `$, leading to four MSSM parameters altogether.
In $`e^+e^{}`$ collisions, for the parameter set (2) almost all chargino and neutralino production channels have cross sections larger than $`0.1\mathrm{fb}`$, and the threshold value for $`\stackrel{~}{\chi }_1^+\stackrel{~}{\chi }_1^{}`$ production is as large as $`1\mathrm{pb}`$ pheno . The NLO electroweak corrections to these production cross sections have been calculated in ILC-NLO . A linear collider with moderate energy and high luminosity would be optimal to probe all these processes, and some kind of fit is the proper method to extract the weak-scale Lagrangian parameters.
We compute the masses and the cross sections for all pair-production processes, with the exception of the $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0`$ channel. To all observables we assign an experimental error, which in our simplified treatment is a relative error of $`0.5\%`$ on all linear-collider mass measurements Aguilar-Saavedra:2001rg , $`5\%`$ on all LHC mass measurements Bachacou:1999zb , and the statistical uncertainty on the number of events at a linear collider corresponding to $`100\mathrm{fb}^1`$ of data at a $`1\mathrm{TeV}`$ collider after all efficiencies.
Around the central parameter point we randomly generate 10000 sets of pseudo-measurements, using a Gaussian smearing. Out of each of these sets we extract the MSSM parameters by a global fit method. The fit results (see Fig. 4, top) show that at the LHC we can extract the Lagrangian mass parameters with reasonable precision. There is sensitivity to one Higgs-sector parameter, which we can take either as $`\mathrm{tan}\beta `$ or as one of the mixing parameters. If we fix $`\mathrm{tan}\beta =4`$, the precision on $`\kappa _u`$ is sufficient to verify consistency with a supersymmetric underlying theory. However, using LHC data alone, a simultaneous fit of all parameters gives only very weak constraints on the anomalous Yukawa couplings (Tab. 1), so no conclusions about the split-supersymmetry renormalization effects can be drawn.
The higher precision of measurements at the ILC, in particular adding cross sections as independent observables, allows us to improve the precision on a five-parameter fit (Fig. 4, bottom) or to simultaneously fit all Lagrangian parameters (Fig. 5). This is the proper treatment, unless we would have reasons to believe that some of the $`\kappa _i^{()}`$ are predicted to be too small to be measured. The results for the precision in determining the anomalous couplings are listed in Tab. 1. (Note that in a complete fit, $`\mathrm{tan}\beta `$ is no longer an independent parameter, so we can fix it to some given value.)
These results for the linear collider indeed indicate that we could not only confirm that the Yukawa couplings and the neutralino and chargino mixing follow the predicted MSSM pattern; for the somewhat larger $`\kappa _i^{}`$ values we can even distinguish the complete weak-scale MSSM from a split supersymmetry spectrum.
While the elements of the neutralino and chargino mixing matrices depend on the Yukawa couplings in a complicated way, the cross sections for chargino/neutralino pair production in association with a Higgs boson are directly proportional to these parameters. Decays of the kind $`\chi _2^\pm \chi _1^\pm H`$ or $`\chi _j^0\chi _i^0H`$ would carry the same information, but typically are kinematically forbidden in split supersymmetry scenarios.
Associated production of charginos and neutralinos with a Higgs boson in the continuum can in principle be observed at a high-luminosity $`e^+e^{}`$ collider. The cross sections for some of these channels exceed $`0.1\mathrm{fb}`$ with little background pheno , so with $`1\mathrm{ab}^1`$ of luminosity we would expect some events of this kind to be detectable. However, due to the small rates for these processes, the achievable precision for parameter determination is limited. The measurement of masses and pair-production cross sections will be the key for establishing split supersymmetry as the underlying physical scenario.
###### Acknowledgements.
W.K. is supported by the German Helmholtz-Gemeinschaft, Contract No. VH–NG–005. |
warning/0507/astro-ph0507252.html | ar5iv | text | # New CO observations and simulations of the NGC 4438/NGC 4435 system Movies of all simulations are available in the electronic version of this article.
## 1 Introduction
The perturbed spiral galaxy NGC 4438 and its companion galaxy NGC 4435 represent the most complicated system in the Virgo cluster (see Fig. 1 and Table 1).
The highly inclined disk of NGC 4438 is heavily perturbed showing a prominent tidal perturbation to the north and with stellar debris displaced to the west of the galaxy’s main disk (Arp 1966). Combes et al. (1988) simulated the NGC 4438/NGC 4435 system using a test-particle code. They showed that the northern tidal tail can be reproduced by a retrograde encounter with an impact parameter of $`6`$ kpc that occurred $`100`$ Myr ago. The relative position and radial velocity between NGC 4438 and NGC 4435 of 730 km s<sup>-1</sup> could be reproduced by their model.
Observations of the interstellar medium (ISM) of NGC 4438 revealed an important extraplanar component that is displaced to the west of the galactic disk (for an overview of the gas distribution at multiple wavelengths see Fig.11 of Kenney et al. 1995). The disk ISM is mainly in the form of molecular gas (Combes et al. 1988, Kenney et al. 1995) characterized by the absence of Hi emission (Cayatte et al. 1990, Hibbard et al. 2001). However, CO and Hi are detected to the west of the center of NGC 4438. The most spectacular gas distribution is that of the dense ionized gas observed in H$`\alpha `$ (Kenney et al. 1995). These observations revealed several filaments which originate from the disk plane and extend out from the disk for 5-10 kpc towards the west and southwest. The radial velocities of these filaments are within $`200`$ km s<sup>-1</sup> of the galaxy’s systemic velocity. Not only the cold and warm phases of the ISM are detected to the west beyond the galactic disk, but also the hot phase (X-rays) and the magnetic field (radio continuum; Kotanyi et al. 1983).
While it is certain that the distortion of the stellar content of NGC 4438 is due to a tidal encounter with NGC 4435, different mechanisms were put forward to explain the displacement of all ISM phases: (i) the tidal interaction which extracted the molecular gas from the center and left it to the west of the galaxy (Combes et al. 1988), (ii) a collision between the ISM of the two galaxies (Kenney et al. 1995), and (iii) ram pressure stripping due to the rapid motion of NGC 4438 through the hot intracluster medium (ICM; Kotanyi et al. 1983, Chincarini & de Souza 1985).
Since NGC 4438 is located in projection only 1$`{}_{}{}^{\mathrm{o}}300`$ kpc from the cluster center<sup>1</sup><sup>1</sup>1We use a distance to the Virgo cluster of 17 Mpc and has a radial velocity of $`1000`$ km s<sup>-1</sup> with respect to the cluster mean, the conditions necessary for strong ram pressure are fulfilled. On the other hand, since Kenney et al. (1995) detected H$`\alpha `$ and Machacek et al. (2004) detected soft X-ray emission in the center of NGC 4435, an ISM-ISM collision is unavoidable if the impact parameter of the encounter with NGC 4438 is smaller than $`10`$ kpc. On the observational side, extraplanar CO and extraplanar, asymmetric, high surface brightness radio continuum emission are very rare phenomena. It is mainly observed in direct galaxy collisions where the two ISM collide (see e.g. Braine et al. 2003). On the other hand, displacements of X-ray and $`H\alpha `$ emission are observed in cluster spiral galaxies (X-rays: Finoguenov et al. 2004, in preparation for the Coma cluster; H$`\alpha `$: Yoshida et al. 2002; NGC 4388 in the Virgo cluster). Thus, there is no clear evidence that rules out or confirms one of the suggested interactions.
In this article we present new CO observations of NGC 4438 and NGC 4435 (Sect. 2) together with numerical simulations of the system (Sect. 3). The latter present the advantage of separating different types of interaction to determine their influence on the ISM. The CO observations show three characteristics that can only be reproduced by a model including ram pressure stripping (Sect. 4.1): (i) the extraplanar gas distribution, (ii) extraplanar double line profiles, and (iii) redshifted line profiles with respect to galactic rotation.
## 2 Observations
The observations were carried out with the 30 meter millimeter-wave telescope on Pico Veleta (Spain) run by the Institut de RadioAstronomie Millimétrique (IRAM) in June 2002. The CO(1–0) and CO(2–1) transitions at 115 and 230 GHz respectively were observed simultaneously and in both polarizations. The observed positions are marked as triangles in Fig. 1. We generally used the $`512\times 1`$ MHz filterbanks at 3mm and the two $`256\times 4`$ MHz filterbanks at 1mm, yielding an instantaneous bandwidth of 1300 km s<sup>-1</sup>.
System temperatures were typically 250–350 K at 3mm and much higher for the CO(2–1) transition (T$`{}_{}{}^{}{}_{A}{}^{}`$ scale), due to the high water vapor content (usually H<sub>2</sub>O $`>`$ 6mm). The forward (main beam) efficiencies at Pico Veleta were taken to be 0.95 (0.74) at 115 GHz and 0.90 (0.54) at 230 GHz. The half-power beamwidths are about 21<sup>′′</sup> and 11<sup>′′</sup>. All observations were done in wobbler-switching mode, usually with a throw of 150<sup>′′</sup> but sometimes more or less depending on the position observed, in order to be sure not to have emission in the reference beam. Pointing was checked on the bright quasar 3C273 roughly every 90 minutes.
The main observational problem was the anomalous refraction that affected pointing measurements and resulted in a widening of the beam until sunset. Much of the badly affected data were rejected and we do not present the CO(2–1) data due to the generally high noise level and refraction.
Data reduction was very simple. After eliminating the bad spectra or bad channels, the spectra for each position were summed. Only zero-order baselines (i.e. continuum levels) were subtracted to obtain the final spectra.
Two features are particularly interesting: the presence of molecular gas in the absence of Hi in the tidal tail to the north (Cayatte et al. 1990, Hibbard et al. 2001) and the molecular gas to the west initially detected by Combes et al. (1988).
### 2.1 NGC 4435
We made one pointing on the center of NGC 4435. For the first time CO emission is detected in this galaxy. The CO(1–0) and CO(2–1) spectra are shown in Fig. 2.
The total line width is $`400`$ km s<sup>-1</sup>, which is about $`100`$ km s<sup>-1</sup> larger than the line width derived from optical spectroscopy (Kenney et al. 1995). This large linewidth justifies our model rotation velocity for NGC 4435 (see Sect. 3). In spiral galaxies the ratio between the CO(2–1) and CO(1–0) lines observed at the same resolution is about 0.9 (Braine et al. 1993). Since the brightness temperature of an unresolved source is inversely proportional to the beamsize, a larger line ratio than 0.9 implies that the source is not resolved at 115 GHz. The fact that the ratio between the CO(2–1) and CO(1–0) lines is about two indicates that the spatial extent of the molecular gas is comparable to the CO(2–1) resolution of $`11^{\prime \prime }900`$ pc, somewhat larger than the extent estimated by Kenney et al. (1995). Using the same conversion factor of $`=2\times 10^{20}`$ H<sub>2</sub> cm<sup>-2</sup> (K km/s)<sup>-1</sup>, the molecular gas mass of NGC 4435 is $`9.5\times 10^7`$ M when the He associated with the molecular gas is included.
### 2.2 NGC 4438
The CO(1–0) spectra of NGC 4438 superimposed on an optical image are shown in Fig. 3.
We have omitted the central spectrum which we show separately in Fig. 14.
Using the same conversion factor as above, the total molecular gas mass (including the He associated with the H<sub>2</sub>) of NGC 4438 is $`1.7\times 10^9`$ M. The conversion ratio is very probably overestimated in the central parts (as for spirals in general) and possibly underestimated in the outer regions. Our observations show weaker CO emission for the western part than reported by Combes et al. (1988) and that it is closely related (spatially) with the Hi and dust lanes.
The CO(1–0) emission extends over almost the entire observed region. The lines are becoming weaker to the west. Whereas we detect only one line at the eastern edge of the galactic disk north-east of the galaxy center, CO emission is detected up to $`80^{\prime \prime }`$ to the west of the galaxy’s major axis. Moreover, all spectra south west of the galaxy center peak at velocities greater than zero, despite the fact that the southern part represents the approaching side of the galaxy. This behaviour is also present in the H$`\alpha `$ spectra of Kenney et al. (1995). We will show in Sect. 4 that these spectra are displaced due to the action of ram pressure. Most striking is the detection of strong, clearly separated double lines west and south west of the galaxy center at the positions ($`20^{\prime \prime }`$,$`20^{\prime \prime }`$), ($`38^{\prime \prime }`$, $`10^{\prime \prime }`$), and ($`40^{\prime \prime }`$, $`10^{\prime \prime }`$). The spectrum south of the galaxy center ($`6^{\prime \prime }`$, $`19^{\prime \prime }`$) also shows a redshifted component which can be better seen in the H$`\alpha `$ Fabry-Perot data of Chemin et al. (2005). Such planar and extraplanar double lines are a unique feature, characteristic of the interaction NGC 4438 underwent and is still undergoing.
The extraplanar gas in NGC 4438, or gas at velocities other than rotation, is at least $`4.7\times 10^8`$ M when the He is included. We took the positive velocity component of all spectra west of a right ascension offset of $`19^{\prime \prime }`$. The spectrum at an offset of $`19.5^{\prime \prime }`$ was not included although its velocity has clearly been increased with respect to normal rotation. Similarly, the spectra along the major axis south of the center have not been included although they too have velocities which are not those of normal rotation (or normal rotation and tides).
### 2.3 The northern tidal tail
As shown in Fig. 1 we also observed several positions on the northern tidal tail (Fig. 4).
Unexpectedly, we detect CO(1–0) emission with narrow line profiles at several positions ($`\mathrm{\Delta }v\stackrel{<}{}50`$ km s<sup>-1</sup>). Using the same conversion factor of $`=2\times 10^{20}`$ H<sub>2</sub> cm<sup>-2</sup> (K km/s)<sup>-1</sup>, the molecular gas mass in this region is about $`2.5\times 10^7`$ M.
## 3 Simulations
We have adopted a model where the ISM is simulated as a collisional component, i.e. as discrete particles which have a mass and a radius and which can have inelastic collisions (sticky particles). In contrast to smoothed particle hydrodynamics (SPH), which is a quasi continuous approach and where the particles cannot penetrate each other, our approach allows a finite penetration length, which is given by the mass-radius relation of the particles. Both methods have their advantages and their limits. The advantage of our approach is that ram pressure can be included easily as an additional acceleration on particles that are not protected by other particles (see Vollmer et al. 2001). In this way we avoid the problem of treating the huge density contrast between the ICM ($`n10^4`$ cm<sup>-3</sup>) and the ISM ($`n>1`$ cm<sup>-3</sup>) of the galaxy.
Since the model is described in detail in Vollmer et al. (2001), we will summarize only its main features. The N-body code consists of two components: a non-collisional component that simulates the stellar bulge/disk and the dark halo, and a collisional component that simulates the ISM. The 25 000 particles of the collisional component represent gas cloud complexes which are evolving in the gravitational potential of NGC 4438. For the ISM-ISM collision models the ISM of NGC 4435 is simulated by 10 000 particles.
The total assumed initial gas masses of NGC 4438 and NGC4435 are $`M_{\mathrm{gas}}^{\mathrm{tot}}=\mathrm{4.9\hspace{0.17em}10}^9`$ M and $`\mathrm{1.3\hspace{0.17em}10}^9`$ M, respectively. For NGC 4438 we begin with a NGC 4501 type galaxy, i.e. the initial Hi diameter is close to the optical diameter. The recent GALEX observations of Boselli et al. (2005) justify this assumption (see also Sect. 4). A larger Hi diameter does not change our results. The total gas mass of NGC 4438 calculated as the sum of the molecular gas mass derived from our observations and the initial Hi mass of NGC 4501 gives $`3.4\times 10^9`$ M. The total gas mass for NGC 4435 represents the upper limit for an S0 galaxy of its luminosity (Welch & Sage 2003). The ensemble of model clouds has a power law mass distribution as described in Vollmer et al. (2001). A radius is attributed to each particle depending on its mass. During the disk evolution the particles can have inelastic collisions, the outcome of which (coalescence, mass exchange, or fragmentation) is simplified following Wiegel (1994). This results in an effective gas viscosity in the disk.
As the galaxy moves through the ICM, its clouds are accelerated by ram pressure. Within the galaxy’s inertial system its clouds are exposed to a wind coming from a direction opposite to that of the galaxy’s motion through the ICM. The temporal ram pressure profile has the form of a Lorentzian, which is realistic for galaxies on highly eccentric orbits within the Virgo cluster (Vollmer et al. 2001). The effect of ram pressure on the clouds is simulated by an additional force on the clouds in the wind direction. Only clouds which are not protected by other clouds against the wind are affected.
The non–collisional component of NGC 4438 consists of 65 536 particles, which simulate the galactic halo, bulge, and disk. NGC 4435 is modeled by 49 152 particles. The characteristics of the different galactic components are listed in Table 2.
The total masses of NGC 4438 and NGC 4435 are $`\mathrm{1.7\hspace{0.17em}10}^{11}`$ M and $`\mathrm{1.3\hspace{0.17em}10}^{11}`$ M, respectively. We have chosen to model NGC 4435 as a bulge dominated system with a bulge to disk mass ratio of 3:1. The bulge has an exponential surface brightness profile with a scale length of $`2`$ kpc, which is the H band exponential scale length given by Gavazzi et al. (2003) ($`R_\mathrm{e}=25^{\prime \prime }`$). The disk scale lengths are 3.5 kpc and 0.9 kpc for NGC 4438 and NGC 4435. The assumed value for disk scale length of NGC 4438 is based on the optical diameter of another Virgo spiral of slightly higher luminosity, NGC 4501 ($`D_{25}=6.9^{}=34`$ kpc, and the relation between the optical radius $`R_{25}`$ and the H band disk scale length $`R_\mathrm{e}`$ derived by Gavazzi et al. (2000): $`R_{25}5\times R_\mathrm{e}`$.
The resulting flat rotation velocities of NGC 4438 and NGC 4435 are $`160`$ km s<sup>-1</sup> and $`250`$ km s<sup>-1</sup>. Our choice of the rotation velocities are close to those derived from our observed CO line profiles (Sect. 2.1 and 4). The disk of NGC 4435 has an inclination angle of $`70^{}`$. NGC 4435 has a higher rotation velocity despite its smaller mass, because its disk scale length is smaller than that of NGC 4438. For comparison, Combes et al. (1988) used a mass ratio between NGC 4438 and NGC 4435 of 2:1 and a ratio between the disk scale lengths of 4:1. The particle trajectories are integrated using an adaptive timestep for each particle. This method is described in Springel et al. (2001). The following criterion for an individual timestep is applied:
$$\mathrm{\Delta }t_\mathrm{i}=\frac{20\mathrm{km}\mathrm{s}^1}{a_\mathrm{i}},$$
(1)
where $`a_i`$ is the acceleration of the particle i. The minimum value of $`t_\mathrm{i}`$ defines the global timestep used for the Burlisch–Stoer integrator that integrates the collisional component.
We present 4 different simulations:
1. tidal interaction with NGC 4435 only,
2. tidal interaction and ISM-ISM collision between NGC 4438 and NGC 4435,
3. tidal interaction and ram pressure stripping,
4. tidal interaction, ISM-ISM collision and ram pressure stripping.
For the simulations without an ISM-ISM collision, NGC 4435 is modeled only by a non-collisional component (halo, bulge and disk).
In a different set of simulation we assumed NGC 4435 to be a disk dominated system before the interaction. These simulations show that (i) NGC 4435 cannot be transformed into a bulge dominated system by the tidal interaction and (ii) the tidal effects on NGC 4438 are qualitatively and quantitatively the same as for the simulations with an initially bulge dominated NGC 4435.
### 3.1 ISM-ISM collision
In a widely accepted picture (see, e.g. Kulkarni & Heiles 1988, Spitzer 1990, McKee 1995), the ISM of the Galaxy consists of 5 different phases: the molecular, cold neutral, warm neutral, warm ionized, and hot ionized gas phases. In our case, the molecular and neutral gas phases are of interest.
We have assumed that in the inner 10 kpc of a spiral galaxy, half of the neutral gas mass is molecular whereas the other half is in atomic form. Moreover, deep Hi observations of local spiral galaxies (Braun 1997) showed that 60-90% of the Hı emission is in form of cold atomic hydrogen. In the calculation of the mass fractions we use a cold to total Hi fraction of 70%. More than 80% of the total gas mass within 10 kpc is neutral and more than 70% has temperatures well below 1000 K.
Two gas phases are important for an ISM-ISM collision: the warm ($`T10^4`$ K), diffuse (ionized or not) and the cold ($`T\stackrel{<}{}100`$ K), dense phases, because they represent more than 90 % of the total gas mass. The diffuse and dense phases have volume filling factors of $`\mathrm{\Phi }_\mathrm{V}`$0.3-0.5 and 0.02, respectively. The area filling factor for the diffuse phase is between 0.5 and 1 and that of the dense phase is about 0.1 (Braun 1997). However, only 20% of the total mass is warm but 75% is cold. Our simulations do not distinguish between the two phases. Therefore, we use average volume and area filling factors of $`\mathrm{\Phi }_\mathrm{V}=\left(\mathrm{\Sigma }_i\frac{4\pi }{3}r_{\mathrm{cl},i}^3\right)/V_{\mathrm{gal}}0.05`$ and $`\mathrm{\Phi }_\mathrm{A}=\left(\mathrm{\Sigma }_i\pi r_{\mathrm{cl},i}^2\right)/A_{\mathrm{gal}}0.2`$ in our ISM-ISM collision simulations, which takes the different volume filling factors and mass fractions into account, where $`r_{\mathrm{cl}}`$ is the cloud radius and $`V_{\mathrm{gal}}`$/$`A_{\mathrm{gal}}`$ are the volume/area occupied by the gas cloud complexes. The collision rate can be controled by changing the mass-radius relation to obtain clouds of larger radii, i.e. larger cross section. The most massive, i.e. the largest clouds determine the volume filling factor. The chosen volume filling factor insures that each cloud of NGC 4435 undergoes at least one collision during the ISM-ISM collision between NGC 4438 and NGC 4435. We have run an additional set of simulations where we increased the cloud-cloud collision rate by a factor of 4. The higher collision rate does not change the distribution and dynamics of the high density gas (CO) significantly. However, the dynamics of the low density gas located in the north-western low surface brightness stellar tidal tail are different. Since we do not have observations in this area, it is not possible to discriminate between the simulations with different collision rates.
In general, SPH simulations which include heating and cooling (Barnes & Hernquist 1996, Struck 1997, Tsuchiya et al. 1998) show that cooling is very important (mainly because it is proportional to the square of the gas density) and that cooling times are so small that in general the gas can be regarded approximately as isothermal. This is what we assume in our model. In Sect. 3.5 we show that our results are similar to those of SPH ISM-ISM collision simulations of Struck (1997) and Tsuchiya (1998), which treat the gas as a continuous phase ($`\mathrm{\Phi }_\mathrm{V}1`$).
We now estimate the mass of the warm phase involved in the ISM-ISM collision between NGC 4438 and NGC 4435: since the fraction of the NGC 4438 disk, which is hit by NGC 4435, is about 10% of the total disk area, the mass of the warm ISM of NGC 4438 involved in the interaction is $`0.1\times 0.2\times `$M$`{}_{}{}^{\mathrm{tot}}{}_{\mathrm{gas}}{}^{}`$. Moreover, the total warm gas mass of NGC 4435 is involved in the interaction: $`0.2\times `$M$`{}_{}{}^{\mathrm{tot}}{}_{\mathrm{gas}}{}^{}`$. Thus even asuming that all clouds in NGC 4435 experience a collision with clouds in NGC 4438, only $`3.4\times 10^8`$ M<sub>/odot</sub> due to the collision of the two warm phases is affected, an unknown fraction of which could be found between the two galaxies. If, in addition, there is an efficient interaction between the warm and the cool phases, this value can increase up to $`1.5\times 10^9`$ M. This interaction between a warm, diffuse phase of one galaxy with the cold, dense phase of the other galaxy is not taken into account in our model. We discuss the implications of this estimate in Sect. 3.5.
### 3.2 Ram pressure
For the simulations including ram pressure stripping we assume that the galaxy is on an eccentric orbit within the cluster. The temporal ram pressure profile can be described by:
$$p_{\mathrm{ram}}=p_{\mathrm{max}}\frac{t_{\mathrm{HW}}^2}{t^2+t_{\mathrm{HW}}^2},$$
(2)
where $`t_{\mathrm{HW}}`$ is the width of the profile (Vollmer et al. 2001). We can estimate the maximum ram pressure using the formula of Gunn & Gott (1972):
$$2\pi G\mathrm{\Sigma }_{}\mathrm{\Sigma }v_{\mathrm{rot}}^2R^1\mathrm{\Sigma }p_{\mathrm{ram}},$$
(3)
where $`G`$ is the gravitational constant, $`\mathrm{\Sigma }_{}`$ the stellar surface density. The truncation radius of the gas disk of NGC 4438 of $`R`$=2.5 kpc together with a rotation velocity of $`v_{\mathrm{rot}}=160`$ km s<sup>-1</sup> and a gas surface density of $`\mathrm{\Sigma }=1.5\times 10^{21}`$ cm<sup>-2</sup> implies a maximum ram pressure of $`p_{\mathrm{max}}`$=5000 cm<sup>-3</sup>(km/s)<sup>2</sup>. A realistic orbit in the Virgo cluster then implies $`t_{\mathrm{HW}}`$=50 Myr. In order to produce extraplanar, high gas column density gas, the inclination angle between the disk and the orbital plane must be higher than $`30^{}`$. On the other hand, the efficiency of ram pressure depends on the inclination angle $`i`$ between the galactic disk and the orbital plane (Vollmer et al. 2001). Too low inclination angles (more face-on stripping) result in a very high column density gas in the south of the galaxy center. The azimuthal angle of the galaxy’s motion within the ICM with respect to the orbit of NGC 4435 is chosen so as to reproduce the asymmetric CO and H$`\alpha `$ emission distributions. In the end, we set $`p_{\mathrm{max}}`$=5000 cm<sup>-3</sup>(km/s)<sup>2</sup>, $`t_{\mathrm{HW}}`$=50 Myr, and an inclination angle of $`i`$=68<sup>o</sup>. This resulted in the best model fit to observations. For ram pressure stripping the model clouds have a column density of $`\mathrm{\Sigma }=1\times 10^{21}`$ cm<sup>-2</sup>.
### 3.3 Relative importance of the interactions
In order to investigate the importance of the different interactions, we estimate the accelerations on the ISM of NGC 4438 due to the galactic potential ($`a_{\mathrm{gal}}`$), the direct gravitational pull of NGC 4435 which is close to the tidal acceleration ($`a_{\mathrm{tid}}`$), and ram pressure ($`a_{\mathrm{ram}}`$):
$$a_{\mathrm{gal}}\frac{M_{\mathrm{N4438}}G}{R^2},$$
(4)
$$a_{\mathrm{tid}}\frac{M_{\mathrm{N4435}}G}{(rR)^2},$$
(5)
$$a_{\mathrm{ram}}\frac{p_{\mathrm{ram}}}{\mathrm{\Sigma }_{\mathrm{ISM}}},$$
(6)
where $`G`$ is the gravitational constant, $`v_{\mathrm{rot},\mathrm{N4438}}`$, $`v_{\mathrm{rot},\mathrm{N4435}}`$ are the rotation velocities of NGC 4438 and NGC 4435, respectively, $`R`$ is the radius of tidal influence, $`r`$ is the distance between the two galaxies, $`p_{\mathrm{ram}}`$ is the ram pressure, and $`\mathrm{\Sigma }_{\mathrm{ISM}}`$ is the gas surface density of NGC 4438 at the radius $`R`$. For the tidal interaction to be important the following condition must be fulfilled:
$$\frac{a_{\mathrm{tid}}}{a_{\mathrm{gal}}}=\frac{M_{\mathrm{N4435}}}{M_{\mathrm{N4438}}}\frac{R^2}{(rR)^2}1.$$
(7)
The mass fraction between the two galaxies is about 0.75. With an impact parameter of $`r=8.5`$ kpc the parts of NGC 4438 at radii $`R>4`$ kpc are affected. For ram pressure to be important we obtain the following condition:
$$\frac{a_{\mathrm{ram}}}{a_{\mathrm{gal}}}=\frac{p_{\mathrm{ram}}R\mathrm{cos}i}{\mathrm{\Sigma }_{\mathrm{ISM}}v_{\mathrm{rot},\mathrm{N4438}}^2}1,$$
(8)
where $`i`$ is the angle between the two vectors $`a_{\mathrm{ram}}`$ and $`a_{\mathrm{gal}}`$. A peak ram pressure of $`p_{\mathrm{ram}}=5000`$ cm<sup>-3</sup>(km s<sup>-1</sup>)<sup>2</sup>, a rotation velocity of $`v_{\mathrm{rot},\mathrm{N4438}}=160`$ km s<sup>-1</sup>, $`i=0^{}`$, and a radius of $`R=4`$ kpc results in a critical gas surface density of $`\mathrm{\Sigma }_{\mathrm{ISM},\mathrm{crit}}2\times 10^{21}`$ cm<sup>-2</sup>. Thus, ISM with surface densities smaller than $`\mathrm{\Sigma }_{\mathrm{ISM},\mathrm{crit}}`$ is stripped by ram pressure at radii greater than 4 kpc. This is close to the observed stripping radius of $`50^{\prime \prime }`$ (Sect. 3.9). On the other hand, giant molecular clouds have column densities much higher than $`10^{21}`$ cm<sup>-2</sup> (several $`10^{22}`$ cm<sup>-2</sup>) and thus should not be stripped.
A single simulation requires about 100 hours on a 1.2 GHz and 1024 MB RAM memory PC. In the following we describe the temporal evolution of the 4 different simulations. The models presented here are those with the best results (i.e. best set of collision parameters, see Sect. 3.2 and 3.1) for each of the four scenarios.
### 3.4 Tidal interaction only
For the tidal interaction we follow Combes et al. (1988) and chose a close, rapid, retrograde encounter between NGC 4438 and NGC 4435. The closest approach occurs at $`t=165`$ Myr. The impact parameter is 9 kpc and the velocity difference is $`\mathrm{\Delta }v=840`$ km s<sup>-1</sup>. We verified their best fit model in varying the impact parameters from 5 to 10 kpc and the inclination angle between the orbital plane of NGC 4435 and the disk plane of NGC 4438 between $`20^{}`$ and $`40^{}`$.
There is no significant change in the distribution of stars and gas of NGC 4438 before the passage of NGC 4435 at $`t=165`$ Myr. Being more extended, NGC 4438 gets most damaged during the interaction. A prominent tidal tail is formed north of the galaxy center. In addition, a second tidal tail in the south forms a loop whose northern/western edge moves towards the north/west while its southern/eastern edge does not change position. There is no stellar or gas bridge between NGC 4438 and NGC 4435. At the end of the simulation ($`t=260`$ Myr) the distribution of gas and stars of NGC 4438 is highly asymmetric with a sharply truncated southern and eastern edge, a prominent tidal arm towards the north, and an extended component to the west (Fig. 5).
NGC 4435 is located close to its observed location. The difference in radial velocities of NGC 4438 and NGC 4435 is $`\mathrm{\Delta }v_\mathrm{r}=670`$ km s<sup>-1</sup>, compatible with the observed value of $`730`$ km s<sup>-1</sup>.
The final stellar distribution at the timestep which we compare to our CO observations is discussed in Sect. 3.8.
### 3.5 Tidal interaction and ISM-ISM collision
Since the collision between the ISM of NGC 4438 and NGC 4435 only concerns the gaseous component whose mass is small compared to that of the non-collisional component, the evolution of the stellar disk is the same as for the simulations of the tidal interaction alone (Sect. 3.4). After the ISM-ISM collision at $`t=165`$ Myr, two gaseous bridges are formed between NGC 4438 and NGC 4435 (Fig. 6).
The southern arm is more prominent and longer lasting than the northern component. The total gas mass in the gas bridge is about $`3.4\times 10^8`$ M which represents $``$25% of the gas mass of NGC 4435. The surface density within the whole bridge is high (5-10 Mpc<sup>-2</sup>) until 45 Myr after the collision ($`t=165`$ Myr). Still $`20`$ Myr later, the surface density in the middle of the bridge drops to $`3`$ Mpc<sup>-2</sup> and higher surface density gas (5 Mpc<sup>-2</sup>) is only found along the minor axis at distances smaller than 4 kpc to the west of the disk of NGC 4438. At $`t=260`$ Myr (95 Myr after the collision) this close, western extraplanar gas has only a mean surface density of about 2 Mpc<sup>-2</sup>. The rest of the gas in the bridge has a lower surface density. At the end of the simulation ($`t=260`$ Myr) the gas distribution of NGC 4438 is very similar to that of the simulations of the tidal interaction alone, i.e. the ISM-ISM collision is neither able to remove the gas from the northern tidal tail, nor can it produce extraplanar, high surface density gas $`100`$ Myr after the collision.
Our simulations can be compared to the SPH ISM-ISM collision simulations of Struck (1997) and Tsuchiya et al. (1998). Since SPH particles cannot interpenetrate by definition, the intruder galaxy makes a hole in the primary galaxy, which is rapidly filled due to differential rotation (Struck 1997). Whereas Struck (1997) explains the effects in more detail, Tsuchiya et al. (1998) give more quantitative information. In particular, Tsuchiya et al. (1998) show the evolution of a head-on collision of a primary galaxy with a massive companion in their Fig. 6. The relative velocity is $`300`$ km s<sup>-1</sup>, thus a factor 3 smaller than in our case and their companion mass is 25% of the primary galaxy. They find 25% of the mass of the companion’s gas mass in the bridge. In addition, the surface density of the gas in the bridge drops drastically within 40 Myr (from $`T=15=195`$ Myr to $`T=18=234`$ Myr) and $`70`$ Myr after the collision the gas surface density in the bridge is small. The column density evolution of the gas in the bridge is qualitative and, as far as we can see, in quantitative agreement with our findings. Interestingly, position-velocity plots of their simulations at $`T=15`$ or $`T=18`$ (their Fig. 9) show that the gas in the bridge near the primary galaxy has negative velocities with respect to the primary’s systemic velocity, because the bridge gas is falling back onto the primary galaxy.
If an ISM-ISM collision were responsible for the observed extraplanar CO, more than $`10^8`$ M of cold and dense gas must have been involved in the collision, since the involved diffuse ISM is not massive enough to account for the observed extraplanar, molecular gas mass (see Sect. 3.1). It is not likely that the extraplanar gas still has a high surface density ($`10`$ Mpc<sup>-2</sup>) as it is observed $`100`$ Myr after the collision. Furthermore, $`100`$ Myr after the collision, the gas which was involved in the collision should be falling back onto NGC 4438 producing lines at negative velocities in the south-west of the galaxy center, which is not observed. We conclude that an ISM-ISM collision is not responsible for the observed molecular gas distribution and kinematics.
### 3.6 Tidal interaction and ram pressure stripping
In this simulation we did not assign gas to NGC 4435 but include ICM ram pressure. The choice of the free parameters of the ISM-ICM interaction are given in Sect. 3.2. Since ram pressure acts only on the ISM of NGC 4438, the evolution of the stellar disk is the same as for the other simulations (Sect. 3.4 and 3.5). For the gas, the situation changes completely. Soon after the passage of NGC 4435, the gas distribution is distorted by the action of ram pressure, which pushes the ISM of NGC 4438 to the west (Fig. 7).
The probability for a close galaxy encounter is highest in the cluster core due to its strongly peaked density of elliptical and S0 galaxies (Schindler et al. 1999). In addition, in the cluster core the galaxy velocity and the ICM density, i.e. ram pressure, are highest. Thus the small time difference between the galaxy collision and maximum ram pressure of $`\mathrm{\Delta }t=85`$ Myr is a natural consequence of the maximized probability for galaxy collisions and maximized ram pressure in the cluster core. The direction of ram pressure is natural for a highly eccentric orbit of NGC 4438 in the Virgo cluster (Vollmer et al. 2001). The trajectory of NGC 4435 can also be naturally explained by a less eccentric orbit.
### 3.7 Tidal interaction, ISM-ISM collision and ram pressure stripping
This simulation contains all three interactions: (i) the tidal interaction, (ii) the ISM-ISM collision and (iii) ram pressure stripping. The evolution of the gaseous component is very similar to that of the previous simulation (tidal interaction and ram pressure stripping), because ram pressure stripping is the most energetic and thus most important effect. Due to the plot mode (each particle is represented by a black dot), the evolutionary plot of this simulation is indistinguishable from that of Fig.7. However, the differences can be seen in the animation or in the final gas distribution (Sect. 3.8).
### 3.8 The final stellar and gas distributions
The final distribution of the stellar content of NGC 4438 and NGC 4435 is shown in Fig. 8. The overall distribution of NGC 4438 is close to the observed one (Fig. 1), i.e. the major characteristics are reproduced: (i) the stellar disk is truncated to the south, (ii) there is a prominent northern tidal arm and (iii) displaced stellar debris to the west of the galaxy’s main disk. In addition, the location and radial velocity of NGC 4435 is close to observations (Sect. 3.4). The reproduction of the stellar component is imperfect in that the model arm to the North is offset. Whereas the observed arm is located on the galaxy’s major axis, the model arm is parallel to the major axis, but $`3`$ kpc offset to the east. The aim here is not to present a perfect model of the tidal interaction but to study the influence of the different kinds of interactions, for which the model stellar distribution is amply sufficient.
We produced a model cube of the gas within the inner $`20\times 20`$ kpc of our final snapshot. From this model cube we extracted column density maps and model spectra at different resolutions. Since our model gas distribution of NGC 4438 does not include the high density core that is observed, we added this component ad hoc. If we had included this component into the numerical simulations, the computational costs would have been too high to carry out our systematic study. The a posteriori addition of this core, does not change our conclusions, because the core does not extend further out than the model gas distribution without the core at $`t=260`$ Myr. The core mainly enhances the surface brightness of the central disk. In addition, the initial model rotation curve of NGC 4438 is not as steep as it is observed because of numerical resolution. The additional core component takes this into account.
An overlay of the stellar component with the gas distribution convolved to a resolution of $`7^{\prime \prime }=580`$ pc for all four simulations is shown in Fig. 9.
The lower panels include ram pressure while the upper panels do not. The final gas distributions of the simulations including ram pressure are very different from those without ram pressure, whereas the effect of the ISM-ISM collision on the gas distribution is minor. This shows that ram pressure is much more important for the overall gas dynamics than the ISM-ISM collision. The gas distributions of the simulations of the tidal interaction with and without an ISM-ISM collision are similar. Both distributions show two tidal arms in the north. Moreover, high column density gas is found to the west of the galactic disk, which belongs to the second spiral arm extending to the south-west. Combes et al. (1988) already observed this effect.
The two simulations including ram pressure stripping also lead to similar distributions. The ISM-ISM collision enhances the velocity dispersion and decreases the volume density of NGC 4438’s gas which is involved in the collision. Due to the decrease of the particle density the ICM penetration length into the ISM of NGC 4438 increases and thus the ram pressure efficiency increases. The net effect is that more gas is pushed to the west by ram pressure when an additional ISM-ISM collision occurs. The gaseous disk of NGC 4438 is truncated at a radius of $`40^{\prime \prime }`$ and the gas of the outer disk is pushed to the west. The displaced gas of highest surface density is found in the south-west and north-west close to the truncated disk (see also Sect. 3.9).
### 3.9 Distribution of the neutral ISM
Since the pointings of our observations form an irregular grid, we used a Delaunay triangulation within IDL (TRIANGULATE and TRIGRID functions, see the IDL reference guide) to interpolate the data and to derive a column density map (Fig. 10).
Within the galactic disk we observe an asymmetry along the major axis, i.e. there is more flux in the north than in the south. The western, extended emission shows the same asymmetry, i.e. there is a local maximum in the north. This emission closely traces the dust absorption features observed towards the west of the galactic disk of NGC 4438 (see e.g. Fig. 1 of Kenney et al. 1995).
The column density maps of our 4 simulations are shown in Fig. 11. Their spatial sampling is complete. For the comparison with the observed column density map one has to keep in mind that these are irregularly gridded and do not cover the entire region.
Only the simulations that include ram pressure stripping reproduce the observed CO gas distribution. In particular, the simulations without ram-pressure cannot account for the observed strong CO(1–0) lines west of the galaxy center. The simulation including all interactions (tidal, ISM-ISM collision and ram pressure stripping) is also consistent with our observations. The main difference with the tidal/ram pressure simulation is that the column density of the western extraplanar gas is lower if an ISM-ISM collision is included, because the gas clouds extracted and “heated” from NGC 4438 during this ISM-ISM collision are efficiently stripped by ram pressure.
The main difference between our CO(1–0) observations and the simulations including ram pressure stripping is that the model gas surface density is greater than observed in the southern, extraplanar gas. Since this is not observed in CO, we suspect that part of this gas is heated and ionized with its molecular gas dissociated. A part of this “missing” gas might be visible in X-rays or in the radio continuum.
## 4 The need for ram pressure
We showed in Sect. 3.9 that only simulations including ram pressure can account for the fact that there is almost no neutral gas detected in the northern tidal tail. Boselli et al. (2005) observed NGC 4438 with the GALEX satellite. They found UV emission in the northern and southern tidal tail where no gas is detected. By fitting a spectral energy distribution locally using NIR, optical and UV data they reconstructed the star formation history of different regions. For the northern and southern tidal tail the data is consistent with a star burst $`100`$ Myr ago. This is entirely consistent with our scenario for the tidal interaction. Since there is UV radiation detected in the tidal tails, gas was associated with these regions before the interaction, which justifies our initial gas distribution of NGC 4438.
In the following we show that the observed CO(1–0) spatial and velocity asymmetry and extraplanar double lines (Fig. 3) can be reproduced by a model that includes ram pressure stripping (Fig. 12). We recall that, in order to reproduce the observed stellar distribution, we adopted the orbital parameters of Combes et al. (1988). We verified their best fit model in the way described in Sect. 3.4. With these parameters, the parameters for the tidal interaction and the ISM-ISM collision are fixed. Concerning ram pressure, the open parameters are the maximum ram pressure, the inclination angle between the disk and the orbital plane and the azimuthal projection angle (Vollmer et al. 2001; see also Sect. 3.2). The maximum ram pressure is chosen such that there is extraplanar high column density gas as observed. The inclination angle is chosen such that there is a a minimum of gas to the south-west. When we decreased the inclination angle, there was far too much extraplanar gas in the SW. The azimuthal projection angle is chosen such that the position and radial velocity of NGC 4435 are close to observations. We call this model our best fitting model.
### 4.1 The best fitting model: tidal interaction and ram pressure stripping
We took the last snapshot of Fig. 7 and produced a model cube out of it (see also Sect. 3.8). CO model spectra with a resolution of $`21^{\prime \prime }`$ were then extracted on a equidistant grid with a cell size of $`20^{\prime \prime }`$. We only show model spectra with a non zero flux (Fig. 12).
The model spectra show the same overall east-west asymmetry as is observed (see also Sect. 3.9). At the eastern edge of the galactic disk the model shows weak lines that are consistent with our CO(1–0) observations. In particular, the observed spectrum about $`30^{\prime \prime }`$ north-east of the galaxy center is well reproduced, even though the maximum of the model spectrum is weaker than the observed one. We find strong model lines to the north-west of the galaxy center as observed. The model spectra in the south-west of the galaxy peak at velocities greater than zero as observed. In addition, the model also reproduces the observed double line profile of the south-western spectrum. Whereas the two observed peaks have the same fluxes, the model peak at high velocities is much stronger than the peak at low velocities. The double line profile is still visible in the model spectra further to the west. The two lines trace two spatially distinct regions in the line of sight. The line at negative velocities traces gas which rotates in the galactic disk, whereas the line at positive velocities traces gas which is accelerated by ram pressure.
For more detailed comparison we show selected observed and model spectra in Fig. 13. The observed spectra (solid line) follow the tidal and ram pressure simulations (dotted line) more closely than the tidal interaction alone (dashed line). Adding the ISM-ISM collision has little effect. On both sides of the galaxy the gas velocities are not well represented by the tidal only simulations whereas they are reasonably reproduced when ram pressure is included. Molecular gas is present in many positions to the West where the tidal forces do not bring gas but is not present on or near the northern or eastern part of the main body of NGC 4438, where strong lines are expected if ram pressure is not acting. NGC 4435 penetrates NGC 4438 in a region which has now rotated to the northern side. A ”smooth” or clumpy ISM-ISM collision $`100`$ Myr ago could not change the rotation velocities in the southern part on or near the major axis (i.e. pos. $`35^{\prime \prime }`$,$`52^{\prime \prime }`$) but the observations show that to the South the CO velocities are much higher than they would be in the absence of ram pressure.
Most of the molecular gas has been efficiently stripped from what is now the northern part of NGC 4438. Only the simulations with ram pressure show this (see also Fig. 11) and a collision with the ISM (dense or diffuse) of NGC 4435 could not empty a large region of its dense gas. While the intensity ratio is not reproduced, it is very interesting that at the ($`24^{\prime \prime }`$, $`13^{\prime \prime }`$) position the ram pressure simulation reproduces the velocities whereas almost no gas is at the velocity of the big (single) peak predicted by the tidal interaction only. The ($`43^{\prime \prime }`$, $`8^{\prime \prime }`$) position is similar. Here as well, the ISM-ISM collision has little effect.
The three spectra to the right of Fig. 13 are all considerably off the plane of NGC 4438 to the west. The simulations without ram pressure show no flux here or, for the ($`54^{\prime \prime }`$, $`46^{\prime \prime }`$) position, at a significantly lower velocity as in the spectra to the right.
The observed central and integrated spectra together with the model spectra are shown in Fig. 14.
In the model spectrum the double horn component at high absolute velocities in the central spectrum is mainly due to the fast rotating core of NGC 4438 (see Sect. 3.8). On the other hand, the bump at $`v140`$ km s<sup>-1</sup> is due to gas that has been accelerated by ram pressure to positive velocity. Thus, we conclude that the observed bump at $`v100`$ km s<sup>-1</sup> is due to ram pressure stripping. The simulated bump is shifted by $`40`$ km s<sup>-1</sup> with respect to the observed bump.
The integrated spectrum shows two main characteristics: (i) a major bump around the galaxy’s systemic velocity $`v=70`$ km s<sup>-1</sup> and (ii) an asymmetry with less gas on the low velocity (approaching) side. The simulations including only a tidal interaction (dashed line) show neither a central bump nor the observed asymmetry in velocity. On the contrary, the model integrated spectrum without ram-pressure shows a reversed asymmetry with more gas at low velocities, because most of the spectra summed are on the approaching side. On the other hand, the integrated spectrum of the simulations including a tidal interaction and ram pressure stripping do reproduce the observed asymmetry and and the central bump. However, the model bump has the same peak value as the high velocity peak and is shifted to smaller velocities.
We conclude that the model including a tidal interaction and ram pressure stripping reproduces the major characteristics of the observed spectra:
* the east west asymmetry,
* the displacement of the peaks of the spectra to higher velocities in the south-west of the galaxy,
* the double line profiles located to the west and south-west of the galaxy center.
Since the detection of extraplanar CO shows that molecular gas is displaced, this implies that either the bulk of the molecular gas is taken along with the atomic gas via ram pressure due to the diffuse ISM or magnetic field coupling, or the molecular clouds are left behind, form stars rapidly and then are destroyed by the energy input due to star formation. This means that the Gunn & Gott criterion might not be applicable to single clouds, but only to the ISM within the galactic disk as a whole.
### 4.2 Alternative scenarios
We have chosen the model including a tidal interaction and ram pressure stripping as the best fitting model, although including an ISM-ISM collision is about as good. In this section we present the alternative scenarios (tidal interaction alone; tidal interaction and ISM-ISM collision; tidal interaction, ISM-ISM collision and ram pressure stripping).
The spectra from the simulation including the tidal interaction alone are shown in Fig. 15.
High column density gas is dragged to the west by the tidal interaction. However, the spectra with non zero flux do not extend further to the west than $`60^{\prime \prime }`$ (see also Fig. 11). Contrary to observations, there are strong lines at the eastern edge of the galactic disk, but not to the north-west of the galaxy center. All spectra located south of the galaxy center peak at velocities smaller than zero in contrast to the observed spectra in this region. At the north-eastern edge of the optical disk multiple line profiles can be seen. These are due to the superposition of the disk gas and gas located in the north-western tidal arm (see Fig. 11). There are no double line profiles in the south-west of the galaxy.
The spectra from the simulation adding ISM-ISM collision are shown in Fig. 16.
The spectra within and east of the optical disk are very similar to those of the simulation with a tidal interaction alone (Fig. 15). The spectral profiles at and beyond the western edge of the optical disk occupy the same velocity range as those of Fig. 15, but their peak fluxes are lower. In contrast to the observed spectra, the model spectra in the south-west of the galaxy center show peaks at negative velocities.
The spectra from the simulation including the tidal interaction, an ISM-ISM collision and ram pressure stripping is shown in Fig. 17.
The spectra are very similar to those of the simulation including a tidal interaction and ram pressure stripping. Again, the effect of the ISM-ISM collision is minor compared to that of the tidal interaction.
We conclude that the simulations without ram pressure do not reproduce the major characteristics of the CO(1–0) observations: (i) the spectra at the eastern edge of the optical disk, (ii) the strong lines to the north-west of the galaxy center, (iii) the double line profile in the south-west of the galaxy center, and (iv) the shift of the line profile in the south-west to positive velocities. Only the simulations with ram pressure stripping are able to reproduce these features.
## 5 The northern tidal tail
We have detected $`2\times 10^7`$ M of molecular gas in the northern tidal tail and the mass limit for atomic gas in this region is a few $`10^7`$ M (Hibbard et al. 2001). The simulations including ram pressure stripping do not show any gas in or near this region (see Fig. 7) nor is any Hi present. The gas mass in the northern tidal arm region of the simulations without ram-pressure (Fig.15) is $`\mathrm{7\hspace{0.17em}10}^8`$ M. Thus, the observed molecular gas mass given the absence of Hi in this region represents a few percent of the gas mass that would be there without ram pressure stripping. This is consistent with a picture where the gas, which was in the form of molecular clouds during the phase of maximum ram pressure ($`t>150`$ Myr), was not affected by ram pressure due to its high gas surface density. This implies that the observed molecular clouds in the northern tidal arm region are rather long-lived (several 10 Myr).
Fig. 18 shows model and observed spectra of two positions within the northern tidal arm region. The model spectra are rescaled to the show a maximum comparable to that of the observed spectra. Note, however, that the mass in the model spectrum is more than 30 times higher than the mass derived from the observed spectra. Since the model northern tidal arm is $`3`$ kpc offset from the galaxy’s major axis to the east, the positions of the model spectra are also shifted by 3 kpc to the east.
Indeed, the model spectra show the same peak velocity and the same line width. The simulation including an ISM-ISM collision fits the observations at ($`38^{\prime \prime }`$, $`134^{\prime \prime }`$) slightly better, because the secondary peak seen in the simulations of the tidal interaction alone is less present.
Thus some giant molecular clouds (a few percent of the total gas mass) seem to have decoupled from the ram pressure wind. They were left behind in the otherwise gas free disk and survived for several 10 Myr. We cannot say if this is a phenomenon exclusively related to the tidal tail or if giant molecular clouds in other parts of the disk also decoupled from the wind. Our observations indicate that the bulk of the molecular gas is stripped by ram pressure. Since only a few percent of the total gas mass decouple, decoupled gas would only represent $`7\times 10^7`$ M.
## 6 Conclusions: the history of NGC 4438
New <sup>12</sup>CO(1–0) observations of the NGC 4438/4435 system are presented. For the first time CO is detected in NGC 4435. As already shown by Combes et al. (1988), the distribution of molecular gas is highly truncated within the disk of NGC 4438 and we find an extraplanar component up to $`1.5^{}`$ to the west of the galaxy center. Within this extraplanar molecular gas we find double line profiles at distances up to $`40^{\prime \prime }`$ to the west and south-west of the galaxy center. In addition the lines in the south of NGC 4438 are all redshifted with respect to galactic rotation. We argue that asymmetry of the molecular gas distribution, the double line profiles and the redshifted lines are characteristic for ram pressure stripping of NGC 4438.
The combination of our new CO(1–0) observations with detailed numerical simulations leads to the following interaction scenario for the NGC 4438/NGC 4435 system: NGC 4435 passed through the disk of NGC 4438 $`100`$ Myr ago at a radial distance of $`510`$ kpc. The encounter was rapid ($`\mathrm{\Delta }v830`$ km s<sup>-1</sup>) and retrograde (see also Combes at al. 1988). With an impact parameter $`<10`$ kpc an ISM-ISM collision is unavoidable. Its importance depends on the initial gas distributions in NGC 4435 and NGC 4438. The estimated extent of NGC 4435’s observed gas disk is $`\stackrel{<}{}1`$ kpc. In our simulations the gas disk of NGC 4438 had an initial extent of $`10`$ kpc. Even with this initial extent the influence of an ISM-ISM collision on the final gas distribution and velocities is small compared to that of ram pressure stripping. NGC 4438 evolves on an eccentric orbit within the Virgo cluster. We observe the galaxy $`10`$ Myr after its closest passage to the cluster center (M87). Our model infers a total velocity of $`v_{\mathrm{tot}}2000`$ of NGC 4438 with respect to the cluster mean. The galaxy is located at a total distance of $`350`$ kpc from the cluster center. With the assumed maximum ram pressure the ICM density at the position of NGC 4438 is $`n_{\mathrm{ICM}}10^3`$ cm<sup>-3</sup> which is high but compatible with the ICM density at the distance of NGC 4438 derived from X-ray observations (Schindler et al. 1999). Ram pressure plays a key role for the evolution of the gaseous component of NGC 4438 together with the tidal interaction. The displacement of the line profiles to higher velocities in the south-western region of the galaxy, the lack of CO emission in the eastern optical disk, and the presence of double line profiles in the south-west of the galaxy center are clear signs of ram pressure stripping. There is evidence that in the northern tidal arm region a few percent of the ISM survived ram-pressure stripping in the form of dense molecular clouds (Sect. 5). These clouds must be stable for several 10 Myr. Vollmer et al. (2001) claimed that gas clouds located within the ICM can be stable and are almost entirely molecular, i.e. they are selfgravitating, heated by the X-ray background and cooled by Cii and Oi line emission. With a typical column density of about $`10^{22}`$ cm<sup>-2</sup> these clouds resist ram pressure. We suggest that the observed molecular clouds in the northern tidal arm regions are of this kind.
NGC 4438 has been greatly affected by both the tidal interaction with NGC 4435 and ram pressure stripping due to the rapid motion of NGC 4438 through the intracluster medium. It is difficult, however, to be more precise than we have been about the relative roles of tides and ram pressure given their partial degeneracy and the uncertainties in the simulations and observations.
###### Acknowledgements.
Based on IRAM observations. IRAM is supported by INSU/CNRS (France), MPG (Germany), and IGN (Spain). We made use of a DSS image. The Digitized Sky Survey was produced at the Space Telescope Science Institute under U.S. Government grant NAG W-2166. The images of these surveys are based on photographic data obtained using the Oschin Schmidt Telescope on Palomar Mountain and the UK Schmidt Telescope. The plates were processed into the present compressed digital form with the permission of these institutions. This research has made use of the GOLD Mine Database (Gavazzi et al. 2003). |
warning/0507/hep-ex0507073.html | ar5iv | text | # Measurement of Transverse Single-Spin Asymmetries for Mid-rapidity Production of Neutral Pions and Charged Hadrons in Polarized p+p Collisions at √𝑠=200 GeV
PHENIX Collaboration
## Abstract
The transverse single-spin asymmetries of neutral pions and non-identified charged hadrons have been measured at mid-rapidity in polarized proton-proton collisions at $`\sqrt{s}=200`$ GeV. The data cover a transverse momentum ($`p_T`$) range 0.5-5.0 GeV/$`c`$ for charged hadrons and 1.0-5.0 GeV/$`c`$ for neutral pions, at a Feynman-$`x`$ ($`x_F`$) value of approximately zero. The asymmetries seen in this previously unexplored kinematic region are consistent with zero within statistical errors of a few percent. In addition, the inclusive charged hadron cross section at mid-rapidity from $`0.5<p_T<7.0`$ GeV/$`c`$ is presented and compared to NLO pQCD calculations. Successful description of the unpolarized cross section above $`2`$ GeV/$`c`$ using NLO pQCD suggests that pQCD is applicable in the interpretation of the asymmetry results in the relevant kinematic range.
proton, spin, polarization, asymmetry
The measurement of transverse single-spin asymmetries (SSAs) in proton-proton collisions and deep-inelastic lepton-nucleon scattering (DIS) probes the quark and gluon structure of transversely polarized nucleons. Large transverse SSAs have been observed in a number of spin-dependent proton-proton scattering experiments at energies ranging from $`\sqrt{s}=510`$ GeV. Asymmetries approaching 30% were observed in inclusive pion production at transverse momentum ($`p_T`$) up to 1.2 GeV/$`c`$ and Feynman-$`x`$ ($`x_F`$) up to 0.8 Dragoset et al. (1978); Allgower et al. (2002). At mid-rapidity and $`x_T=\frac{2p_T}{\sqrt{s}}`$ up to 0.8, asymmetries were also observed in inclusive $`\pi ^0`$ and $`\pi ^+`$ production but not in $`\pi ^{}`$ production Antille et al. (1980); Saroff et al. (1990); Apokin et al. (1990). At higher center-of-mass energies of 20 and 200 GeV, $`\pi ^+`$, $`\pi ^{}`$, and $`\pi ^0`$ asymmetries were found to persist at large $`x_F`$ Adams et al. (1991a, b, 2004) while the asymmetry in $`\pi ^0`$ production at mid-rapidity was found to be consistent with zero at $`\sqrt{s}=20`$ GeV and for $`p_T<4`$ GeV/$`c`$ Adams et al. (1996). Non-zero transverse asymmetries have also been observed in semi-inclusive DIS experiments Airapetian et al. (2000, 2002, 2004).
Three different mechanisms have been studied as the possible origin of transverse SSAs in hadron collisions at high energies: (1) Transversity distributions, the quark spin distributions in a transversely polarized proton, can give rise to SSAs in combination with spin-dependent fragmentation functions (FFs), e.g. the Collins function Collins (1993). Spin-dependent FFs serve as analyzers for the transverse spin of the struck quark. (2) Quark and gluon distributions that are asymmetric in the transverse intrinsic parton momentum, $`k_T`$, first suggested by Sivers Sivers (1990), can lead to SSAs. (3) Alternatively, interference between quark and gluon fields in the initial or final state can also generate SSAs Qiu and Sterman (1999); Kanazawa and Koike (2000). Sivers parton distributions can exist both for quarks and gluons, and a possible connection to orbital angular momentum of partons in the nucleon has been suggested Sivers (1990); Burkardt and Hwang (2004).
It is expected that SSAs measured at the Relativistic Heavy Ion Collider (RHIC) result from a combination of these three effects (see Adams et al. (2004) and references therein). Model calculations leading to predictions for the Sivers and transversity distributions have been performed to describe existing data at forward rapidities. Precision measurements of SSAs in different regions of $`x_F`$ and $`p_T`$ and their QCD analysis may serve to quantify contributions from the competing mechanisms. In this Letter we present first measurements of transverse single-spin asymmetries at mid-rapidity and collider energies.
These data were collected during the 2001-2 polarized proton run at RHIC, in which approximately 0.15 $`pb^1`$ of integrated luminosity were collected using the PHENIX detector. Two beams of 55 bunches of polarized protons, with approximately $`5\times 10^{10}`$ protons per bunch, were injected into RHIC and accelerated to 100 GeV each.
Measurements of the unpolarized production of charged hadrons and of the spin-dependent production of both neutral pions and charged hadrons were made in the central arms of the PHENIX detector. These cover a pseudorapidity range of $`|\eta |<0.35`$ and two azimuthal angle intervals of $`\mathrm{\Delta }\varphi =90^{}`$, offset $`33.75^{}`$ from vertical Adcox et al. (2003a).
A minimum-bias (MB) collision trigger and the vertex position in the beam direction are provided by two beam-beam counters (BBCs) Allen et al. (2003). The BBCs, which cover 2$`\pi `$ in azimuth and $`3.0<|\eta |<3.9`$, are sensitive to charged particles and select approximately half of the total inelastic proton-proton cross section. A $`\pm 30`$ cm event vertex cut was applied for all analyses, corresponding to the central arm acceptance. The approximate vertex resolution was 2 cm in the beam direction.
Charged-particle tracks from MB events were reconstructed using a drift chamber and pad chambers Adcox et al. (2003b) as well as the collision vertex, which is the assumed point of origin because the tracking chambers are placed outside the magnetic field. Thus charged particles that do not originate at the vertex have incorrectly reconstructed momentum, leading to low-momentum, long-lived particle decays (e.g. $`K^\pm `$, $`K_L^0`$) and conversion electrons as the two main sources of background.
For the charged hadron cross section, approximately 17 million MB events were analyzed. The luminosity was measured as $`N_{\mathrm{BBC}}/\sigma _{\mathrm{BBC}}`$ with $`\sigma _{\mathrm{BBC}}=21.8`$ mb Adler et al. (2003), accounting for the fraction of the yield for which the MB condition was satisfied. Backgrounds were estimated and subtracted statistically following the method of Adler et al. (2004): conversion electrons were estimated using the different response of the ring-imaging $`\stackrel{ˇ}{\mathrm{C}}`$erenkov detector (RICH) Aizawa et al. (2003) to electrons and charged pions, and decay particles were estimated using the track bend in the residual magnetic field in the tracking detectors. Weak decays of short-lived particles, mainly $`K_S^0`$, $`\mathrm{\Lambda }`$, and $`\overline{\mathrm{\Lambda }}`$, remain, especially when they decay close to the vertex. Based on a Monte Carlo simulation, the reported cross section was reduced by 7% over the entire $`p_T`$ range to correct for these decays.
The unpolarized cross section for inclusive charged hadron production at mid-rapidity is presented in Fig. 1 and Table 1. The dominant systematic uncertainty for $`p_T>5`$ GeV/$`c`$ is from the background subtraction, while for $`p_T<5`$ GeV/$`c`$ it is due to the weak-decay correction. There is a 9.6% normalization uncertainty due to the luminosity measurement. In Fig. 1 the cross section is compared to a next-to-leading-order (NLO) pQCD calculation using the CTEQ6M Pumplin et al. (2002) parton distribution functions and Kniel-Kramer-Pötter FFs Kniehl et al. (2001) and found to be consistent above $`p_T2`$ GeV/$`c`$. The unpolarized cross sections for mid-rapidity and forward production of neutral pions have also been measured in 200-GeV proton-proton collisions at RHIC Adler et al. (2003); Adams et al. (2004) and have been found to agree well with NLO pQCD calculations Aversa et al. (1989); de Florian (2003); Jäger et al. (2003). The agreement between all of these unpolarized measurements and the theoretical calculations indicates that NLO pQCD is applicable in interpreting polarized data at $`\sqrt{s}=200`$ GeV and provides a solid theoretical foundation for the study of the spin structure of the proton at RHIC.
The stable spin direction of the protons through acceleration and storage is vertical, and there is an approximately equal number of bunches filled with the spin of the protons up as there is down.. With both beams polarized, single-spin analyses were performed by taking into account the spin states of one beam, averaging over those of the other. The beam polarization at 100 GeV was obtained using the same analyzing power ($`A_N^{pC}`$) in proton-carbon elastic scattering in the Coulomb-nuclear interference region measured at 22 GeV (see Jinnouchi et al. (2003) and references therein), near RHIC injection energy. The average beam polarization was 15$`\pm `$5%.
The left-right transverse single-spin asymmetry, $`A_N`$, can be extracted using
$$A_N=\frac{1}{P_\mathrm{b}}\left(\frac{\sigma ^{}\sigma ^{}}{\sigma ^{}+\sigma ^{}}\right)=\frac{1}{P_\mathrm{b}}\left(\frac{N^{}N^{}}{N^{}+N^{}}\right),$$
(1)
where $`P_\mathrm{b}`$ is the beam polarization, $`\sigma ^{}`$ ($`\sigma ^{}`$) the production cross section when the protons in the bunch are polarized up (down), $`N^{}`$ ($`N^{}`$) the experimental yield from up- (down-) polarized bunches, and $`=^{}/^{}`$ the relative integrated luminosity of bunches of opposite polarization sign. The above formula as written applies to yields observed to the left of the polarized beam. An overall minus sign is required for yields observed to the right of the polarized beam. Alternatively, we derive the asymmetry using
$$A_N=\frac{1}{P_\mathrm{b}}\left(\frac{\sqrt{N_L^{}N_R^{}}\sqrt{N_L^{}N_R^{}}}{\sqrt{N_L^{}N_R^{}}+\sqrt{N_L^{}N_R^{}}}\right),$$
(2)
which calculates a single value for the asymmetry taking into account yields from both the left ($`N_L`$) and right ($`N_R`$) sides of the polarized beam and provides a consistency check on the relative luminosity Spinka (1999).
The BBCs were used to determine the relative luminosity ($``$ in Eq. 1) between bunches of opposite polarization sign fill-by-fill. A typical $``$ for the data sample analyzed here was approximately 1.09, measured to better than $`10^3`$. In the asymmetry analysis of charged hadrons, which utilized $`13`$M minimum-bias events, it was required that there be no hits in the RICH in order to eliminate electrons from photon conversions which mimic high-$`p_T`$ charged tracks. The momentum threshold for production of $`\stackrel{ˇ}{\mathrm{C}}`$erenkov radiation by pions was 4.7 GeV/$`c`$, allowing the RICH veto to preserve nearly all charged pions. The electron contamination in the final data sample was less than 1%. The decay background from long-lived particles was less than 5%.
Neutral pions were reconstructed via their decay to two photons using finely segmented ($`\mathrm{\Delta }\varphi \times \mathrm{\Delta }\eta 0.01\times 0.01`$) electromagnetic calorimeters (EMCal) Aphecetche et al. (2003). Photon clusters were selected by their shower shape and a charged track veto. Approximately 18M events recorded by an EMCal-based high-energy photon trigger in coincidence with the BBC collision trigger were analyzed Adler et al. (2003). The trigger efficiency for neutral pions varied from $`24`$% in the 1-2 GeV/$`c`$ bin to $`78`$% in the 4-5 GeV/$`c`$ bin. Only triggered events were used in this analysis. The $`\pi ^0`$ peak widths varied from 13.2 MeV/$`c^2`$ in the 1-2 GeV/$`c`$ bin to 10.6 MeV/$`c^2`$ in the 4-5 GeV/$`c`$ bin. The contribution from combinatorial background ranged from 34% to 5% across these bins; in order to avoid errors associated with peak extraction it was not subtracted.
The asymmetry for neutral pions and charged hadrons was determined for each fill using Eq. 1, then averaged over all fills. The contribution to the $`\pi ^0`$ asymmetry by the background under the peak was estimated by calculating the asymmetry of 50-MeV/$`c^2`$ regions on both sides of the signal, from 60-110 MeV/$`c^2`$ and 170-220 MeV/$`c^2`$ (see Table 2). The asymmetry of the signal region and its uncertainty were then corrected using
$$A_N^{\pi ^0}=\frac{A_N^{\mathrm{peak}}rA_N^{\mathrm{bg}}}{1r},\sigma _{A_N^{\pi ^0}}=\frac{\sqrt{\sigma _{A_N^{\mathrm{peak}}}^2+r^2\sigma _{A_N^{\mathrm{bg}}}^2}}{1r}$$
(3)
where $`r`$ is the fraction of background under the peak.
As the dominant systematic uncertainty is expected to be from the determination of the relative luminosity, systematic errors were evaluated by direct comparison of the asymmetry values calculated using Eq. 1 and Eq. 2. Any potential effect should be the same for both the charged hadron and neutral pion analyses. No $`p_T`$ dependence was expected or observed; therefore, we take the weighted average of the systematic uncertainties calculated for each bin, 0.002, as the overall, uniform systematic uncertainty.
The resulting asymmetries are plotted vs. $`p_T`$ in Fig. 2 and shown in Tables 2 and 3. The asymmetries are consistent with zero over the entire transverse momentum range.
In this Letter we have presented the first measurement of transverse-spin asymmetries $`A_N`$ at mid-rapidity and high $`p_T`$ at collider energies and the cross section for inclusive charged hadrons at mid-rapidity. NLO pQCD calculations have been found to reproduce experimental results for $`p_T>2`$ GeV/$`c`$ not only for the cross section presented here but also for inclusive neutral pion and production, indicating that pQCD can be used to interpret the high-$`p_T`$ asymmetries. The transverse SSAs observed for mid-rapidity production of both neutral pions and charged hadrons are consistent with zero within statistical errors of a few percent, measured over $`0.5<p_T<5`$ GeV/$`c`$. The result is consistent with the mid-rapidity results for neutral pions at $`\sqrt{s}=20`$ GeV Adams et al. (1996). The present measurement is complementary to that of Adams et al. (2004). The large asymmetries observed in neutral pion production at forward rapidity at $`\sqrt{s}=200`$ GeV Adams et al. (2004) are expected to originate from partonic processes involving valence quarks ($`x>0.1`$), whereas the particle production at mid-rapidity presented here is dominated by gluon-gluon and quark-gluon processes ($`x<0.1`$). Our results are consistent with the pQCD expectation that quark-gluon correlations are suppressed at high $`p_T`$ and mid-rapidity Kane et al. (1978); Qiu and Sterman (1999). A QCD analysis of the presented $`A_N`$ may lead to constraints on gluon-Sivers contributions to observed transverse-spin phenomena. The present transverse single-spin asymmetries represent an early measurement in a rigorous program to study transverse proton spin structure at hard scales using a pQCD framework at RHIC.
We thank the staff of the Collider-Accelerator Department, Magnet Division, and Physics Department at BNL and the RHIC polarimetry group for their vital contributions. We thank W. Vogelsang for calculations as well as numerous useful discussions. We acknowledge support from the Department of Energy and NSF (U.S.A.), MEXT and JSPS (Japan), CNPq and FAPESP (Brazil), NSFC (China), CNRS-IN2P3 and CEA (France), BMBF, DAAD, and AvH (Germany), OTKA (Hungary), DAE and DST (India), ISF (Israel), KRF and CHEP (Korea), RMIST, RAS, and RMAE (Russia), VR and KAW (Sweden), U.S. CRDF for the FSU, US-Hungarian NSF-OTKA-MTA, and US-Israel BSF. |
warning/0507/math0507455.html | ar5iv | text | # Local rigidity of group actions: past, present, future
## 1. Prologue
Let $`\mathrm{\Gamma }`$ be a finitely generated group, $`D`$ a topological group, and $`\pi :\mathrm{\Gamma }D`$ a homomorphism. We wish to study the space of deformations or perturbations of $`\pi `$. Certain trivial perturbations are always possible as soon as $`D`$ is not discrete, namely we can take $`d\pi d^1`$ where $`d`$ is a small element of $`D`$. This motivates the following definition:
###### Definition 1.1.
Given a homomorphism $`\pi :\mathrm{\Gamma }D`$, we say $`\pi `$ is locally rigid if any other homomorphism $`\pi ^{}`$ which is close to $`\pi `$ is conjugate to $`\pi `$ by a small element of $`D`$.
We topologize $`\mathrm{Hom}(\mathrm{\Gamma },D)`$ with the compact open topology which means that two homomorphisms are close if and only if they are close on a generating set for $`\mathrm{\Gamma }`$. If $`D`$ is path connected, then we can define deformation rigidity instead, meaning that any continuous path of representations $`\pi _t`$ starting at $`\pi `$ is conjugate to the trivial path $`\pi _t=\pi `$ by a continuous path $`d_t`$ in $`D`$ with $`d_0`$ being the identity in $`D`$. If $`D`$ is an algebraic group over $``$ or $``$, it is possible to prove that deformation rigidity and local rigidity are equivalent since $`\mathrm{Hom}(\mathrm{\Gamma },D)`$ is an algebraic variety and the action of $`D`$ by conjugation is algebraic; see \[Mu\], for example. For $`D`$ infinite dimensional and path-connected, this equivalence is no longer clear.
The study of local rigidity of lattices in semi-simple Lie groups is probably the beginning of the general study of rigidity in geometry and dynamics, a subject that is by now far too large for a single survey. See \[Sp1\] for the last attempt at a comprehensive survey and \[Sp2\] for a more narrowly focused updating of that survey. Here we abuse language slightly by saying a subgroup is locally rigid if the defining embedding is locally rigid as a homomorphism. See subsection 3.1 for a brief history of local rigidity of lattices and some discussion of subsequent developments that are of particular interest in the study of rigidity of group actions.
In this article we will focus on a survey of local rigidity when $`D=\mathrm{Diff}^{\mathrm{}}(M)`$ or occasionally $`D=\mathrm{Diff}^k(M)`$ for some finite $`k`$. Here we often refer to a homomorphism $`\pi :\mathrm{\Gamma }\mathrm{Diff}^{\mathrm{}}(M)`$ as an action, since it can clearly be thought of as $`C^{\mathrm{}}`$ action $`\pi :\mathrm{\Gamma }\times MM`$. We will consistently use the same letter $`\pi `$ to denote either the action or the homomorphism to $`\mathrm{Diff}^{\mathrm{}}(M)`$. The title of this article refers to this interpretation of $`\pi `$ as defining a group action. In this context, one considers rigidity of actions of connected groups as well as of discrete groups. In cases where $`\mathrm{\Gamma }`$ has any topology, we will always only study continuous actions, i.e. ones for which the defining homomorphism $`\pi `$ is a continuous map.
One can, in this context, develop more refined notions of local rigidity, since the topology on $`\mathrm{Diff}^{\mathrm{}}(M)`$ is an inverse limit topology. This means that two $`C^{\mathrm{}}`$ diffeomorphisms of $`M`$ are close precisely when they are $`C^k`$ close for some large value of $`k`$. The most exhaustive definition of local rigidity is probably the following:
###### Definition 1.2.
Let $`\mathrm{\Gamma }`$ be a discrete group and $`\pi :\mathrm{\Gamma }\mathrm{Diff}^k(M)`$ a homomorphism where $`k`$ is either a positive integer or $`\mathrm{}`$. We say that $`\pi `$ is $`C^{k,l,i,j,m}`$ locally rigid if any $`\pi ^{}:\mathrm{\Gamma }\mathrm{Diff}^l(M)`$ which is close to $`\pi `$ in the $`C^i`$ topology is conjugate to $`\pi `$ by a $`C^j`$ diffeomorphism $`\varphi `$ which is $`C^m`$ small. Here $`l,i,j,m`$ are all either non-negative integers or $`\mathrm{}`$ and the only a priori constraint are $`i\mathrm{min}(k,l)`$ and $`mj`$. When $`j=0`$, we will call the action stable or structurally stable. When $`j>0`$, we will call the action locally rigid or simply rigid.
We will avoid using this cumbersome notation when at all possible. There is a classical, dynamical notion of structural stability which is equivalent to $`C^{1,1,1,0,0}`$ local rigidity. I.e. a $`C^1`$ action $`\pi `$ of a group $`\mathrm{\Gamma }`$ is structurally stable if any $`C^1`$ close $`C^1`$ action of $`\mathrm{\Gamma }`$ is conjugate to $`\pi `$ by a small homeomorphism. For actions of $``$ this notion arose in hyperbolic dynamics in the work of Anosov and Smale \[An, Sm\]. From a dynamical point of view structural stability is important since it allows one to control dynamical properties of an open set of actions in $`\mathrm{Diff}^1(M)`$. Local rigidity can be viewed as a strengthening of this property in that it shows that an open set of actions is exhausted by smooth conjugates of a single action.
Though actions of $``$ and free groups on $`k`$ generators are often structurally stable, they are never locally rigid, and it is an interesting question as to how “large” a group needs to be in order to have actions which are locally rigid. Many of the original questions and theorems concerning local rigidity were for lattices in higher rank semi-simple Lie groups, where here higher rank means that all simple factors have real rank at least $`2`$. (See subsection 2.1 for a definition of rank.) Fairly early in the theory it became clear that local rigidity often held, and was in fact easier to prove, for certain actions of higher rank abelian groups, i.e. $`^k`$ for $`k2`$, see \[KL1\]. In addition, local rigidity results have been proven for actions of a wider variety of groups, including
1. certain non-volume preserving actions of lattices in $`SO(1,n)`$ in \[Kan1\]
2. all isometric actions of groups with property $`(T)`$ in \[FM2\],
3. certain affine actions of lattices in $`SP(1,n)`$ in \[Hi\].
There is extremely interesting related work of Ghys, older than the work just mentioned, which shows that the space of deformations of certain actions of surface groups on the circle is finite dimensional \[Gh1, Gh2, Gh3\]. Ghys also proved some very early results on local and global rigidity of very particular actions of connected solvable groups, see \[GhS, Gh1\] and subsection 4.2.
The study of local rigidity of group actions has had three primary historical motivations, one from the theory of lattices in Lie groups, one from dynamical systems theory and a third from the theory of foliations. (This statement is perhaps a bit coarse, and there is heavy overlap between these motivations, particularly the second and the third.) The first is the general study of rigidity of actions of large groups, as discussed in \[Z3, Z4\], see \[La, FK\] for more up to date surveys. This area is motivated by the study of rigidity of lattices in semi-simple Lie groups, particularly by Margulis’ super-rigidity theorem and it’s non-linear generalization by Zimmer to a cocycle super-rigidity theorem, see subsection 3.1 and \[Z4\] for more discussion. This motivation also stems from an analogy between semi-simple Lie groups and diffeomorphism groups. When $`M`$ is a compact manifold, not only is $`\mathrm{Diff}^{\mathrm{}}(M)`$ an infinite dimensional Lie group, but its connected component is simple. Simplicity of the connected component of $`\mathrm{Diff}^{\mathrm{}}(M)`$ was proven by Thurston using results of Epstein and Herman \[Th2, Ep, Hr\]. Herman had used Epstein’s work to see that the connected component of $`\mathrm{Diff}^{\mathrm{}}(𝕋^n)`$ is simple and Thurston’s proof of the general case uses this. See also further work on the topic by Banyaga and Mather \[Ba1, Mt1, Mt2, Mt3\], as well as Banyaga’s book \[Ba2\].
The dynamical motivation for studying rigidity of group actions comes from the study of structural stability of diffeomorphisms and flows in hyperbolic dynamics, see the introduction of \[KS2\]. This area provides many of the basic techniques by which results in the area have been proven, particularly early in the history of the field. Philosophically, hyperbolic diffeomorphisms are structurally stable, group actions generated by structurally stable diffeomorphisms are quite often structurally stable, and the presence of a large group action frequently allows one to improve the regularity of the conjugacy. See subsection 3.2 for a brief history of relevant results on structural stability and subsections 4.1, 4.2, and 5.1 for some applications of these ideas.
The third motivation for studying rigidity of group actions comes from the theory of foliations. Many techniques and ideas in this area are also related to work on hyperbolic dynamics, and many of the foliations of interest are dynamical foliations of hyperbolic dynamical systems. A primary impetus in this area is the theory of codimension one foliations, and so many of the ideas here were first developed either for groups acting on the circle or for actions of connected groups on manifolds only one dimension larger then the acting group. See particularly \[GhS, Gh1\] for the early history of these developments.
### Some remarks on biases and omissions.
Like any survey of this kind, this work is informed by it’s authors biases and experiences. The most obvious of these is that my point of view is primarily motivated by the study of rigidity properties of semi-simple Lie groups and their lattices, rather than primarily motivated by hyperbolic dynamics or foliation theory. This informs the biases of this article and a very different article would result from different biases.
There are two large omissions in this article. The first omission is that it is primarily occupied with local rigidity of discrete group actions. When similar results are known for actions of Lie groups, they are mentioned, though frequently only special cases are stated. This is partially because results in this context are often complicated by the need to consider time changes, and I did not want to dwell on that issue here. The second omission is that little to no care is taken to state optimal results relating the various constants in $`C^{k,l,i,j,m}`$ local rigidity. Dwelling on issues of regularity seemed likely to obscure the main line of the developments, so many results are stated without any explicit mention of regularity. Usually this is done only when the action can be shown to be locally rigid in $`\mathrm{Diff}^{\mathrm{}}(M)`$ in the sense of Definition 1.1. This implicitly omits both the degree of regularity to which the perturbed and unperturbed actions are close and the degree of regularity with which the size of the conjugacy is small. In other words local rigidity is $`C^{\mathrm{},\mathrm{},i,\mathrm{},m}`$ local rigidity for some unspecified $`i`$ and $`m`$, and I always fail to specify $`i`$ and $`m`$ even when they are known. Occasionally a result is stated that only produces a finite regularity conjugacy, with this issue only remarked on following the statement of the result. It seems quite likely that most results of this kind can be improved to produce $`C^{\mathrm{}}`$ conjugacies using the techniques of \[FM2, FM3\], see discussion at the end of Section 5.1.
Lastly we remark that the study of local rigidity of group actions is often closely intertwined with the study of global rigidity of group actions. The meaning of the phrase global rigidity is not entirely precise, but it is typically used to cover settings in which one can classify all group actions satisfying certain hypotheses on a certain manifold or class of manifolds. The study of global rigidity is too broad and interesting to summarize briefly, but some examples are mentioned below when they are closely related to work on local rigidity. See \[FK, La\] for recent surveys concerning both local and global rigidity.
## 2. A brief digression: some examples of groups and actions
In this section we briefly describe some of the groups that will play important roles in the results discussed here. The reader already familiar with semi-simple Lie groups and their lattices may want to skip to the second subsection where we give descriptions of group actions.
### 2.1. Semi-simple groups and their lattices.
By a simple Lie group, we mean a connected Lie group all of whose normal subgroups are discrete, though we make the additional convention that $``$ and $`S^1`$ are not simple. By a semi-simple Lie group we mean the quotient of a product of simple Lie groups by some subgroup of the product of their centers. Note that with our conventions, the center of a simple Lie group is discrete and is in fact the maximal normal subgroup. There is an elaborate structure theory of semi-simple Lie groups and the groups are competely classified, see \[He\] or \[Kn\] for details. Here we merely describe some examples, all of which are matrix groups. All connected semisimple Lie groups are discrete central extensions of matrix groups, so the reader will lose very little by always thinking of matrix groups.
1. The groups $`SL(n,),SL(n,)`$ and $`SL(n,)`$ of $`n`$ by $`n`$ matrices of determinant one over the real numbers, the complex numbers or the quaternions.
2. The group $`SP(2n,)`$ of $`2n`$ by $`2n`$ matrices of determinant one which preserve a real symplectic form on $`^{2n}`$.
3. The groups $`SO(p,q),SU(p,q)`$ and $`SP(p,q)`$ of matrices which preserve inner products of signature $`(p,q)`$ where the inner product is real linear on $`^{p+q}`$, hermitian on $`^{p+q}`$ or quaternionic hermitian on $`^{p+q}`$ respectively.
Let $`G`$ be a semi-simple Lie group which is a subgroup of $`GL(n,)`$. We say that $`G`$ has real rank $`k`$ if $`G`$ has a $`k`$ dimensional abelian subgroup which is conjugate to a subgroup of the real diagonal matrices and no $`k+1`$ dimensional abelian subgroups with the same property. The groups in $`(1)`$ have rank $`n1`$, the groups in $`(2)`$ have rank $`n`$ and the groups in $`(3)`$ have rank $`\mathrm{min}(p,q)`$.
Since this article focuses primarily on finitely generated groups, we are more interested in discrete subgroups of Lie groups than in the Lie groups themselves. A discrete subgroup $`\mathrm{\Gamma }`$ in a Lie group $`G`$ is called a lattice if $`G/\mathrm{\Gamma }`$ has finite Haar measure. The lattice is called cocompact or uniform if $`G/\mathrm{\Gamma }`$ is compact and non-uniform or simply not cocompact otherwise. If $`G=G_1\times \mathrm{}\times G_n`$ is a product then we say a lattice $`\mathrm{\Gamma }<G`$ is irreducible if it’s projection to each $`G_i`$ is dense. More generally we make the same definition for an almost direct product, by which we mean a direct product $`G`$ modulo some subgroup of the center $`Z(G)`$. Lattices in semi-simple Lie groups can always be constructed by arithmetic methods, see \[Bo\] and also \[Mr\] for more discussion. In fact, one of the most important results in the theory of semi-simple Lie groups is that if $`G`$ is a semi-simple Lie group without compact factors, then all irreducible lattices in $`G`$ are arithmetic unless $`G`$ is locally isomorphic to $`SO(1,n)`$ or $`SU(1,n)`$. For $`G`$ of real rank at least two, this is Margulis’ arithmeticity theorem, which he deduced from his super-rigidity theorems \[Ma2, Ma3, Ma4\]. For non-uniform lattices, Margulis had an earlier proof which does not use the superrigidity theorems, see \[Ma1, Ma2\]. This earlier proof depends on the study of dynamics of unipotent elements on the space $`G/\mathrm{\Gamma }`$, and particularly on what is now known as the “non-divergence of unipotent flows”. Special cases of the super-rigidity theorems were then proven for $`Sp(1,n)`$ and $`F_4^{20}`$ by Corlette and Gromov-Schoen, which sufficed to imply the statement on arithmeticity given above \[Co2, GS\]. As we will be almost exclusively concerned with arithmetic lattices, we do not give examples of non-arithmetic lattices here, but refer the reader to \[Ma4\] and \[Mr\] for more discussion. A formal definition of arithmeticity, at least when $`G`$ is algebraic is:
###### Definition 2.1.
Let $`G`$ be a semisimple algebraic Lie group and $`\mathrm{\Gamma }<G`$ a lattice. Then $`\mathrm{\Gamma }`$ is arithmetic if there exists a semi-simple algebraic Lie group $`H`$ defined over $``$ such that
1. there is a homomorphism $`\pi :H^0G`$ with compact kernel,
2. there is a rational structure on $`H`$ such that the projection of the integer points of $`H`$ to $`G`$ are commensurable to $`\mathrm{\Gamma }`$, i.e. $`\pi (H())\mathrm{\Gamma }`$ is of finite index in both $`H()`$ and $`\mathrm{\Gamma }`$.
We now give some examples of arithmetic lattices. The simplest is to take the integer points in a simple (or semi-simple) group $`G`$ which is a matrix group, e.g. $`SL(n,)`$ or $`Sp(n,)`$. This exact construction always yields lattices, but also always yields non-uniform lattices. In fact the lattices one can construct in this way have very special properties because they will contain many unipotent matrices. If a lattice is cocompact, it will necessarily contain no unipotent matrices. The standard trick for understanding the structure of lattices in $`G`$ which become integral points after passing to a compact extension is called change of base. For much more discussion see \[Ma4, Mr, Z2\]. We give one example to illustrate the process. Let $`G=SO(m,n)`$ which we view as the set of matrices in $`SL(n+m,)`$ which preserve the inner product
$$v,w=\left(\sqrt{2}\underset{i=1}{\overset{m}{}}v_iw_i\right)+\left(\underset{i=m+1}{\overset{n+m}{}}v_iw_i\right)$$
where $`v_i`$ and $`w_i`$ are the $`i`$th components of $`v`$ and $`w`$. This form, and therefore $`G`$, are defined over the field $`(\sqrt{2})`$ which has a Galois conjugation $`\sigma `$ defined by $`\sigma (\sqrt{2})=\sqrt{2}`$. If we looks at the points $`\mathrm{\Gamma }=G([\sqrt{2}])`$, we can define an embedding of $`\mathrm{\Gamma }`$ in $`SO(m,n)\times SO(m+n)`$ by taking $`\gamma `$ to $`(\gamma ,\sigma (\gamma ))`$. It is straightforward to check that this embedding is discrete. In fact, this embeds $`\mathrm{\Gamma }`$ in $`H=SO(m,n)\times SO(m+n)`$ as integral points for the rational structure on $`H`$ where the rational points are exactly the points $`(m,\sigma (m))`$ where $`mG((\sqrt{2}))`$. This makes $`\mathrm{\Gamma }`$ a lattice in $`H`$ and it is easy to see that $`\mathrm{\Gamma }`$ projects to a lattice in $`G`$, since $`G`$ is cocompact in $`H`$. What is somewhat harder to verify is that $`\mathrm{\Gamma }`$ is cocompact in $`H`$, for which we refer the reader to the list of references above.
Similar constructions are possible with $`SU(m,n)`$ or $`SP(m,n)`$ in place of $`SO(m,n)`$ and also with more simple factors and fields with more Galois automorphisms. There are also a number of other constructions of arithmetic lattices using division algebras. See \[Mr\] for a comprehensive treatment.
We end this section by defining a key property of many semisimple groups and their lattices. This is property $`(T)`$ of Kazhdan, and was introduced by Kazhdan in \[Ka1\] in order to prove that non-uniform lattices in higher rank semi-simple Lie groups are finitely generated and have finite abelianization. It has played a fundamental role in many subsequent developments. We do not give Kazhdan’s original definition, but one which was shown to be equivalent by work of Delorme and Guichardet \[De, Gu\].
###### Definition 2.2.
A group $`\mathrm{\Gamma }`$ has property $`(T)`$ of Kazhdan if $`H^1(\mathrm{\Gamma },\pi )=0`$ for every continuous unitary representation $`\pi `$ of $`\mathrm{\Gamma }`$ on a Hilbert space. This is equivalent to saying that any continuous isometric action of $`\mathrm{\Gamma }`$ on a Hilbert space has a fixed point.
###### Remarks 1.
1. Kazhdan’s definition is that the trivial representation is isolated in the unitary dual of $`\mathrm{\Gamma }`$.
2. If a continuous group $`G`$ has property $`(T)`$ so does any lattice in $`G`$. This result was proved in \[Ka1\].
3. Any semi-simple Lie group has property $`(T)`$ if and only if it has no simple factors locally isomorphic to $`SO(1,n)`$ or $`SU(1,n)`$. For a discussion of this fact and attributions, see \[HV\]. For groups with all simple factors of real rank at least three, this is proven in \[Ka1\].
4. No noncompact amenable group, and in particular no noncompact abelian group, has property $`(T)`$. An easy averaging argument shows that all compact groups have property $`(T)`$.
Groups with property $`(T)`$ play an important role in many areas of mathematics and computer science.
### 2.2. Some actions of groups and lattices.
Here we define and give examples of the general classes of actions for which local rigidity results have been proven. Let $`H`$ be a Lie group and $`L<H`$ a closed subgroup. Then a diffeomorphism $`f`$ of $`H/L`$ is called affine if there is a diffeomorphism $`\stackrel{~}{f}`$ of $`H`$ such that $`f([h])=\stackrel{~}{f}(h)`$ where $`\stackrel{~}{f}=A\tau _h`$ with $`A`$ an automorphism of $`H`$ with $`A(L)=L`$ and $`\tau _h`$ is left translation by some $`h`$ in $`H`$. Two obvious classes of affine diffeomorphisms are left translations on any homogeneous space and either linear automorphisms of tori or more generally automorphisms of nilmanifolds. A group action is called affine if every element of the group acts by an affine diffeomorphism. It is easy to check that the full group of affine diffeomorphisms $`\mathrm{Aff}(H/L)`$ is a finite dimensional Lie group and an affine action of a group $`D`$ is a homomorphism $`\pi :D\mathrm{Aff}(H/L)`$. The structure of $`\mathrm{Aff}(H/L)`$ is surprisingly complicated in general, it is a quotient of a subgroup of the group $`\mathrm{Aut}(H)H`$ where $`\mathrm{Aut}(H)`$ is a the group of automorphisms of $`H`$. For a more detailed discussion of this relationship, see \[FM1, Section 6\]. While it is not always the case that any affine action of a group $`D`$ on $`H/L`$ can be described by a homomorphism $`\pi :D\mathrm{Aut}(H)H`$, this is true for two important special cases:
1. $`D`$ is a connected semi-simple Lie group and $`L`$ is a cocompact lattice in $`H`$,
2. $`D`$ is a lattice in a semi-simple Lie group $`G`$ where $`G`$ has no compact factors and no simple factors locally isomorphic to $`SO(1,n)`$ or $`SU(1,n)`$, and $`L`$ is a cocompact lattice in $`H`$.
These facts are \[FM1, Theorem 6.4 and 6.5\] where affine actions as in $`(1)`$ and $`(2)`$ above are classified.
The most obvious examples of affine actions of large groups are of the following forms, which are frequently referred to as standard actions:
1. Actions of groups by automorphisms of nilmanifolds. I.e. let $`N`$ be a simply connected nilpotent group, $`\mathrm{\Lambda }<N`$ a lattice (which is necessarily cocompact) and assume a finitely generated group $`\mathrm{\Gamma }`$ acts by automorphisms of $`N`$ preserving $`\mathrm{\Lambda }`$. The most obvious examples of this are when $`N=^n`$, $`\mathrm{\Lambda }=^n`$ and $`\mathrm{\Gamma }<SL(n,)`$, in which case we have a linear action of $`\mathrm{\Gamma }`$ on $`𝕋^n`$.
2. Actions by left translations. I.e. let $`H`$ be a Lie group and $`\mathrm{\Lambda }<H`$ a cocompact lattice and $`\mathrm{\Gamma }<H`$ some subgroup. Then $`\mathrm{\Gamma }`$ acts on $`H/\mathrm{\Lambda }`$ by left translations. Note that in this case $`\mathrm{\Gamma }`$ need not be discrete.
3. Actions by isometries. Here $`K`$ is a compact group which acts by isometries on some compact manifold $`M`$ and $`\mathrm{\Gamma }<K`$ is a subgroup. Note that here $`\mathrm{\Gamma }`$ is either discrete or a discrete extension of a compact group.
We now briefly define a few more general classes of actions, for which local rigidity results are either known or conjectured. We first fix some notations. Let $`A`$ and $`D`$ be topological groups, and $`B<A`$ a closed subgroup. Let $`\rho :D\times A/BA/B`$ be a continuous affine action.
###### Definition 2.3.
1. Let $`A,B,D`$ and $`\rho `$ be as above. Let $`C`$ be a compact group of affine diffeomorphisms of $`A/B`$ that commute with the $`D`$ action. We call the action of $`D`$ on $`C\backslash A/B`$ a generalized affine action.
2. Let $`A`$, $`B`$, $`D`$ and $`\rho `$ be as in $`1`$ above. Let $`M`$ be a compact Riemannian manifold and $`\iota :D\times A/B\mathrm{Isom}(M)`$ a $`C^1`$ cocycle. We call the resulting skew product $`D`$ action on $`A/B\times M`$ a quasi-affine action. If $`C`$ and $`D`$ are as in $`2`$, and $`\alpha :D\times C\backslash A/B\mathrm{Isom}(M)`$ is a $`C^1`$ cocycle, then we call the resulting skew product $`D`$ action on $`C\backslash A/B\times M`$ a generalized quasi-affine action.
For many of the actions we consider here, there will be a foliation of particular importance. If $`\rho `$ is an action of a group $`D`$ on a manifold $`N`$, and $`\rho `$ preserves a foliation $`𝔉`$ and a Riemannian metric along the leaves of $`𝔉`$, we call $`𝔉`$ a central foliation for $`\rho `$. For quasi-affine and generalized quasi-affine actions on manifolds of the form $`C\backslash A/B\times M`$ the foliation by leaves of the $`\{[a]\}\times M`$ is always a central foliation. There are also actions with more complicated central foliations. For example if $`H`$ is a Lie group, $`\mathrm{\Lambda }<H`$ is discrete and a subgroup $`G<H`$ acts on $`H/\mathrm{\Lambda }`$ by left translations, then the foliation of $`H/\mathrm{\Lambda }`$ by orbits of the centralizer $`Z_H(G)`$ of $`G`$ in $`H`$ is a central foliation. It is relatively easy to construct examples where this foliation has dense leaves. Another example of an action which has a foliation with dense leaves is to embed the $`[\sqrt{2}]`$ points of $`SO(m,n)`$ into $`SL(2(m+n),)`$ as described in the preceding subsection and then let this group act on $`𝕋^{2(m+n)}`$ linearly. It is easy to see in this case that the maximal central foliation for the action is a foliation by densely embedded leaves none of which are compact.
We end this section by describing briefly the standard construction of an induced or suspended action. This notion can be seen as a generalization of the construction of a flow under a function or as an analogue of the more algebraic notion of inducing a representation. Given a group $`H`$, a (usually closed) subgroup $`L`$, and an action $`\rho `$ of $`L`$ on a space $`X`$, we can form the space $`(H\times X)/L`$ where $`L`$ acts on $`H\times X`$ by $`h(l,x)=(lh{}_{}{}^{1},\rho (h)x)`$. This space now has a natural $`H`$ action by left multiplication on the first coordinate. Many properties of the $`L`$ action on $`X`$ can be studied more easily in terms of properties of the $`H`$ action on $`(H\times X)/L`$. This construction is particularly useful when $`L`$ is a lattice in $`H`$.
## 3. Pre-history
### 3.1. Local and global rigidity of homomorphisms into finite dimensional groups.
The earliest work on local rigidity in the context of Definition 1.1 was contained in series of works by Calabi–Vesentini, Selberg, Calabi and Weil, which resulted in the following:
###### Theorem 3.1.
Let $`G`$ be a semi-simple Lie group and assume that $`G`$ is not locally isomorphic to $`SL(2,)`$. Let $`\mathrm{\Gamma }<G`$ be an irreducible cocompact lattice, then the defining embedding of $`\mathrm{\Gamma }`$ in $`G`$ is locally rigid.
###### Remarks 2.
1. If $`G=SL(2,)`$ the theorem is false and there is a large, well studied space of deformation of $`\mathrm{\Gamma }`$ in $`G`$, known as the Teichmuller space.
2. There is an analogue of this theorem for lattices that are not cocompact. This result was proven later and has a more complicated history which we omit here. In this case it is also necessary to exclude $`G`$ locally isomorphic to $`SL(2,)`$.
This theorem was originally proven in special cases by Calabi, Calabi–Vesentini and Selberg. In particular, Selberg gives a proof for cocompact lattices in $`SL(n,)`$ for $`n3`$ in \[S\], Calabi–Vesentini give a proof when the associated symmetric space $`X=G/K`$ is Kähler in \[CV\] and Calabi gives a proof for $`G=SO(1,n)`$ where $`n3`$ in \[C\]. Shortly afterwards, Weil gave a complete proof of Theorem 3.1 in \[We1, We2\].
In all of the original proofs, the first step was to show that any perturbation of $`\mathrm{\Gamma }`$ was discrete and therefore a cocompact lattice. This is shown in special cases in \[C, CV, S\] and proven in a somewhat broader context than Theorem 3.1 in \[W1\].
The different proofs of cases of Theorem 3.1 are also interesting in that there are two fundamentally different sets of techniques employed and this dichotomy continues to play a role in the history of rigidity. Selberg’s proof essentially combines algebraic facts with a study of the dynamics of iterates of matrices. He makes systematic use of the existence of singular directions, or Weyl chamber walls, in maximal diagonalizable subgroups of $`SL(n,)`$. Exploiting these singular directions is essential to much later work on rigidity, both of lattices in higher rank groups and of actions of abelian groups. It seems possible to generalize Selberg’s proof to the case of $`G`$ an $``$-split semi-simple Lie group with rank at least $`2`$. Selberg’s proof, which depended on asymptotics at infinity of iterates of matrices, inspired Mostow’s explicit use of boundaries in his proof of strong rigidity \[Mo2\]. Mostow’s work in turn provided inspiration for the use of boundaries in later work of Margulis, Zimmer and others on rigidity properties of higher rank groups.
The proofs of Calabi, Calabi–Vesentini and Weil involve studying variations of geometric structures on the associated locally symmetric space. The techniques are analytic and use a variational argument to show that all variations of the geometric structure are trivial. This work is a precursor to much work in geometric analysis studying variations of geometric structures and also informs later work on proving rigidity/vanishing of harmonic forms and maps. The dichotomy between approaches based on algebra/dynamics and approaches that are in the spirit of geometric analysis continues through much of the history of rigidity and the history of local rigidity of group actions in particular.
Shortly after completing this work, Weil discovered a new criterion for local rigidity \[We3\]. In the context of Theorem 3.1, this allows one to avoid the step of showing that a perturbation of $`\mathrm{\Gamma }`$ remains discrete. In addition, this result opened the way for understanding local rigidity of more general representations of discrete groups.
###### Theorem 3.2.
Let $`\mathrm{\Gamma }`$ be a finitely generated group, $`G`$ a Lie group and $`\pi :\mathrm{\Gamma }G`$ a homomorphism. Then $`\pi `$ is locally rigid if $`H^1(\mathrm{\Gamma },𝔤)=0`$. Here $`𝔤`$ is the Lie algebra of $`G`$ and $`\mathrm{\Gamma }`$ acts on $`𝔤`$ by $`Ad_G\pi `$.
Weil’s proof of this result uses only the implicit function theorem and elementary properties of the Lie group exponential map. The same theorem is true if $`G`$ is an algebraic group over any local field of characteristic zero. In \[We3\], Weil remarks that if $`\mathrm{\Gamma }<G`$ is a cocompact lattice and $`G`$ satisfies the hypothesis of Theorem 3.1, then the vanishing of $`H^1(\mathrm{\Gamma },𝔤)`$ can be deduced from the computations in \[We2\]. The vanishing of $`H^1(\mathrm{\Gamma },𝔤)`$ is proven explicitly by Matsushima and Murakami in \[MM\].
Motivated by Weil’s work and other work of Matsushima, conditions for vanishing of $`H^1(\mathrm{\Gamma },𝔤)`$ were then studied by many authors. See particularly \[MM\] and \[Rg1\]. The results in these papers imply local rigidity of many linear representations of lattices.
To close this section, I will briefly discuss some subsequent developments concerning rigidity of lattices in Lie groups that motivated the study of both local and global rigidity of group actions.
The first remarkable result in this direction is Mostow’s rigidity theorem, see \[Mo1\] and references there. Given $`G`$ as in Theorem 3.1, and two irreducible cocompact lattices $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ in $`G`$, Mostow proves that any isomorphism from $`\mathrm{\Gamma }_1`$ to $`\mathrm{\Gamma }_2`$ extends to an isomorphism of $`G`$ with itself. Combined with the principal theorem of \[We1\] which shows that a perturbation of a lattice is again a lattice, this gives a remarkable and different proof of Theorem 3.1, and Mostow was motivated by the desire for a “more geometric understanding” of Theorem 3.1 \[Mo1\]. Mostow’s theorem is in fact a good deal stronger, and controls not only homomorphisms $`\mathrm{\Gamma }G`$ near the defining homomorphism, but any homomorphism into any other simple Lie group $`G^{}`$ where the image is lattice. As mentioned above, Mostow’s approach was partially inspired by Selberg’s proof of certain cases of Theorem 3.1, \[Mo2\]. A key step in Mostow’s proof is the construction of a continuous map between the geometric boundaries of the symmetric spaces associated to $`G`$ and $`G^{}`$. Boundary maps continue to play a key role in many developments in rigidity theory. A new proof of Mostow rigidity, at least for $`G_i`$ of real rank one, was provided by Besson, Courtois and Gallot. Their approach is quite different and has had many other applications concerning rigidity in geometry and dynamics, see e.g. \[BCG, CF\].
The next remarkable result in this direction is Margulis’ superrigidity theorem. Margulis proved this theorem as a tool to prove arithmeticity of irredudicible uniform lattices in groups of real rank at least $`2`$. For irreducible lattices in semi-simple Lie groups of real rank at least $`2`$, the superrigidity theorems classifies all finite dimensional linear representations. Margulis’ theorem holds for irreducible lattices in semi-simple Lie groups of real rank at least two. Given a lattice $`\mathrm{\Gamma }<G`$ where $`G`$ is simply connected, one precise statement of some of Margulis results is to say that any linear representation $`\sigma `$ of $`\mathrm{\Gamma }`$ almost extends to a linear representation of $`G`$. By this we mean that there is a linear representation $`\stackrel{~}{\sigma }`$ of $`G`$ and a bounded image representation $`\overline{\sigma }`$ of $`\mathrm{\Gamma }`$ such that $`\sigma (\gamma )=\stackrel{~}{\sigma }(\gamma )\overline{\sigma }(\gamma )`$ for all $`\gamma `$ in $`G`$. Margulis’ theorems also give an essentially complete description of the representations $`\overline{\sigma }`$, up to some issues concerning finite image representations. The proof here is partially inspired by Mostow’s work: a key step is the construction of a measurable “boundary map”. However the methods for producing the boundary map in this case are very dynamical. Margulis’ original proof used Oseledec Multiplicative Ergodic Theorem. Later proofs were given by both Furstenberg and Margulis using the theory of group boundaries as developed by Furstenberg from his study of random walks on groups \[Fu1, Fu2\]. Furstenberg’s probabalistic version of boundary theory has had a profound influence on many subsequent developments in rigidity theory. For more discussion of Margulis’ superrigidity theorem, see \[Ma2, Ma3, Ma4\].
A main impetus for studying rigidity of group actions on manifolds came from Zimmer’s theorem on superrigidity for cocycles. This theorem and it’s proof were strongly motivated by Margulis’ work. In fact, Margulis’ theorem is Zimmer’s theorem for a certain coccyle $`\alpha :G\times G/\mathrm{\Gamma }\mathrm{\Gamma }`$ and the proof of Zimmer’s theorem is quite similar to the proof of Margulis’. In order to avoid technicalities, we describe only a special case of this result. Let $`M`$ be a compact manifold, $`H`$ a matrix group and $`P=M\times H`$. Now let a group $`\mathrm{\Gamma }`$ act on $`M`$ and $`P`$ continuously, so that the projection from $`P`$ to $`M`$ is equivariant. Further assume that the action on $`M`$ is measure preserving and ergodic. If $`\mathrm{\Gamma }`$ is a lattice in a simply connected, semi-simple Lie group $`G`$ all of whose simple factors have real rank at least two then there is a measurable map $`s:MH`$, a representation $`\pi :GH`$, a compact subgroup $`K<H`$ which commutes with $`\pi (G)`$ and a measurable map $`\mathrm{\Gamma }\times MK`$ such that
(1)
$$\gamma s(m)=k(m,\gamma )\pi (\gamma )s(\gamma m).$$
It is easy to check from this equation that the map $`K`$ satisfies a certain equation that makes it into a cocycle over the action of $`\mathrm{\Gamma }`$. One should view $`s`$ as providing coordinates on $`P`$ in which the $`\mathrm{\Gamma }`$ action is almost a product. For more discussion of this theorem the reader should see any of \[Fe1, Fe2, FM1, Fu3, Z2\]. (The version stated here is only proven in \[FM1\], previous proofs all yielded somewhat more complicated statements.) As a sample application, let $`M=𝕋^n`$ and let $`P`$ be the frame bundle of $`M`$, i.e. the space of frames in the tangent bundle of $`M`$. Since $`𝕋^n`$ is parallelizable, we have $`P=𝕋^n\times GL(n,^n)`$. The cocycle super-rigidity theorem then says that “up to compact noise”, the derivative of any measure preserving $`\mathrm{\Gamma }`$ action on $`𝕋^n`$ looks measurably like a constant linear map. In fact, the cocycle superrigidity theorems apply more generally to continuous actions on any principal bundle $`P`$ over $`M`$ with fiber $`H`$, an algebraic group, and in this context produces a measurable section $`s:MP`$ satisfying equation $`(\text{1})`$. So in fact, cocycle superrigidity implies that for any action preserving a finite measure on any manifold the derivative cocycle looks measurably like a constant cocycle, up to compact noise. That cocycle superrigidity provides information about actions of groups on manifolds through the derivative cocycle was first observed in \[Fu3\]. Zimmer originally proved cocycle superrigidity in order to study orbit equivalence of group actions. For a recent survey of subsequent developments concerning orbit equivalence rigidity and other forms of superrigidity for cocycles, see \[Sl2\].
### 3.2. Stability in hyperbolic dynamics.
A diffeomorphism $`f`$ of a manifold $`X`$ is said to be Anosov if there exists a continuous $`f`$ invariant splitting of the tangent bundle $`TX=E_f^uE_f^s`$ and constants $`a>1`$ and $`C,C^{}>0`$ such that for every $`xX`$,
1. $`Df^n(v^u)Ca^nv^u`$ for all $`v^uE_f^u(x)`$ and,
2. $`Df^n(v^s)C^{}a^nv^s`$ for all $`v^sE_f^s(x)`$.
We note that the constants $`C`$ and $`C^{}`$ depend on the choice of metric, and that a metric can always be chosen so that $`C=C^{}=1`$. There is an analogous notion for a flow $`f_t`$, where $`TX=T𝒪E_{f_t}^uE_{f_t}^s`$ where $`T𝒪`$ is the tangent space to the flow direction and vectors in $`E_{f_t}^u`$ (resp. $`E_{f_t}^s`$) are uniformly expanded (resp. uniformly contracted) by the flow. This notion was introduced by Anosov and named after Anosov by Smale, who popularized the notion in the United States \[An, Sm\]. One of the earliest results in the subject is Anosov’s proof that Anosov diffeomorphisms are structurally stable, or, in our language $`C^{1,1,1,0,0}`$ locally rigid. There is an analgous result for flows, though this requires that one introduce a notion of time change that we will not consider here. Since Anosov also showed that $`C^2`$ Anosov flows and diffeomorphisms are ergodic, structural stability implies that the existence of an open set of “chaotic” dynamical systems.
The notion of an Anosov diffeomorphism has had many interesting generalizations, for example: Axiom A diffeomorphisms, non-uniformly hyperbolic diffeomorphisms, and diffeomorphisms admitting a dominated splitting. The notion that has been most useful in the study of local rigidity is the notion of a partially hyperbolic diffeomorphism as introduced by Hirsch, Pugh and Shub. Under strong enough hypotheses, these diffeomorphisms have a weaker stability property similar to structural stability. More or less, the diffeomorphisms are hyperbolic relative to some foliation, and any nearby action is hyperbolic to some nearby foliation. To describe more precisely the class of diffeomorphisms we consider and the stability property they enjoy, we require some definitions.
The use of the word foliation varies with context. Here a foliation by $`C^k`$ leaves will be a continuous foliation whose leaves are $`C^k`$ injectively immersed submanifolds that vary continuously in the $`C^k`$ topology in the transverse direction. To specify transverse regularity we will say that a foliation is transversely $`C^r`$. A foliation by $`C^k`$ leaves which is tranversely $`C^k`$ is called simply a $`C^k`$ foliation. (Note our language does not agree with that in the reference \[HPS\].)
Given an automorphism $`f`$ of a vector bundle $`EX`$ and constants $`a>b1`$, we say $`f`$ is $`(a,b)`$-partially hyperbolic or simply partially hyperbolic if there is a metric on $`E`$, a constant and $`C1`$ and a continuous $`f`$ invariant, non-trivial splitting $`E=E_f^uE_f^cE_f^s`$ such that for every $`x`$ in $`X`$:
1. $`f^n(v^u)Ca^nv^u`$ for all $`v^uE_f^u(x)`$,
2. $`f^n(v^s)C{}_{}{}^{1}a_{}^{n}v^s`$ for all $`v^sE_f^s(x)`$ and
3. $`C{}_{}{}^{1}b_{}^{n}v^0<f^n(v^0)Cb^nv^0`$ for all $`v^0E_f^c(x)`$ and all integers $`n`$.
A $`C^1`$ diffeomorphism $`f`$ of a manifold $`X`$ is $`(a,b)`$-partially hyperbolic if the derivative action $`Df`$ is $`(a,b)`$-partially hyperbolic on $`TX`$. We remark that for any partially hyperbolic diffeomorphism, there always exists an adapted metric for which $`C=1`$. Note that $`E_f^c`$ is called the central distribution of $`f`$, $`E_f^u`$ is called the unstable distribution of $`f`$ and $`E_f^s`$ the stable distribution of $`f`$.
Integrability of various distributions for partially hyperbolic dynamical systems is the subject of much research. The stable and unstable distributions are always tangent to invariant foliations which we call the stable and unstable foliations and denote by $`𝒲_f^s`$ and $`𝒲_f^u`$. If the central distribution is tangent to an $`f`$ invariant foliation, we call that foliation a central foliation and denote it by $`𝒲_f^c`$. If there is a unique foliation tangent to the central distribution we call the central distribution uniquely integrable. For smooth distributions unique integrability is a consequence of integrability, but the central distribution is usually not smooth. If the central distribution of an $`(a,b)`$-partially hyperbolic diffeomorphism $`f`$ is tangent to an invariant foliation $`𝒲_f^c`$, then we say $`f`$ is $`r`$-normally hyperbolic to $`𝒲_f^c`$ for any $`r`$ such that $`a>b^r`$. This is a special case of the definition of $`r`$-normally hyperbolic given in \[HPS\].
Before stating a version of one of the main results of \[HPS\], we need one more definition. Given a group $`G`$, a manifold $`X`$, two foliations $`𝔉`$ and $`𝔉^{}`$ of $`X`$, and two actions $`\rho `$ and $`\rho ^{}`$ of $`G`$ on $`X`$, such that $`\rho `$ preserves $`𝔉`$ and $`\rho ^{}`$ preserves $`𝔉^{}`$, following \[HPS\] we call $`\rho `$ and $`\rho ^{}`$ leaf conjugate if there is a homeomorphism $`h`$ of $`X`$ such that:
1. $`h(𝔉)=𝔉^{}`$ and
2. for every leaf $`𝔏`$ of $`𝔉`$ and every $`gG`$, we have $`h(\rho (g)𝔏)=\rho ^{}(g)h(𝔏)`$.
The map $`h`$ is then referred to as a leaf conjugacy between $`(X,𝔉,\rho )`$ and $`(X,𝔉^{},\rho ^{})`$. This essentially means that the actions are conjugate modulo the central foliations.
We state a special case of some the results of Hirsch-Pugh-Shub on perturbations of partially hyperbolic actions of $``$, see \[HPS\]. There are also analogous definitions and results for flows. As these are less important in the study of local rigidity, we do not discuss them here.
###### Theorem 3.3.
Let $`f`$ be an $`(a,b)`$-partially hyperbolic $`C^k`$ diffeomorphism of a compact manifold $`M`$ which is $`k`$-normally hyperbolic to a $`C^k`$ central foliation $`𝒲_f^c`$. Then for any $`\delta >0`$, if $`f^{}`$ is a $`C^k`$ diffeomorphism of $`M`$ which is sufficiently $`C^1`$ close to $`f`$ we have the following:
1. $`f^{}`$ is $`(a^{},b^{})`$-partially hyperbolic, where $`|aa^{}|<\delta `$ and $`|bb^{}|<\delta `$, and the splitting $`TM=E_f^{}^uE_f^{}^cE_f^{}^s`$ for $`f^{}`$ is $`C^0`$ close to the splitting for $`f`$;
2. there exist $`f^{}`$ invariant foliations by $`C^k`$ leaves $`𝒲_f^{}^c`$ tangent to $`E_f^{}^c`$, which is close in the natural topology on foliations by $`C^k`$ leaves to $`𝒲_f^c`$,
3. there exists a (non-unique) homeomorphism $`h`$ of $`M`$ with $`h(𝒲_f^c)=𝒲_f^{}^c`$, and $`h`$ is $`C^k`$ along leaves of $`𝒲_f^c`$, furthermore $`h`$ can be chosen to be $`C^0`$ small and $`C^k`$ small along leaves of $`𝒲_f^c`$
4. the homeomorphism $`h`$ is a leaf conjugacy between the actions $`(M,𝒲_f^c,f)`$ and $`(M,𝒲_f^{}^c,f^{})`$.
Conclusion $`(1)`$ is easy and probably older than \[HPS\]. One motivation for Theorem 3.3 is to study stability of dynamical properties of partially hyperbolic diffeomorphisms. See the survey, \[BPSW\], for more discussion of that and related issues.
## 4. History
In this section, we describe the history of the subject roughly to the year 2000. More recent developments will be discussed below. Here we do not treat the subject entirely chronologically, but break the discussion into four subjects: first, the study of local rigidity of volume preserving actions, second the study of local rigidity of certain non-volume preserving actions called boundary actions, third the existence of (many) deformations of (many) actions of groups that are typically quite rigid, and lastly a brief discussion of infinitesimal rigidity. This is somewhat ahistorical as the first results on smooth conjugacy of perturbations of group actions appear in \[Gh1\], which we describe in subsection 4.2. While those results are not precisely local rigidity results, they are clearly related and the techniques involved inform some later approaches to local rigidity.
### 4.1. Volume preserving actions.
In this subsection, we discuss local rigidity of volume preserving actions. The acting groups will usually be lattices in higher rank semi-simple Lie groups or higher rank free abelian groups. Many of the results discussed here were motivated by conjectures of Zimmer in \[Z4, Z5\]. The first result we mention, due to Zimmer, does not prove local rigidity, but did motivate much later work on the subject.
###### Theorem 4.1.
Let $`\mathrm{\Gamma }`$ be a group with property $`(T)`$ of Kazhdan and let $`\rho `$ be a Riemannian isometric action of $`\mathrm{\Gamma }`$ on a compact manifold $`M`$. Further assume the action is ergodic. Then any $`C^k`$ action $`\rho ^{}`$ which is $`C^k`$ close to $`\rho `$, volume preserving and ergodic, preserves a $`C^{k3}`$ Riemannian metric.
###### Remarks 3.
1. Zimmer first proved this theorem in \[Z1\], but only for $`\mathrm{\Gamma }`$ a lattice in a semi-simple group, all of whose simple factors have real rank at least two, and only producing a $`C^0`$ invariant metric for $`\rho ^{}`$. In \[Z2\], he gave the proof of the regularity stated here and in \[Z4\] he extended the theorem to all Kazhdan groups.
2. In this theorem if $`\rho ^{}`$ is $`C^{\mathrm{}}`$, the invariant metric for $`\rho ^{}`$ can also be chosen $`C^{\mathrm{}}`$.
The first major result that actually produced a conjugacy between the perturbed and unperturbed actions was due to Hurder, \[H1\]. Again, this result is not quite a local rigidity theorem, but only a deformation rigidity theorem. Hurder’s work is the first place where hyperbolic dynamics is used in the theory, and is the beginning of a long development in which hyperbolic dynamics play a key role.
###### Theorem 4.2.
The standard action of any finite index subgroup of $`SL(n,)`$ on the $`n`$ dimensional torus is deformation rigid when $`n3`$.
Hurder actually proves a much more general result. His theorem proves deformation rigidity of any group $`\mathrm{\Gamma }`$ acting on the $`n`$ dimensional torus by linear transformations provided:
1. the set of periodic points for the $`\mathrm{\Gamma }`$ action is dense, and
2. the first cohomology of any finite index subgroup of $`\mathrm{\Gamma }`$ in any $`n`$ dimensional representation vanishes, and
3. the action contains “enough” Anosov elements.
Here we intentionally leave the meaning of $`(3)`$ vague, as the precise notion needed by Hurder is quite involved. To produce a continuous path of continuous conjugacies, Hurder only need conditions $`(1)`$ and $`(2)`$ and the existence of an Anosov diffeomorphism in the stabilizer of every periodic point for the action. The additional Anosov elements needed in $`(3)`$ are used to improve regularity of the conjugacy, and better techniques for this are now available. Hurder’s work has some applications to actions of irreducible lattices in products of rank $`1`$ Lie groups, which we discuss below in subsection 6.2. These applications do not appear to be accessible by later techniques.
A key element in Hurder’s argument is to use results of Stowe on persistence of fixed points under perturbations of actions \[St1, St2\]. To use Stowe’s result one requires that the cohomology in the derivative representation at the fixed point vanishes. This is where $`(2)`$ above is used. Hurder constructs his conjugacy by using the theorem of Anosov, mentioned above in subsection 3.2, that any Anosov diffeomorphism is structurally stable. This produces a conjugacy $`h`$ for an Anosov element $`\rho (\gamma _0)`$ which one then needs to see is a conjugacy for the entire group action. Hurder uses Stowe’s results to show that $`h`$ is conjugacy for the $`\mathrm{\Gamma }`$ actions at all of the periodic points for the $`\mathrm{\Gamma }`$ action, and since periodic points are dense it is then a conjugacy for the full actions. The precise argument using Stowe’s theorem is quite delicate and we do not attempt to summarize it here. This argument applies much more generally, see \[H1, Theorem 2.9\]. That the conjugacy depends continuously, and in fact even smoothly, on the original action is deduced from results on hyperbolic dynamics in \[dlLMM\].
The first major development after Hurder’s theorem was a theorem of Katok and Lewis \[KL1\] of which we state a special case:
###### Theorem 4.3.
Let $`\mathrm{\Gamma }<SL(n,)`$ be a finite index subgroup, $`n>3`$. Then the linear action of $`\mathrm{\Gamma }`$ on $`𝕋^n`$ is locally rigid.
It is worth noting that this theorem does not cover the case of $`n=3`$. A major ingredient in the proof is studying conjugacies produced by hyperbolic dynamics for certain $``$ actions generated by hyperbolic and partially hyperbolic diffeomorphisms in both the original action and the perturbation. The strategy is to find a hyperbolic generating set and to show that the conjugacies produced by the stability of those diffeomorphisms agree. A key ingredient idea is to show that they agree on the set of periodic orbits. Periodic orbits are then studied via Theorem 3.3 for elements of $`\mathrm{\Gamma }`$ with large centralizers and large central foliations. Periodic orbits are detected as intersections of central foliations for different elements, and this allows the authors to show that the structure of the periodic set persists under deformations. Elements with large centralizers had previously been exploited by Lewis for studying infinitesimal rigidity of similar actions. The authors also exploit their methods to prove the following remarkable result:
###### Theorem 4.4.
Let $`^n`$ be a maximal diagonalizable (over $``$) subgroup of $`SL(n+1,)`$ where $`n2`$. The linear action of $`^n`$ on $`𝕋^{n+1}`$ is locally rigid.
This result is the first in a long series of results showing that many actions of higher rank abelian groups are locally rigid. See Theorems 4.6 and 5.8 below for more instances of this remarkable behavior.
The next major development occurs in a paper of Katok, Lewis and Zimmer where Theorem 4.3 is extended to cover the case of $`n=3`$ as well as some more general groups acting on tori. Though this does not seem, on the face of it, to be a very dramatic development, an important idea is introduced in this paper. The authors proceed by comparing the measurable data coming from cocycle super-rigidity to the continuous data provided by hyperbolic dynamics. In this context, this essentially allows the authors to show that the map $`s`$ in equation $`(\text{1})`$, described in the statement of cocycle superrigidity given in subsection 3.1, is continuous. This idea of comparing the output of cocycle super-rigidity to information provided by hyperbolic dynamics has played a major role in the development of both local and global rigidity of group actions.
The results in the papers \[H1, KL1, KLZ\] are all proven for particular actions of particular groups, and in particular are all proven for actions on tori. The next sequence of developments was a generalization of the ideas and methods contained in these papers to fairly general Anosov actions of higher rank lattices on nilmanifolds. Part of this development takes place in the works \[Q1, Q2, QY\]. A key difficulty in generalizing the early approaches to rigidity of groups of toral automorphims is in adapting the methods from hyperbolic dynamics which are used to improve the regularity of the conjugacy. In \[KS1, KS2\], Katok and Spatzier developed a broadly applicable method for smoothing conjugacies which depends on the theory of non-stationary normal forms as developed by Guysinsky and Katok in \[GK, G\]. Two main consequences of this method are:
###### Theorem 4.5.
Let $`G`$ be a semi-simple Lie group with all simple factors of real rank at least two, $`\mathrm{\Gamma }<G`$ a lattice, $`N`$ a nilpotent Lie group and $`\mathrm{\Lambda }<N`$ a lattice. Then any affine Anosov action of $`\mathrm{\Gamma }`$ on $`N/\mathrm{\Lambda }`$ is locally rigid.
Here by Anosov action, we mean that some element of $`\mathrm{\Gamma }`$ acts on $`N/\mathrm{\Lambda }`$ as an Anosov diffeomorphism.
###### Theorem 4.6.
Let $`N`$ a nilpotent Lie group and $`\mathrm{\Lambda }<N`$ a lattice. Let $`^d`$ be a group of affine transformations of $`N/\mathrm{\Lambda }`$ such that the derivative action (on a subgroup of finite index) is simultaneously diagonalizable over $``$ with no eigenvalues on the unit circle (i.e. on subgroup of finite index, each element of $`^d`$ is an Anosov diffeomorphism which has semi-simple derivative). Then the $`^d`$ action on $`N/\mathrm{\Lambda }`$ is locally rigid.
Katok and Spatzier also apply their method to show that for certain standard Anosov $`^d`$ actions the orbit foliation is locally rigid. I.e. any nearby action has conjugate orbit foliation. This yields interesting applications to rigidity of boundary actions, see Theorem 4.16 below. Also, combined with other results of the same authors on rigidity of cocycles over actions of Abelian groups, this yields local rigidity of certain algebraic actions of $`^d`$, \[KS3, KS4\]. We state a special case of these results here:
###### Theorem 4.7.
Let $`G`$ be an $``$-split semi-simple Lie group of real rank at least two. Let $`\mathrm{\Lambda }<G`$ be a cocompact lattice and let $`^d<G`$ be a maximal $``$-split subgroup. Then the $`^d`$ action on $`G/\mathrm{\Lambda }`$ is locally rigid.
###### Remarks 4.
1. Here “local rigidity” has a slightly different meaning than above. Since the automorphism group of $`^d`$ has non-trivial connected component, it is possible to perturb the action by taking a small automorphism of $`^d`$. What is proven in this theorem is that any small enough perturbation is conjugate to one obtained in this way.
2. Another approach to related cocycle rigidity results is developed in the paper \[KNT\].
3. The actual theorem in \[KS2\] is much more general.
A key ingredient in the Katok–Spatzier method is to find foliations which are orbits of transitive, isometric, smooth group actions for both the perturbed and unperturbed action. To show smoothness of the conjugacy, one constructs such group actions that
1. are intertwined by a continuous conjugacy and
2. exist on enough foliations to span all directions in the space.
This proves that the conjugacy is “smooth along many directions” and one then uses a variety of analytic methods to prove that the conjugacy is actually globally smooth. The fact that the transitive group exists and acts smoothly on the leaves of some foliation for the unperturbed action is typically obvious. One then uses the continuous conjugacy to define the group action along leaves for the perturbed action and the fact that the resulting action is smooth along leaves is verified using the normal form theory. The foliations along which one builds transitive group actions are typically central foliations for certain special elements of the suspension of the action. If the original action is a $`^k`$ action by automorphisms on some nilmanifold $`N/\mathrm{\Lambda }`$, the suspension of the action is the left action of $`^k`$ on the solv-manifold $`M=(^kN)/(^k\mathrm{\Lambda })`$. A typical one parameter subgroup of $`^k`$ acts hyperbolically $`M`$, but certain special directions in $`^k`$, those in so-called Weyl chamber walls give rise to one parameter subgroups with non-trivial central direction. A key fact used in the argument is that one can find another subgroup of $`^k`$ for which the central foliation for some one parameter subgroup $`\rho (t)`$ is also a dynamically defined, contracting foliation for the some other element $`a^k`$.
All results quoted so far have strong assumptions on hyperbolicity of the action. For actions of semi-simple groups and their lattices, the ultimate result on local rigidity in hyperbolic context was proven by Margulis and Qian in \[MQ\]. This result is for so-called weakly hyperbolic actions, which we define below. This work proceeds by first using a comparison between hyperbolic data and data from cocycle superrigidity to produce a continuous conjugacy between $`C^1`$ close actions and then uses an adaptation of the Katok–Spatzier smoothing method mentioned above. A key technical innovation in this work is the choice of cocycle to which cocycle super-rigidity is applied. In all work to this point, it was applied to the derivative cocycle. Here it is applied to a cocycle that measures the difference between the action and the perturbation. To illustrate the idea, we give the definition of this cocycle, which we refer to as the Margulis–Qian cocycle, in the special case of actions by left translations. As this construction is quite general, we will let $`D`$ be the acting group. Let the $`D`$ action $`\rho `$ on $`H/\mathrm{\Lambda }`$ be defined via a homomorphism $`\pi _0:DH`$. Let $`\rho ^{}`$ be a perturbation of $`\rho `$. If $`D`$ is connected it is clear that the action lifts to $`\stackrel{~}{H}`$ and therefore to $`H`$. If $`D`$ is discrete, this lifting still occurs, since the obstacle to lifting is a cohomology class in $`H^2(D,\pi _1(H/\mathrm{\Lambda }))`$ which does not change under a small perturbation of the action. (A direct justification without reference to group cohomology can be found in \[MQ\] section 2.3.) Write the lifted actions of $`D`$ on $`H`$ by $`\stackrel{~}{\rho }`$ and $`\stackrel{~}{\rho }^{}`$ respectively. We can now define a cocycle $`\alpha :D\times HH`$ by
$$\stackrel{~}{\rho }^{}(g)x=\alpha (g,x)x$$
for any $`g`$ in $`D`$ and any $`x`$ in $`H`$. It is easy to check that this is a cocycle and that it is right $`\mathrm{\Lambda }`$ invariant, and so defines a cocycle $`\alpha :D\times H/\mathrm{\Lambda }H`$. See \[MQ\] section 2 or \[FM1\] section 6 for more discussion as well as for more general variants on this definition. We remark that the use of this cocycle allowed Margulis and Qian to prove the first local rigidity results for volume preserving actions of lattices that have no global fixed point. The construction of this cocycle is inspired by a cocycle used by Margulis in his first proof of superrigidity. This is the cocycle $`G\times G/\mathrm{\Gamma }\mathrm{\Gamma }`$ defined by the choice of a fundamental domain for $`\mathrm{\Gamma }`$ in $`G`$. See Example $`4`$ in subsection 6.2 for a more explicit description.
The work of Margulis and Qian applies to actions which satisfy the following condition. This condition essentially says that the action is hyperbolic in all possible directions, at least for some element of the acting group. It is easy to construct weakly hyperbolic actions of lattices, in particular $`\mathrm{\Gamma }`$ acting on $`G/\mathrm{\Lambda }`$ where $`G`$ is a simple Lie group and $`\mathrm{\Lambda }<G`$ is a cocompact lattice and $`\mathrm{\Gamma }`$ is any other lattice in $`G`$. It is important to note that for this example, no element acts as an Anosov diffeomorphism and, with an appropriate choice of $`\mathrm{\Gamma }`$ and $`\mathrm{\Lambda }`$, there are no finite $`\mathrm{\Gamma }`$ orbits.
###### Definition 4.8.
An action $`\rho `$ of a group $`D`$ on a manifold $`M`$ is called weakly hyperbolic if there exist elements $`d_1,\mathrm{},d_k`$ and constants $`a_i>b_i1`$ for $`i=1,\mathrm{},k`$ such that each $`\rho (d_i)`$ is $`(a_i,b_i)`$-partially hyperbolic in the sense of subsection 3.2 and we have $`TM=E_{\rho (d_i)}^s`$. I.e. there are partially hyperbolic elements whose stable (or unstable) directions span the tangent space at any point.
###### Theorem 4.9.
Let $`H`$ be a real algebraic Lie group and $`\mathrm{\Lambda }<H`$ a cocompact lattice. Assume $`G`$ is a semisimple Lie group with all simple factors of real rank at least two and $`\mathrm{\Gamma }<G`$ is a lattice. Then any weakly hyperbolic, affine algebraic action of $`\mathrm{\Gamma }`$ or $`G`$ on $`H/\mathrm{\Lambda }`$ is locally rigid.
###### Remark 5.
This is somewhat more general than the result claimed in \[MQ\], as they only work with certain special classes of affine actions, which they call standard. This result can proven by the methods of \[MQ\], and a proof in precisely this generality can be read out of \[FM1, FM3\], simply by assuming that the common central direction for the acting group is trivial.
The next major result was a remarkable theorem of Benveniste concerning isometric actions. This is a stronger result than Theorem 4.1 because it actually produces a conjugacy, but is weaker in that it requires much stronger assumptions on the acting group.
###### Theorem 4.10.
Let $`\mathrm{\Gamma }`$ be a cocompact lattice in a semi-simple Lie group with all simple factors of real rank at least two. Let $`\rho `$ be an isometric $`\mathrm{\Gamma }`$ action on a compact manifold $`M`$. Then $`\rho `$ is locally rigid.
The proof of this theorem is inspired by the work of Calabi, Vesentini and Weil in the original proofs of Theorem 3.1 and is based on showing that certain deformations of foliated geometric structures are trivial. The argument is much more difficult than the classical case and uses Hamilton’s implicit function theorem. This is the first occasion on which analytic methods like KAM theory or hard implicit function theorems appear in work on local rigidity of group actions. More recently these kinds of methods have been applied more systematically, see subsection 5.2.
The theorems described so far concern actions that are either isometric or weakly hyperbolic. There are many affine actions which satisfy neither of these dynamical hypotheses, but are genuinely partially hyperbolic. Local rigidity results for actions of this kind first arise in work of Nitica and Torok. We state special cases of two of their theorems:
###### Theorem 4.11.
Let $`\mathrm{\Gamma }<SL(n,)`$ be a finite index subgroup with $`n3`$. Let $`\rho _1`$ be the standard $`\mathrm{\Gamma }`$ action on $`𝕋^n`$ and let $`\rho `$ be the diagonal $`\mathrm{\Gamma }`$ action on $`𝕋^n\times 𝕋^m`$ defined by $`\rho _1`$ on the first factor and the trivial action on the second factor. The action $`\rho `$ is deformation rigid.
###### Theorem 4.12.
Let $`\mathrm{\Gamma },\rho _1,\rho `$ be as above and further assume that $`m=1`$. The action $`\rho `$ is locally rigid.
###### Remarks 6.
1. Nitica and Torok prove more general theorems in which both $`\mathrm{\Gamma }`$ and $`\rho _1`$ can be more general. The exact hypotheses required are different in the two theorems.
2. We are being somewhat ahistorical here, Theorem 4.11 predates the work of Margulis and Qian.
3. The conjugacy produced in the papers \[NT1, NT2, T\] is never $`C^{\mathrm{}}`$, but only $`C^k`$ for some choice of $`k`$. The choice of $`k`$ is essentially free and determines the size of perturbations or deformations that can be considered. It should be possible to produce a $`C^{\mathrm{}}`$ conjugacy by combining the arguments in these papers with arguments in \[FM2, FM3\], see the end of subsection 5.1 for some discussion.
The work of Nitica and Torok is quite complex, using several different ideas. The most novel is to study rigidity of cocycles over hyperbolic dynamical systems taking values in diffeomorphism groups. The dynamical system is either the action $`\rho _1`$ in Theorem 4.11 or 4.12 or it’s restriction to any sufficiently generic subgroup containing an Anosov diffeomorphism of $`𝕋^n`$, and the target group is the group of diffeomorphisms of $`𝕋^m`$. This part of the work is inspired by a classical theorem of Livsic and the proof his modelled on his proof. To reduce the rigidity question to the cocycle question is quite difficult and depends on an adaptation of the work of \[HPS\] discussed in subsection 3.2 as well as use of results of Stowe \[St1, St2\]. Regrettably, the technology seems to limit the applicability of the ideas to diagonal actions $`\rho _1\times \mathrm{Id}`$ on products $`M\times N`$ where the action on $`M`$ has many periodic points and the action on $`N`$ is trivial. Theorem 4.12 also depends on the acting group having property $`(T)`$ of Kazhdan. The method of proof of Theorem 4.11 has additional applications, see particularly subsection 6.2 below.
To close this section, we remark that local rigidity is often considerably easier in the analytic setting. Not much work has been done in this direction, but there is an interesting note of Zeghib \[Zg\]. A sample result is the following:
###### Theorem 4.13.
Let $`\mathrm{\Gamma }<SL(n,)`$ be a subgroup of finite index and let $`\rho `$ be the standard action of $`\mathrm{\Gamma }`$ on $`𝕋^n`$. Then any analytic action close enough to $`\rho `$ is analytically conjugate to $`\rho `$. Furthermore, if $`M`$ is a compact analytic manifold on which $`\mathrm{\Gamma }`$ acts trivially and we let $`\stackrel{~}{\rho }`$ be the diagonal action of $`\mathrm{\Gamma }`$ on $`𝕋^n\times M`$, then $`\stackrel{~}{\rho }`$ is also locally rigid in the analytic category.
Zeghib also proves a number of other interesting results for both volume preserving and non-volume preserving actions and it is clear that his method has applications not stated in his note. The key point for all of his arguments is a theorem of Ghys and Cairns that says that one can linearize an analytic action of a higher rank lattice in a neighborhood of a fixed point. Zeghib proves his results by using results of Stowe \[St1, St2\] to find fixed points for the perturbed action and then studying the largest possible set to which the linearization around this point can be extended. We end this section by stating the theorem of Cairns and Ghys from \[CGh\] which Zeghib uses.
###### Theorem 4.14.
Let $`G`$ be a semisimple Lie group of real rank at least two with no compact factors and finite center and let $`\mathrm{\Gamma }<G`$ be a lattice. Then every analytic action of $`\mathrm{\Gamma }`$ with a fixed point $`p`$ is analytically linearizable in a neighborhood of $`p`$.
###### Remarks 7.
1. By analytically linearizable in a neighborhood of $`p`$, we mean that there exists a neighborhood $`U`$ of $`p`$ and an analytic diffeomorphism $`\varphi `$ of $`U`$ into the ambient manifold $`M`$ such that the action of $`\mathrm{\Gamma }`$, conjugated by $`\varphi `$ is the restriction of a linear action to $`\varphi (U)`$.
2. In the same paper, Ghys and Cairns give an example of a $`C^{\mathrm{}}`$ action of $`SL(3,)`$ on $`^8`$ fixing the origin, which is not $`C^0`$ linearizable in any neighborhood of the origin. So the assumption of analyticity in the theorem is necessary.
### 4.2. Actions on boundaries.
In this subsection we discuss rigidity results for groups acting on homogeneous spaces known as “boundaries”. In contrast to the last section, the actions we describe here never preserve a volume form, or even a Borel measure. We will not discuss here all the geometric, function theoretic or probabalistic reasons why these spaces are termed boundaries, but merely describe examples. For us, if $`G`$ is a semisimple Lie group, then a boundary of $`G`$ is a space of the form $`G/P`$ where $`P`$ is a connected Lie subgroup of $`G`$ such that the quotient $`G/P`$ is compact. The groups $`P`$ having this property are often called parabolic subgroups. The space $`G/P`$ is also considered a boundary for any lattice $`\mathrm{\Gamma }<G`$. For more precise motivation for this terminology, see \[Fu1, Fu2, Mo1, Ma4\].
The simplest example of a boundary is for the group $`SL(2,)`$ in which case the only choice of $`P`$ resulting in a non-trivial boundary is the group of upper (or lower) triangular matrices. The resulting quotient is naturally diffeomorphic to the circle and $`SL(2,)`$ acts on this circle by the action on rays through the origin in $`^2`$. We can restrict this $`SL(2,)`$ action to any lattice $`\mathrm{\Gamma }`$ in $`SL(2,)`$. The following remarkable theorem was first proved by Ghys in \[Gh1\]:
###### Theorem 4.15.
Let $`\mathrm{\Gamma }<SL(2,)`$ be a cocompact lattice and let $`\rho `$ be the action of $`\mathrm{\Gamma }`$ on $`S^1`$ described above. If $`\rho ^{}`$ is any perturbation of $`\rho `$, then $`\rho ^{}`$ is smoothly conjugate to an action defined by another embedding $`\pi ^{}`$ of $`\mathrm{\Gamma }`$ in $`SL(2,)`$ close to the original embedding. In particular $`\pi ^{}(\mathrm{\Gamma })`$ is a cocompact lattice in $`SL(2,)`$.
Ghys gives two proofs of this fact, one in \[Gh1\] and another different one in \[Gh2\]. A third and also different proof is in later work of Kononenko and Yue \[KY\]. Ghys’ first proof derives from a remarkable global rigidity result for actions on certain three dimensional manifolds by the affine group of the line, while his second derives from rigidity results concerning certain Anosov flows on three dimensional manifolds. We remark that the fact that $`\rho ^{}`$ is continuously conjugate to an action defined by a nearby embedding into $`SL(2,)`$ was known and so Theorem 4.15 can be viewed as a regularity theorem though this is not how the proof proceeds.
Both of Ghys’ proofs pass through a statement concerning local rigidity of foliations. This uses the following variant on the construction of the induced action. Let $`\mathrm{\Gamma }`$ be a cocompact lattice in $`SL(2,)`$, and let $`\rho `$ be the $`\mathrm{\Gamma }`$ action on $`S^1`$ defined by the action of $`SL(2,)`$ there. There is also a $`\mathrm{\Gamma }`$ action on the hyperbolic plane $`^2`$. We form the manifold $`(^2\times S^1)/\mathrm{\Gamma }_\rho `$, where $`\mathrm{\Gamma }`$ acts diagonally. This manifold is diffeomorphic to the unit tangent bundle of $`^2/\mathrm{\Gamma }`$ which is also diffeomorphic to $`SL(2,)/\mathrm{\Gamma }`$, and the foliation by planes of the form $`^2\times \{point\}`$ is the weak stable foliation for the geodesic flow and also the orbit foliation for the action of the affine group. Given a $`C^r`$ perturbation $`\rho ^{}`$ of the $`\mathrm{\Gamma }`$ action on $`S^1`$, we can form the corresponding bundle $`(^2\times S^1)_\rho ^{}/\mathrm{\Gamma }_\rho ^{}`$, and the foliation by planes of the form $`^2\times \{point\}`$ is $`C^r`$ close to analogous foliation in $`(^2\times S^1)\mathrm{\Gamma }_\rho `$. To show that $`\rho `$ and $`\rho ^{}`$ are conjugate, it suffices to find a diffeomorphism of $`(^2\times S^1)/\mathrm{\Gamma }_\rho `$ conjugating the two foliations. Both of Ghys’ proofs proceed by constructing such a conjugacy of foliations. This reduction to studying local rigidity of foliations has further applications in slightly different settings, see Theorems 4.16 and 4.17 below.
In later work, Ghys proved a remarkable result which characterized an entire connected component of the space of actions of $`\mathrm{\Gamma }`$ on $`S^1`$. Let $`X`$ be the component of $`\mathrm{Hom}(\mathrm{\Gamma },\mathrm{Diff}(S^1))`$ containing the actions described in Theorem 4.15. In \[Gh3\], Ghys showed that this component consisted entirely of actions conjugate to actions defined by embeddings $`\pi ^{}`$ of $`\mathrm{\Gamma }`$ into $`SL(2,)`$ where $`\pi ^{}(\mathrm{\Gamma })`$ is a cocompact lattice. This result builds on earlier work of Ghys where a similar result was proven concerning $`\mathrm{Hom}(\mathrm{\Gamma },\mathrm{Homeo}(S^1))`$. A key ingredient is the use of the Euler class of the action, viewed as a bounded cocycle.
As mentioned above, Ghys’ method of reducing local rigidity of an action to local rigidity of a foliation has had two more applications. The first of these is due to Katok and Spatzier \[KS2\].
###### Theorem 4.16.
Let $`G`$ be a semisimple Lie group with no compact factors and real rank at least two. Let $`\mathrm{\Gamma }<G`$ be a cocompact lattice and $`B=G/P`$ a boundary for $`\mathrm{\Gamma }`$. Then the $`\mathrm{\Gamma }`$ action on $`G/P`$ is locally rigid.
The proof of this result uses an argument similar to Ghys’ to reduce to a need to study regularity of foliations for perturbations of the action of certain connected abelian subgroups of $`G`$ on $`G/\mathrm{\Gamma }`$. The result used in the proof here is the same as the one used in the proof of Theorem 4.7. For $`\mathrm{\Gamma }`$ a lattice but not cocompact, some partial results are obtained by Yaskolko in his Ph.D. thesis \[Yk\].
Following a similar outline, Kanai proved the following:
###### Theorem 4.17.
Let $`G=SO(n,1)`$ and $`\mathrm{\Gamma }<G`$ be a cocompact lattice. Then the action of $`\mathrm{\Gamma }`$ on the boundary $`G/P`$ is locally rigid.
Partial results in this direction were proven earlier by Chengbo Yue. Yue also proves partial results in the case where $`SO(1,n)`$ is replaced by any rank $`1`$ non-compact simple Lie group.
In somewhat earlier work, Kanai had also proven a special case of Theorem 4.16. More precisely:
###### Theorem 4.18.
Let $`\mathrm{\Gamma }<SL(n,)`$ be a cocompact lattice where $`n21`$ and let $`\rho `$ be the $`\mathrm{\Gamma }`$ action on $`S^{n1}`$ by acting on the space of rays in $`^n`$. Then $`\rho `$ is locally rigid.
Kanai’s proof proceeds in two steps. In the first step, he uses Thomas’ notion of a projective connection to reduce the question to one concerning vanishing of certain cohomology groups. In the second step, he uses stochastic calculus to prove a vanishing theorem for the relevant cohomology groups. The first step is rather special, and is what restricts Kanai’s attention to spheres, rather than other boundaries, which are Grassmanians. The method in the second step seems a good deal more general and should have further applications, perhaps in the context of Theorem 5.12 below. While the approach here is similar in spirit to the work of Benveniste in \[Be1\], it should be noted that Kanai does not use a hard implicit function theorem.
We end this subsection by recalling a construction due to Stuck, which shows that much less rigidity should be expected from actions which do not preserve volume \[Sk\]. Let $`G`$ be a semi-simple Lie group and $`P<G`$ a minimal parabolic. Then there always exists a homomorphism $`\pi :P`$. Given any manifold $`M`$ and any action $`s`$ of $``$ on $`M`$, we can then form the induced $`G`$ action $`\rho _s`$ on $`(G\times M)/P`$ where $`P`$ acts on $`G`$ on the left and on $`M`$ by $`\pi s`$. Varying the action $`s`$ varies the action $`\rho _s`$. It is easy to see that if $`\rho _s`$ and $`\rho _s^{}`$ are two such actions, then they are conjugate as $`G`$ actions if and only if $`s`$ and $`s^{}`$ are conjugate. If one picks an irreducible lattice $`\mathrm{\Gamma }`$ in $`G\times G`$, project $`\mathrm{\Gamma }`$ to $`G`$ and restricts the actions to $`\mathrm{\Gamma }`$, then it is also easy to see that the restriction of $`\rho _s`$ and $`\rho _s^{}`$ to $`\mathrm{\Gamma }`$ are conjugate if and only if $`s`$ and $`s^{}`$ are conjugate. The author does not know a proof that this is also true if one simply takes a lattice $`\mathrm{\Gamma }`$ in $`G`$, but believes that this is also true and may even be known.
### 4.3. “Flexible” actions of rigid groups.
In this subsection, I discuss a sequence of results concerning flexible actions of large groups. More or less, the sequence of examples provides counter-examples to most naive conjectures of the form “all of actions of some lattice $`\mathrm{\Gamma }`$ are locally rigid.” There are some groups, for example compact groups and finite groups, all of whose smooth actions are locally rigid. It seems likely that there should be infinite discrete groups with this property as well, but the constructions in this subsection show that one must look beyond lattices in Lie groups for examples.
Essentially all of the examples given here derive from the simple construction of “blowing up” a point or a closed orbit, which was introduced to this subject in \[KL2\]. The further developments after that result are all further elaborations on one basic construction. The idea is to use the “blow up” construction to introduce distinguished closed invariant sets which can be varied in some manner to produce deformations of the action. The “blow up” construction is a classical tool from algebraic geometry which takes a manifold $`N`$ and a point $`p`$ and constructs from it a new manifold $`N^{}`$ by replacing $`p`$ by the space of directions at $`p`$. Let $`P^l`$ be the $`l`$ dimensional projective space. To blow up a point, we take the product of $`N\times P^{dim(N)}`$ and then find a submanifold where the projection to $`N`$ is a diffeomorphism off of $`p`$ and the fiber of the projection over $`p`$ is $`P^{dim(N)}`$. For detailed discussion of this construction we refer the reader to any reasonable book on algebraic geometry.
The easiest example to consider is to take the action of $`SL(n,)`$, or any subgroup $`\mathrm{\Gamma }<SL(n,)`$ on the torus $`𝕋^n`$ and blow up the fixed point, in this case the equivalence class of the origin in $`^n`$. Call the resulting manifold $`M`$. Provided $`\mathrm{\Gamma }`$ is large enough, e.g. Zariski dense in $`SL(n,)`$, this action of $`\mathrm{\Gamma }`$ does not preserve the measure defined by any volume form on $`M`$. A clever construction introduced in \[KL2\] shows that one can alter the standard blowing up procedure in order to produce a one parameter family of $`SL(n,)`$ actions on $`M`$, only one of which preserves a volume form. This immediately shows that this action on $`M`$ admits perturbations, since it cannot be conjugate to the nearby, non-volume preserving actions. Essentially, one constructs different differentiable structures on $`M`$ which are diffeomorphic but not equivariantly diffeomorphic.
After noticing this construction, one can proceed to build more complicated examples by passing to a subgroup of finite index, and then blowing up several fixed points. One can also glue together the “blown up” fixed points to obtain an action on a manifold with more complicated topology. See \[KL2, FW\] for discussion of the topological complications one can introduce.
In \[Be2\] it is observed that a similar construction can be used for the action of a simple group $`G`$ by left translations on a homogeneous space $`H/\mathrm{\Lambda }`$ where $`H`$ is a Lie group containing $`G`$ and $`\mathrm{\Lambda }<H`$ is a cocompact lattice. Here we use a slightly more involved construction from algebraic geometry, and “blow up” the directions normal to a closed submanifold. I.e. we replace some closed submanifold $`N`$ in $`H/\mathrm{\Lambda }`$ by the projectived normal bundle to $`N`$. In all cases we consider here, this normal bundle is trivial and so is just $`N\times P^l`$ where $`l=dim(H)dim(N)`$.
Benveniste used his construction to produce more interesting perturbations of actions of higher rank simple Lie group $`G`$ or a lattice $`\mathrm{\Gamma }`$ in $`G`$. In particular, he produced volume preserving actions which admit volume preserving perturbations. He does this by choosing $`G<H`$ such that not only are there closed $`G`$ orbits but so that the centralizer $`Z=Z_H(G)`$ of $`G`$ in $`H`$ has no-trivial connected component. If we take a closed $`G`$ orbit $`N`$, then any translate $`zN`$ for $`z`$ in $`Z`$ is also closed and so we have a continuum of closed $`G`$ orbits. Benveniste shows that if we choose two closed orbits $`N`$ and $`zN`$ to blow up and glue, and then vary $`z`$ in a small open set, the resulting actions can only be conjugate for a countable set of choices of $`z`$.
This construction is further elaborated in \[F1\]. Benveniste’s construction is not optimal in several senses, nor is his proof of rigidity. In \[F1\], I give a modification of the construction that produces non-conjugate actions for every choice of $`z`$ in a small enough neighborhood. By blowing up and gluing more pairs of closed orbits, this allows me to produce actions where the space of deformations contains a submanifold of arbitrarily high, finite dimension. Further, Benveniste’s proof that the deformation are non-trivial is quite involved and applies only to higher rank groups. In \[F1\], I give a different proof of non-triviality of the deformations, using consequences of Ratner’s theorem to due Witte and Shah \[R, Sh, W2\]. This shows that the construction produces non-trivial perturbations for any semisimple $`G`$ and any lattice $`\mathrm{\Gamma }`$ in $`G`$.
In \[BF\] we show that none of these actions preserve any rigid geometric structure in the sense of Gromov. It is possible that any action of a higher rank lattice which preserves a rigid geometric structure is locally rigid. It is also possible that any such action is generalized quasi-affine.
### 4.4. Infinitesimal rigidity.
In \[Z4\], Zimmer introduced a notion of infinitesimal rigidity motivated by Weil’s Theorem 3.2 and the analogy between finite dimensional Lie algebras and vector fields. Let $`\rho `$ be a smooth action of a group $`\mathrm{\Gamma }`$ on a manifold $`M`$, then $`\rho `$ is infinitesimally rigid if $`H^1(\mathrm{\Gamma },\mathrm{Vect}^{\mathrm{}}(M))=0`$. Here the $`\mathrm{\Gamma }`$ action on $`\mathrm{Vect}^{\mathrm{}}(M)`$ is given by the derivative of $`\rho `$. The notion of infinitesimal rigidity was introduced with the hope that one could prove an analogue of Weil’s Theorem 3.2 and then results concerning infinitesimal rigidity would imply results concerning local rigidity. Many infinitesimal rigidity results were then proven, see \[H2, Ko, L, LZ, Q3, Z6\]. For some more results on infinitesimal rigidity, see Theorems 5.10 and 5.11. Also see subsection 5.2 for a discussion of known results on the relation between infinitesimal and local rigidity.
## 5. Recent developments
In this section we discuss the most recent dramatic developments in the field. The first subsection discusses work of the author and Margulis on rigidity of actions of higher rank groups and lattices. Our main result is that if $`H`$ is the real points of an algebraic group defined over $``$ and $`\mathrm{\Lambda }<H`$ is a cocompact lattice, then any affine action of $`G`$ or $`\mathrm{\Gamma }`$ on $`H/\mathrm{\Lambda }`$ is locally rigid. This work is quite involved and spans a sequence of three long papers \[FM1, FM2, FM3\]. One of the main goals of subsection 5.1 is to provide something of a “reader’s guide” to those papers.
The second subsection discusses some recent developments involving more geometric and analytic approaches to local rigidity. Till this point, the study of local rigidity of group actions has been dominated by algebraic ideas and hyperbolic dynamics with the exception of \[Ka1\] and \[Be2\]. The results described in subsection 5.2 represent (the beginning of) a dramatic development in analytic and geometric techniques. The first of these is the work of Damjanovich and Katok on local rigidity of certain partially hyperbolic affine actions of abelian groups on tori using a KAM approach \[DK1, DK2\]. The second is the author’s proof of a criterion for local rigidity of groups actions modelled on Weil’s Theorem 3.2 and proven using Hamilton’s implicit function theorem \[F2\]. This result currently has an unfortunate “side condition” on second cohomology that makes it difficult to apply.
The final subsection concerns a few other very recent results and developments that the author feels point towards the future development of the field.
### 5.1. The work of Margulis and the author.
Let $`G`$ be a (connected) semi-simple Lie group with all simple factors of real rank at least two, and $`\mathrm{\Gamma }<G`$ is a lattice. The main result of the papers \[FM1, FM2, FM3\] is:
###### Theorem 5.1.
Let $`\rho `$ be a volume preserving quasi-affine action of $`G`$ or $`\mathrm{\Gamma }`$ on a compact manifold $`X`$. Then the action locally rigid.
###### Remarks 8.
1. This result subsumes essentially all of the theorems in subsection 4.1, excepting those concerning actions of abelian groups.
2. In \[FM3\] we also achieve some remarkable results for perturbations of very low regularity. In particular, we prove that any perturbation which is a $`C^3`$ close $`C^3`$ action is conjugate back to the original action by a $`C^2`$ diffeomorphism.
3. The statement here is slightly different than that in \[FM3\]. Here $`X=H/L\times M`$ with $`L`$ cocompact, while there $`X=H/\mathrm{\Lambda }\times M`$ with $`\mathrm{\Lambda }`$ discrete and cocompact. An essentially algebraic argument using results in \[W1\], shows that possibly after changing $`H`$ and $`M`$, these hypotheses are equivalent.
Another main result of the research resulting in Theorem 5.1 is the following:
###### Theorem 5.2.
Let $`\mathrm{\Gamma }`$ be a discrete group with property $`(T)`$. Let $`X`$ be a compact smooth manifold, and let $`\rho `$ be a smooth action of $`\mathrm{\Gamma }`$ on $`X`$ by Riemannian isometries. Then $`\rho `$ is locally rigid.
###### Remarks 9.
1. A key step in the proof of Theorem 5.1 is a foliated version of Theorem 5.2.
2. As in Theorem 5.1, there is a finite regularity version of Theorem 5.2 and it’s foliated generalization, we refer the reader to \[FM2\] for details.
The remainder of this subsection will consist of a sketch of the proof of Theorem 5.1. The intention is essentially to provide a reader’s guide to the three papers \[FM1, FM2, FM3\]. Throughout the remainder of this subsection, to simplify notation, we will discuss only the case of affine $`\mathrm{\Gamma }`$ actions on $`H/\mathrm{\Lambda }`$ with $`\mathrm{\Lambda }`$ a cocompact lattice. The proof for connected groups and quasi-affine action on $`X=H/\mathrm{\Lambda }\times M`$ is similar. To further simplify the discussion, we assume that $`\rho ^{}`$ is a $`C^{\mathrm{}}`$ perturbation of $`\rho `$.
### Step 1: An invariant “central” foliation for the perturbed action and leaf conjugacy.
To begin the discussion and the proof, we need some knowledge of the structure of the affine actions considered. By \[FM1, Theorem 6.4\], there is a finite index subgroup $`\mathrm{\Gamma }^{}<\mathrm{\Gamma }`$ such that the action of $`\mathrm{\Gamma }^{}`$ on $`H/\mathrm{\Lambda }`$ is given by a homomorphism $`\sigma :\mathrm{\Gamma }^{}\mathrm{Aut}(H)H`$. We simplify the discussion by assuming $`\mathrm{\Gamma }=\mathrm{\Gamma }^{}`$ throughout. Using Margulis superrigidity theorems, which are also used in the proof of \[FM1, Theorem 6.4\], it is relatively easy to understand the maximal central foliation $`𝔉`$ for $`\rho `$: there is a subgroup $`Z<H`$ whose orbit foliation is exactly the central foliation. For example, if $`G<H`$ acts on $`H/\mathrm{\Lambda }`$ by left translations and $`\rho `$ is restriction of that action to $`\mathrm{\Gamma }`$, then $`Z=Z_H(G)`$. For details on what $`Z`$ is more generally, see \[FM1\].
Given a perturbation $`\rho ^{}`$ of $`\rho `$, we begin by finding a $`\rho ^{}`$ invariant foliation $`𝔉^{}`$ and a leaf conjugacy $`\varphi `$ from $`(H/\mathrm{\Lambda },\rho ,𝔉)`$ to $`(H/\mathrm{\Lambda },\rho ^{},𝔉^{})`$. To do this, we apply a result concerning local rigidity of cocycles over actions of higher rank groups and lattices to the Margulis–Qian cocycle defined by the perturbation. As the statements of the local rigidity results for cocycles are somewhat technical, we refer the reader to \[FM1, Theorems 1.1 and 5.1\]. Those Theorems are proven in Section $`5`$ of that paper using results in Section $`4`$ concerning orbits in representation varieties as well as the cocycle superrigidity theorems. The construction of the leaf conjugacy is completed in \[FM3, Section 2.2\] using \[FM1, Theorem 1.8\]. We remark that we actually construct $`\varphi ^1`$ rather than $`\varphi `$. The paper \[FM1\] also contains a proof of superrigidity for cocycles that results in many technical improvements to that result.
### Step 2: Smoothness of the central foliation, reduction to a foliated perturbation.
The next step in the proof is to show that $`𝔉^{}`$ is a foliation by smooth leaves. In fact, it is only possible to show at this point that it is a foliation by $`C^k`$ leaves for some $`k`$ depending on $`\rho `$ and $`\rho ^{}`$ and particularly on the $`C^1`$ size of the perturbation. This is done using the work of Hirsch, Pugh and Shub described in subsection 3.2. If the central foliation for $`\rho `$ is the central foliation for $`\rho (\gamma )`$ for some single element $`\gamma `$ in $`\mathrm{\Gamma }`$, this amounts to showing that the foliation $`𝔉^{}`$ constructed in step one is the same foliation as the central foliation for $`\rho ^{}(\gamma )`$ constructed in the proof of Theorem 3.3. To prove this, one needs to analyze the proof of Theorem 3.3. More generally, we show, in \[FM3, Section 3.2\] that there is a finite collection of elements $`\gamma _1,\mathrm{},\gamma _k`$ in $`\mathrm{\Gamma }`$ such that each leaf of the foliation $`𝔉`$ is a transverse intersection of central leaves of $`\rho (\gamma _1),\mathrm{},\rho (\gamma _k)`$. One then needs to combine an analysis of the proof of Theorem 3.3 with some arguments concerning persistence of transversality under certain kinds of perturbations. This argument is carried out in \[FM3, Section 3.3\].
Once we know that $`𝔉^{}`$ is $`C^k`$, it is easy to see that the leaf conjugacy $`\varphi `$ is $`C^k`$ and $`C^k`$ small along leaves of $`𝔉`$ though all derivatives are only continuous in the transverse direction. Conjugating the action $`\rho ^{}`$ by $`\varphi `$, we obtain a new action $`\rho ^{\prime \prime }`$ on $`H/\mathrm{\Lambda }`$ which preserves $`𝔉`$. This action is only continuous, but it is $`C^0`$ close to $`\rho `$ and $`C^k`$ and $`C^k`$ close to $`\rho `$ along leaves of $`𝔉`$. In \[FM2, FM3\] we refer to perturbations of this type as foliated perturbations.
### Step 3: Conjugacy along the central foliation.
The next step is to apply a foliated generalization of Theorem 5.2 to the actions $`\rho `$ and $`\rho ^{\prime \prime }`$. The exact result we apply is \[FM2, Theorem 2.11\] which produces a semi-conjugacy $`\psi `$ between $`\rho `$ and $`\rho ^{}`$. This result is somewhat involved to state and the regularity of $`\psi `$ is hard to describe. The map $`\psi `$ is $`C^{k1\epsilon }`$ along the leaves of $`𝔉`$ at almost all points in $`H/\mathrm{\Lambda }`$, for $`\epsilon `$ depending on the size of the perturbation, but only transversely measurable. In addition, the map $`\psi `$ satisfies a certain Sobolev estimate, that implies that it is $`C^{k1\epsilon }`$ small in a small ball in $`𝔉`$ at most points, and that the $`C^{k1\epsilon }`$ norm is only large on very small sets. Rather than try to make this precise here, we include a sketch of the proof of Theorem 5.2. Before doing so we remark that the map $`\phi =\varphi \psi `$ is a semiconjugacy from between the $`\mathrm{\Gamma }`$ action $`\rho `$ and the $`\mathrm{\Gamma }`$ action $`\rho ^{}`$. The last step in the argument is to improve the regularity of $`\phi `$ which we will discuss following the sketch of the proof of Theorem 5.2
We recall two definitions and another theorem from \[FM2\].
###### Definition 5.3.
Let $`\epsilon 0`$ and $`Z`$ and $`Y`$ be metric spaces. Then a map $`h:ZY`$ is an $`\epsilon `$-almost isometry if
$$(1\epsilon )d_Z(x,y)d_Y(h(x),h(y))(1+\epsilon )d_Z(x,y)$$
for all $`x,yZ`$.
The reader should note that an $`\epsilon `$-almost isometry is a bilipschitz map. We prefer this notation and vocabulary since it emphasizes the relationship to isometries.
###### Definition 5.4.
Given a group $`\mathrm{\Gamma }`$ acting on a metric space $`X`$, a compact subset $`S`$ of $`\mathrm{\Gamma }`$ and a point $`xX`$. The number $`sup_{kS}d(x,kx)`$ is called the $`S`$-displacement of $`x`$ and is denoted $`\mathrm{disp}_S(x)`$.
###### Theorem 5.5.
Let $`\mathrm{\Gamma }`$ be a locally compact, $`\sigma `$-compact group with property $`(T)`$ and $`S`$ a compact generating set. There exist positive constants $`\epsilon `$ and $`D`$, depending only on $`\mathrm{\Gamma }`$ and $`S`$, such that for any continuous action of $`\mathrm{\Gamma }`$ on a Hilbert space $``$ where $`S`$ acts by $`\epsilon `$-almost isometries there is a fixed point $`x`$; furthermore for any $`y`$ in $`X`$, the distance from $`y`$ to the fixed set is not more than $`D\mathrm{disp}_S(y)`$.
We now sketch the proof of Theorem 5.2 for Theorem 5.5. Given a compact Riemannian manifold $`X`$, there is a canonical construction of a Sobolev inner product on $`C^k(X)`$ such that the Sobolev inner product is invariant under isometries of the Riemannian metric, see \[FM2, Section 4\]. We call the completion of $`C^k(X)`$ with respect to the metric induced by the Sobolev structure $`L^{2,k}(X)`$. Given an isometric $`\mathrm{\Gamma }`$ action $`\rho `$ on a manifold $`M`$ there may be no non-constant $`\mathrm{\Gamma }`$ invariant functions in $`L^{2,k}(X)`$. However, if we pass to the diagonal $`\mathrm{\Gamma }`$ action on $`X\times X`$, then any function of the distance to the diagonal is $`\mathrm{\Gamma }`$ invariant and, if $`C^k`$, is in $`L^{2,k}(X\times X)`$.
We choose a smooth function $`f`$ of the distance to the diagonal in $`X\times X`$ which has a unique global minimum at $`x`$ on $`\{x\}\times X`$ for each $`x`$, and such that any function $`C^2`$ close to $`f`$ also has a unique minimum on each $`\{x\}\times X`$. This is guaranteed by a condition on the Hessian and the function is obtained from $`d(x,y)^2`$ by renormalizing and smoothing the function away from the diagonal. This implies $`f`$ is invariant under the diagonal $`\mathrm{\Gamma }`$ action defined by $`\rho `$. Let $`\rho ^{}`$ be another action $`C^k`$ close to $`\rho `$. We define a $`\mathrm{\Gamma }`$ action on $`X\times X`$ by acting on the first factor by $`\rho `$ and on the second factor by $`\rho ^{}`$. For the resulting action $`\overline{\rho }^{}`$ of $`\mathrm{\Gamma }`$ on $`L^{2,k}(X\times X)`$ and every $`\gamma S`$, we show that $`\overline{\rho }^{}(k)`$ is an $`\epsilon `$-almost isometry and that the $`S`$-displacement of $`f`$ is a small number $`\delta `$, where both $`\epsilon `$ and $`\delta `$ can be made arbitrarily small by choosing $`\rho ^{}`$ close enough to $`\rho `$. Theorem 5.5 produces a $`\overline{\rho }^{}`$ invariant function $`f^{}`$ close to $`f`$ in the $`L^{2,k}`$ topology. Then $`f^{}`$ is $`C^{kdim(X)}`$ close to $`f`$ by the Sobolev embedding theorems and if $`kdim(X)2`$, then $`f`$ has a unique minimum on each fiber $`\{x\}\times X`$ which is close to $`(x,x)`$. We verify that the set of minima is a $`C^{kdim(X)1}`$ submanifold and, in fact, the graph of a conjugacy between the $`\mathrm{\Gamma }`$ actions on $`X`$ defined by $`\rho `$ and $`\rho ^{}`$.
Note that this argument yields worse regularity than we discussed in the foliated context or than is stated in Theorem 5.2. There are considerable difficulties involved in achieving lower loss of regularity or in producing a $`C^{\mathrm{}}`$ conjugacy and we do not dwell on these here, but refer the reader to \[FM2\] and also to some discussion in the next step.
### Step 4: Regularity of the conjugacy.
We improve the regularity of $`\phi `$ in three stages. First we show it is a homeomorphism in \[FM3, Section 5.2\]. The key to this argument is proving that there is a finite collection $`\gamma _1,\mathrm{},\gamma _k`$ of elements of $`\mathrm{\Gamma }`$ such that $`\phi `$ takes stable foliations for $`\rho (\gamma _i)`$ to stable foliations for $`\rho ^{}(\gamma _i)`$. If $`\phi `$ were continuous, this is both easy and classical. In our context, we require a density point argument to prove this, which is given in \[FM3, Section 5.1\]. Once we have that stable foliations go to stable foliations, we use this to show that $`\phi `$ is actually uniformly continuous along central foliations and then patch together continuity along various foliations. (In \[FM3, Section 5.3\], we show how to remove the assumption, made above, that the $`\mathrm{\Gamma }`$ action lifts from $`H/\mathrm{\Lambda }`$ to $`H`$. This cannot be done until we have produced a continuous conjugacy.)
The next stage is to show that $`\phi `$ is a finite regularity diffeomorphism. To show this, we show that $`\phi `$ is smooth (with estimates) along certain foliations which span the tangent space. This step is essentially an implementation of the method of Katok and Spatzier described in subsection 4.1. A few technical difficulties occur as we need to keep careful track of estimates in the method for use at later steps and because we need to identify ergodic components of the measure for certain elements in the unperturbed action. After applying the Katok–Spatzier method, we have that $`\phi `$ is smooth along many contracting directions and smooth along the central foliation, and then use a fairly standard argument involving elliptic operators to show that it is actually smooth, and even $`C^k^{}`$ small, for $`k^{}=k1\epsilon \frac{1}{2}dim(H)`$. It would be interesting to see if one could lose less regularity at this step, for example by a method like Journé’s. A key difficulty in adapting the method of \[Jn\] is that we only have a Sobolev estimate along central leaves and not a uniform one.
The last stage of the argument is to show that $`\phi `$ is smooth. There are two parts to this argument. The first is to show that if $`\rho `$ and $`\rho ^{}`$ are $`C^k`$ close, we can actually show that $`\phi `$ is $`C^l`$ for some $`lk`$. The main difficulty here is obtaining better regularity in the foliated version of Theorem 5.2. This requires the use of estimates on convexity of derivatives and estimates on compositions of diffeomorphisms, as well as an iterative method of constructing the semi-conjugacy $`\psi `$, for this we refer the reader to \[FM2, Sections 6 and 7.3\]. Once we know we can produce a conjugacy of greater regularity, we can then approximate $`\phi `$ in the $`C^l`$ topology by a $`C^{\mathrm{}}`$ map and smoothly conjugate $`\rho ^{}`$ to a very small $`C^l`$ perturbation perturbation of $`\rho `$. The point is to iterate this procedure while obtaining estimates on the size of the conjugacy produced at each step. We then show that the iteration converges to produce a smooth conjugacy. We give here a general theorem whose proof follows from arguments in \[FM2, FM3\].
It is convenient to fix right invariant metrics $`d_l`$ on the connected components of $`\mathrm{Diff}^l(X)`$ with the additional property that if $`\phi `$ is in the connected component of $`\mathrm{Diff}^{\mathrm{}}(X)`$, then $`d_l(\phi ,\mathrm{Id})d_{l+1}(\phi ,\mathrm{Id})`$. To fix $`d_l`$, it suffices to define inner products $`,_l`$ on $`\mathrm{Vect}^l(X)`$ which satisfy $`V,V_lV,V_{l+1}`$ for $`V\mathrm{Vect}^{\mathrm{}}(X)`$. As remarked in \[FM2, Section 6\], after fixing a Riemannian metric $`g`$ on $`X`$, it is straightforward to introduce such metrics using the methods of \[FM2, Section 4\].
###### Definition 5.6.
Let $`\mathrm{\Gamma }`$ be a group, $`M`$ a compact manifold and assume $`\rho :\mathrm{\Gamma }\mathrm{Diff}^{\mathrm{}}(M)`$. We say $`\rho `$ is strongly $`C^{k,l,n}`$ locally rigid if for every $`\epsilon >0`$ there exists $`\delta >0`$ such that if $`\rho ^{}`$ is an action of $`\mathrm{\Gamma }`$ on $`M`$ with $`d_k(\rho ^{}(g)\rho (g){}_{}{}^{1},\mathrm{Id})<\delta `$ for all $`gK`$ then there exists a $`C^l`$ conjugacy $`\phi `$ between $`\rho `$ and $`\rho ^{}`$ such that $`d_{kn}(\phi ,\mathrm{Id})<ϵ`$.
We are mainly interested in the case where $`l>k`$.
###### Theorem 5.7.
Let $`\mathrm{\Gamma }`$ be a group, $`M`$ a compact manifold and assume $`\rho :\mathrm{\Gamma }\mathrm{Diff}^{\mathrm{}}(M)`$. Assume that there are constants $`n>0`$ and $`k_00`$ and that for every $`k>k_0`$ there exists an $`l>k`$ such that $`\rho `$ is strongly $`C^{k,l,n}`$ locally rigid. Then $`\rho `$ is $`C^{\mathrm{}}`$ locally rigid.
The proof of Theorem 5.7 follows \[FM3, Corollary 7.2\], though the result is not stated in this generality there.
### 5.2. KAM, implicit function theorems: work of Damjanovich–Katok and the author.
In this subsection, we discuss some new results that use more geometric/analytic methods to approach questions of local rigidity. These methods are entirely independent of methods using hyperbolic dynamics and appear likely to be robustly applicable.
We begin with a theorem of Katok and Damjanovich concerning abelian groups of toral automorphisms. Here we consider actions $`\pi :^n\mathrm{Diff}^{\mathrm{}}(𝕋^m)`$ where $`\pi (^n)`$ lies in $`GL(m,Z)`$ acting on $`𝕋^m`$ by linear automorphisms or more generally where $`\pi (^n)`$ acts affinely on $`𝕋^m`$. An affine factor $`\pi ^{}`$ of $`\pi `$ is another affine action $`\pi ^{}:^n𝕋^l`$ and there is an affine map $`\varphi :𝕋^n𝕋^l`$ such that $`\varphi \pi (v)=\pi ^{}(v)\varphi `$ for every $`v^n`$. We say a factor $`\pi ^{}`$ has rank one if $`\pi ^{}(^n)`$ has a finite index cyclic subgroup.
###### Theorem 5.8.
Let $`\pi :^n\mathrm{Aff}(𝕋^m)`$ have no rank one factors. Then $`\pi `$ is locally rigid.
This theorem is proven by a KAM method. One should note that all the theorems on actions of abelian groups by toral automorphisms stated in subsection 4.1 are special cases of this theorem. (This is not quite literally true. Those theorems apply to perturbations that are only $`C^1`$ close, while the result currently under discussion only applies to actions that are close to very high order.) It is also worthwhile to note that Theorem 5.8 can be proven using no techniques of hyperbolic dynamics.
We now briefly describe the KAM method. Let $`\mathrm{\Gamma }`$ be a finitely generated group and $`\pi :\mathrm{\Gamma }\mathrm{Diff}^{\mathrm{}}(M)`$ a homomorphism. To apply a KAM-type argument, define $`L:\mathrm{Diff}^{\mathrm{}}(M)^k\times \mathrm{Diff}^{\mathrm{}}(M)\mathrm{Diff}^{\mathrm{}}(M)^k`$ by taking
$$L(\varphi _1,\mathrm{},\varphi _k,f)=(\pi (\gamma _1)f\varphi _1f{}_{}{}^{1},\mathrm{},\pi (\gamma _1)f\varphi _1f{}_{}{}^{1})$$
where $`\pi :\mathrm{\Gamma }\mathrm{Diff}^{\mathrm{}}(M)`$ is a homomorphism. If $`\pi ^{}`$ is another $`\mathrm{\Gamma }`$ action on $`M`$ then $`L(\pi ^{}(\gamma _1),\mathrm{},\pi ^{}(\gamma _k),f)=(\mathrm{Id},\mathrm{},\mathrm{Id})`$ implies $`f`$ is a conjugacy between $`\pi `$ and $`\pi ^{}`$, so the problem of finding a conjugacy is the same as finding a diffeomorphism $`f`$ which solves $`L(\pi ^{}(\gamma _1),\mathrm{},\pi ^{}(\gamma _k),f)=(\mathrm{Id},\mathrm{},\mathrm{Id})`$ subject to the constraint that $`\pi ^{}`$ is a $`\mathrm{\Gamma }`$ action.
The $`KAM`$ method proceeds by taking the derivative $`DL`$ of $`L`$ at $`(\pi ,f)`$ and solving the resulting linear equation instead subject to a linear constraint that is the derivative of the condition that $`\pi ^{}`$ is a $`\mathrm{\Gamma }`$ action. This produces an “approximate solution” to the non-linear problem and one proceeds by an iteration. If the original perturbation is of size $`\epsilon `$ then the perturbation obtained after one step in the iteration is of size $`\epsilon ^2`$ at least to whatever order one can control the size of the solution of $`DL`$. This allows one to show that the iteration converges even under conditions where there is some dramatic “loss of regularity”, i.e. when solutions of $`DL`$ are only small to much lower order than the initial data. A standard technique used to combat this loss of regularity is to alter the equation given by $`DL`$ by inserting smoothing operators. That one can solve the linearized equation modified by smoothing operators in place of the original linearized equation and still expect to prove convergence of the iterative procedure depends heavily on the quadratic convergence just described. The main difficulty in applying this outline is obtaining so-called tame estimates on inverses of linearized operators. For a definition of a tame estimate, see following Theorem 5.12.
The KAM method is often presented as a method for proving hard implicit function theorems. The paradigmatic theorem of this kind is due to Hamilton \[Ha1, Ha2\], and is used by the author in the proof of the following theorem. For a brief discussion of the relation between this work and that of Katok and Damjanovich, see the end of this subsection.
###### Theorem 5.9.
Let $`\mathrm{\Gamma }`$ be a finitely presented group, $`(M,g)`$ a compact Riemannian manifold and $`\pi :\mathrm{\Gamma }\mathrm{Isom}(M,g)\mathrm{Diff}^{\mathrm{}}(M)`$ a homomorphism. If $`H^1(\mathrm{\Gamma },\mathrm{Vect}^{\mathrm{}}(M))=0`$ and $`H^2(\mathrm{\Gamma },\mathrm{Vect}^{\mathrm{}}(M))`$ is Hausdorff in the tame topology, the homomorphism $`\pi `$ is locally rigid as a homomorphism into $`\mathrm{Diff}^{\mathrm{}}(M)`$.
I believe the condition on $`H^2(\mathrm{\Gamma },\mathrm{Vect}^{\mathrm{}}(M))`$ should hold automatically under the other hypotheses of the theorem. If this is true, then one has a new proof of Theorem 5.2 using a result in \[LZ\]. There are some other infinitesimal rigidity results that would then yield more novel applications. For example:
###### Theorem 5.10.
Let $`\mathrm{\Gamma }`$ be an irreducible lattice in a semisimple Lie group $`G`$ with real rank at least two. Then for any Riemannian isometric action of $`\mathrm{\Gamma }`$ on a compact manifold $`H^1(\mathrm{\Gamma },\mathrm{Vect}^{\mathrm{}}(M))=0`$.
Theorem 5.10 naturally applies in greater generality, in particular to irreducible $`S`$-arithmetic lattices and to irreducible lattices in products of more general locally compact groups.
To give another result on infinitesimal rigidity, we require a definition. For certain cocompact arithmetic lattices $`\mathrm{\Gamma }`$ in a simple group $`G`$, the arithmetic structure of $`\mathrm{\Gamma }`$ comes from a realization of $`\mathrm{\Gamma }`$ as the integer points in $`G\times K`$ where $`K`$ is a compact Lie group. In this case it always true that the projection to $`G`$ is a lattice and the projection to $`K`$ is dense. We say a $`\mathrm{\Gamma }`$ action is arithmetic isometric if it is defined by projecting $`\mathrm{\Gamma }`$ to $`K`$, letting $`K`$ act by $`C^{\mathrm{}}`$ diffeomorphisms on a compact manifold $`M`$ and restricting the $`K`$ action to $`\mathrm{\Gamma }`$.
###### Theorem 5.11.
For certain congruence lattices $`\mathrm{\Gamma }<SU(1,n)`$, any arithmetic isometric action of $`\mathrm{\Gamma }`$ has $`H^1(\mathrm{\Gamma },\mathrm{Vect}^{\mathrm{}}(M))=0`$.
For a description of which lattices this applies to, we refer the reader to \[F2\]. The theorem depends on deep results of Clozel \[Cl\].
Theorem 5.9 actually follows from the following, more general result.
###### Theorem 5.12.
Let $`\mathrm{\Gamma }`$ be a finitely presented group, $`M`$ a compact manifold, and $`\pi :\mathrm{\Gamma }\mathrm{Diff}^{\mathrm{}}(M)`$ a homomorphism. If $`H^1(\mathrm{\Gamma },\mathrm{Vect}^{\mathrm{}}(M))=0`$ and the sequence
admits a tame splitting then the homomorphism $`\pi `$ is locally rigid. I.e. $`\pi `$ is locally rigid provided there exist tame linear maps
$$V_1:C^1(\mathrm{\Gamma },\mathrm{Vect}^{\mathrm{}}(M))C^0(\mathrm{\Gamma },\mathrm{Vect}^{\mathrm{}}(M))$$
and
$$V_2:C^2(\mathrm{\Gamma },\mathrm{Vect}^{\mathrm{}}(M))C^1(\mathrm{\Gamma },\mathrm{Vect}^{\mathrm{}}(M))$$
such that $`d_1V_1+V_2d_2`$ is the identity on $`C^1(\mathrm{\Gamma },\mathrm{Vect}^{\mathrm{}}(M))`$.
Here $`C^i(\mathrm{\Gamma },Vect^{\mathrm{}})`$ is the group of $`i`$-cochains with values in $`\mathrm{Vect}^{\mathrm{}}(M)`$ and $`d_i`$ are the standard coboundary maps where we have identified the cohomology of $`\mathrm{\Gamma }`$ with the cohomology of a $`K(\mathrm{\Gamma },1)`$ space with one vertex, one edge for each generator in our presentation of $`\mathrm{\Gamma }`$ and one $`2`$ cell for each relator in our presentation of $`\mathrm{\Gamma }`$. A map $`L`$ is called tame if there is an estimate of the type $`Lv_kC_kv_{k+r}`$ for a fixed choice of $`r`$. Here the $`_l`$ can be taken to be the $`C^l`$ norm on cochains with values in $`\mathrm{Vect}^{\mathrm{}}(M)`$. This notion clearly formalizes the notion of being able to solve an equation with some loss of regularity.
The proof of Theorem 5.12 proceeds by reducing the question to Hamilton’s implicit function theorem for short exact sequences and is similar in outline to Weil’s proof of Theorem 3.2.
Fix a finitely presented group $`\mathrm{\Gamma }`$ and a presentation of $`\mathrm{\Gamma }`$. This is a finite collection $`S`$ of generators $`\gamma _1,\mathrm{},\gamma _k`$ and finite collection $`R`$ of relators $`w_1,\mathrm{},w_r`$ where each $`w_i`$ is a finite word in the $`\gamma _j`$ and their inverses. More formally, we can view each $`w_i`$ as a word in an alphabet on $`k`$ letters. Let $`\pi :\mathrm{\Gamma }\mathrm{Diff}^{\mathrm{}}(M)`$ be a homomorphism, which we can identify with a point in $`\mathrm{Diff}^{\mathrm{}}(M)^k`$ by taking the images of the generators. We have a complex:
(2)
Where $`P`$ is defined by taking $`\psi `$ to $`(\psi \pi (\gamma _1)\psi {}_{}{}^{1},\mathrm{},\psi \pi (\gamma _k)\psi {}_{}{}^{1})`$ and $`Q`$ is defined by viewing each $`w_i`$ as a word in $`k`$ letters and taking $`(\psi _1,\mathrm{},\psi _k)`$ to $`(w_1(\psi _1,\mathrm{},\psi _k),\mathrm{},w_r(\psi _1,\mathrm{},\psi _k))`$. To this point this is simply Weil’s proof where $`\mathrm{Diff}^{\mathrm{}}(M)`$ is replacing a finite dimensional Lie group $`H`$. Letting $`\mathrm{Id}`$ be the identity map on $`M`$, it follows that $`P(\mathrm{Id})=\pi `$ and $`Q(\pi )=(\mathrm{Id},\mathrm{},\mathrm{Id})`$. Also note that $`Q{}_{}{}^{1}(\mathrm{Id}_M,\mathrm{},\mathrm{Id}_M)`$ is exactly the space of $`\mathrm{\Gamma }`$ actions. Note that $`P`$ and $`Q`$ are $`\mathrm{Diff}^{\mathrm{}}(M)`$ equivariant where $`\mathrm{Diff}^{\mathrm{}}(M)`$ acts on itself by left translations and on $`\mathrm{Diff}^{\mathrm{}}(M)^k`$ and $`\mathrm{Diff}^{\mathrm{}}(M)^r`$ by conjugation. Combining this equivariance with Hamilton’s implicit function theorem, I show that local rigidity is equivalent to producing a tame splitting of the sequence
(3)
To complete the proof of Theorem 5.12 requires that one compute $`DP_{\mathrm{Id}}`$ and $`DQ_\pi `$ in order to relate the sequence in equation $`(\text{3})`$ to the cohomology sequence in Theorem 5.12.
We remark here that the information needed to split the sequence in Theorem 5.12 is quite similar to the information one would need to apply a KAM method. This is not particularly surprising as Hamilton’s implicit function theorem is a formalization of the KAM method. In particular, to prove Theorem 5.8, one can apply Theorem 5.12 using estimates and constructions from \[DK2, Section 3\] to produce the required tame splitting. This avoids the use of the explicit KAM argument in \[DK2, Section 4\].
Finally, we remark that there is a theorem of Fleming that is an analogue of Theorems 5.9 and 5.12 in the setting of finite, or finite Sobolev, regularity \[Fl\]. This is proven using an infinite dimensional variant of Stowe’s fixed point theorems, though it has recently been reproven by An and Neeb using a new implicit function theorem \[AN\]. With either proof, this result also has a similar condition on second cohomology. We remark here that due to the nature of the respective topologies on spaces of vector fields, the condition on second cohomology in the work of Fleming or An-Neeb is considerably stronger than what is needed in Theorems 5.9 and 5.12. As an illustration, no version of Theorem 5.8 can be proven using these results. This is because a cohomological equation can have solutions with tame estimates, i.e. with some loss of regularity, without having solutions with an estimate at any fixed regularity.
### 5.3. Further results.
In this subsection, we describe a few more recent developments related to the results discussed so far. These results are either very recent or somewhat removed from the main stream of research.
The first result we discuss concerns actions of lattices in $`Sp(1,n)`$ or $`F_4^{20}`$ and is due to T.J.Hitchman. We state here only a special case of his results. For this result, we assume that $`\mathrm{\Gamma }`$ is an arithmetic subgroup of $`Sp(1,n)`$ or $`F_4^{20}`$ in the standard $``$ structures on those groups. This means that $`\mathrm{\Gamma }`$ is a finite index subgroup of the integer points in the standard matrix representation of these groups. This means that in the defining representation of $`Sp(1,n)`$ or $`F_4^{20}`$ on $`^m`$ the action of $`\mathrm{\Gamma }`$ preserves the integer lattice $`^m`$ and therefore defines an action $`\rho `$ of $`\mathrm{\Gamma }`$ on $`𝕋^m`$.
###### Theorem 5.13.
The action $`\rho `$ defined in the last paragraph is deformation rigid.
The proof proceeds in two steps. Building a path of $`C^0`$ conjugacies follows more or less as in \[H1\], see subsection 4.1 above. The main novelty in \[Hi\] is the proof that these conjugacies are in fact smooth. Theorem 5.13 is a special case of the results obtained in \[Hi\].
Another recent development should lead to a common generalization of Theorems 5.8 and 4.7. For example, one can consider actions of an abelian subgroup $`^k`$ of the full diagonal group $`^{n1}`$ in $`SL(n,)`$ on $`SL(n,)/\mathrm{\Lambda }`$ where $`\mathrm{\Lambda }`$ is a cocompact lattice. In this context, Damjanovich and Katok are developing a more geometric approach in contrast to the analytic method of \[DK1, DK2\]. In \[DK3\], under a natural non-degeneracy condition on the subgroup $`^k<^{n1}`$, the authors prove the cocycle rigidity result required to generalize Theorem 4.7 to the natural result for the $`^k`$ action. Here there are “additional trivial perturbations” of the action arising from $`\mathrm{Hom}(^k,^{n1})`$. A rigidity theorem in this context is work in progress, see \[DK3\] for some discussion.
To close this section, we mention two other recent works. The first is a paper by Burslem and Wilkinson which investigates local and global rigidity questions for actions of certain solvable groups on the circle. Particularly striking is their construction of group actions which admit $`C^r`$ perturbations but no $`C^{r+1}`$ perturbations for every integer $`r`$. The second is a paper by M.Einsiedler and T.Fisher in which the method of proof of Theorem 4.6 is extended to affine actions of $`^d`$ where the matrices generating the group action have non-trivial Jordan blocks. For perturbations of the group action which are close to very high order, this result follows from Theorem 5.8, but in \[EF\] the result only requires that the perturbation be $`C^1`$ close to the original action.
## 6. Directions for future research and conjectures
In this section, I mention a few conjectures and point a few directions for future research. These are particularly informed by my taste.
### 6.1. Actions of groups with property $`(T)`$.
Lattices in $`SP(1,n)`$ and $`F_4^{20}`$ share many of the rigidity properties of higher rank lattices. In light of Theorems 5.1, 5.2 and 5.13, it seems natural to conjecture:
###### Conjecture 6.1.
Let $`G`$ be a semi-simple Lie group with no compact factors and no simple factors isomorphic to $`SO(1,n)`$ or $`SU(1,n)`$, and let $`\mathrm{\Gamma }<G`$ be a lattice. Then any volume preserving generalized quasi-affine action of $`G`$ or $`\mathrm{\Gamma }`$ on a compact manifold is locally rigid.
Note that for $`G`$ with no rank one factors and quasi-affine actions, this is just Theorem 5.1. To apply the methods of \[FM1, FM2, FM3\] in the setting of Conjecture 6.1 there are essentially three difficulties:
1. If the action is generalized quasi-affine and not quasi-affine, then one cannot use the construction of the Margulis–Qian cocycle described above. This is easiest to see for a generalized affine action $`\rho `$ on some $`K\backslash H/\mathrm{\Lambda }`$. The action $`\rho `$ lifts to $`H/\mathrm{\Lambda }`$, but a perturbation $`\rho ^{}`$ need not. If $`K`$ is finite, this difficulty can be overcome by passing to a subgroup of finite index $`\mathrm{\Gamma }^{}<\mathrm{\Gamma }`$ and arguments in \[FM3\] can be used prove rigidity of $`\mathrm{\Gamma }`$ from rigidity of $`\mathrm{\Gamma }^{}`$. If $`K`$ is compact and connected, this is a genuine and surprisingly intractable difficulty.
2. Proving a version of Zimmer’s cocycle super-rigidity for groups as in the assumptions of the conjecture. Partial results in this direction were obtained by Corlette–Zimmer and Korevaar–Schoen, but their results all require hypotheses that are obviously restrictive or simply difficult to verify. Very recently, the author and T.J.Hitchman have proven a complete version of cocycle super-rigidity, at least for so-called $`L^2`$-cocycles. In light of our work, this difficulty is already overcome.
3. Replace the method of Katok–Spatzier in the proof that the conjugacy is actually smooth. The proof of Theorem 5.13 gives some progress in this direction, but there are significant technical difficulties to overcome in applying Hitchman’s methods at this level of generality. There is some progress on this question by Gorodnik, Hitchman and Spatzier.
The author and T.J.Hitchman have another approach to Conjecture 6.1 based on Theorem 5.12, some estimates proven in \[FH1\], and using heat flow and those estimates to produce a tame splitting of the short exact sequence in Theorem 5.12. It is not yet clear how generally applicable this method will be.
The following question is also interesting in this context:
###### Question 6.2.
Let $`G`$ and $`\mathrm{\Gamma }`$ be as in Conjecture 6.1. Let $`\rho `$ be a non-volume preserving affine action of $`\mathrm{\Gamma }`$ on a compact manifold $`M`$. When is $`\rho `$ locally rigid?
By the work of Stuck discussed at the end of Section 4.2, it is clear that local rigidity will not hold in full generality here. In particular for the product of the action of $`\mathrm{\Gamma }`$ on the boundary $`G/P`$ with the trivial action of $`\mathrm{\Gamma }`$ on any manifold, there is already strong evidence against local rigidity. We give two particularly interesting special cases where we expect local rigidity to occur. The first is to take a lattice $`\mathrm{\Gamma }`$ in $`SP(1,n)`$ and ask if the $`\mathrm{\Gamma }`$ action on the boundary $`SP(1,n)/P`$ is locally rigid. This is analogous to Theorems 4.16 and 4.18. Another direction worth pursuing is to see if the action of say $`Sl(n,)`$ on $`P^{n+k}`$ is locally rigid for $`n3`$ and any $`k0`$. For cocompact lattices instead of $`SL(n,)`$ and $`k=1`$ it is Theorem 4.16. One can ask a wide variety of similar questions for both compact and non-cocompact lattices acting on homogeneous spaces that are “larger than” any natural boundary for the group as long as one avoids settings in which Stuck’s examples can occur. We remark that in some instances, partial results for analytic perturbations can be obtained by Zeghib’s method \[Zg\].
We remark here that the action of $`SL(n,)`$ on $`𝕋^n`$ is not locally rigid in $`\mathrm{Homeo}(𝕋^n)`$ by a construction of Weinberger. It would be interesting to understand local rigidity in low regularity for other actions.
### 6.2. Actions of irreducible lattices in products.
In this subsection, we formulate a general conjecture concerning local rigidity of actions of irreducible lattices in products. We begin by making a few remarks on other rigidity properties of irreducible lattices and by describing a few examples where rigidity might hold, as well as some examples where it does not.
Rigidity properties of irreducible lattices have traditionally been studied together with rigidity of lattices in simple groups, and irreducible lattices enjoy many of the same rigidity properties. We list a few here to motivate our conjectures on rigidity of actions of these lattices. The properties we list also rule out certain trivial constructions of perturbations and deformations of actions. In the following, $`\mathrm{\Gamma }`$ is an irreducible lattice in $`G=(_{i=1}^kG_i)/Z`$ where each $`G_i`$ is a noncompact semi-simple Lie group and $`Z`$ is a subgroup of the center of $`_{i=1}^kG_i`$. Many of these results hold more generally, see below.
Properties of irreducible lattices:
1. There are no non-trivial homomorphisms $`\mathrm{\Gamma }`$ (and therefore no non-trivial homomorphisms to any abelian or non-abelian free group).
2. All linear representations of $`\mathrm{\Gamma }`$ are classified. In particular, given any representation $`\rho `$ of $`\mathrm{\Gamma }`$ into $`GL(V)`$, where $`V`$ is a finite dimensional vector space, then $`H^1(\mathrm{\Gamma },V)=0`$.
3. All normal subgroups of $`\mathrm{\Gamma }`$ are either finite or finite index.
The first property, for $`\mathrm{\Gamma }`$ cocompact, was originally proven by Bernstein and Kazhdan \[BK\]. Some special cases of this result were proven earlier by Matsushima and Shimura \[MS\]. All other properties are originally due to Margulis, see \[Ma4\] and historical references there. The original proofs all go through with only minor adaptations if some of the $`G_i`$ are replaced with the $`k`$-points of a $`k`$-algebraic group over some other local field $`k`$. These properties have been shown to hold for appropriate classes of lattices in even more general products of locally compact groups by Bader, Monod and Shalom, \[Sm, Md1, Md2, BSh\].
Example $`1`$: Let $`\mathrm{\Gamma }=SL(2,[\sqrt{2}])`$. We embed $`\mathrm{\Gamma }`$ in $`SL(2,)\times SL(2,)`$ by taking $`\gamma `$ to $`(\gamma ,\sigma (\gamma ))`$ where $`\sigma `$ is the Galois automorphism of $`(\sqrt{2})`$ taking $`\sqrt{2}`$ to $`\sqrt{2}`$. This embedding defines an action of $`\mathrm{\Gamma }`$ on $`𝕋^4`$ where $`𝕋^4=^4/^4`$ where we identity $`^4`$ with the image in $`^4`$ of $`[\sqrt{2}]^2`$ via the embedding $`v(v,\sigma (v))`$. We first note that the list of properties given above imply one can prove deformation rigidity of this action using Hurder’s argument from \[H1\] to produce a continuous conjugacy and using the method of Katok–Spatzier \[KS1, KS2\] to show that the conjugacy is smooth. Anatole Katok has suggested one might be able to show local rigidity of this action by using the methods in \[KL1\]. This example is just the first in a large class of Anosov actions of irreducible lattices, all of which should be locally rigid. We leave the general construction to the interested reader.
Example $`2`$: We take the action of $`\mathrm{\Gamma }=SL(2,[\sqrt{2}]`$ and let $`\mathrm{\Gamma }`$ act on $`𝕋^5=𝕋^4\times 𝕋^1`$ by a diagonal action where the action on $`𝕋^4`$ is as defined in Example $`1`$ and the action on $`𝕋^1`$ is trivial. Once again, this is merely the first example of a large class of partially hyperbolic actions of irreducible lattices. In this instance, the central foliation for the $`\mathrm{\Gamma }`$ action consists of compact tori. For this type of example, many of the argument of \[NT1, NT2, T\] carry over, but the fact that $`\mathrm{\Gamma }`$ does not have property $`(T)`$ prevents one from using those outlines to prove local rigidity. On the other hand, the methods of \[NT1\] can be adapted to prove deformation rigidity again replacing their argument for smoothness of the conjugacy by the Katok–Spatzier method. We remark that it is also possible to give many examples of actions of irreducible lattices in products where the central foliation is by dense leaves, and the methods of Nitica and Torok cannot be applied. See discussion below for the obstructions to applying the methods of \[FM1, FM2, FM3\].
Example $`3`$: We now give an example of a family of actions which extend to an action of $`G=G_1\times G_2`$. Let $`H`$ be a simple Lie group with $`G<H`$, for example, $`H=SL(4,)`$ or $`H=Sp(4,)`$ and $`\mathrm{\Lambda }<H`$ a cocompact lattice. Then both $`G`$ and $`\mathrm{\Gamma }`$ act by left translations on $`H/\mathrm{\Lambda }`$.
Example $`4`$:We end with a family of examples for which there exists a large family of deformations. Let $`\mathrm{\Gamma }`$ be as above and let $`\mathrm{\Lambda }<SL(2,)=G_1`$ be an irreducible lattice. Then $`\mathrm{\Gamma }`$ acts by left translations on $`M=SL(2,)/\mathrm{\Lambda }`$. Call this action $`\overline{\rho }`$. I do not currently know whether it is possible to deform this action, but one can use this action to build actions with perturbations on a slightly larger manifold. Let $`\mathrm{\Gamma }`$ act trivially on any manifold $`N`$ and take the diagonal action $`\rho `$ on $`M\times N`$. It is well-known that there exists a non-trivial homomorphism $`\sigma :\mathrm{\Lambda }`$. There is also a standard construction of a cocycle $`\alpha :G_1\times G_1/\mathrm{\Lambda }\mathrm{\Lambda }`$ over the $`G_1`$ action on $`G_1/\mathrm{\Lambda }`$. The cocycle is defined by taking a fundamental domain $`X`$ for $`\mathrm{\Lambda }`$ in $`G_1`$, identifying $`G_1/\mathrm{\Lambda }`$ with $`X`$ and letting $`\alpha (g,x)`$ be the element of $`\mathrm{\Lambda }`$ such that $`gx\alpha (g,x)^1`$, as an element of $`G_1`$, is in $`X`$. Taking any vector field $`V`$ on $`N`$ and let $`s_\epsilon `$ be the $``$ action on $`N`$ defined by having $`1`$ act by flowing to time $`\epsilon `$ along $`V`$. We then define a $`\mathrm{\Gamma }`$ action on $`M\times N`$ by taking
$$\rho _\epsilon (\gamma )(m,n)=(\overline{\rho }(\gamma )m,s_\epsilon (\sigma (\alpha (g,m))(n)).$$
We leave it to the interested reader to show that the actions $`\rho _\epsilon `$ are not conjugate to $`\rho `$ for essentially any choice of $`V`$.
The key point in example $`4`$, which is not present in the first three examples, is that the action factors through a projection of $`\mathrm{\Gamma }`$ into a simple factor of $`G`$. Motivated by the examples so far, by the results in Theorem 5.10, and by analogy with results on actions of higher rank abelian groups, we make the following definition and conjecture:
###### Definition 6.3.
Let $`G=G_1\times \mathrm{}\times G_k`$ be a semisimple Lie group where all the $`G_i`$ are non-compact. Let $`\mathrm{\Gamma }<G`$ be an irreducible lattice. Let $`\rho `$ be an affine action of $`\mathrm{\Gamma }`$ on some $`H/\mathrm{\Lambda }`$ where $`H`$ is a Lie group and $`\mathrm{\Lambda }<H`$ is a cocompact lattice. Then we say $`\rho `$ has rank one factors if there exists
1. an action $`\overline{\rho }`$ of $`\mathrm{\Gamma }`$ on a some space $`X`$ which is a factor of $`\rho `$
2. and a rank one factor $`G_i`$ of $`G`$
such that $`\overline{\rho }`$ is the restriction of a $`G_i`$ action. I.e. $`\mathrm{\Gamma }`$ acts on $`X`$ by projecting $`\mathrm{\Gamma }`$ to $`G_i`$ and restricting a $`G_i`$ action.
###### Conjecture 6.4.
Let $`G,\mathrm{\Gamma },H,\mathrm{\Lambda }`$ and $`\rho `$ be as in Definition 6.3. Then if $`\rho `$ has no rank one factors $`\rho `$ is locally rigid.
###### Remark 10.
It is a consequence of Ratner’s measure rigidity theorem, see \[R, Sh, W2\], that any rank one factor of an affine action for these groups is in fact affine. This implies that any rank one factor is a left translation action on some $`H^{}/\mathrm{\Lambda }^{}`$. So a special case of the conjecture is that any affine $`\mathrm{\Gamma }`$ action on a torus or nilmanifold is locally rigid.
There is a variant of Conjecture 6.4 for $`G`$ actions. We recall that a measure preserving action of $`G=G_1\times \mathrm{}\times G_k`$ action is irreducible if each $`G_i`$ acts ergodically. We extend this notion to non-ergodic $`G`$ actions by saying that the action is weakly irreducible if every ergodic component of the volume measure for the action of any $`G_i`$ is an ergodic component of the volume measure for the action of $`G`$.
###### Conjecture 6.5.
Let $`G,H,\mathrm{\Lambda }`$ be as in Definition 6.3 and let $`G<H`$, so that we have a left translation action $`\rho `$ of $`G`$ on $`H/\mathrm{\Lambda }`$. Then if $`\rho `$ is weakly irreducible, $`\rho `$ is locally rigid.
The relation of the two conjectures follows from the following easy lemma.
###### Lemma 6.6.
Let $`G,\mathrm{\Gamma },H,\mathrm{\Lambda }`$ and $`\rho `$ be as in Definition 6.3, then the $`\mathrm{\Gamma }`$ action on $`H/\mathrm{\Lambda }`$ has no rank one factors if and only if the induced $`G`$ action on $`(G\times H/\mathrm{\Lambda })/\mathrm{\Gamma }`$ is weakly irreducible.
To prove the lemma requires both some algebraic untangling of induced actions and a use of Margulis’ super-rigidity theorem to describe affine actions of $`G`$ and $`\mathrm{\Gamma }`$ along the lines of \[FM1, Theorems 6.4 and 6.5\]. We leave this as an exercise for the interested reader. It is fairly easy to check the lemma for any particular affine $`\mathrm{\Gamma }`$ action.
We end this subsection by pointing out the difficulty in approaching this conjecture by means of the methods of \[FM1, FM2, FM3\]. A central difficulty is that the lattices in question do not have property $`(T)`$ and so the foliated generalization of Theorem 5.2 does not apply. However, even in the case of weakly hyperbolic actions, there are significant difficulties. To begin the argument, one would like to apply cocycle superrigidity to the Margulis–Qian cocycle. To do this requires the existence of an invariant measure which is usually established using property $`(T)`$ by an argument of Seydoux \[Sy\]. In this setting, where property $`(T)`$ does not hold, one might try instead to use the work of Nevo and Zimmer, \[NZ1, NZ2, NZ3\], but there are non-trivial difficulties here as well. One cannot apply their theorems without first showing that the perturbed action satisfies some irreducibility assumption. Even if one were to obtain an invariant measure, the precise form of cocycle superrigidity required is not known for products of rank one groups or their irreducible lattices. And the strongest possible forms of cocycle superrigidity in this context again require a kind of irreducibility of the perturbed action. So to proceed by this method one would need to show that perturbations of the actions in Conjecture 6.4 and 6.5 still satisfied some irreducibility conditions. This seems quite difficult. It may also be possible to approach these questions by using Theorem 5.12, but even proving that the relevant cohomology groups vanish seems subtle.
### 6.3. Other questions and conjectures.
We end this article by discussing some other questions and conjectures.
In the context of Theorem 5.11 it is interesting to ask if isometric actions of lattices in $`SU(1,n)`$ are locally rigid. For some choices of lattice, the answer is trivially no. Namely some cocompact lattices in $`SU(1,n)`$ have homomorphisms $`\rho `$ to $``$ \[Ka2, BW\], and so have arithmetic actions with deformations provided the centralizer $`Z`$ of $`K`$ in $`\mathrm{Diff}^{\mathrm{}}(M)`$ is non-trivial. Having centralizer allows one to deform the action along the image of the homomorphism $`\rho \sigma _t:FZ`$ where $`\sigma _t:Z`$ is any one parameter family of homomorphisms. It seems reasonable to conjecture:
###### Conjecture 6.7.
Let $`\rho `$ be an arithmetic isometric action of a lattice in $`SU(1,n)`$. Then any sufficiently small perturbation of $`\rho `$ is of the form described in the previous paragraph.
This conjecture is in a certain sense an infinite dimensional analogue of work of Goldman–Millson and Corlette \[Co1, GM\]. Another conjecture concerning complex hyperbolic lattices, for which work of Yue provides significant evidence \[Yu\], is:
###### Conjecture 6.8.
Is the action of any lattice in $`SU(1,n)`$ on the boundary of complex hyperbolic space locally rigid?
There are also many interesting questions concerning the failure of local rigidity for lattices in $`SO(1,n)`$. The only rigidity theorem we know of in this context is Kanai’s, Theorem 4.17, and it would be interesting to extend Kanai’s theorem to non-uniform lattices. In \[F1, F3\] various deformations of lattices in $`SO(1,n)`$ are constructing for affine and isometric actions. These constructions both adapt the bending construction of Johnson and Millson, \[JM\]. It seems likely that in some cases one should be able to prove results concerning the structure of the representation space and, in particular, to show that it is “singular” in an appropriate sense. See \[F3\] for more discussion.
Two other paradigmatic examples of large groups are the outer automorphism group of the free group, $`\mathrm{Out}(F_n)`$, and the mapping class group of a surface $`S`$, $`MCG(S)`$. These groups do not admit many natural actions on compact manifolds, but there are some natural interesting actions quite analogous to those we have already discussed. For $`MCG(S)`$, the question we raise here is already raised in \[La\]. The actions we consider are “non-linear” analogues of the standard actions of $`SL(n,)`$ on $`𝕋^n`$ and $`SP(2n,)`$ on $`𝕋^{2n}`$. The spaces acted upon are moduli spaces of representations of either the free group or the fundamental group of a surface $`S`$, where the representations take values in compact groups. More precisely, we have an action of $`\mathrm{Out}(F_n)`$ on $`\mathrm{Hom}(F_n,K)/K`$ and an action of $`MCG(S)`$ on $`\mathrm{Hom}(\pi (S),K)/K`$ where $`K`$ is a compact group. It is natural to ask whether these actions are locally rigid, though the meaning of the question is somewhat obscured by the fact that the representation varieties are not smooth. For $`K=S^1`$, one obtains actions on manifolds, and in fact tori, and one might begin by considering that case.
We end with a question motivated by the recent work of Damjanovic and Katok. We only give a special case here. Let $`G`$ be a real split, simple Lie group of real rank at least two. Let $`\mathrm{\Gamma }<G\times G`$ be an irreducible lattice. Let $`K<G`$ be a maximal compact subgroup and view $`K`$ as a subgroup of $`G\times G`$ by viewing it as a subgroup of the second factor. The quotient $`K\backslash (G\times G)/\mathrm{\Gamma }`$ has a natural $`G`$ action on the left on the first factor. We can restrict this action to the action of a maximal split torus $`A`$ in $`G`$. Note that $`A`$ is isomorphic to $`^d`$ for some $`d2`$.
###### Question 6.9.
Is the action of $`^d`$ described in the paragraph above locally rigid? |
warning/0507/hep-th0507187.html | ar5iv | text | # 𝐴-Model Correlators from the Coulomb Branch
## 1 Introduction
Topological Field Theory (TFT) is a powerful tool for studying the RG-invariant properties of non-trivial quantum field theories. A particularly important class of examples is provided by the $`A`$ and $`B`$ twists of an $`𝒩=(2,2)`$ SUSY Non-Linear Sigma Model (NLSM) defined on a Riemann surface $`\mathrm{\Sigma }`$ . These TFTs provide rich examples of solvable quantum field theories, and they have important applications to compactification in string theory. In addition, these TFTs can be used to study enumerative geometry and refined topological invariants, such as the Gromov-Witten invariants, of many manifolds. Finally, these theories provide a natural setting for the study of mirror symmetry of Calabi-Yau manifolds. Although the TFT perspective on the NLSM is immediately conceptually useful, aside from particularly tractable examples such as the classic work of Candelas et al and its various generalizations, direct study of these models is difficult.
Remarkably, a large class of NLSMs, including those corresponding to Calabi-Yau manifolds constructed as hypersurfaces or complete intersections in Fano toric varieties, may be constructed as IR limits of certain $`𝒩=(2,2)`$ SUSY abelian gauge theories termed Gauged Linear Sigma Models (GLSMs). The RG-invariant observables of the $`A`$ and $`B`$ models may be computed in the massive theory. The $`A`$-model in particular has been studied extensively in . It was found that many properties of the GLSM, including the correlators in the topologically twisted models, are constrained by toric geometry. For example, toric methods allow an explicit and general formulation of the instanton sum for the “toric” subset of $`A`$-model correlators. The study of these instanton sums is enlightening: it allows for a careful definition of the monomial-divisor mirror map, and it reduces mirror symmetry to the mirror map, a non-trivial renormalization of the GLSM versus NLSM parameters. Furthermore, for a genus zero Riemann surface, the quantum restriction formula of reduces the computation of topological correlators for a Calabi-Yau hypersurface to the computation of correlators on the ambient Fano toric variety. Where the mirror map is known, this allows for an explicit verification of mirror symmetry at the level of TFT. These are powerful results. Still, the instanton sums are unwieldy. In general they may be quite intricate, and simple but important properties of the correlators, like the quantum cohomology relations, seem to follow from obscure relations between intersection numbers on various toric varieties . Finally, the toric methods for performing the gauge instanton sums have only been developed for genus zero correlators, and their extension to $`g>0`$ correlators is non-trivial.
In this work we reduce the computation of $`A`$-model correlators in a wide class of GLSMs to a well-studied algebraic problem. Our result is simple to state. Consider a GLSM with a set of chiral matter $`𝒩=(2,2)`$ multiplets $`\mathrm{\Phi }^i`$, $`i=1,\mathrm{},n`$, charged under a gauge group $`\left[\mathrm{U}(1)\right]^r`$ with charges $`Q_i^a`$, $`a=1,\mathrm{},r`$, and zero superpotential for the matter fields. Such models will be referred to as toric GLSMs. Upon twisting, the local $`A`$-model observables are found to be functions of the $`\sigma _a`$, the bosonic super-partners of the gauge fields. The most general $`A`$-model correlator on $`\mathrm{\Sigma }_g`$, a Riemann surface of genus $`g`$, may be obtained from linear combinations of $`\sigma _{a_1}(z_1)\mathrm{}\sigma _{a_s}(z_s)_g`$ and derivatives of these with respect to the GLSM parameters. Since this is a TFT computation, the correlator is independent of the generic points $`z_1,\mathrm{},z_s`$ on $`\mathrm{\Sigma }_g`$. Following a standard notation, we will denote these correlators by $`F(\sigma )_g`$, where $`F(\sigma )`$ is to be understood as a power series in $`\sigma _{a_k}(z_k)`$ with a generic choice of the $`z_k`$. In a sense made more precise below, “most” toric GLSMs possess a region of parameter space where the theory has a number of discrete Coulomb vacua determined as solutions to the equations of motion for a certain effective twisted superpotential $`\stackrel{~}{W}_{\text{eff}}(\sigma )`$. This class includes all GLSMs corresponding to compact toric varieties. For these compact, toric GLSMs there is an additional simplification: there exists a region of the parameter space where these discrete Coulomb vacua are the only vacua. In this case, the $`A`$-model correlators at genus $`g`$ are given by
$$F(\sigma )_g=\underset{\widehat{\sigma }|d\stackrel{~}{W}_{\text{eff}}(\widehat{\sigma })=0}{}\stackrel{~}{H}(\widehat{\sigma })^{g1}F(\widehat{\sigma }),$$
(1)
where $`\stackrel{~}{H}(\sigma )=H_{i=1}^n\xi _i`$, $`H`$ is the Hessian of $`\stackrel{~}{W}_{\text{eff}}`$, and $`\xi _i=_aQ_i^a\sigma _a.`$ The equations of motion that follow from $`d\stackrel{~}{W}_{\text{eff}}=0`$ are polynomial in the $`\sigma _a`$, so the computation of $`A`$-model correlators is now reduced to an algebraic problem.
This form is convenient for obtaining explicit correlators in many toric examples. Furthermore, by the use of the quantum restriction formula of , it becomes a useful tool to compute genus zero $`A`$-model correlators on GLSMs corresponding to Calabi-Yau surfaces. In addition, eqn.(1) manifestly satisfies the quantum cohomology relations, and, as we will see in more detail below, it is a useful probe for the physics of the Coulomb branch of the GLSM.
Since we have not coupled the TFT to topological gravity, our result is of limited use for the computation of general Gromov-Witten invariants—factorization of the correlators implies that our higher genus results do not generate “new” invariants for $`g>0`$. However, we believe that for the purposes of enumerative geometry, eqn.(1) is a neat packaging of the requisite combinatorics. Essentially, while it is true that the $`g>0`$ correlators may be obtained by factorization from the $`g=0`$ correlators, if one interested in explicit numbers, eqn.(1) may eliminate much algebraic suffering.
It is no accident that the form we find is reminiscent of the correlators in topological Landau-Ginzburg models studied by Vafa . In fact, our analysis is a simple extension of those techniques to include the zero modes of the matter fields. These additional zero modes are the source of the factor of $`_i\xi _i`$ in our expression. We will discuss this further below.
While eqn.(1) computes the correlators in compact toric GLSMs, in non-compact toric GLSMs the above expression is only a part of the story. In general, only a subset of the solutions to $`d\stackrel{~}{W}_{\text{eff}}=0`$ correspond to Coulomb vacua, and eqn.(1) provides the correct measure for the contributions to the correlators due to these vacua. In addition, other, non-Coulomb vacua may also contribute, and these contributions may invalidate the quantum cohomology relations.
The rest of this note is organized as follows. In section 2 we briefly review the toric GLSM and the corresponding $`A`$-model, and we take care to distinguish the GLSM phases according to the properties of the Coulomb vacua. In section 3 we prove eqn.(1) by studying the $`A`$-model localization onto the Coulomb vacua. We provide some applications of the result to compact toric GLSMs in section 4. In section 5 we turn to a study of a non-compact example, and we conclude with a discussion in section 6.
## 2 A GLSM Overview
### 2.1 Some Superspace Details
The GLSM is a $`d=2`$ abelian gauge theory with $`𝒩=(2,2)`$ supersymmetry. The field content is neatly summarized in terms of $`𝒩=(2,2)`$ multiplets. The matter fields belong to chiral multiplets $`\mathrm{\Phi }^i=(\varphi ^i,\psi _\pm ^i,\overline{\psi }_\pm ^i,F^i)`$, with $`\varphi ^i`$ a complex scalar, $`\psi _\pm ^i`$ left/right-moving Weyl fermions, $`F^i`$ a complex auxiliary field, and $`i=1,\mathrm{},n`$. These fields are charged under the gauge group $`G=\left[\mathrm{U}(1)\right]^r`$ with integral charges $`Q_i^a`$, $`a=1,\mathrm{},r`$. The gauge fields reside in real vector supermultiplets $`V_a`$, and the gauge-invariant field-strengths are to be found in twisted chiral multiplets $`\mathrm{\Sigma }_a=(\sigma _a,\lambda _{\pm ,a},\overline{\lambda }_{\pm ,a},D_aif_{01,a})`$, where $`\sigma _a`$ is a complex scalar, $`\lambda _{\pm ,a}`$ are left/right-moving Weyl fermions, $`D_a`$ is a real auxiliary field, and $`f_{01,a}`$ is the abelian gauge field-strength.
We define the GLSM at a scale $`\mu `$ by a Lagrange density $`\mathrm{L}^\mu `$ given by a sum of two terms, the Kähler term $`\mathrm{L}_K^\mu `$ and the twisted superpotential term $`\mathrm{L}_{\stackrel{~}{W}}^\mu `$. We take the Kähler term to be
$$\mathrm{L}_K^\mu =d^4\theta \left(\frac{1}{4}\underset{i=1}{\overset{n}{}}\overline{\mathrm{\Phi }}^i\mathrm{exp}\left(2\underset{a=1}{\overset{r}{}}Q_i^aV_a\right)\mathrm{\Phi }^i+\frac{1}{4\mu ^2g(\mu )^2}\underset{a=1}{\overset{r}{}}\overline{\mathrm{\Sigma }}_a\mathrm{\Sigma }_a\right),$$
(2)
where $`g(\mu )`$ is the dimensionless coupling of the gauge theory. The tree-level twisted superpotential is given by
$$\mathrm{L}_{\stackrel{~}{W}}^\mu =\left[\frac{i}{2\sqrt{2}}𝑑\theta ^+𝑑\overline{\theta }^{}\underset{a=1}{\overset{r}{}}\mathrm{\Sigma }_a\tau ^a(\mu )\right]+\text{c.c.}.$$
(3)
The $`\tau ^a(\mu )=ir^a(\mu )+\frac{\theta ^a}{2\pi }`$ are the parameters of the model. Each $`\tau ^a`$ is a combination of a Fayet-Iliopoulos (F-I) term $`r^a`$ and a $`\theta `$-angle $`\theta ^a`$. It is useful to define single-valued parameters $`q_a=e^{2\pi i\tau ^a}.`$ For generic values of these parameters the moduli space of classical vacua of the GLSM so defined is a toric variety. The GLSM may be generalized by including a superpotential $`W(\mathrm{\Phi })`$ which serves to restrict the moduli space to a hypersurface or a complete intersection in the ambient toric variety. We will restrict attention to toric GLSMs, those with $`W(\mathrm{\Phi })=0`$.
### 2.2 Basic GLSM Properties
Let us begin with a brief review of the Higgs vacua of the GLSM.<sup>1</sup><sup>1</sup>1There are no photons and no Higgs mechanism in two dimensions, and a description based on Higgs vacua is only valid at weak coupling. Fortunately, this is just where we will use it, and so we will ignore this subtlety in what follows. This material is well known, and we refer the reader to for further details. The classical moduli space of a toric GLSM is obtained by solving the $`D`$-terms modulo the gauge group as functions of the F-I parameters $`r^a`$. One finds that there exists a cone $`𝒦_c^r`$ where the space of solutions to the $`D`$-terms is non-empty.<sup>2</sup><sup>2</sup>2The D-term equations are $`D^a=_iQ_i^a|\varphi ^i|^2r^a,`$ whence it follows that $`𝒦_c`$ is indeed a cone, the space positively generated by the $`n`$ vectors $`𝐐_i^r`$. For generic values of the $`r^a𝒦_c`$, the gauge group is completely broken, the $`\sigma _a`$ are massive, and the moduli space is a toric variety of complex dimension $`d=nr`$. The geometric properties of this toric variety vary smoothly with the $`r^a`$ away from co-dimension one sub-cones of $`\overline{𝒦_c}`$, where the gauge group is un-Higgsed and some or all of the $`\sigma _a`$ become massless. These boundaries subdivide $`𝒦_c`$ into a set of cones $`𝒦_V`$, indexed by a set of birationally equivalent toric varieties. There is a natural association between the varieties and the cones: for $`r^a𝒦_V`$ the moduli space of the GLSM’s classical vacua is the variety $`V`$. The cones $`𝒦_V`$ are termed phases of the GLSM. By choosing the F-I parameters deep in the interior of any such phase, we obtain a weakly coupled theory whose low energy theory is that of a NLSM with target-space $`V`$. For reasons that will become clear below, we will also refer to the complement of $`𝒦_c`$ as a phase.
As the F-I terms are tuned to approach a lower dimensional face of $`𝒦_V`$, the low energy description seems to break down as $`V`$ becomes singular, or equivalently, there appear new massless degrees of freedom corresponding to an un-Higgsed subgroup of the gauge group. Quantum effects lift the corresponding singularities when the un-Higgsed gauge group satisfies $`_iQ_i0`$, and even for gauge groups with $`_iQ_i=0`$, the singularities are lifted for generic values of the corresponding $`\theta `$ angle. Thus, all phases are smoothly connected, and the low-energy NLSM description is smooth away from a complex co-dimension one subvariety in the space of the $`q_a`$—the singular locus. Of course, from the point of view of the GLSM there is no real singularity on the singular locus. However, we do expect that the theory is strongly coupled on the singular locus, and strong coupling effects may invalidate results based on the weakly coupled description.
In addition to the Higgs vacua, the GLSM possesses Coulomb vacua. These are obtained when some of the $`\sigma _a`$ acquire non-zero expectation values and give masses to some or all of the matter fields. Integrating out these massive $`\mathrm{\Phi }^i`$ multiplets leads to an effective interaction for the $`\mathrm{\Sigma }_a`$ fields, which can be expressed in terms of an effective twisted superpotential $`\stackrel{~}{W}_{\text{eff}}(\mathrm{\Sigma })`$ . The solutions to $`d\stackrel{~}{W}_{\text{eff}}(\sigma )=0`$ are continuous if $`_iQ_i^a=0`$ for all $`\sigma _a`$ with non-zero expectation values, and they are discrete otherwise. The former exist only on the singular locus of the model, while the latter vary smoothly with the parameters. When all of the matter fields are massive, $`\stackrel{~}{W}_{\text{eff}}(\mathrm{\Sigma })`$ is given by<sup>3</sup><sup>3</sup>3We have left off a conventional over-all factor of $`\frac{1}{4\pi \sqrt{2}}`$. As far as our results are concerned, this factor can be absorbed in the definition of the string coupling.
$$\stackrel{~}{W}_{\text{eff}}=\underset{a=1}{\overset{r}{}}\mathrm{\Sigma }_a\mathrm{log}\left[\underset{i=1}{\overset{n}{}}\left(\frac{1}{\mathrm{exp}(1)\mu }\underset{b=1}{\overset{r}{}}Q_i^b\mathrm{\Sigma }_b\right)^{Q_i^a}/q_a\right].$$
(4)
The vacua corresponding to $`d\stackrel{~}{W}_{\text{eff}}=0`$ will occupy us for most of this note. For future reference, we give the equations of motion which follow from $`d\stackrel{~}{W}_{\text{eff}}=0`$:
$$\underset{i|Q_i^a>0}{}\left(\frac{\xi _i}{\mu }\right)^{Q_i^a}=q_a\underset{i|Q_i^a<0}{}\left(\frac{\xi _i}{\mu }\right)^{Q_i^a},a=1,\mathrm{},r,$$
(5)
where $`\xi _i=_aQ_i^a\sigma _a`$. We will also have use for
$$\mathrm{H}^{\mathrm{a}\mathrm{b}}:=\frac{^2\stackrel{~}{\mathrm{W}}_{\text{eff}}}{\sigma _\mathrm{a}\sigma _\mathrm{b}}=\underset{\mathrm{i}}{}\frac{\mathrm{Q}_\mathrm{i}^\mathrm{a}\mathrm{Q}_\mathrm{i}^\mathrm{b}}{\xi _\mathrm{i}}.$$
(6)
Of course, the Hessian of $`\stackrel{~}{W}_{\text{eff}}`$ is given by $`H=det\mathrm{H}`$.
The Coulomb vacua are derived by integrating out massive matter fields, and thus are only reliable in the regions of the parameter space where these fields are indeed massive. In principle, this may depend on the renormalization of the Kähler terms, but at least in the weak coupling regimes (i.e deep in the interior of some $`𝒦_V`$) one may discern which Coulomb vacua are reliable. At weak coupling, the $`\varphi `$ mass term has the contribution $`2_{i,a,b}|\varphi ^i|^2Q_i^aQ_i^b\sigma _a\overline{\sigma }_b`$, so that a critical point of $`\stackrel{~}{W}_{\text{eff}}`$ is not reliable if all the $`\sigma _a`$ are small.
When do these discrete Coulomb vacua arise? By examining the equations of motion in eqn.(5) it is clear that these are homogeneous in the $`\sigma _a`$ whenever $`_iQ_i^a=0`$ for all $`a`$, leading to a continuous set of solutions for the $`\sigma _a`$. We have not shown it, but it seems likely that if $`\mathrm{rank}(Q)=r`$ then this homogeneity is the only way to obtain a continuum of solutions. So, we expect that $`\stackrel{~}{W}_{\text{eff}}`$ describes discrete Coulomb vacua whenever $`_iQ_i^a0`$ for some $`a`$. It is always possible to choose a basis for the action of the gauge group so that $`Q_i^a`$ satisfy $`_iQ_i^a=0`$ for $`a>1`$. We will work in this basis, taking $`\mathrm{\Delta }=_iQ_i^1`$. The condition for $`\stackrel{~}{W}_{\text{eff}}`$ to describe discrete Coulomb vacua is then just $`\mathrm{\Delta }0`$. When $`\mathrm{\Delta }=0`$, the continuous solutions to $`d\stackrel{~}{W}_{\text{eff}}=0`$ emerge on the principal component of the singular locus .
The $`\stackrel{~}{W}_{\text{eff}}`$ above describes the vacua where all of the $`\mathrm{\Phi }^i`$ are massive. Of course, there may also be Coulomb-Higgs vacua, where the gauge group is partially Higgsed. Just as the Coulomb vacua described by $`\stackrel{~}{W}_{\text{eff}}`$, these may be labelled according to the space of $`\sigma `$ vevs as either continuous or discrete. The former are found on various non-principal components of the singular locus of the model, while the latter, like their discrete Coulomb cousins, may be found in various phases. We will not study the discrete Coulomb-Higgs vacua in this note. However, when analyzing a particular phase one should be careful to check that the results are not invalidated by the presence of these vacua.
We will find it useful to characterize the phases of the GLSM by the types of vacua found at weak coupling. Whenever a phase does not have any discrete Coulomb-Higgs vacua, we will refer to it as:
* a Geometric Phase if its weak coupling limit has no reliable Coulomb vacua, and the vacua are purely Higgs;
* a Non-Geometric Phase if the situation is reversed and there are no Higgs vacua;
* a Mixed Phase if both Coulomb and Higgs vacua are present at weak coupling.
These distinctions are important. For example, in a model with a Non-Geometric Phase, eqn.(1) yields the correlators at any genus, while in a model with a Geometric Phase, we may be able to compute the genus zero correlators by gauge instanton sums. As we will see below, in a Mixed Phase the correlators may be obtained by simply adding the Higgs and Coulomb contributions. Note that a Non-Geometric Phase may only exist outside of $`𝒦_c`$, so it is only if $`𝒦_c\simeq ̸^r`$, that our strongest results hold. Happily, this holds for compact toric GLSMs.
## 3 $`A`$-Model Localization on the Discrete Coulomb Branch
### 3.1 Twisting and Localization: Generalities
The topological twisting of an $`𝒩=(2,2)`$ theory may be accomplished by shifting the spin connection on the world-sheet by the (ultraviolet) $`R`$-symmetry of the model. In effect, this produces a new theory by modifying the spins of the fields. Let us now point out some basic consequences of this twisting in the context of the GLSM. The twisted theory possesses a world-sheet anti-commuting scalar operator $`𝒬`$ which can be used to project the theory onto the $`𝒬`$-cohomology. In the GLSM, this leads to topological observables parametrized by powers of the $`\sigma _a`$. Another consequence of the twisting is that the path integral localizes onto the $`𝒬`$-invariant configurations—the SUSY vacua of the untwisted theory. This property of localization plays a crucial role in the study of TFT. Roughly, this is the statement that the path integral will localize onto the vacua of the theory, and, under certain conditions, the contribution of a particular vacuum may be computed semi-classically. This still leaves a difficult problem, especially when quantum vacua are involved. Of course, this is the case when we wish to study the correlators in a Mixed or Non-Geometric Phases of the GLSM. Fortunately, in that case the vacua are controlled by $`\stackrel{~}{W}_{\text{eff}}(\mathrm{\Sigma })`$, a quantity determined by holomorphy and ’t Hooft anomaly matching.
A further simplification makes the TFT computations tractable: since the singular locus is a complex co-dimension one variety in the space of the $`q_a`$, we expect that we should be able to compute correlators at weak coupling in any phase, and then unambiguously obtain the result for generic $`q_a`$ by analytic continuation. Indeed, this has been demonstrated for the Geometric Phases in . As we will see below, the result also holds in more general situations involving the Coulomb vacua.
### 3.2 $`A`$-Twist Details
On a Euclidean signature world-sheet $`\mathrm{\Sigma }_g`$ the spin connection may be thought of as a $`\mathrm{U}(1)`$ connection, and twisting amounts to shifting the Lorentz $`\mathrm{U}(1)`$ charges of the fields by a combination of the $`R`$-charges. In the toric GLSM the classical $`\mathrm{U}(1)_+\times \mathrm{U}(1)_{}`$ $`R`$-symmetry group leaves the superfields $`\mathrm{\Phi }^i`$, $`V_a`$ invariant while acting on the $`\theta ^\pm `$ with charges $`Q_\pm (\theta ^\pm )=+1`$, $`Q_{}(\theta ^\pm )=0`$. All other charges are determined by this choice, and, in particular, $`Q_\pm (\sigma _a)=\pm 1`$. When $`\mathrm{\Delta }0`$, this classical $`R`$-symmetry suffers from an anomaly in the presence of gauge instantons. That means that only the vector combination may be used to obtain a consistent twisted theory. This is the $`A`$-model, obtained by twisting with $`Q_V=\frac{1}{2}(Q_++Q_{})`$. If we designate the Lorentz charges of the fields by $`Q_L`$, the new Lorentz charges are $`Q_L^{}=Q_LQ_V`$. Applying this to the fields of the GLSM we find that the $`\varphi ^i`$, $`\sigma _a`$ remain world-sheet scalars, but the spins of the fermions are shifted. The $`\psi _+,\overline{\psi }_{},\lambda _{},\overline{\lambda }_+`$ become world-sheet one-forms, while $`\psi _{},\overline{\psi }_+,\lambda _+,\overline{\lambda }_{}`$ become world-sheet scalars :
$$\begin{array}{cccccc}\psi _+& & \psi _z& \lambda _+& & \eta \\ \psi _{}& & \chi & \lambda _{}& & \rho _{\overline{z}}\\ \overline{\psi }_+& & \overline{\chi }& \overline{\lambda }_+& & \overline{\rho }_z\\ \overline{\psi }_{}& & \overline{\psi }_{\overline{z}}& \overline{\lambda }_{}& & \overline{\eta }.\end{array}$$
(7)
### 3.3 Localization
We are interested in computing expectation values of the form $`F(\sigma )_g`$ in the twisted theory. As described above, to perform these computations we deform the model to weak coupling, where we have a reasonable handle on the vacua of the theory. Since the TFT path integral localizes onto the vacua, we may compute the correlators by summing contributions from the vacua. Let us suppose we are working in a phase of the GLSM without discrete Coulomb-Higgs vacua. At weak coupling, the path integral then receives contributions from the Higgs vacua—the gauge instantons, and the discrete Coulomb vacua:
$$Z=Z^{\text{Higgs}}+Z^{\text{Coulomb}}.$$
(8)
At $`g=0`$, the contribution from the gauge instantons may be determined by the methods of Morrison and Plesser , and we will now show how to compute $`Z^{\text{Coulomb}}`$ at any genus.
Since we wish to compute $`F(\sigma )_g`$, we may first perform the integration over the $`\mathrm{\Phi }^i`$ multiplets. As we have argued above, in the untwisted theory this leads to a factor of $`\mathrm{exp}\left(d^2z\mathrm{L}_{\stackrel{~}{W}_{\text{eff}}}\right)`$ in the path integral over the $`\mathrm{\Sigma }_a`$ multiplets. When we perform this integration in the twisted theory, we find a similar result, but we must be careful of one subtlety: the zero modes of the $`\mathrm{\Phi }^i`$ multiplets. In the untwisted theory the $`\psi _\pm ,\overline{\psi }_\pm `$ had no zero modes, and, aside from the factor of $`\mathrm{exp}\left(d^2z\mathrm{L}_{\stackrel{~}{W}_{\text{eff}}}\right)`$ and a deformation of the (irrelevant) Kähler term of the $`\mathrm{\Sigma }_a`$, the one-loop fermion determinants cancelled the bosonic ones as a consequence of the $`𝒩=(2,2)`$ SUSY. In the twisted theory, while the non-zero modes of the $`\mathrm{\Phi }^i`$ multiplets continue to be paired up just as in the untwisted theory, the zero modes of the fermions no longer pair up with the $`\varphi ^i`$ zero modes. To deal with this subtlety, we will separate out the integral over the $`\mathrm{\Phi }^i`$ zero modes. Of course, integrating out the non-zero modes will still lead to the factor of $`\mathrm{exp}\left(d^2z\mathrm{L}_{\stackrel{~}{W}_{\text{eff}}}\right)`$.
We are now in a position to apply the standard localization arguments to the Coulomb vacua. The corresponding $`𝒬`$-invariant configurations are a subset of the field configurations
$$\varphi ^i=0,f_a=0,_z\sigma _a=_{\overline{z}}\sigma _a=0,d\stackrel{~}{W}_{\text{eff}}=0.$$
(9)
A solution to $`d\stackrel{~}{W}_{\text{eff}}(\sigma )=0`$ does not necessarily correspond to a Coulomb vacuum, and at generic $`q_a`$ it is difficult to determine which of the solutions to $`d\stackrel{~}{W}_{\text{eff}}(\sigma )=0`$ are reliable. However, at weak coupling, i.e. deep in the interior of a phase, we may answer this question unambiguously: a $`\sigma `$ vacuum is trustworthy only if in the weak coupling limit the corresponding $`|\sigma _a|`$ grow in such a way that all of the $`\mathrm{\Phi }^i`$ multiplets may be consistently integrated out. We will label these reliable vacua by $`\widehat{\sigma }`$. The contribution to the path integral is then
$$Z^{\text{Coulomb}}=\underset{\widehat{\sigma }}{}Z(\widehat{\sigma }),$$
(10)
and we can compute $`Z(\widehat{\sigma })`$ at weak coupling by expanding in fluctuations about the $`\widehat{\sigma }`$ vacuum. Integration over the massive modes of the $`\mathrm{\Sigma }`$ multiplets leads to determinants that exactly cancel between the bosons and fermions (this familiar fact may be traced back to the primordial $`𝒩=(2,2)`$ SUSY), and we are left with an integral over the zero-modes, which factorizes into an integral over the $`\mathrm{\Phi }`$ fluctuations and an integral over $`\mathrm{\Sigma }`$ fluctuations:
$$Z(\widehat{\sigma })=[D\mathrm{\Phi }][D\mathrm{\Sigma }]\mathrm{exp}(S_\mathrm{\Phi }S_\mathrm{\Sigma })=[D\mathrm{\Phi }]\mathrm{exp}(S_\mathrm{\Phi })[D\mathrm{\Sigma }]\mathrm{exp}(S_\mathrm{\Sigma })=Z_\mathrm{\Phi }(\widehat{\sigma })Z_\mathrm{\Sigma }(\widehat{\sigma }).$$
(11)
The terms in the action are given by
$`S_\mathrm{\Phi }`$ $`=`$ $`{\displaystyle \underset{i}{}}\left\{V_\mathrm{\Sigma }\left[2|\xi _i|^2|\varphi ^i|^2+\sqrt{2}\chi ^i\overline{\xi }_i\overline{\chi }^i\right]+\sqrt{2}{\displaystyle _{\mathrm{\Sigma }_g}}\psi _z^i\xi _i\overline{\psi }_{\overline{z}}^i\right\},`$
$`S_\mathrm{\Sigma }`$ $`=`$ $`{\displaystyle \underset{a,b}{}}\left\{V_\mathrm{\Sigma }\left[4\mu ^2\overline{\sigma }_a\left(\mathrm{H}^{}\mathrm{H}\right)^{ab}\sigma _b+\overline{\eta }_a2\overline{\mathrm{H}}^{ab}\eta _b\right]+{\displaystyle _{\mathrm{\Sigma }_g}}\overline{\rho }_{z,a}2\mathrm{H}^{\mathrm{a}\mathrm{b}}\rho _{\overline{\mathrm{z}},\mathrm{b}}\right\},`$ (12)
where $`V_\mathrm{\Sigma }`$ is the volume of $`\mathrm{\Sigma }_g`$, and $`\xi _i`$ and $`\mathrm{H}^{\mathrm{a}\mathrm{b}}`$—the latter defined in eqn.(6)— are to be evaluated at $`\sigma =\widehat{\sigma }`$. We can now evaluate the Gaussian integrals. For $`Z_\mathrm{\Phi }`$ we find
$$Z_\mathrm{\Phi }=Z_\varphi Z_\chi Z_\psi ,$$
(13)
$$Z_\varphi ^1=V_\mathrm{\Sigma }^n\underset{i}{}\left|\xi _i\right|^2,Z_\chi =V_\mathrm{\Sigma }^n\underset{i}{}\overline{\xi }_i,Z_\psi =\underset{i}{}\xi _i^g.$$
(14)
We have counted the $`1`$ zero mode of $`\chi `$ (there is just one constant function on a compact $`\mathrm{\Sigma }_g`$) and the $`g`$ zero modes of the $`\psi _z^i`$ (there are $`g`$ holomorphic one-forms on $`\mathrm{\Sigma }_g`$). Similarly,
$$Z_\mathrm{\Sigma }=Z_\sigma Z_\eta Z_\rho ,$$
(15)
with
$$Z_\sigma ^1=V_\mathrm{\Sigma }^r\left|H\right|^2,Z_\eta =V_\mathrm{\Sigma }^r\overline{H},Z_\rho =H^g,$$
(16)
where $`H=det\mathrm{H}`$. Putting it all together, we find the measure in eqn.(1).
The careful reader will note that we have ignored any subtleties associated to gauge invariance. This simplification follows because the computations are performed on the Coulomb branch, where the condition $`f_a=0`$ and the Riemann-Roch theorem ensure that our computation of $`Z_\mathrm{\Phi }`$ is correct. We have also neglected various constants which may be absorbed into an overall normalization of the correlators or a re-definition of the string coupling constant. Furthermore, the sign of the fermion integration measure has been chosen to match results from the Higgs Phase computations at genus zero.
In what follows, we will work in units of the scale $`\mu `$. This scale plays an important role in the untwisted theory, but upon twisting it becomes superfluous, essentially because the TFT is a theory at zero energy. The scale can be important if one wants to make connections with the untwisted theory, in which case it is easy to restore in our formulas.
## 4 A Few Applications to Compact Toric GLSMs
As discussed above, we expect that whenever $`𝒦_c\simeq ̸^r`$, eqn.(1) directly gives the correlators at arbitrary genus. This makes it useful for elucidating various properties of the compact toric GLSM $`A`$-model correlators, as well as actual computations. We illustrate this in this section.
### 4.1 Some Properties of the Correlators
We can easily demonstrate some important properties of these GLSM $`A`$-model correlators from the explicit form. Perhaps the simplest observation is that the result presents the correlators as a sum over all the solutions to a system of polynomial equations with finitely many common zeroes. This finite sum has a natural expansion in terms of symmetric functions, and, thus, it is clear that the correlators are meromorphic functions of the $`q_a`$. This is not obvious from the form of the instanton sum in a Geometric Phase. Another equally simple but important observation is that the quantum cohomology relations, which are just the equations of motion in eqn.(5) considered as operator relations, obviously hold. This should also be compared with the Geometric phase computation, where this is a non-trivial combinatorics result. Below, we give a few more technical observations.
#### 4.1.1 TFT Factorization
The $`A`$-twist reduces the Hilbert space of the GLSM to a vector space of dimension $`N_v`$, where $`N_v`$ is the number of discrete Coulomb vacua. The operators $`\sigma _a`$ are now simply $`N_v\times N_v`$ matrices, and correlators are obtained by taking a matrix trace:
$$F(\sigma )_g=\mathrm{Tr}\left[F(\sigma )\left(\stackrel{~}{H}(\sigma )\right)^{g1}\right].$$
(17)
These correlators are easily shown to satisfy the factorization axioms of topological field theory. These axioms state that if we choose a complete basis of states $`|i`$, and the corresponding operators $`\varphi _i`$ have the metric $`\eta _{ij}=\varphi _i\varphi _j_0,`$ with inverse $`\eta ^{ij}`$, then
1. if $`F(\sigma )=f_1(\sigma )f_2(\sigma )`$, then
$$F(\sigma )_g=\underset{ij}{}f_1(\sigma )\varphi _i_g^{}\eta ^{ij}\varphi _jf_2(\sigma )_{gg^{}},$$
(18)
and
2. for any $`F(\sigma )`$
$$F(\sigma )_g=\underset{ij}{}\eta ^{ij}\varphi _i\varphi _jF(\sigma )_{g1}.$$
(19)
These properties are apparent in a basis of states corresponding to the $`N_v`$ $`\sigma `$-vacua. The state operator correspondence is $`|i\varphi _i=\delta _{\sigma ,\sigma _i}`$, where $`\sigma _i`$ is the value of $`\sigma `$ in the $`i`$-th vacuum. In this basis the operator $`\sigma `$ is diagonal, and $`\eta ^{ij}=\stackrel{~}{H}(\sigma _i)\delta ^{ij}`$. Factorization follows immediately. Thus, as expected, any genus correlator may be obtained from the $`g=0`$ results.
#### 4.1.2 The Ghost Number Selection Rule
These $`A`$-model correlators obey a simple selection rule. Working in the basis where
$$\underset{i}{}Q_i^1=\mathrm{\Delta }\text{and}\underset{i}{}Q_i^a=0\text{for}a>1,$$
(20)
we may write the $`\sigma `$ equations of motion in the form $`\sigma _a=\omega _a\sigma _1`$ for $`a>2`$, where $`\omega _a`$ are now determined by solving $`r1`$ polynomial equations, and $`\sigma _1`$ satisfies $`\sigma _1^\mathrm{\Delta }=q_1s(\omega )`$ for some $`s(\omega )`$. Thus, the sum over the vacua includes a sum over the $`\mathrm{\Delta }`$-th roots of unity. Since $`\stackrel{~}{H}`$ has degree $`d=nr`$, if $`F(\sigma )`$ has degree $`s`$, then $`F(\sigma )_g`$ is non-zero only if $`s+d(g1)=m\mathrm{\Delta }`$ for some integer $`m`$, in which case $`F(\sigma )_gq_1^m.`$ This selection rule is just the ghost number selection rule familiar from TFTs in general and GLSMs in particular.
#### 4.1.3 The All Genus Correlation Function
Although factorization makes this exercise purely one of convenience, we can easily sum over the genera to obtain
$$F(\sigma )=\underset{g0}{}\lambda ^{2g2}F(\sigma )_g=\mathrm{Tr}\left[\frac{F(\sigma )}{\lambda ^2\stackrel{~}{H}}\frac{1}{1\lambda ^2\stackrel{~}{H}}\right].$$
(21)
From the selection rule above it follows that if $`F(\sigma )`$ has degree $`s`$ then
$$F(\sigma )=q_1^{s/\mathrm{\Delta }}f\left(\left(q_1^d\lambda ^{2\mathrm{\Delta }}\right)^{1/gcd(d,\mathrm{\Delta })}\right).$$
(22)
#### 4.1.4 The Quantum Restriction Formula
Given a Calabi-Yau hypersurface in a toric variety $`V`$, there exists a simple method for obtaining the “toric” subset of the $`g=0`$ $`A`$-model correlators for the Calabi-Yau model. These correlators can be computed by the quantum restriction formula of , which expresses a hypersurface toric $`A`$-model correlator, denoted by $`F(\sigma )`$, to a sum over the $`A`$-model correlators for $`V`$:
$$F(\sigma )_{g=0}=F(\sigma )\frac{K}{1K}_{g=0},$$
(23)
where $`K`$ is the operator corresponding to the anti-canonical divisor on $`V`$, given by $`K=_i\xi _i`$. Using our form of the correlators on $`V`$, it follows that
$$F(\sigma )_{g=0}=\mathrm{Tr}\left[\frac{F(\sigma )}{\stackrel{~}{H}}\frac{K}{1K}\right].$$
(24)
### 4.2 Two Examples
In this section we will apply our simple result and the observations above to two examples. These models are not difficult to solve, but they illustrate some techniques and ideas that should be useful even in much more intricate examples.
#### 4.2.1 $`A`$-model correlators for $`^4`$.
Let us start with the canonical GLSM example: $`^4`$. This is a one parameter model with $`Q=(1,1,1,1,1).`$ The equation of motion from eqn.(5) is just $`\sigma ^5=q,`$ and $`\stackrel{~}{H}=5\sigma ^4,`$ yielding
$$\sigma ^a_g=5^{g1}\underset{\sigma |\sigma ^5=q}{}\sigma ^{a+4(g1)}.$$
(25)
The correlators satisfy the selection rule discussed earlier:
$$\sigma ^a_g=0\text{unless}a+4(g1)=5n\text{for some integer}n,$$
in which case
$$\sigma ^{5n+4(1g)}_g=5^gq^n.$$
(26)
The all-genus correlation function is given by
$$\sigma ^a=\underset{g0}{}\lambda ^{2g2}\sigma ^a_g.$$
(27)
Evaluating this for $`0a4`$, we find
$`\sigma ^a`$ $`=`$ $`{\displaystyle \frac{5(5q\lambda ^2)^a}{15^5q^4\lambda ^{10}}},a=0,\mathrm{},3,`$
$`\sigma ^4`$ $`=`$ $`{\displaystyle \frac{\lambda ^2}{15^5q^4\lambda ^{10}}}.`$ (28)
The intriguing pole at $`q^4\lambda ^{10}=5^5`$ agrees with the findings of . The interpretation of this pole is far from clear. While we might expect such a pole in a topological string theory, where it could be a manifestation of non-perturbative effects in $`\lambda `$, we have not coupled the model to $`d=2`$ gravity, and thus any string-based interpretation does not seem appropriate.
Finally, we can use the quantum restriction formula to compute the unique $`A`$-model correlator on the quintic in $`^4`$. The anti-canonical divisor corresponds to $`K=5\sigma `$, and we find
$$\sigma ^3_{g=0}=\mathrm{Tr}\frac{\sigma ^3}{5\sigma ^4}\frac{5\sigma }{1+5\sigma }=\mathrm{Tr}\frac{1}{1+(5\sigma )^5}=\frac{5}{1+5^5q}.$$
(29)
#### 4.2.2 A Two Parameter Example
This is another example that has been studied in detail in . This GLSM corresponds to the toric variety obtained by resolving the curve of $`_2`$ singularities in the weighted projective space $`_{1,1,2,2,2}^4`$. The GLSM has $`n=6`$, $`r=2`$ and charges
$$Q=\left(\begin{array}{cccccc}0& 0& 1& 1& 1& 1\\ 1& 1& 0& 0& 0& 2\end{array}\right).$$
(30)
Obviously, $`\mathrm{\Delta }=4`$, and $`\stackrel{~}{H}=8\sigma _1^3\sigma _2`$. Letting $`\sigma _2=\omega \sigma _1`$, the equations of motion $`d\stackrel{~}{W}_{\text{eff}}=0`$ may be written as
$$\sigma _1^4=\frac{q_1}{12\omega },$$
(31)
and
$$P(\omega )=\omega ^2q_2(12\omega )^2=0.$$
(32)
The selection rule implies that $`\sigma _1^a\sigma _2^b_g`$ is zero unless $`a+b=4(m+1)`$, and if $`m0`$, we have
$$\sigma _1^{4(m+1)b}\sigma _2^b=\underset{g0}{}\mathrm{Tr}\left[\sigma _1^{4(m+g)}\omega ^{b+g1}(8\lambda ^2)^{g1}\right],$$
(33)
which we can reduce to a trace on the roots of $`P(\omega )`$, denoted by $`\mathrm{Tr}^{}`$:
$$\sigma _1^{4(m+1)b}\sigma _2^b=\frac{q_1^m}{2\lambda ^2(1(8\lambda ^2q_1)^2q_2)}\mathrm{Tr}^{}\left[\frac{\omega ^{b1}(1+(8\lambda ^2q_12)\omega )}{(12\omega )^{m+1}}\right].$$
(34)
Again, we observe the interesting $`\lambda `$-dependent pole.
At genus zero the above expression simplifies to
$$\sigma _1^{4(m+1)b}\sigma _2^b_{g=0}=\frac{q_1^m}{2}\mathrm{Tr}^{}\frac{\omega ^{b1}}{(12\omega )^m}.$$
(35)
We can again use the quantum restriction formula to compute correlators on the anti-canonical hypersurface. We have $`K=4\sigma _1`$, and
$$\sigma _1^{3j}\sigma _2^j_{g=0}=4\frac{\sigma ^{4j}\sigma _2^j}{1+(4\sigma _1)^4}_{g=0}=2\mathrm{Tr}^{}\frac{\omega ^{j1}(12\omega )}{14^4q_12\omega }.$$
(36)
In this and other two-parameter models it is convenient to rewrite the $`\mathrm{Tr}^{}`$ as a contour integral:
$$\mathrm{Tr}^{}f(\omega )=\underset{\widehat{\omega }|P(\widehat{\omega })=0}{}_{C(\widehat{\omega })}\frac{d\omega }{2\pi i}\frac{f(\omega )P^{}(\omega )}{P(\omega )},$$
(37)
where $`C(\widehat{\omega })`$ is a small contour about $`\omega =\widehat{\omega }`$. This form makes it easy to evaluate the traces. In the case of more than two parameters, more sophisticated residue techniques may be applied . Applying this to the case at hand,
$$\sigma _1^{3j}\sigma _2^j_{g=0}=\underset{\widehat{\omega }|P(\widehat{\omega })=0}{}_{C(\widehat{\omega })}\frac{d\omega }{2\pi i}\frac{4\omega ^j}{(14^4q_12\omega )P(\omega )}.$$
(38)
Pulling the contour off the roots of $`P(\omega )`$, the correlators are written as
$$\sigma _1^{3j}\sigma _2^j_{g=0}=2\frac{\omega ^j}{P(\omega )}|_{\omega =\frac{14^4q_1}{2}}+2\mathrm{Res}\left\{\frac{\omega ^j}{(\omega \frac{14^4q_1}{2})P(\omega )}\right\}|_{\omega =\mathrm{}}.$$
(39)
Straightforward algebra yields
$`\sigma _1^3_{g=0}`$ $`=`$ $`{\displaystyle \frac{8}{D}},`$
$`\sigma _1^2\sigma _2_{g=0}`$ $`=`$ $`{\displaystyle \frac{4(12^8q_1)}{D}},`$
$`\sigma _1\sigma _2^2_{g=0}`$ $`=`$ $`{\displaystyle \frac{8q_2(2^9q_11)}{(14q_2)D}},`$
$`\sigma _2^3_{g=0}`$ $`=`$ $`{\displaystyle \frac{4q_2(1+4q_22^8q_13072q_1q_2)}{(14q_2)^2D}},`$ (40)
and $`D=(12^8q_1)^22^{18}q_1^2q_2`$. This reproduces the results of eqn.(4.28) of .<sup>4</sup><sup>4</sup>4Our expression corrects a sign error in the last correlator in eqn.(4.28) of .
## 5 A Non-Compact Example
Having examined the properties of models with a Non-Geometric Phase, we now turn to models where Higgs vacua are present in every phase. Sadly, this means that with the current technology we will need to restrict attention to genus zero correlators, but, nevertheless, we will be able to uncover some surprises.
We will work with the example studied at length in . This GLSM has $`n=5`$, $`r=2`$ and charges
$$Q=\left(\begin{array}{ccccc}1& 1& 1& N& 1\\ 0& 0& 1& 1& 2\end{array}\right).$$
(41)
There are four classical phases, each corresponding to a triangulation of a toric fan. The fan without any subdivisions is an orbifold, $`^3/_{(2N+1)(2,2,1)}`$, the partially subdivided fans correspond to partial resolutions, and the completely subdivided fan is the smooth phase. These phases are depicted in Fig. (1). We will assume $`N>2`$, and we have labelled the phases according to the Geometric-Mixed terminology defined above. This model has a continuous Coulomb branch which emerges for small $`|q_1|`$ and $`q_2=1/4`$, and, naively, one would expect that some observables in the TFT will be sensitive to this singularity.
The $`A`$-model correlators of interest are the $`Y_{a,b}=\sigma _1^a\sigma _2^b_{g=0}`$. The ghost number selection rule requires that for a non-zero correlator $`a+b=3+(2N)n`$. In , we were able to compute the $`Y_{3+(2N)nb,b}`$ correlators for $`n<0`$ by summing the instantons in the Geometric Phase, and we found that these correlators could be put into a “Coulomb” form:<sup>5</sup><sup>5</sup>5We call this the “Coulomb” form because it is the answer that one would get by naive application of eqn.(1).
$$Y_{3+(2N)nb,b}^{\text{Geometric}}=q_1^n\mathrm{Tr}^{}\frac{\omega ^bs(\omega )^n}{3N+1+2N\omega },$$
(42)
where $`\mathrm{Tr}^{}`$ is to be taken over the roots of
$$P(\omega )=(1+\omega )(N+\omega )q_2(1+2\omega )^2,$$
(43)
and $`s(\omega )`$ is given by
$$s(\omega )=\frac{(12\omega )(N+\omega )^N}{1+\omega }.$$
(44)
For later convenience, we will re-write these as a contour integral:
$$Y_{3+(2N)nb,b}=\underset{\widehat{\omega }=\omega _+,\omega _{}}{}_{C(\widehat{\omega })}\frac{d\omega }{2\pi i}\frac{\omega ^bs(\omega )^n}{(1+2\omega )P(\omega )},$$
(45)
where $`C(\widehat{\omega })`$ is a small contour about $`\omega =\widehat{\omega }`$, and $`\omega _\pm `$ are the roots of $`P(\omega )`$.
The computation of the $`n=0`$ correlators is complicated by the non-compactness of the orbifold, and in we circumvented that problem by using the quantum cohomology relations to determine the $`Y_{3b,b}`$. As we saw above, these relations are powerful, and it is easy to show that if the $`Y_{3b,b}`$ are determined from the $`Y_{3+(2N)nb,b}`$ by the relations, they must be of given by eqn.(42) with $`n=0`$. Upon computing the trace, one finds that the $`Y_{3b,b}`$ so determined are sensitive to the $`q_2=1/4`$ singularity. As we will show below, our basic assumption was incorrect. The quantum cohomology relations simply do not hold!
Even without further computations, there are several reasons to suspect the validity of this result. First, the $`Y_{3b,b}`$ are independent of $`q_1`$ and thus, if they are sensitive to the $`q_2=1/4`$ singularity at small $`|q_1|`$, they are equally singular at $`q_2=1/4`$ for arbitrary $`q_1`$. However, we are hard pressed to explain the singularity at large $`|q_1|`$ and $`q_2=1/4`$. After all, this is deep in the weakly coupled regime of the Geometric Phase, where a classical analysis is reliable and does not reveal any singularities. In addition, we know that the gauge instantons are labelled by sets of integers $`𝐧\left(^d\right)^{}`$, and each such instanton contributes a term $`Y^𝐧_aq_a^{n_a}`$ to the correlator. In a particular phase $`𝒦_V`$, the instanton numbers corresponding to non-zero $`Y^𝐧`$ must lie in the dual cone defined by
$$𝒦_V^{}=\left\{𝐧\left(^d\right)^{}\right|𝐧,𝐫0𝐫𝒦_V\}.$$
(46)
It is easy to see that in the Geometric Phase the only instantons that can contribute to $`Y_{3b,b}`$ have $`𝐧=0`$. Hence, one would expect $`Y_{3b,b}`$ to be constants, and any $`q_2`$ dependence, let alone a singular one, is strange indeed. This would seem to indicate that the quantum cohomology relations are violated whenever the $`Y_{3b,b}`$ correlators are involved.
To explore this further, let us now work out the correlators in one of the Mixed Phases. We will choose the phase $`A`$, but the computation may be easily repeated for other phases. First, let us compute the contribution from the Coulomb vacua. For this model, $`\stackrel{~}{H}`$ is given by
$$\stackrel{~}{H}=(N2)\sigma _1^2((3N+1)\sigma _1+2N\sigma _2),$$
(47)
and the equations of motion that follow from $`d\stackrel{~}{W}_{\text{eff}}=0`$ are
$`\left(\sigma _1+\sigma _2\right)\left(\sigma _2N\sigma _1\right)`$ $`=`$ $`q_2\left(\sigma _1+2\sigma _2\right)^2,`$
$`\sigma _1^2\left(\sigma _1+\sigma _2\right)`$ $`=`$ $`q_1\left(\sigma _2N\sigma _1\right)^N\left(\sigma _1+2\sigma _2\right).`$ (48)
As in the previous section, we may parametrize the solutions by $`\sigma _1`$ and the ratio $`\omega =\sigma _2/\sigma _1`$:
$`P(\omega _\pm )`$ $`=`$ $`0`$
$`\sigma _{1,\pm ;p}`$ $`=`$ $`\zeta ^p\left(q_1s(\omega _\pm )\right)^{\frac{1}{2N}},`$
$`\sigma _{2,\pm ;p}`$ $`=`$ $`\omega _\pm \sigma _{1,\pm ;p},`$ (49)
where $`\zeta =e^{\frac{2\pi i}{N2}}`$, $`p=0,\mathrm{},N1`$, and $`P(\omega )`$ and $`s(\omega )`$ are as in eqns.(43,44).
Let us consider the solutions $`(\sigma _{1,\pm ;p},\sigma _{2,\pm ;p})`$ in the weak coupling limit. Weak coupling in the Mixed Phase A corresponds to
$$|q_2|^N|q_1||q_2|,$$
(50)
and in particular, $`|q_2|0`$. In this limit the $`\sigma `$-vacua have a simple structure:
$`\omega _+`$ $``$ $`1{\displaystyle \frac{1}{N+1}}q_2,`$
$`\omega _{}`$ $``$ $`N+{\displaystyle \frac{(2N+1)^2}{N+1}}q_2,`$ (51)
and hence
$`\sigma _{1;+}^{2N}`$ $``$ $`q_1q_2^1(N1)^{N+1},`$
$`\sigma _{2;}^{2N}`$ $``$ $`q_1q_2^N{\displaystyle \frac{(2N+1)^{2N+1}}{(N+1)^{N+1}}}.`$ (52)
Thus, we see that in the weak coupling limit of the Mixed phase the $`N2`$$``$” critical points of $`\stackrel{~}{W}_{\text{eff}}`$ have growing $`\sigma `$ vevs, while the $`N2`$$`+`$” critical points have decreasing $`\sigma `$ vevs. Thus, only the “$``$” solutions correspond to actual Coulomb vacua, and their contribution, $`Y_{3+(2N)nb,b}^{\text{Coulomb}}`$, is given by the $`\omega _{}`$ contribution in eqn.(45):
$$Y_{3+(2N)nb,b}^{\text{Coulomb}}=q_1^n_{C(\omega _{})}\frac{\omega ^bs(\omega )^n}{(1+2\omega )P(\omega )}.$$
(53)
Next, we consider the Higgs contribution. Unlike the Coulomb computation, which gives the same form regardless of whether $`n=0`$ or $`n<0`$, here this distinction makes a crucial difference. First, let us consider the situation when $`n<0`$. Using the standard toric techniques of Morrison and Plesser, we can perform the instanton sum and evaluate the Higgs branch contribution to the correlators. Performing the requisite toric intersection computations, we reduce the correlators to a single sum:
$$Y_{3+(2N)nb,b}^{\text{Higgs}}=q_1^n\underset{m=n}{\overset{\mathrm{}}{}}_{C(1)}\frac{d\omega }{2\pi i}\frac{\omega ^bs(\omega )^nR^m}{(1+\omega )(N+\omega )(12\omega )},$$
(54)
where
$$R=q_2\frac{(1+2\omega )^2}{(1+\omega )(N+\omega )},$$
(55)
and $`C(1)`$ is a small contour about $`\omega =1`$: $`\omega =1+ϵe^{i\theta }`$. For uniform convergence we must have
$$|q_2|<\frac{ϵ(N+1ϵ)}{1+2ϵ}.$$
(56)
Provided that this condition holds, we can exchange the integral and the sum to obtain
$$Y_{3+(2N)nb,b}^{\text{Higgs}}=q_1^n_{C(1)}\frac{\omega ^bs(\omega )^nR^n}{(1+2\omega )P(\omega )}.$$
(57)
And now comes a pleasant surprise: the condition for convergence ensures that $`\omega =\omega _+`$ is enclosed by $`C(1)`$, while $`\omega =\omega _{}`$ remains outside of it, and so, since $`R(\omega _\pm )=1`$,
$$Y_{3+(2N)nb,b}^{\text{Higgs}}=q_1^n_{C(\omega _+)}\frac{d\omega }{2\pi i}\frac{\omega ^bs(\omega )^n}{(1+2\omega )P(\omega )},$$
(58)
and for $`n<0`$ we precisely have the desired form for the correlators:
$$Y_{a,b}^{\text{Geometric}}=Y_{a,b}^{\text{Higgs}}+Y_{a,b}^{\text{Coulomb}}.$$
(59)
The contribution to the $`n=0`$ correlators is even more remarkable. The standard manipulation of the instanton sum yields
$$Y_{3b,b}^{\text{Higgs}}=Y_{3b,b}^0+\underset{m=0}{\overset{\mathrm{}}{}}_{C(N)}\frac{d\omega }{2\pi i}\frac{\omega ^bR^m}{(1+\omega )(N+\omega )(12\omega )},$$
(60)
with $`R`$ as above, and $`C(N)`$ a small contour about $`\omega =N`$: $`\omega =N+ϵe^{i\theta }`$. The constants $`Y_{3b,b}^0`$ parametrize our ignorance of how to compute intersection numbers on a non-compact variety. Presumably, these are computed by an appropriate cohomolgy theory. For uniform convergence we must have
$$|q_2|<\frac{ϵ(1ϵ)}{1+2N+ϵ}.$$
(61)
Carrying out the sum, we have
$$Y_{3b,b}^{\text{Higgs}}=Y_{3b,b}^0_{C(N)}\frac{d\omega }{2\pi i}\frac{\omega ^b}{(1+2\omega )P(\omega )}.$$
(62)
The condition for convergence ensures that, this time, $`\omega _{}`$ is enclosed by $`C(N)`$, while $`\omega _+`$ remains outside. The crucial overall minus sign means that putting this together with the Coulomb contribution, we find that the $`Y_{3b,b}`$ are just constants, as predicted by our discussion above.
A similar analysis may be carried out in the other Mixed Phases. In those phases all of the Coulomb vacua are reliable and contribute. For $`n<0`$ there are no contributions from the gauge instantons of the Higgs branch, while for $`n=0`$ the instanton sums cancel the Coulomb contribution up to the constants $`Y_{3b,b}^0`$. It would be interesting to examine these constants in more detail, but whatever they are, the resulting correlators are incompatible with quantum cohomology relations.
One final aspect of this example deserves mention—the disappearance of the semi-classical singularity at $`q_2=1/4.`$ We know that semi-classical analysis of the Higgs branch in Mixed Phase B or C shows this singularity. We expect that analysis to be valid for small $`|q_1|`$. Of course, the discrete Coulomb vacua also exist in this limit, and it is possible that the presence of these additional Coulomb vacua washes out the singularity. It appears that this is so, at least in the topological theory.
## 6 Discussion
We have found a simple algebraic formula for the Coulomb contribution to the $`A`$-model correlators in toric GLSMs. We hope to have convinced the reader that this expression is conceptually satisfying and computationally useful. We will now conclude with an outlook on some interesting questions that remain.
### 6.1 Some Observations on the Coulomb Vacua and the GLSM
Our work is a simple application of the general principle of localization in TFTs. It has been known for a long time that in the Geometric Phases the path integral localizes onto the gauge instantons. We have merely extended this result to the Mixed and Non-Geometric Phases. In models with a Non-Geometric Phase our result gives a surprisingly complete form for the $`A`$-model correlators. In models without such a phase we are restricted to genus zero by our inability to compute the contribution from the Higgs vacua.
Perhaps the most surprising finding of our work is that for models without a Non-Geometric Phase the quantum cohomology relations may fail for a subset of the correlators which, in some Mixed Phase, receive contributions from both Higgs and Coulomb vacua. There are two simple “proofs” of quantum cohomology relations: the first is a Geometric Phase analysis of the intersection numbers on the gauge instanton moduli spaces, and the second is a Non-Geometric Phase analysis given above. The first argument is subtle when the phase corresponds to a non-compact variety, and the second does not apply in the absence of a Non-Geometric Phase. The example of the last section illustrates that these problems are manifestations of the same failure of the quantum cohomology relations in different phases.
We have not addressed computations of the correlators in phases where discrete Higgs-Coulomb vacua are present. It would be interesting to understand the contributions from these vacua. This exercise may well provide some new insights into the more mysterious aspects of GLSM physics, and it may provide us with another set of phases where the computation of the correlators is made tractable.
Finally, we have derived our results in the context of traditional GLSMs. It would be interesting to extend our treatment to the recently much-discussed GLSMs corresponding to supermanifolds.
### 6.2 A Pure Landau-Ginzburg Description?
The contribution to the $`A`$-model correlators bears a striking resemblance to Vafa’s result on correlators in topological Landau-Ginzburg models . Vafa studied a topological Landau-Ginzburg model with a superpotential $`W(X)`$, and he showed that the correlators are given by
$$F(x)=\mathrm{Tr}F(x)H^{g1},$$
(63)
where the trace is taken over the critical points of $`W`$, and each contribution is weighted by the Hessian of $`W`$, $`H`$, evaluated at that point.
While it is certainly not true that the GLSM model correlators are computed by a Landau-Ginzburg theory with $`\mathrm{\Sigma }_a`$ as the fields and $`\stackrel{~}{W}_{\text{eff}}`$ as the superpotential, because of the remarkable similarity between eqn.(1) and eqn.(63), it is natural to wonder if some Landau-Ginzburg theory computes the same correlators. It is easy to see that this is possible for at least some models. For example, a Landau-Ginzburg theory of a single field $`\mathrm{\Sigma }`$ and superpotential $`W=\frac{1}{6}\mathrm{\Sigma }^6q\mathrm{\Sigma }`$ has the equation of motion $`\sigma ^5=q`$ and Hessian of $`5\sigma ^4`$. Hence, according to eqn.(63), this model will compute the same correlators as the GLSM with target space $`^4`$.
It is clear that, in general, the requisite Landau-Ginzburg model will be much more complicated. After all, the theories studied by Vafa can be constructed as relevant deformations of a free theory with an ultraviolet $`R`$-symmetry with charges $`Q_+(X)=Q_{}(X)=0`$, and all of the solutions to the classical equations of motion $`dW(x)=0`$ are on the same footing. It is tempting to suggest that the method of Hori and Vafa, which relies on dualizing the matter fields, yields this Landau-Ginzburg description. It would be interesting to check whether this is the case.
### 6.3 Coupling to Topological Gravity
Another direction to pursue is to couple the model to topological gravity in the spirit of . The resulting theory would be an interesting topological string theory, where perhaps the $`\lambda `$-dependent poles we have found would find a natural interpretation. Furthermore, this model would, in principle, compute a much larger subset of non-trivial Gromov-Witten invariants. Hopefully, the simplicity of our result for the correlators would persist to some extent in the string theory.
## Acknowledgments
It is a pleasure to thank C. Haase , E. Katz, G. Moore, and S. Rinke for useful comments and conversations. M.R.P. would like to thank the Perimeter Institute and the Theoretical High Energy Physics group at Rutgers University for hospitality while some of this work was completed. This article is based upon work supported in part by the National Science Foundation under Grants DMS-0074072 and DMS-0301476. Any opinions, findings, and conclusions or recommendations expressed in this article are those of the authors and do not necessarily reflect the views of the National Science Foundation. |
warning/0507/math-ph0507028.html | ar5iv | text | # MICZ-Kepler Problems in All Dimensions
## 1 Introduction
The Kepler problem is the physics problem about two bodies which attract each other by a force proportional to the inverse square of the distance. By solving this problem in classical mechanics, Newton gave a satisfactory explanation for Kepler’s laws for the planetary motion. The Kepler problem plays a significant role in the development of quantum mechanics, too; in fact, the solution of this problem in the Schrödinger’s wave mechanics firmly puts the Schrödinger equation right at the center of quantum mechanics.
After more than three centuries, the Kepler problem still plays an important role in mathematics and physics. There has been a continuous interest in this problem; in particular, in the last three decades we have witnessed an explosion of its interactions with quantum mechanics, celestial mechanics and mathematics. For a recent comprehensive treatment of the Kepler problem, the interested readers may consult Ref. .
The MICZ-Kepler problems are natural cousins of the Kepler problem, and they were independently discovered by McIntosh-Cisneros and Zwanziger more than thirty years ago. Roughly speaking, a MICZ-Kepler problem is the Kepler problem in the case when the nucleus of a hypothetic hydrogen atom also carries a magnetic charge. These generalized problems share the following characteristic beauty with the Kepler problem: the existence of the Runge-Lenz vector and the dynamical $`\text{Spin}(4)`$ symmetry for the bound states; therefore, they provide a rich family of examples for the exploration of extra hidden dynamic symmetry.
The hamiltonian of a MICZ-Kepler problem is constructed from that of the Kepler problem by adding the vector potential of a Dirac monopole and a repulsive centrifugal potential; explicitly, we have
$`H={\displaystyle \frac{1}{2m}}(\stackrel{}{p}+e\stackrel{}{A})^2+{\displaystyle \frac{\mu ^2}{2mr^2}}{\displaystyle \frac{e^2}{r}}`$ (1)
where $`\stackrel{}{p}`$ is the canonical momentum of the electron, $`\stackrel{}{A}`$ is the vector potential of a Dirac monopole, $`r`$ is the distance from the electron to the hydrogen nucleus, $`m`$ is the (reduced) mass of the electron, $`e`$ is the fundamental unit of the electric charge, and $`\mu `$ is the magnetic charge of the Dirac monopole measured in unit $`\frac{c}{e}`$, i.e., $`\mu \frac{c}{e}`$ is the magnetic charge of the Dirac monopole, here $`c`$ is the speed of light in vacuum<sup>1</sup><sup>1</sup>1The Dirac quantization condition becomes $`\frac{\mu }{\mathrm{}}=`$ a half integer..
Quantum mechanically, via rescaling
$`r{\displaystyle \frac{\mathrm{}^2}{me^2}}r,\mu \mathrm{}\mu ,`$
we arrive at the following hamiltonian operator:
$`\widehat{H}={\displaystyle \frac{me^4}{\mathrm{}^2}}\left({\displaystyle \frac{1}{2}}\mathrm{\Delta }_A+{\displaystyle \frac{\mu ^2}{2r^2}}{\displaystyle \frac{1}{r}}\right):={\displaystyle \frac{me^4}{\mathrm{}^2}}\widehat{h}`$
where<sup>2</sup><sup>2</sup>2Remark that $`\widehat{h}`$ is dimensionless and is expressed in terms of dimensionless quantities. It is $`\widehat{h}`$ (not $`\widehat{H}`$) that will be generalized later.
$`\widehat{h}={\displaystyle \frac{1}{2}}\mathrm{\Delta }_A+{\displaystyle \frac{\mu ^2}{2r^2}}{\displaystyle \frac{1}{r}}.`$ (2)
Here $`\mathrm{\Delta }_A`$ is the Laplace operator twisted by the gauge potential $`A`$ of a Dirac monopole, and $`\mu =0`$, $`\pm \frac{1}{2}`$, $`\pm 1`$, $`\mathrm{}`$ is the magnetic charge of the Dirac monopole measured in terms of the fundamental unit<sup>3</sup><sup>3</sup>3The case $`\mu =0`$ corresponds to the Kepler problem.. Locally, with a gauge chosen, we have
$`\mathrm{\Delta }_A=_a_a`$
where the repeated index $`a`$ is summed up. Here $`_a`$ is the $`a`$-th covariant partial derivative and is written as $`_a+iA_a`$ by physicists with $`\stackrel{}{A}=(A_1,A_2,A_3)`$ being the gauge potential of the Dirac magnetic monopole. Mathematically $`_a=_a+\omega _a`$ where $`\omega =\omega _adx_a`$ has been previously identified with the Levi-Civita spin connection form of the cylindrical metric
$$ds^2=\frac{1}{r^2}(dx_1^2+dx_2^2+dx_3^2)$$
on the punctured $`3`$-space, see Ref. for the details.
The MICZ-Kepler problems exist in higher dimensions just as the Kepler problem does, and that is a main observation here. In fact, the existence in dimension five has been previously observed ; however, the existence in all dimensions greater than two, though very straightforward from a canonical geometric point of view, was probably not expected by the community. This overlook is very likely due to a general belief in the literature: the existence of Dirac monopoles and its five dimensional analogue (the Yang monopoles ) has to do with the existence of the division algebras or Hopf bundles.
In section 2, we construct the MICZ-Kepler problems in all dimensions and then state the main results. The construction is geometric and canonical. A key ingredient in the construction is the higher dimensional generalization of the Dirac monopoles — a canonical geometric object that has been used in Ref. . In section 3, we first introduce the explicit formulas for the gauge potential of these generalized Dirac monopoles; then we list and prove some crucial identities necessary for the exhibition of the extra large hidden dynamical symmetry. In section 4, we introduce the angular momentum and Rung-Lenz vector for our MICZ-Kepler problems, and derive the symmetry algebra. In section 5, we obtain the energy spectrum and the energy eigenspaces for bound states by using Painlevé analysis plus representation theory, and then show that the Hilbert space of bound states has a hidden dynamical $`\text{Spin}(D+1)`$-symmetry for a $`D`$-dimensional MICZ-Kepler problem even though in general the Runge-Lenz vector fails to be conserved when $`D`$ is even.
Remark that the MICZ-Kepler problems in higher dimensions constructed here are based on modern geometry, but they are solved by classical analytic method with the help of the representation theory for Lie groups. The solution of these new MICZ-Kepler problems can in principle be solved by the modern geometric quantization approach pioneered by Simms , Mladenov and Tsanov , but that will be reserved for the future for the following reasons: 1) the primary objective of this paper is to inform the experts in the fields that the MICZ-Kepler problems do exist in higher dimensions, 2) the classical analytic approach is more elementary and easier to understand, 3) the modern geometric quantization approach is a bit more involved and deserves an independent research.
### Acknowledgment
I would like to thank SiXia Yu for a conversation on the Kepler problem. This work is supported by the Hong Kong Research Grants Council under the RGC project no. 602504. I would also like to thank the referee for his or her careful reading of the manuscript and for his or her valuable suggestions.
## 2 The main results
From the physics point of view, a MICZ Kepler problem is obtained from the Kepler problem by adding a suitable background magnetic field, while at the same time making a suitable adjustment to the scalar Coulomb potential so that the problem is still integrable. The background magnetic field is just the spin connection<sup>4</sup><sup>4</sup>4For readers without sufficient background in modern geometry, just take our explicit formulas for gauge potential in Eq. (5) for granted. of the cylindrical metric on the configuration space that we have mentioned in the introduction. The configuration space is the punctured Euclidean space. With this in mind, we are now ready to give the detailed presentation of our generalized MICZ Kepler problems.
Let $`D3`$ be an integer, $`_{}^D`$ be the punctured $`D`$-space, i.e., $`^D`$ with the origin removed. Let $`ds^2`$ be the cylindrical metric on $`_{}^D`$. Then $`(_{}^D,ds^2)`$ is the product of the straight line $``$ with the round sphere $`\mathrm{S}^{D1}`$. When $`D`$ is odd, we let $`𝒮_\pm `$ be the positive/negative spinor bundle of $`(_{}^D,ds^2)`$, and when $`D`$ is even, we let $`𝒮`$ be the spinor bundle of $`(_{}^D,ds^2)`$. Note that, these bundles correspond to the fundamental spin representations $`𝐬_\pm `$ of $`\mathrm{so}(even)`$ and $`𝐬`$ of $`\mathrm{so}(odd)`$ respectively.
The above spinor bundles come with a natural $`\mathrm{SO}(D)`$ invariant connection — the Levi-Civita spin connection of $`(_{}^D,ds^2)`$. As a result, the Young product of $`I`$ copies of these bundles, denoted by $`𝒮_+^I`$, $`𝒮_{}^I`$ (when $`D`$ is odd) and $`𝒮^I`$ (when $`D`$ is even) respectively, come with a natural connection, too.
For the sake of notational sanity, from here on, when $`D`$ is odd and $`\mu `$ is a half integer, we rewrite $`𝒮_+^{2\mu }`$ as $`𝒮^{2\mu }`$ if $`\mu 0`$ and rewrite $`𝒮_{}^{2\mu }`$ as $`𝒮^{2\mu }`$ if $`\mu 0`$; moreover, we adopt this convention for $`\mu =0`$: $`𝒮^0`$ is the product complex line bundle with the product connection. When $`D`$ is odd, $`𝒮^{2\mu }`$ is the product complex line bundle with the product connection in the case $`\mu =0`$, and is the fundamental spinor bundle $`𝒮`$ in the case $`\mu =1/2`$.
Note that $`𝒮^{2\mu }`$ is our analogue of the Dirac monopole with magnetic charge $`\mu `$, and its corresponding representation of $`\mathrm{so}(D1)`$ will be denoted by $`𝐬^{2\mu }`$. We are now ready to present our definitions.
###### Definition 1.
Let $`n1`$ be an integer, $`\mu `$ a half integer. The $`(2n+1)`$-dimensional MICZ-Kepler problem with magnetic charge $`\mu `$ is defined to be the quantum mechanical system on $`_{}^{2n+1}`$ for which the wave-functions are sections of $`𝒮^{2\mu }`$, and the hamiltonian is
$`\widehat{h}={\displaystyle \frac{1}{2}}\mathrm{\Delta }_\mu +{\displaystyle \frac{(n1)|\mu |+\mu ^2}{2r^2}}{\displaystyle \frac{1}{r}}`$ (3)
where $`\mathrm{\Delta }_\mu `$ is the Laplace operator twisted by $`𝒮^{2\mu }`$.
###### Definition 2.
Let $`n>1`$ be an integer, $`\mu =0`$ or $`1/2`$. The $`2n`$-dimensional MICZ-Kepler problem with magnetic charge $`\mu `$ is defined to be the quantum mechanical system on $`_{}^{2n}`$ for which the wave-functions are sections of $`𝒮^{2\mu }`$, and the hamiltonian is
$`\widehat{h}={\displaystyle \frac{1}{2}}\mathrm{\Delta }_\mu +{\displaystyle \frac{(n1)\mu }{2r^2}}{\displaystyle \frac{1}{r}}`$ (4)
where $`\mathrm{\Delta }_\mu `$ is the Laplace operator twisted by $`𝒮^{2\mu }`$.
Note that we require $`\mu =0`$ or $`1/2`$ in the even dimensional case. There is both an analytic and an algebraic reason for this requirement, which shall be pointed out in appropriate places. Remark also that, upon a choice of a local gauge, the background magnetic potential $`A_\alpha `$ can be explicitly written down, then $`\mathrm{\Delta }_\mu =_\alpha (_\alpha +iA_\alpha )^2`$ can be explicitly written down, too. We are now ready to state our main results.
###### Theorem 1.
Let $`n1`$ be an integer and $`\mu `$ be a half integer. For the $`(2n+1)`$-dimensional MICZ-Kepler problem with magnetic charge $`\mu `$, the following statements are true:
1) The negative energy spectrum is
$$E_I=\frac{1/2}{(I+n+|\mu |)^2}$$
where $`I=0`$, $`1`$, $`2`$, …;
2) The Hilbert space $``$ of negative-energy states admits a linear $`\mathrm{Spin}(2n+2)`$-action under which there is a decomposition
$$=\widehat{}_{I=0}^{\mathrm{}}_I$$
where $`_I`$ is a model for the irreducible $`\mathrm{Spin}(2n+2)`$-representation with highest weight $`(I+|\mu |,|\mu |,\mathrm{},|\mu |,\mu )`$;
3) $`\mathrm{Spin}(2n+1,1)`$ acts linearly on the positive-energy states and $`\mathrm{Spin}(2n+1)^{2n+1}`$ acts linearly on the zero-energy states;
4) The linear action in either part 2) or part 3) extends the manifest linear action of $`\mathrm{Spin}(2n+1)`$, and $`_I`$ in part 2) is the energy eigenspace with eigenvalue $`E_I`$ in part 1).
###### Theorem 2.
Let $`n>1`$ be an integer and $`\mu =0`$ or $`1/2`$. For the $`2n`$-dimensional MICZ-Kepler problem with magnetic charge $`\mu `$, the following statements are true:
1) The negative energy spectrum is
$$E_I=\frac{1/2}{(I+n+\mu \frac{1}{2})^2}$$
where $`I=0`$, $`1`$, $`2`$, …;
2) The Hilbert space $``$ of negative-energy states admits a linear $`\mathrm{Spin}(2n+1)`$-action under which there is a decomposition
$$=\widehat{}_{I=0}^{\mathrm{}}_I$$
where $`_I`$ is a model for the irreducible $`\mathrm{Spin}(2n+1)`$-representation with highest weight $`(I+\mu ,\mu ,\mathrm{},\mu )`$;
3) $`\mathrm{Spin}(2n,1)`$ acts linearly on the positive-energy states and $`\mathrm{Spin}(2n)^{2n}`$ acts linearly on the zero-energy states;
4) The linear action in part 2) extends the manifest linear action of $`\mathrm{Spin}(2n)`$, and $`_I`$ in part 2) is the energy eigenspace with eigenvalue $`E_I`$ in part 1).
Remark that, based on the analysis done in later sections, we know that bound eigen-states are always the ones with negative energy eigenvalues.
## 3 Generalized Dirac monopoles
We write $`\stackrel{}{r}=(x_1,x_2,\mathrm{},x_{D1},x_0)`$ for a point in $`^D`$ and $`r`$ for the length of $`\stackrel{}{r}`$. The small Greek letters $`\mu `$, $`\nu `$, etc run from $`0`$ to $`D1`$ and the small Latin letters $`a`$, $`b`$ etc run from $`1`$ to $`D1`$. We use the Einstein convention: the repeated index is always summed up.
To do computations, we just need to choose a gauge on $`^D`$ minus the negative $`0`$-th axis and then write down the gauge potential explicitly. We have done that before in Eq. (10) of Ref. . Note that, if we use the rectangular coordinates $`\stackrel{}{r}=(\stackrel{}{x},x_0)`$, then the gauge potential $`A=A_\mu dx_\mu `$ from Eq. (10) of Ref. can be written as<sup>5</sup><sup>5</sup>5In Ref. we only consider the case $`D`$ is odd — the topological nontrivial case, but the basic construction there is valid in any dimension, see appendix A in Ref. .
$`A_0=0,A_b={\displaystyle \frac{1}{r(r+x_0)}}x_a\gamma _{ab}`$ (5)
where $`\gamma _{ab}=\frac{i}{4}[\gamma _a,\gamma _b]`$ with $`\gamma _a`$ being the “gamma matrix” for physicists. Note that $`\gamma _a=ie_a`$ with $`\stackrel{}{e}_a`$ being the element in the Clifford algebra that corresponds to the $`a`$-th standard coordinate vector of $`^{D1}`$.
It is straightforward to calculate the gauge field strength $`F_{\mu \nu }=_\mu A_\nu _\nu A_\mu +i[A_\mu ,A_\nu ]`$ and get
$`F_{0b}`$ $`=`$ $`{\displaystyle \frac{1}{r^3}}x_a\gamma _{ab}`$ (6)
$`F_{ab}`$ $`=`$ $`{\displaystyle \frac{2\gamma _{ab}}{r(r+x_0)}}+{\displaystyle \frac{1}{r^2(r+x_0)^2}}`$ (8)
$`\left((2+{\displaystyle \frac{x_0}{r}})x_c(x_a\gamma _{cb}x_b\gamma _{ca})+ix_dx_c[\gamma _{da},\gamma _{cb}]\right)`$
The following lemma is crucially used when we check the dynamical symmetry of our models.
###### Lemma 1.
For the gauge potential defined in Eq. (5), we have
1) Let $`_\alpha =_\alpha +iA_\alpha `$, then the following identities are valid in any representation:
$`F_{\mu \nu }F^{\mu \nu }={\displaystyle \frac{2}{r^4}}c_2\text{where }c_2=c_2[\mathrm{so}(D1)]=\frac{1}{2}\gamma _{ab}\gamma _{ab}`$ (9)
$`[_\kappa ,F_{\mu \nu }]={\displaystyle \frac{1}{r^2}}\left(x_\mu F_{\nu \kappa }+x_\nu F_{\kappa \mu }2x_\kappa F_{\mu \nu }\right)`$ (10)
$`x_\mu A_\mu =0,x_\mu F_{\mu \nu }=0,[_\mu ,F_{\mu \nu }]=0`$ (11)
$`r^2[F_{\mu \nu },F_{\alpha \beta }]+iF_{\mu \beta }\delta _{\alpha \nu }iF_{\nu \beta }\delta _{\alpha \mu }+iF_{\alpha \mu }\delta _{\beta \nu }iF_{\alpha \nu }\delta _{\beta \mu }`$ (12)
$`={\displaystyle \frac{i}{r^2}}\left(x_\mu x_\alpha F_{\beta \nu }+x_\mu x_\beta F_{\nu \alpha }x_\nu x_\alpha F_{\beta \mu }x_\nu x_\beta F_{\mu \alpha }\right)`$ (13)
2) When $`D=2n+1`$, identity
$`r^2F_{\lambda \alpha }F_{\lambda \beta }={\displaystyle \frac{c_2}{n}}\left({\displaystyle \frac{1}{r^2}}\delta _{\alpha \beta }{\displaystyle \frac{x_\alpha x_\beta }{r^4}}\right)+i(n1)F_{\alpha \beta }`$ (14)
holds in the irreducible representation $`𝐬^{2\mu }`$ of $`\mathrm{so}(2n)`$ whose highest weight is of the form $`(|\mu |,\mathrm{},|\mu |,\mu )`$.
3) When $`D=2n`$, identity
$`r^2F_{\lambda \alpha }F_{\lambda \beta }={\displaystyle \frac{n1}{2}}\left({\displaystyle \frac{1}{r^2}}\delta _{\alpha \beta }{\displaystyle \frac{x_\alpha x_\beta }{r^4}}\right)+i(n{\displaystyle \frac{3}{2}})F_{\alpha \beta }`$ (15)
holds in the fundamental spin representation $`𝐬`$ of $`\mathrm{so}(2n1)`$.
One can show that the Eq. (15) is valid only when $`\mu =0`$ or $`1/2`$. That is the algebraic reason for requiring $`\mu =0`$ or $`1/2`$.
### 3.1 Proof of Lemma 1
The verification of these identities is just a direct and lengthy calculation. However, if we exploit the symmetry, we just need to check the identities at point $`\stackrel{}{r}_0=(0,\mathrm{},0,r)`$, a much easier task. For example, since
$`A_\mu =0,F_{0a}=0,F_{ab}={\displaystyle \frac{1}{r^2}}\gamma _{ab}`$ (16)
at $`\stackrel{}{r}_0`$, identity (9) is obvious.
Proof of part 1). We have just remarked that identity (9) is obvious. Also,
$$x_\mu F_{\mu \nu }|_{\stackrel{}{r}_0}=x_0F_{0\nu }|_{\stackrel{}{r}_0}=0.$$
It is also easy to see that $`x_\mu A_\mu |_{\stackrel{}{r}_0}=0`$ and
$$[_\mu ,F_{\mu \nu }]|_{\stackrel{}{r}_0}=_\mu F_{\mu \nu }|_{\stackrel{}{r}_0}=0.$$
Therefore, identity (10) is checked.
To check identity (9), first we assume $`\mu =0`$, $`\nu =b`$, then we need to check that
$$_\kappa F_{0b}=\frac{1}{r}F_{b\kappa }$$
at $`\stackrel{}{r}_0`$, and that can be easily seen to be true whether $`\kappa =0`$ or $`a`$. Next we assume that $`\mu =a`$ and $`\nu =b`$, then we need to check that
$$_\kappa F_{ab}=\frac{2}{r^2}x_\kappa F_{ab}$$
at $`\stackrel{}{r}_0`$, and that can be easily verified, too.
We divide the checking of identity (11) at $`\stackrel{}{r}_0`$ into two cases: 1) one of indices is zero, easy to check; 2) none of the indices is zero, then the identity becomes
$`[\gamma _{ab},\gamma _{cd}]=i\gamma _{ad}\delta _{bc}+i\gamma _{bd}\delta _{ac}i\gamma _{ca}\delta _{bd}+i\gamma _{cb}\delta _{ad}`$ (17)
which is of course true because $`\gamma _{ab}`$’s are the generators of $`\mathrm{so}(D1)`$.
Proof of part 2). Write
$$r^2F_{\lambda \alpha }F_{\lambda \beta }=\frac{r^2}{2}\{F_{\lambda \alpha },F_{\lambda \beta }\}+\frac{r^2}{2}[F_{\lambda \alpha },F_{\lambda \beta }].$$
Using identity (12), we have
$$r^2F_{\lambda \alpha }F_{\lambda \beta }=\frac{r^2}{2}\{F_{\lambda \alpha },F_{\lambda \beta }\}+i(n1)F_{\alpha \beta }.$$
Therefore, by checking at $`\stackrel{}{r}_0`$, we just need to verify that identity
$`{\displaystyle \underset{k}{}}\{\gamma _{ki},\gamma _{kj}\}={\displaystyle \frac{\delta _{ij}}{n}}{\displaystyle \underset{a,b}{}}(\gamma _{ab})^2`$ (18)
holds in the irreducible representation $`𝐬^{2\mu }`$ of $`\mathrm{so}(2n)`$ whose highest weight is of the form $`(|\mu |,\mathrm{},|\mu |,\mu )`$. Since this checking is a bit involved, we do it in the appendix.
Proof of part 3). Write
$$r^2F_{\lambda \alpha }F_{\lambda \beta }=\frac{r^2}{2}\{F_{\lambda \alpha },F_{\lambda \beta }\}+\frac{r^2}{2}[F_{\lambda \alpha },F_{\lambda \beta }].$$
Using identity (12), we have
$$r^2F_{\lambda \alpha }F_{\lambda \beta }=\frac{r^2}{2}\{F_{\lambda \alpha },F_{\lambda \beta }\}+i(n\frac{3}{2})F_{\alpha \beta }.$$
Therefore, by checking at $`\stackrel{}{r}_0`$, we just need to verify that identity
$`{\displaystyle \underset{k}{}}\{\gamma _{ki},\gamma _{kj}\}=(n1)\delta _{ij}`$ (19)
holds in the spin representation $`𝐬`$ of $`\mathrm{so}(2n1)`$, but this is easy to check by using the Clifford algebra.
## 4 The hidden dynamical symmetry
To exhibit the dynamical symmetry for our MICZ-Kepler problems, as usual, we introduce the angular momentum tensor
$`\widehat{L}_{\alpha \beta }=i(x_\alpha _\beta x_\beta _\alpha )+r^2F_{\alpha \beta }`$ (20)
and the Runge–Lenz vector
$`\widehat{L}_\beta ={\displaystyle \frac{i}{2}}\left(_\alpha \widehat{L}_{\alpha \beta }+\widehat{L}_{\alpha \beta }_\alpha \right)+{\displaystyle \frac{x_\beta }{r}}`$ (21)
With the help of the identities stated in lemma 1, a lengthy calculation yields the following commutation relations:
$`\overline{)\begin{array}{ccc}[\widehat{L}_{\mu \nu },\widehat{h}]\hfill & =& 0\hfill \\ [\widehat{L}_{\mu \nu },\widehat{L}_{\alpha \beta }]\hfill & =& i\delta _{\mu \alpha }\widehat{L}_{\nu \beta }i\delta _{\nu \alpha }\widehat{L}_{\mu \beta }i\delta _{\mu \beta }\widehat{L}_{\nu \alpha }+i\delta _{\nu \beta }\widehat{L}_{\mu \alpha }\hfill \\ [\widehat{L}_{\mu \nu },\widehat{L}_\lambda ]\hfill & =& i\delta _{\mu \lambda }\widehat{L}_\nu i\delta _{\nu \lambda }\widehat{L}_\mu \hfill \\ [\widehat{L}_\mu ,\widehat{h}]\hfill & =& 0\hfill \\ [\widehat{L}_\mu ,\widehat{L}_\nu ]\hfill & =& 2i\widehat{h}\widehat{L}_{\mu \nu }.\hfill \end{array}}`$ (27)
On the Hilbert space of negative-energy states, we can introduce $`\widehat{J}_{MN}`$ where the capital Latin letters $`M`$, $`N`$ run from $`0`$ to $`D`$:
$`\widehat{J}_{MN}=\{\begin{array}{cc}\widehat{L}_{\mu \nu }\hfill & \text{if }M=\mu \text{}N=\nu \hfill \\ (2\widehat{h})^{\frac{1}{2}}\widehat{L}_\mu \hfill & \text{if }M=\mu \text{}N=D\hfill \\ (2\widehat{h})^{\frac{1}{2}}\widehat{L}_\nu \hfill & \text{if }M=D\text{}N=\nu \hfill \\ 0\hfill & \text{if }M=N\hfill \end{array}`$
Then the commutation relation (27) says that a $`D`$-dimensional MICZ-Kepler problem has a dynamical $`\text{SO}(D+1)`$-symmetry on the Hilbert space of negative-energy states. Actually, the dynamical symmetry group should be $`\text{Spin}(D+1)`$, rather than $`\text{SO}(D+1)`$. It is also clear that a $`D`$-dimensional MICZ-Kepler problem has a dynamical $`\text{Spin}(D,1)`$-symmetry on the positive-energy states and a dynamical $`\text{Spin}(D)^D`$-symmetry on the zero-energy states.
It also follows from (27) that $`\widehat{h}`$ must be in the center of Lie algebra $`\text{so}(D+1)`$; in fact, it is a function of the quadratic Casimir operator of $`\text{so}(D+1)`$:
$`\widehat{h}={\displaystyle \frac{1/2}{c_2[\mathrm{so}(D+1)]+(\frac{D1}{2})^2\overline{c}_2}}`$ (29)
where $`\overline{c}_2`$ is the value of $`c_2[\mathrm{so}(D1)]`$ in representation $`𝐬^{2\mu }`$.
To prove Eq. (29), we first note that $`\widehat{L}_\mu =(2\widehat{h})^{\frac{1}{2}}\widehat{J}_{\mu D}`$, so $`\widehat{L}_\mu \widehat{L}_\mu =2\widehat{h}_\mu \widehat{J}_{\mu D}\widehat{J}_{\mu D}`$. On the other hand, based on the definition of $`L_\mu `$ given in Eq. (21), a direct computation yields
$`\widehat{L}_\mu \widehat{L}_\mu =1+\left({\displaystyle \frac{1}{2}}(D1)^22\overline{c}_2+\widehat{J}_{\mu \nu }\widehat{J}_{\mu \nu }\right)\widehat{h}.`$
Therefore,
$$1+\left(\frac{1}{2}(D1)^22\overline{c}_2+\widehat{J}_{MN}\widehat{J}_{MN}\right)\widehat{h}=0,$$
then we have Eq. (29).
It is clear now that in order to determine the spectrum of $`\widehat{h}`$, we just need to find out which irreducible representation of $`\text{Spin}(D+1)`$ enters into the Hilbert space $``$ of negative-energy states. However, we shall find the discrete spectrum by solving the Schrödinger equation directly and then figure out the decomposition of the Hilbert space of negative-energy states into the irreducible representations of $`\text{Spin}(D+1)`$ via representation theory.
## 5 The spectrum analysis
The Schrödinger equation for the stationary states, in terms of the polar coordinates, is
$`\left({\displaystyle \frac{1}{2r^{D1}}}_rr^{D1}_r+{\displaystyle \frac{c_2[\text{so}(D)]\overline{c}_2+\delta _D}{2r^2}}{\displaystyle \frac{1}{r}}\right)\psi =E\psi `$ (30)
where $`E`$ is the energy, $`c_2[\text{so}(D)]=\frac{1}{2}\widehat{L}_{\mu \nu }\widehat{L}_{\mu \nu }`$, and $`\delta _D`$ is equal to $`(n1)\mu `$ if $`D=2n`$ and is equal to $`(n1)|\mu |+\mu ^2`$ if $`D=2n+1`$. Under the action of $`\text{Spin}(D)`$, the Hilbert space of negative-energy states splits into the direct sum of irreducible components. These irreducible components are essentially labeled by a nonnegative integer $`l`$, and we shall be able to see that shortly. On the irreducible component labeled by $`l`$, the Schrödinger equation becomes an equation for the radial part:
$`\left({\displaystyle \frac{1}{2r^{D1}}}_rr^{D1}_r+{\displaystyle \frac{c_2[l]\overline{c}_2+\delta _D}{2r^2}}{\displaystyle \frac{1}{r}}\right)R_{kl}=E_{kl}R_{kl}`$ (31)
where $`c_2[l]`$ is the value of the quadratic Casimir operator of $`\text{so}(D)`$ in the irreducible component labeled by $`l`$, and the additional label $`k`$ is introduced for the purpose of listing the radial eigenfunctions, just as in the Kepler problem.
Let $`E_{kl}=\frac{1}{2}\lambda _{kl}^2`$ and $`R_{kl}(r)=e^{\lambda _{kl}r}u_{kl}`$, then the preceding radial Schrödinger equation becomes
$`\left({\displaystyle \frac{1}{2r^{D1}}}_rr^{D1}_r+\lambda _{kl}{\displaystyle \frac{1}{r^{\frac{D1}{2}}}}_rr^{\frac{D1}{2}}+{\displaystyle \frac{c_2[l]\overline{c}_2+\delta _D}{2r^2}}{\displaystyle \frac{1}{r}}\right)u_{kl}=0`$ (32)
Let $`y_{kl}=r^{\frac{D1}{2}}u_{kl}`$, then the above equation becomes
$`\left({\displaystyle \frac{d^2}{dr^2}}2\lambda _{kl}{\displaystyle \frac{d}{dr}}+\left[{\displaystyle \frac{2}{r}}{\displaystyle \frac{c_2[l]\overline{c}_2+\delta _D+\frac{(D1)(D3)}{4}}{r^2}}\right]\right)y_{kl}(r)=0`$ (33)
Assume that $`y_{kl}(r)r^s`$ as $`r0^+`$, then we must have the following indicial equation:
$`\overline{)s(s1)=c_2[l]\overline{c}_2+\delta _D+\frac{(D1)(D3)}{4}.}`$ (34)
The further analysis is divided into two cases: 1) $`D`$ is odd, 2) $`D`$ is even.
### 5.1 The odd dimensional cases
Let $`D=2n+1`$. Let $`L^2(𝒮^{2\mu }|_{\mathrm{S}^{2n}})`$ be the $`L^2`$-sections of vector bundle $`𝒮^{2\mu }`$ restricted to the unit sphere $`\mathrm{S}^{2n}`$. From the representation theory, we know that
$`L^2(𝒮^{2\mu }|_{\mathrm{S}^{2n}})=\widehat{{\displaystyle }}_{l0}_l`$ (35)
where $`_l`$ is the irreducible representation space of $`\mathrm{Spin}(2n+1)`$ with highest weight $`(l+|\mu |,|\mu |,\mathrm{},|\mu |)`$. It is then clear that the Hilbert spaces of bound states is
$`=\widehat{{\displaystyle }}_{l0}_l`$ (36)
with $`_l`$ being a subspace of $`L^2(_+,r^{2n}dr)_l`$. Here $`L^2(_+,r^{2n}dr)`$ is the $`L^2`$-space of complex-valued functions on half-line $`_+`$ with measure $`r^{2n}dr`$.
The value of the quadratic Casimir operator of $`\mathrm{so}(2n)`$ on representation $`𝐬^{2\mu }`$ is
$$\overline{c}_2=n\mu ^2+n(n1)|\mu |.$$
The value of the quadratic Casimir operator of $`\mathrm{so}(2n+1)`$ on $`_l`$ is
$$c_2[l]=l^2+2l(n+|\mu |\frac{1}{2})+n\mu ^2+n^2|\mu |.$$
Plugging the values for $`\overline{c}_2`$ and $`c_2[l]`$ into Eq. (34), we get
$$s(s1)=(l+n+|\mu |)(l+n+|\mu |1).$$
Therefore, $`s=l+n+|\mu |`$ or $`s=1ln|\mu |`$. The solution $`s=1ln|\mu |`$ must be rejected; otherwise, the wave-functions cannot be square integrable near $`r=0`$. Just as in solving the hydrogen atom problem, with $`s=l+n+|\mu |`$, we continue the analysis by setting
$$y_{kl}=r^s\underset{m=0}{\overset{\mathrm{}}{}}a_mr^m$$
with $`a_0=1`$ and then get the recursive relation: for $`m1`$, one has
$`a_m\left((m+s)(m+s1)s(s1)\right)=\left(1\lambda _{kl}(m+s1)\right)a_{m1}.`$ (37)
As it has been demonstrated in Ref. , the power series solution must be a polynomial solution; otherwise, the wave-function will not be square integrable for $`r`$ near infinity. Therefore, we must have
$$\lambda _{kl}=\frac{1}{k+s1}=\frac{1}{k+l+n+|\mu |1}$$
and that leads to the energy spectrum
$`E_{kl}={\displaystyle \frac{1/2}{(k+l+n+|\mu |1)^2}}`$ (38)
where $`k`$ must be a positive integer, and an orthogonal decomposition
$$_l=\widehat{}_{k=1}^{\mathrm{}}_{kl}$$
with each of $`_{kl}`$ being isomorphic to $`_l`$ as $`\mathrm{Spin}(2n+1)`$-modules. Therefore, in view of Eq. (36), we have an orthogonal decomposition of $``$ into energy eigen-states:
$`=\widehat{{\displaystyle }}_{I=0}^{\mathrm{}}_I`$ (39)
where
$`_I={\displaystyle \underset{k+l=I+1}{}}_{kl}.`$
Since linear action of $`\mathrm{Spin}(2n+2)`$ on $``$ commutes with the hamiltonian, this linear action must leave the energy eigen-states $`_I`$ invariant. On the other hand, from representation theory, as a $`\mathrm{Spin}(2n+1)`$-module, being isomorphic to
$$\underset{l=0}{\overset{I}{}}_l,$$
$`_I`$ must be the irreducible representation of $`\text{Spin}(2n+2)`$ with highest weight $`(I+|\mu |,|\mu |,\mathrm{},|\mu |,\pm |\mu |)`$. As a consistency check, one can see that Eq. (29) yields
$`E_I={\displaystyle \frac{1/2}{(I+n+|\mu |)^2}}`$ (40)
on such representation, in complete agreement with Eq. (38) because $`I=k+l1`$. One can show that, as a $`\mathrm{Spin}(2n+2)`$-module, $`_I`$ has the highest weight $`(I+|\mu |,|\mu |,\mathrm{},|\mu |,\mu )`$.
In summary, the energy spectrum is
$`E_I={\displaystyle \frac{1/2}{(I+n+|\mu |)^2}}`$ (41)
where $`I=0`$, $`1`$, $`2`$, …; and $``$ furnishes a representation for $`\mathrm{Spin}(2n+2)`$ and has the following decomposition into energy eigenstates:
$`=\widehat{{\displaystyle }}_{I=0}^{\mathrm{}}_I`$ (42)
where $`_I`$ is the irreducible component of $``$ with highest weight $`(I+|\mu |,|\mu |,\mathrm{},|\mu |,\mu )`$.
### 5.2 The even dimensional cases
Let $`D=2n`$. Let $`L^2(𝒮^{2\mu }|_{\mathrm{S}^{2n1}})`$ be the $`L^2`$-sections of vector bundle $`𝒮^{2\mu }`$ restricted to the unit sphere $`\mathrm{S}^{2n1}`$. From the representation theory, we know that<sup>6</sup><sup>6</sup>6For a fixed $`l`$, there are $`(2\mu +1)`$ many of $`_l`$’s. When $`\mu >1/2`$, Eq. (51) below is no longer valid for some $`_l`$ and the subsequent analysis fails. That is the analytic reason for requiring $`\mu =0`$ or $`1/2`$.
$`L^2(𝒮^{2\mu }|_{\mathrm{S}^{2n1}})=\{\begin{array}{cc}\widehat{}_{l0}(_l^+_l^{})\hfill & \text{if }\mu =1/2\hfill \\ & \\ \widehat{}_{l0}_l^0\hfill & \text{if }\mu =0\hfill \end{array}`$ (46)
where $`_l^\pm `$ is the irreducible representation of $`\mathrm{Spin}(2n)`$ with highest weight $`(l+1/2,1/2,\mathrm{},1/2,\pm 1/2)`$ and $`_l^0`$ is the irreducible representation of $`\mathrm{Spin}(2n)`$ with highest weight $`(l,0,\mathrm{},0)`$. It is then clear that the Hilbert spaces of bound states is
$`=\{\begin{array}{cc}\widehat{}_{l0}(_l^+_l^{})\hfill & \text{if }\mu =1/2\hfill \\ & \\ \widehat{}_{l0}_l^0\hfill & \text{if }\mu =0\hfill \end{array}`$ (50)
with $`_l^\sigma `$ being a subspace of $`L^2(_+,r^{2n1}dr)_l^\sigma `$. Here $`L^2(_+,r^{2n1}dr)`$ is the $`L^2`$-space of complex-valued functions on half-line $`_+`$ with measure $`r^{2n1}dr`$.
The value of the quadratic Casimir operator of $`\mathrm{so}(2n1)`$ on representation $`𝐬^{2\mu }`$ is
$$\overline{c}_2=(n1)\mu ^2+(n1)^2\mu .$$
The value of the quadratic Casimir operator of $`\mathrm{so}(2n)`$ on $`_l^\sigma `$ is
$`c_2[l]=l^2+2l(n+\mu 1)+n\mu ^2+(n^2n)\mu .`$ (51)
Plugging the values for $`\overline{c}_2`$ and $`c_2[l]`$ into Eq. (34), we get
$$s(s1)=(l+n+\mu \frac{1}{2})(l+n+\mu \frac{3}{2}).$$
Therefore, $`s=l+n+\mu \frac{1}{2}`$ or $`s=\frac{3}{2}ln\mu `$. The solution $`s=\frac{3}{2}ln\mu `$ must be rejected; otherwise, the wave-functions cannot be square integrable near $`r=0`$. Just as in solving the hydrogen atom problem, with $`s=l+n+\mu \frac{1}{2}`$, we continue the analysis by setting
$$y_{kl}=r^s\underset{m=0}{\overset{\mathrm{}}{}}a_mr^m$$
with $`a_0=1`$ and then get the recursive relation: for $`m1`$, one has
$`a_m\left((m+s)(m+s1)s(s1)\right)=\left(1\lambda _{kl}(m+s1)\right)a_{m1}.`$ (52)
As it has been demonstrated in Ref. , the power series solution must be a polynomial solution; otherwise, the wave-function will not be square integrable for $`r`$ near infinity. Therefore, we must have
$$\lambda _{kl}=\frac{1}{k+s1}=\frac{1}{k+l+n+\mu \frac{3}{2}}$$
and that leads to the energy spectrum
$`E_{kl}={\displaystyle \frac{1/2}{(k+l+n+\mu \frac{3}{2})^2}}`$ (53)
where $`k`$ must be a positive integer, and an orthogonal decomposition
$$_l^\sigma =\widehat{}_{k=1}^{\mathrm{}}_{kl}^\sigma $$
with each $`_{kl}^\sigma `$ being isomorphic to $`_l^\sigma `$ as $`\mathrm{Spin}(2n)`$-modules. Therefore, in view of Eq. (50), we have an orthogonal decomposition of $``$ into energy eigen-states:
$`=\widehat{{\displaystyle }}_{I=0}^{\mathrm{}}_I`$ (54)
where
$`_I=\{\begin{array}{cc}_{k+l=I+1}(_{kl}^+_{kl}^{})\hfill & \text{if }\mu =1/2\hfill \\ & \\ _{k+l=I+1}_{kl}^0\hfill & \text{if }\mu =0\text{ .}\hfill \end{array}`$
Note that $`_I`$ is isomorphic to
$$\{\begin{array}{cc}_{l=0}^I(_l^+_l^{})\hfill & \text{if }\mu =1/2\hfill \\ & \\ _{l=0}^I_l^0\hfill & \text{if }\mu =0\hfill \end{array}$$
as a $`\mathrm{Spin}(2n)`$-module, from the representation theory, the manifest $`\mathrm{Spin}(2n)`$ linear action on $`_I`$ can be extended to a linear action of $`\mathrm{Spin}(2n+1)`$ such that $`_I`$ is the irreducible representation of $`\text{Spin}(2n+1)`$ with highest weight $`(I+\mu ,\mu ,\mathrm{},\mu )`$ and Eq. (29) is valid. As a consistency check, one can see that Eq. (29) yields
$`E_I={\displaystyle \frac{1/2}{(I+n+\mu \frac{1}{2})^2}}`$ (56)
on such representation, in complete agreement with Eq. (53) because $`I=k+l1`$.
In summary, the energy spectrum is
$`E_I={\displaystyle \frac{1/2}{(I+n+\mu \frac{1}{2})^2}}`$ (57)
where $`I=0`$, $`1`$, $`2`$, …; and $``$ furnishes a representation for $`\mathrm{Spin}(2n+1)`$ and has the following decomposition into energy eigenstates:
$`=\widehat{{\displaystyle }}_{I=0}^{\mathrm{}}_I`$ (58)
where $`_I`$ is the irreducible component of $``$ with highest weight $`(I+\mu ,\mu ,\mathrm{},\mu )`$.
## Appendix A Proof of the remaining part of Lemma 1
The case $`n=1`$ is trivial. So we assume that $`n2`$. To prove identity (18), we first note that we just need to prove that identities
$`{\displaystyle \underset{k}{}}(\gamma _{1,k})^2`$ $`=`$ $`{\displaystyle \frac{1}{n}}c_2`$ (59)
and
$`{\displaystyle \underset{k}{}}\{\gamma _{1,k},\gamma _{2,k}\}`$ $`=`$ $`0`$ (60)
hold in the representation $`𝐬_+^{2\mu }`$ for any non-negative half integer $`\mu `$.
To continue, a digression on Lie algebra $`\mathrm{so}(2n)`$ is needed. Recall that the root space of $`\mathrm{so}(2n)`$ is $`^n`$. Let $`e^i`$ be the vector in $`^n`$ whose $`i`$-th entry is $`1`$ and all other entries are zero. The positive roots are $`e^i\pm e^j`$ with $`1i<jn`$. The simple roots are $`\alpha ^i=e^ie^{i+1}`$, $`i=1`$ to $`n1`$ and $`\alpha ^n=e^{n1}+e^n`$. For the Cartan basis, we make the following choice: The commuting generators in the Cartan subalgebra are taken to be
$$H_j=\gamma _{2j1,2j}j=1\text{ to }n,$$
and the $`E`$ generators are taken to be
$`E_{\eta e^j+\eta ^{}e^k}={\displaystyle \frac{1}{2}}\left(\gamma _{2j1,2k1}+i\eta \gamma _{2j,2k1}+i\eta ^{}\gamma _{2j1,2k}\eta \eta ^{}\gamma _{2j,2k}\right)`$ (61)
where $`j<k`$, and $`\eta `$, $`\eta ^{}`$ $`\{1,1\}`$. Note that, the fact that
$$[E_\alpha ,E_\beta ]=0\text{if }\alpha +\beta \text{ is neither a root nor zero}$$
is frequently used in all subsequent calculations. All of these are standard materials taken from a textbook such as Ref. .
Let
$`𝒪`$ $`=`$ $`{\displaystyle \underset{ni2}{}}E_{e^1e^i}E_{e^1+e^i}`$ (62)
$`𝒪_1`$ $`=`$ $`H_1^2+{\displaystyle \frac{1}{2}}{\displaystyle \underset{ni2}{}}\left(\{E_{e^1e^i},E_{e^1+e^i}\}+\{E_{e^1+e^i},E_{e^1e^i}\}\right)`$ (63)
By simple computations, we have
$`{\displaystyle \underset{k}{}}(\gamma _{1,k})^2`$ $`=`$ $`𝒪_1+𝒪^{}+𝒪`$ (64)
$`{\displaystyle \underset{k}{}}\{\gamma _{1,k},\gamma _{2,k}\}`$ $`=`$ $`{\displaystyle \frac{2}{i}}(𝒪^{}𝒪)`$ (65)
It is then clear from the above calculations that identities (59) and (60) are valid modulo the following claim:
###### Claim 1.
Let $`|\mathrm{\Lambda }`$ be an element of the $`\mathrm{so}(2n)`$-module $`𝐬_+^{2\mu }`$. Then
$`𝒪|\mathrm{\Lambda }`$ $`=`$ $`0,`$ (66)
$`𝒪^{}|\mathrm{\Lambda }`$ $`=`$ $`0,`$ (67)
$`𝒪_1|\mathrm{\Lambda }`$ $`=`$ $`\mu (n+\mu 1)|\mathrm{\Lambda }`$ (68)
$`=`$ $`{\displaystyle \frac{1}{n}}c_2|\mathrm{\Lambda }.`$ (69)
Proof of the claim. We first remark that, among the four equalities in the claim, we just need to prove the first and the third, that is because the 2nd is a consequence of the first and the last is true because $`c_2=n\mu (n+\mu 1)`$. We also remark that we may assume that $`|\mathrm{\Lambda }`$ is a state created from $`|\mu ,\mathrm{},\mu `$ by applying a bunch of lowing operators of the form $`E_{\alpha ^j}`$’s, that is because a general state is always a linear combination of states of this kind.
Next, we observe that
$`E_{e^1+e^i}|\mu ,\mathrm{},\mu =0;`$ (70)
consequently, $`𝒪|\mu ,\mathrm{},\mu =0`$. This observation can be shown by the following trick: Let $`e=E_{e^1e^i}`$, $`f=E_{e^1+e^i}`$, $`h=H_1H_i`$, then
$`[h,e]`$ $`=`$ $`2e,`$
$`[h,f]`$ $`=`$ $`2f,`$
$`[e,f]`$ $`=`$ $`h.`$
I.e., $`\{e,f,h\}`$ forms the standard Cartan basis for $`\mathrm{su}(2)`$. It follows from the following computation
$`f|\mu ,\mathrm{},\mu ^2=\mu ,\mathrm{},\mu |ef|\mu ,\mathrm{},\mu =\mu ,\mathrm{},\mu |h+fe|\mu ,\mathrm{},\mu =0`$
that $`f|\mu ,\mathrm{},\mu =0`$. Moreover, when $`|\mathrm{\Lambda }=|\mu ,\mathrm{},\mu `$, the third identity of the claim is just the consequence of a direct computation. Therefore, the claim is true when $`|\mathrm{\Lambda }=|\mu ,\mathrm{},\mu `$.
Finally, we need to reduce the general case to the special case discussed in the previous paragraph. Combining the computational fact that
$`[E_{e^1+e^j},E_{e^j+e^{j+1}}]=iE_{e^1+e^{j+1}},[E_{e^1e^{j+1}},E_{e^j+e^{j+1}}]=iE_{e^1e^j}`$ (71)
(where $`1<j<n`$) and the computational fact that
$`[E_{e^1+e^{n1}},E_{e^{n1}e^n}]=iE_{e^1e^n},[E_{e^1+e^n},E_{e^{n1}e^n}]=iE_{e^1e^{n1}},`$ (72)
one can show that $`[𝒪,E_{\alpha ^j}]=0`$ for any $`1jn`$. Since $`|\mathrm{\Lambda }`$ can be assumed to be a state created from $`|\mu ,\mathrm{},\mu `$ by applying a bunch of lowing operators of the form $`E_{\alpha ^j}`$’s, in view of the fact that $`𝒪|\mu ,\mathrm{},\mu =0`$, we have $`𝒪|\mathrm{\Lambda }=0`$.
The reduction to the special case for the third identity is a bit involved. The key is to introduce a series of operators:
$`𝒪_2`$ $`=`$ $`2[𝒪_1,E_{\alpha ^1}]`$ (73)
$`𝒪_k`$ $`=`$ $`i[𝒪_{k1},E_{\alpha ^{k1}}]\text{for }3kn`$ (74)
$`𝒪^{n1}`$ $`=`$ $`i[𝒪_{n1},E_{\alpha ^n}]`$ (75)
$`𝒪^{k1}`$ $`=`$ $`i[𝒪^k,E_{\alpha ^k}]\text{for }1kn1`$
and then observe that
$`[𝒪_k,E_{\alpha ^j}]`$ $`=`$ $`0\text{ if }kn1\text{ and }jk`$ (76)
$`[𝒪_{n1},E_{\alpha ^j}]`$ $`=`$ $`0\text{ if }jn1,n`$ (77)
$`[𝒪^k,E_{\alpha ^j}]`$ $`=`$ $`0\text{ if }jk`$ (78)
$`𝒪^0`$ $`=`$ $`4i𝒪`$ (79)
$`𝒪_k|\mu ,\mathrm{},\mu `$ $`=`$ $`0\text{for }2kn`$ (80)
$`𝒪^k|\mu ,\mathrm{},\mu `$ $`=`$ $`0\text{for }0kn1`$ (81)
With the help of the above equations starting from (73) and ending at (76), an induction argument finishes the proof. Here the induction is done on the number of lowing operators of the form $`E_{\alpha ^j}`$’s which are used to create state $`|\mathrm{\Lambda }`$.
We end this appendix with the following conjecture: *identity (18) holds if and only if when the representation is a Young power of $`𝐬_+`$ or $`𝐬_{}`$.* For our purpose, we have just proved the “if” part in this appendix. |
warning/0507/gr-qc0507098.html | ar5iv | text | # Length Uncertainty in a Gravity’s Rainbow Formalism
## I Introduction
A standard result in Quantum Mechanics is that the measurement of the position of a quantum state is affected by an uncertainty that satisfies the Heisenberg relations GP . In order to diminish the position uncertainty one is thus forced to consider states with increasing momentum uncertainty, achieving an infinite spatial resolution only at the cost of completely delocalizing the momentum. In the presence of gravity, however, the situation becomes more complicated. Via Einstein equations, an uncertainty in the (energy-)momentum of the system results in one in the geometry, which implies an additional uncertainty in the position. The total position uncertainty will therefore consist in the combined effect of a purely quantum mechanical contribution and a contribution of gravitational origin Luis . In these circumstances, one should not expect that an infinite spatial resolution can be reached, unless there exists a very specific relation between these types of contributions. Similar conclusions apply to the measurements of length of spatial intervals, determined by the positions of their endpoints.
The most common approach to analyze the emergence of a minimum spatial (or time) uncertainty when gravity comes into the scene consists in adopting a perturbative scheme. The starting point is a flat background where the matter is inserted. This matter curves the spacetime, producing a deformation of the geometry which in turn modifies the expression of the physical energy and momentum of the system (usually defined in terms of normalized -asymptotic- Killing vectors). The process continues with successive corrections that one assumes to be less and less important. The studies in the literature indicate that a minimum uncertainty is ineluctable in this kind of perturbative quantization (at least in the next-to-leading-order approximation) Luis ; Pad ; unc ; unc2 . A different issue, which is still open to debate, is whether the same result holds as well in the context of a non-perturbative quantum description PG ; BMV .
A suitable arena to test some of these issues is provided by doubly special relativity (DSR) Amelino ; Amelino2 . In this kind of theories, the definition of the physical energy and momentum of particles is modified with respect to the standard relativistic one in order to encode, at least to some extent, the possible effects of the gravitational interactions, without necessarily adhering to any perturbative interpretation. The modification is such that the system presents an energy and/or momentum scale which is invariant under Lorentz transformations. This is possible because the action of the Lorentz group becomes nonlinear on the physical energy-momentum space Amelino ; Amelino2 ; DSR1 ; DSR2 ; DSR12 .
Several proposals have been put forward for the realization of DSR in position space kappa ; position ; mignemi . In a previous paper PG we suggested that this realization should be determined by completing into a canonical transformation the nonlinear mapping that relates the original energy-momentum variables of standard relativity in Minkowski spacetime (that we will call pseudo variables from now on) with the physical energy-momentum of the system in DSR judes . In this framework, the background Minkowski coordinates are mapped to a new set of spacetime coordinates that can be regarded as canonically conjugate to the physical energy-momentum. Those coordinates are linear in the Minkowski ones, but depend in a non-trivial way on the energy and momentum of the particle. Owing to this dependence of the spacetime description, the formalism can be considered a kind of gravity’s rainbow rainbow .
Our discussion in Ref. PG was focused on the existence of a minimum time uncertainty in quantum theories derived from DSR. In particular, we considered the different possibilities of describing the quantum evolution in terms of a parameter that corresponds either to the original time of the Minkowski background or to the physical time of the system. According to our comments above, we will respectively refer to these two types of quantization as perturbative and non-perturbative ones, given the distinct philosophy in the use of background structures. Our analysis proved that, while there always exists a non-vanishing uncertainty in the physical time when a perturbative quantization is adopted, an infinite time resolution can be achieved in certain theories when the quantization is non-perturbative. More precisely, no minimum time uncertainty arises non-perturbatively in DSR theories whose physical energy is unbounded from above. The aim of the present work is to extend this study of the uncertainty from time lapses to the case of spatial intervals.
A particular class of spacetimes in which the commented analysis of the time uncertainty has been carried out in detail is that of the Einstein-Rosen waves BMV . These linearly polarized waves are described by cylindrically symmetric spacetimes in 3+1 dimensions, but can equivalently be described in terms of a massless scalar field coupled to gravity in 2+1 dimensions with axial symmetry ERscalar ; ERash ; ERBMV . In this dimensionally reduced formulation, the system can in fact be viewed as an example of DSR theories, with a physical energy that is bounded from above 2+1 ; ERbound . Therefore, for Einstein-Rosen waves, a non-vanishing quantum time uncertainty emerges both in the perturbative and in the non-perturbative approaches. The study of the spatial uncertainty is not specially interesting in this case, because the associated DSR theory involves no modification in the definition of the momenta nor in the canonically conjugate position variables.
The rest of the paper is organized as follows. In the following section, we review some aspects of the formulation of DSR theories in momentum space and introduce our canonical proposal for their realization in position space. We obtain spacetime coordinates that are conjugate to the physical energy-momentum, arriving at a gravity’s rainbow formalism. Next, we study the quantization of this formalism, restricting our considerations to free systems that can be described within a Hamiltonian scheme. Adopting a perturbative approach to the quantization, we analyze in Sec. III the length uncertainty, i.e. the uncertainty in the difference of spatial positions. We show in Sec. IV that this uncertainty cannot vanish in the perturbative case under quite generic assumptions. Furthermore, in Sec. V we prove that the appearance of a minimum length uncertainty persists when the quantum evolution is described in terms of the physical time, i.e., in a non-perturbative quantization. However, we comment the possibility that in some DSR models one could construct a different type of non-perturbative quantum theory where the physical position operator became explicitly time independent. In this scenario, the resolution in the spatial position could in principle be made as large as desired if the DSR theory does not possess an invariant momentum scale. The uncertainty in the physical length (as well as in the physical time lapse) is studied in Sec. VI in the low-energy sector, approximating the results of the perturbative quantization up to first order corrections. In Sec. VII we consider the *massless* case in this approximation for large values of the Minkowski time $`T`$. We show that the uncertainty increases then like the square root of $`T`$, just as it occurs in Salecker and Wigner devices SalWig . We present our conclusions in Sec. VIII. Finally, two appendices are added. In the following, we will adopt units in which $`\mathrm{}=c=1`$ (with $`\mathrm{}`$ being Planck constant and $`c`$ the speed of light).
## II DSR in momentum and position spaces
A characteristic feature of DSR theories is that they possess a Lorentz invariant energy and/or momentum scale, apart from the scale provided in standard relativity by the speed of light Amelino ; Amelino2 ; DSR1 ; DSR2 ; DSR12 . The invariance of such a scale is possible only thanks to a nonlinear realization of the Lorentz group in momentum space. A simple way to construct a realization of this kind is by introducing an invertible map $`U`$ between the physical energy-momentum $`P^a=(E,p^i)`$ and a standard Lorentz 4-vector $`\mathrm{\Pi }^a=(ϵ,\pi ^i)`$, which we call the pseudo energy-momentum judes (lowercase Latin indices from the beginning and the middle of the alphabet represent Lorentz and flat spatial indices, respectively). Denoting the usual linear action of the Lorentz group by $``$, the nonlinear Lorentz transformations are then given by $`L(P)=(U^1U)(P)`$ judes ; MS .
The map $`U`$ must reduce to the identity when energies and momenta are negligibly small compared to the DSR scale, so that the physical and pseudo variables coincide in this limit. In addition, a simplifying assumption that is generally accepted is that the standard action of rotations is preserved; only boosts are modified in DSR MS ; kappa . So, with the notation $`p:=|\stackrel{}{p}|`$ and $`\pi :=|\stackrel{}{\pi }|`$, the most general expression for the map $`U`$ becomes kappa ; PG
$`\mathrm{\Pi }=U(P)`$ $``$ $`\{\begin{array}{c}ϵ=\stackrel{~}{g}(E,p),\hfill \\ \pi ^i=\stackrel{~}{f}(E,p)\frac{p^i}{p},\hfill \end{array}`$ (2.3)
$`P=U^1(\mathrm{\Pi })`$ $``$ $`\{\begin{array}{c}E=g(ϵ,\pi ),\hfill \\ p^i=f(ϵ,\pi )\frac{\pi ^i}{\pi }.\hfill \end{array}`$ (2.6)
Since the only invariant energy-momentum scale in standard special relativity is at infinity, the DSR theory admits a Lorentz invariant scale at a finite value of the energy and/or momentum only if the map $`U`$ has a singularity there MS . The domain of definition of $`U`$ (which is assumed to contain the low energy-momentum sector) is therefore bounded from above by that scale. Consequently, DSR theories can be classified in three types: DSR1 if it is only the physical momentum that is bounded from above, DSR3 if it is the physical energy what is bounded, and DSR2 if both the physical energy and momentum are bounded.
As it is implicit in our discussion, DSR theories are usually formulated in momentum space, mainly owing to the increasing interest in investigating the observational implications of deformed dispersion relations Amelino ; phenomen . There are different proposals to determine what is the modified spacetime geometry and the corresponding transformation rules in position space that should complement this formulation kappa ; position . Among them, one of the most popular consists in abandoning the commutativity of the spacetime coordinates, as it happens e.g. in $`\kappa `$-deformed Minkowski spacetime DSR12 ; kappa .
However, noncommutative geometries are by no means the only way to obtain a consistent realization in position space. The same goal can be achieved without renouncing the conventional framework of commutative spacetimes. In fact, the literature contains several suggestions for realizations of this kind position ; mignemi ; PG ; hinterleitner . A particular example was put forward by Magueijo and Smolin rainbow , who required that the contraction between the energy-momentum and an infinitesimal spacetime displacement were a linear invariant in DSR. This requirement leads to new spacetime coordinates that depend on the energy-momentum. Ultimately, the system adopts a spacetime metric that directly depends on the energy and momentum of its particle content. This explains the name of *gravity’s rainbow* that has been given to this class of DSR implementations.
In this work, we will follow a suggestion for the realization of DSR in position space that differs from that of Magueijo and Smolin, although it leads as well to a gravity’s rainbow formalism in the sense of the energy dependence of the geometry. We will adopt the proposal of Ref. PG , namely, we will specify the realization by demanding the invariance of the symplectic form $`𝐝q^a𝐝\mathrm{\Pi }_a`$ (where the wedge denotes the exterior product and Lorentz indices are lowered with the Minkowski metric). This assigns to the system new, modified spacetime coordinates $`x^a`$ that are conjugate to the physical energy-momentum $`P_a`$, so that the relation between $`(q^a,\mathrm{\Pi }_a)`$ and $`(x^a,P_a)`$ is given by a canonical transformation. Similar proposals for a canonical implementation of DSR theories have been analyzed by other authors mignemi ; hinterleitner .
By completing the map $`U`$ into a canonical transformation, one easily derives the following expressions for the new spacetime coordinates PG :
$`x^i`$ $`=`$ $`{\displaystyle \frac{1}{J}}\left[_\pi g{\displaystyle \frac{\pi ^i}{\pi }}q^0+_ϵg{\displaystyle \frac{\pi ^i\pi _j}{\pi ^2}}q^j\right]+{\displaystyle \frac{\pi }{f}}(q^i{\displaystyle \frac{\pi ^i\pi _j}{\pi ^2}}q^j),`$
$`x^0`$ $`=`$ $`{\displaystyle \frac{1}{J}}\left[_\pi fq^0+_ϵf{\displaystyle \frac{\pi _i}{\pi }}q^i\right].`$ (2.7)
Here, $`J=_ϵg_\pi f_\pi g_ϵf`$ is the determinant of the Jacobian of the transformation $`U^1`$ between $`(ϵ,\pi )`$ and $`(E,P)`$, and the functions $`f`$ and $`g`$ (and therefore $`J`$) depend on $`(ϵ,\pi )`$. We point out that the transformation (II) is linear in the coordinates $`q^a`$, but generally depends non-trivially on the energy-momentum.
We will refer to $`(x^a,P_a)`$ and $`(q^a,\mathrm{\Pi }_a)`$ as physical and background (or pseudo) variables, respectively, and will denote $`q^0`$ by $`T`$ and $`x^0`$ by $`t`$ to emphasize the role played by the evolution parameter in our discussion. In addition, we assume in the following that the system admits a Hamiltonian description, so that the value of the physical and pseudo energies are respectively given by a physical Hamiltonian $`H`$ and a background Hamiltonian $`H_0`$. Together with Eq. (2.3), we then get $`EH=g(H_0,\pi )`$ and $`ϵH_0=\stackrel{~}{g}(H,p)`$. Finally, since DSR theories are essentially conceived as effective descriptions of free particles that incorporate quantum gravitational phenomena, we will concentrate our analysis on free systems. For such systems, the energy and momentum are constants of motion. The Hamiltonian is hence time independent and commutes with the momentum under Poisson brackets, both for the physical and the background variables.
## III Physical Length Uncertainty: Perturbative Case
In this section, we will consider the perturbative approach to the quantization of the system in which one adopts the background time coordinate $`q^0=T`$ as evolution parameter, so that the evolution is generated by the Hamiltonian $`H_0`$. We assume that a quantization of this kind is feasible. In such a quantum description, the physical time is represented by a genuine operator $`\widehat{t}`$ BMV ; PG . We want to study whether the spatial position and length determined by the physical coordinates $`x^i`$ is affected in this case by a non-vanishing quantum uncertainty. In order to simplify the analysis and deal only with scalar quantities (circumventing the kind of problems derived from the use of vector components and their dependence on choices of fixed background structures, choices which are questionable both from the viewpoint of general relativity and of the fluctuations inherent to quantum mechanics) we will focus our attention exclusively on the projection of the position vector along the direction of motion:
$$X:=x^i\frac{p_i}{p}=x^i\frac{\pi _i}{\pi }=\frac{1}{J}\left[_\pi gT+_{H_0}g\frac{\pi _j}{\pi }q^j\right].$$
(3.1)
We recall that $`g`$, $`f`$, and $`J`$ are functions of only $`H_0`$ and $`\pi `$. Remarkably, this expression is similar to that given in (II) for the time coordinate $`x^0=t`$ with the exchange of the function $`f`$ for $`g`$ and a flip of global sign (so that the determinant of the Jacobian $`J`$ is preserved under the commented exchange).
Given our restriction to free systems, where the energy and momentum are conserved, the only variable in the expression for $`X`$ that evolves in time (apart from the parameter $`T`$) is
$$s_T:=\pi _jq^j.$$
(3.2)
The subscript $`T`$ emphasizes this time dependence. Moreover, since the system is free, the background Hamiltonian $`H_0`$ is a function of only the pseudo momentum. Then, from the Hamiltonian equations of motion, the time derivative of $`s_T`$ equals $`\pi H_0^{}`$, which is a constant of motion. Here, the prime denotes the derivative with respect to $`\pi `$. Thus, we conclude that $`s_T=s_0+T\pi H_0^{}`$, where $`s_0`$ is the value of $`s_T`$ at the initial instant of time.
For our quantum analysis we will only consider differences between position variables, avoiding in this way the arbitrariness in the choice of an origin and the conceptual tensions that arise from fixing it classically while allowing quantum fluctuations in the spatial position. The physics of the problem suggests two possible elections of reference for the position, namely, either the physical or the background initial value (of the projection along the direction of motion) of the position vector. In the first case, the position difference determines the physical interval covered by the particle in the background lapse $`T`$. In the second case the difference includes as well the effective corrections to the initial background position contained in DSR. We will study both possibilities to show that our conclusions do not depend on the specific choice adopted. To distinguish between the two cases, we introduce a parameter $`\eta `$, with $`\eta =0`$ corresponding to the initial physical position and $`\eta =1`$ to the background one. Explicitly, the former of these positions is given by Eq. (3.1) with $`T=0`$ and $`\pi _jq^j`$ replaced with $`s_0`$, whereas the latter is equal to $`s_0/\pi `$.
From the difference between $`X`$ and any of these reference positions, we obtain the following length:
$`L_\eta `$ $`:=`$ $`{\displaystyle \frac{1}{J}}\left[_\pi gT+_{H_0}gS_T+\eta {\displaystyle \frac{(_{H_0}gJ)}{\pi }}s_0\right],`$
$`S_T`$ $`:=`$ $`{\displaystyle \frac{s_Ts_0}{\pi }}.`$ (3.3)
We will refer to it as the physical length. To represent it as an operator, we write
$`\widehat{L}_\eta `$ $`:=`$ $`\widehat{M}(H_0,\pi )T+\widehat{R}_{T,\eta },`$ (3.4)
$`\widehat{R}_{T,\eta }`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\widehat{N}(H_0,\pi )\widehat{S}_T+\widehat{S}_T\widehat{N}(H_0,\pi )\right)`$ (3.5)
$`+`$ $`{\displaystyle \frac{\eta }{2}}\left(\widehat{O}(H_0,\pi )\widehat{s}_0+\widehat{s}_0\widehat{O}(H_0,\pi )\right),`$
where
$$M:=\frac{_\pi g}{J},N:=\frac{_{H_0}g}{J},O:=\frac{_{H_0}gJ}{\pi J}.$$
(3.6)
The subscript $`T`$ denotes again dependence on time. In Eqs. (3.4) and (3.5), we have symmetrized the products of $`\widehat{N}`$ with $`\widehat{S}_T`$ and $`\widehat{O}`$ with $`\widehat{s}_0`$, and displayed explicitly the arguments of the functions $`M`$, $`N`$, and $`O`$. As we have commented, these functions correspond to constants of motion. Their respective operators can be defined in terms of those for $`H_0`$ and $`\pi `$ employing the spectral theorem. As for the operator representing $`s_T`$ (and hence $`S_T`$), we will comment on its definition later in this section.
It is worth pointing out that our expressions are to some extent similar to those introduced in Ref. PG for the physical time operator $`\widehat{t}`$. The differences come from the fact that in the latter case the role of the initial background position variable $`s_0/\pi `$ is played by the initial background time ($`T=0`$), and that in that work we only analyzed the choice $`\eta =1`$ (initial time identified with that of the background time parameter). Our analysis here can be easily applied to the resulting time lapse, $`t_\eta `$, the precise correspondence being the disappearance of the contribution $`1/\pi `$ in the function $`O(H_0,\pi )`$ (and therefore in $`\widehat{R}_{T,\eta }`$), the exchange of the function $`f`$ for $`g`$ in the resulting formulas, and a flip of global sign.
In order to calculate the uncertainty in the physical length operator $`\widehat{L}_\eta `$, we will follow the same procedure that was explained in Ref. PG . Given a quantum state, one can measure the probability densities of any set of observables at any instant of time note0 . In this way, one can determine e.g. the expectation value of those operators. In addition, one can estimate the value of the parameter $`T`$ at that instant of time by analyzing the evolution of the probability densities of observables in the considered state. This procedure allows to derive a statistical distribution for $`T`$ with probability density $`\rho (T)`$ (and mean value $`\overline{T}`$). Heisenberg relations imply that the uncertainty $`\mathrm{\Delta }T`$ of this distribution satisfies the inequality $`\mathrm{\Delta }T\mathrm{\Delta }H_01/2`$ (usually called the fourth Heisenberg relation) GP ; PG . The double average process involved by the quantum expectation value $``$ and by the estimation of the time parameter leads to the following uncertainty:
$$(\mathrm{\Delta }L_\eta )^2=𝑑T\rho (T)\left(\widehat{M}T+\widehat{R}_{T,\eta }\widehat{M}\overline{T}\widehat{R}_{\overline{T},\eta }\right)^2.$$
(3.7)
Here, $`\widehat{R}_{\overline{T},\eta }`$ is the mean value of the operator $`\widehat{R}_{T,\eta }`$ computed with the commented double average PG .
At this stage, some remarks are in order about the precise operator representation adopted for $`s_T`$ when defining $`\widehat{R}_{T,\eta }`$ and how this affects the measurements that are necessary to determine the mean value of this observable. Two cases are worth commenting. On the one hand, one can represent $`s_T`$ as an explicitly $`T`$-independent operator by simply adopting a symmetrized factor ordering in Eq. (3.2) and directly promoting the canonical background variables $`(q^i,\pi _i)`$ to operators. Similarly, we can define $`\widehat{S}_T`$ from its symmetrized classical expression. By performing quantum measurements at the fixed instant of time in which the system is analyzed, one can then determine the probability distribution for $`s_T`$ at that instant. No estimation of the value of the evolution parameter is needed, so that the average over $`T`$ becomes spurious. Similar arguments apply to the products of $`s_T`$ with constants of motion that appear in $`\widehat{R}_{T,\eta }`$. At least in principle, one may hence identify $`\widehat{R}_{\overline{T},\eta }`$ and $`\widehat{R}_{T,\eta }`$ in Eq. (3.7), even if the exact value of $`T`$ in which the measurements are made is not known.
On the other hand, one can instead reflect explicitly all the $`T`$-dependence of $`s_T`$ in the definition of its associated operator. Starting with the solution to its evolution equation, one arrives at $`\widehat{s}_T:=\widehat{s}_0+T\widehat{\pi }\widehat{H_0^{}}`$. So $`\widehat{S}_T:=T\widehat{H_0^{}}`$. Here, $`\widehat{H_0^{}}`$ can be defined in terms of the pseudo momentum using the spectral theorem. Since the operator $`\widehat{H_0^{}}`$ corresponds to a constant of motion, its probability density does not evolve in time. Actually, the same happens with $`\widehat{s}_0`$, $`\widehat{M}`$, $`\widehat{N}`$ and $`\widehat{O}`$, appearing in Eqs. (3.4) and (3.5). In particular, the measurements of all of their densities can be performed at an initial instant of time, identified with $`T=0`$. For all other instants, the only missing piece of information is the probability density $`\rho (T)`$, obtained through measurements of distributions of observables that track the passage of time. In this case, obviously, the average with $`\rho (T)`$ cannot be obviated when calculating the mean value of $`\widehat{R}_{T,\eta }`$.
The two cases can nevertheless be studied in exactly the same way by simply combining all the explicit linear $`T`$-dependence of $`\widehat{X}`$. In the latter case, one gets
$`\widehat{L}_\eta `$ $`=`$ $`\widehat{Y}(H_0,\pi )T+\widehat{Z}_\eta (H_0,\pi ,s_0),`$ (3.8)
$`\widehat{Y}(H_0,\pi )`$ $`=`$ $`\widehat{M}(H_0,\pi )+\widehat{N}(H_0,\pi )\widehat{H_0^{}}(\pi ),`$ (3.9)
$`\widehat{Z}_\eta (H_0,\pi ,s_0)`$ $`=`$ $`{\displaystyle \frac{\eta }{2}}\left(\widehat{O}(H_0,\pi )\widehat{s}_0+\widehat{s}_0\widehat{O}(H_0,\pi )\right).`$ (3.10)
For computational purposes, expression (3.4) can be considered a particular example of formula Eq. (3.8) with $`\widehat{Y}=\widehat{M}`$ and $`\widehat{Z}_\eta =\widehat{R}_{T,\eta }`$. With the same substitutions in Eq (3.7), the physical length uncertainty can then be rewritten:
$$(\mathrm{\Delta }L_\eta )^2=[\mathrm{\Delta }(Y\overline{T}+Z_\eta )]^2+\widehat{Y}^2(\mathrm{\Delta }T)^2+(\mathrm{\Delta }T\mathrm{\Delta }Y)^2.$$
(3.11)
The case of the physical time lapse can be treated in a completely similar way PG , removing the contribution $`1/\pi `$ to $`O`$ in the definition of $`Z_\eta `$, interchanging the functions $`f`$ and $`g`$, and introducing a global change of sign (to preserve that of $`J`$).
## IV Existence of a Minimum Uncertainty in the Perturbative Case
The physical length uncertainty vanishes if and only if the three positive terms that form the r.h.s. of equation (3.11) are equal to zero. We will show in this section that this cannot generally occur.
In order for the uncertainty to vanish, it must in particular do so at large $`T`$, times for which the contribution $`(\overline{T}\mathrm{\Delta }Y)^2`$ dominates in (3.11). Therefore, $`\mathrm{\Delta }Y`$ (which is independent of time) must vanish. Let us assume that the expression of the background Hamiltonian $`H_0`$ as a function of $`\pi `$ is invertible for the whole range of pseudo energies, i.e. $`\pi =\pi (H_0)`$ PG . One can then define the function $`𝒴(H_0):=Y[H_0,\pi (H_0)]`$. In these circumstances, it suffices that the system satisfies, e.g., one of the following generic sets of hypotheses to prove that the physical length uncertainty is strictly positive.
i) We first assume that the function $`𝒴(H_0)`$ is strictly monotonic, namely $`d𝒴/dH_00`$, so that it provides a one-to-one map. Then, via the spectral theorem, the eigenstates of the operators $`𝒴`$ and $`H_0`$ coincide, and the demand $`\mathrm{\Delta }Y=\mathrm{\Delta }𝒴=0`$ implies that $`\mathrm{\Delta }H_0=0`$. The fourth Heisenberg relation leads to $`\mathrm{\Delta }T\mathrm{}`$. Let us then prove that the third term in Eq. (3.11) does not vanish when $`\mathrm{\Delta }H_0`$ tends to zero. Expanding $`𝒴`$ around the mean value of $`H_0`$ nota , we find
$$(\mathrm{\Delta }𝒴)^2=\widehat{𝒴}^2\widehat{𝒴}^2\left(\frac{d𝒴}{dH_0}|_{\widehat{H}_0}\mathrm{\Delta }H_0\right)^2,$$
(4.1)
$$\underset{\mathrm{\Delta }H_00}{lim}2\mathrm{\Delta }T\mathrm{\Delta }𝒴\underset{\mathrm{\Delta }H_00}{lim}\frac{\mathrm{\Delta }𝒴}{\mathrm{\Delta }H_0}=|\frac{d𝒴}{dH_0}|_{\widehat{H}_0}|0.$$
(4.2)
We hence conclude that the physical length uncertainty cannot vanish in this case.
ii) We suppose instead that $`𝒴(H_0)`$ is positive and, for large pseudo energies, grows at least like $`H_0`$ multiplied by a constant. We analyze first the case in which $`𝒴`$ is strictly positive. Since $`\widehat{Y}=\widehat{𝒴}`$ is then different from zero, the vanishing of the second term in Eq. (3.11) requires $`\mathrm{\Delta }T=0`$. So, the fourth Heisenberg relation implies that $`\mathrm{\Delta }H_0\mathrm{}`$. Let us consider again the third term in Eq. (3.11). Our condition on the behavior of $`𝒴`$ for large $`H_0`$ can be rephrased by saying that $`lim_{H_0\mathrm{}}(𝒴/H_0)>r`$ for a certain number $`r>0`$. As a consequence, one can see that $`lim_{\mathrm{\Delta }H_0\mathrm{}}(\mathrm{\Delta }𝒴/\mathrm{\Delta }H_0)>r`$. Therefore, the product $`\mathrm{\Delta }T\mathrm{\Delta }Y=\mathrm{\Delta }T\mathrm{\Delta }𝒴`$ cannot vanish when $`\mathrm{\Delta }H_0`$ tends to infinity, and the physical length uncertainty is strictly positive. On the other hand, in the case that $`𝒴`$ can also take the zero value, $`\widehat{Y}=\widehat{𝒴}`$ may occasionally vanish, but this may only happen if the quantum state is in the kernel of the operator $`\widehat{𝒴}`$. We then introduce the additional assumption that this kernel is formed exclusively by the eigenvectors corresponding to a unique eigenvalue $`\overline{H}_0`$ of $`\widehat{H}_0`$, a result that holds when $`𝒴(H_0)`$ vanishes only at that value of the pseudo energy. If the system approaches such an eigenvector, the uncertainty of $`H_0`$ tends to zero and $`\mathrm{\Delta }T\mathrm{}`$. Assuming finally that $`(d𝒴/dH_0)|_{\overline{H}_0}0`$, one arrives at the same conclusion about the third term in Eq. (3.11) that was obtained in inequality (4.2) nota . Therefore, under this set of hypotheses, it is impossible to achieve an infinite resolution in the physical length.
An important class of DSR theories in which the positivity of $`𝒴(H_0)`$ is satisfied when $`s_T`$ is represented by an explicitly time-dependent operator is when the physical energy does not depend on the pseudo momentum, i.e., when the function $`g`$ depends only on $`H_0`$. In this case,
$$M=\frac{_\pi g}{J}=0,N=\frac{_{H_0}g}{J}=\frac{1}{_\pi f},Y=\frac{H_0^{}}{_\pi f}.$$
As a consequence, $`𝒴(H_0)`$ is non-zero, because both the map $`U`$ and $`H_0(\pi )`$ are invertible by assumption (this guarantees that $`_\pi f0`$ and $`H_0^{}0`$). Since $`𝒴(H_0)`$ has a definite sign, and $`_\pi f1`$ in the sector of small pseudo energy-momentum, in the standard situation with a pseudo energy that increases with $`\pi `$ in that sector we conclude that $`𝒴(H_0)`$ is strictly positive nota2 .
In conclusion, a non-vanishing uncertainty generically affects the physical length in the perturbative quantization of the system. The above discussion can also be applied to the study of the physical time uncertainty considered in Ref. PG . All the hypotheses can be easily generalized to that case with the due substitution of $`𝒴`$ by the function $`𝒱`$ defined in that reference.
## V Physical Position Uncertainty: Non-Perturbative Case
We turn now to the analysis of the physical length uncertainty when one adopts what we have called a non-perturbative quantization, i.e., when the quantum evolution is described in terms of the physical time.
In principle, one can always construct a non-perturbative quantum theory (in the sense indicated above) starting with the perturbative one, which has been assumed to exist. Employing the spectral decomposition of the pseudo momentum $`\pi `$ and recalling that $`H_0=H_0(\pi )`$, one can define the physical Hamiltonian $`H=g(H_0,\pi )`$ as an operator. The parameter of the evolution generated by this Hamiltonian can be identified with the physical time $`t`$. By contrast, the background time gets now promoted to an operator. This fact changes the expression of the observable $`\widehat{L}_\eta `$ when regarded as an explicitly time dependent operator. From Eqs. (3.4) and (II), one gets
$`\widehat{L}_\eta ^{[2]}`$ $`=`$ $`\widehat{M}^{[2]}(H_0,\pi )t+\widehat{R}_{t,\eta }^{[2]},`$ (5.1)
$`\widehat{R}_{t,\eta }^{[2]}`$ $`:=`$ $`{\displaystyle \frac{1}{2}}\left(\widehat{N}^{[2]}(H_0,\pi )\widehat{S}_t+\widehat{S}_t\widehat{N}^{[2]}(H_0,\pi )\right)`$ (5.2)
$`+`$ $`{\displaystyle \frac{1}{2}}\left(\widehat{O}_\eta ^{[2]}(H_0,\pi )\widehat{s}_0+\widehat{s}_0\widehat{O}_\eta ^{[2]}(H_0,\pi )\right),`$
where
$`M^{[2]}`$ $`:=`$ $`{\displaystyle \frac{_\pi g}{_\pi f}},N^{[2]}:={\displaystyle \frac{1}{_\pi f}}`$ (5.3)
$`O_\eta ^{[2]}`$ $`:=`$ $`\eta {\displaystyle \frac{_{H_0}gJ}{\pi J}}{\displaystyle \frac{_\pi g_{H_0}f}{\pi J_\pi f}}.`$ (5.4)
The analysis is parallel to that followed in Sec. III and Sec. IV, with the caveat that $`s_t:=\pi _jq^j`$ \[and therefore $`S_t:=(s_ts_0)/\pi `$\] must now be considered a variable that evolves in the physical time $`t`$, rather than in the background time. In particular, by extracting explicitly all the time dependence of $`s_t`$ when defining its operator counterpart, one arrives at
$$\widehat{L}_\eta ^{[2]}=\widehat{Y}^{[2]}(H_0,\pi )t+\widehat{Z}_\eta ^{[2]}(H_0,\pi ,s_0),$$
(5.5)
with
$`\widehat{Y}^{[2]}(H_0,\pi )`$ $`=`$ $`(\widehat{H_0^{}}\widehat{_{H_0}g}+\widehat{_\pi g})\widehat{N}^{[2]}(H_0,\pi )`$ (5.6)
$`+`$ $`\widehat{M}^{[2]}(H_0,\pi ),`$
$`\widehat{Z}_\eta ^{[2]}(H_0,\pi ,s_0)`$ $`=`$ $`{\displaystyle \frac{\widehat{O}_\eta ^{[2]}(H_0,\pi )\widehat{s}_0+\widehat{s}_0\widehat{O}_\eta ^{[2]}(H_0,\pi )}{2}}.`$ (5.7)
Here, the observable $`\widehat{s}_0`$ represents the value of $`s_t`$ at the initial physical time, which is a constant of motion.
In order to calculate the physical length uncertainty, one has to average now over the time parameter $`t`$, instead of averaging over $`T`$, as we did in Eq. (3.7). This leads to
$`\left(\mathrm{\Delta }L_\eta ^{[2]}\right)^2`$ $`=`$ $`\left[\mathrm{\Delta }\left(Y^{[2]}\overline{t}+Z_\eta ^{[2]}\right)\right]^2+\left(\widehat{Y}^{[2]}\mathrm{\Delta }t\right)^2`$ (5.8)
$`+`$ $`\left(\mathrm{\Delta }t\mathrm{\Delta }Y^{[2]}\right)^2,`$
where $`\overline{t}`$ and $`\mathrm{\Delta }t`$ are the mean value and the uncertainty of the distribution deduced for the parameter $`t`$ by analyzing the evolution of the probability densities of observables in our quantum state. Obviously, the time uncertainty satisfies the fourth Heisenberg relation $`\mathrm{\Delta }t\mathrm{\Delta }H1/2`$.
Notice that the physical length uncertainty is again given by the sum of three positive terms. The analysis of the previous section can be easily extended to the case considered here. From the behavior of $`\mathrm{\Delta }L_\eta ^{[2]}`$ at large times we conclude that $`\mathrm{\Delta }Y^{[2]}`$ must vanish. Moreover, taking into account the assumption that the function $`H_0(\pi )`$ be invertible, remembering that $`H=g(H_0,\pi )`$, and using the implicit function theorem, it is possible to define $`Y^{[2]}`$ as a function of only $`H`$ -that we denote $`𝒴^{[2]}(H)`$\- provided that $`H_0^{}_{H_0}g+_\pi g0`$. One can then introduce the same two sets of hypotheses that were discussed in Sec. IV, but with the role of $`𝒴(H_0)`$ played by $`𝒴^{[2]}(H)`$. In this way one concludes that, under quite generic assumptions, an infinite resolution cannot be reached for the physical length in a non-perturbative quantization of the system constructed from the perturbative quantum theory.
Finally, we want to comment on the possibility that the system might admit a different non-perturbative quantization (with evolution still generated by the physical Hamiltonian) in which the canonically conjugate physical variables $`(X,p)`$ were promoted to explicitly time-independent operators and such that the quantum spectrum of the physical momentum $`p`$ were contained in its corresponding classical domain. This is non-trivial in general, and the viability of such a quantization cannot be taken for granted starting from the only assumption of the existence of a perturbative quantum description with the properties that we have discussed. From Eq. (3.1), we see that a situation in which this possibility is realized is when the physical energy does not depend on the pseudo momentum, $`_\pi g=0`$. In this case (which includes the example of the Einstein-Rosen waves), the physical position $`X`$ is independent of the background time. It may then be promoted to an operator that does not display any explicit time dependence, in terms of those for $`\pi ^i`$ and for the background coordinates $`q^i`$, the latter evolving only implicitly in the time parameter. Strictly speaking, nonetheless, the discussion presented in the paragraphs above cannot be applied in these circumstances because, with such an operator representation, $`Y^{[2]}(H_0,\pi )`$ must be identified with $`M^{[2]}(H_0,\pi )`$, the latter being identically zero when so is $`_\pi g`$ \[see Eqs. (5.1) and (5.3)\]. This vanishing invalidates the sets of hypotheses under which our study was carried out.
When a non-perturbative quantization with those characteristics exists, the Heisenberg uncertainty principle implies that $`\mathrm{\Delta }X\mathrm{\Delta }p1/2`$. As a consequence, the resolution in the physical position is limited if and only if the physical momentum is bounded from above. This happens in DSR1 and DSR2 theories, but not in DSR3. The same phenomenon occurs with the physical length if it is determined by the difference of two uncorrelated position observables. In conclusion, we see that the emergence of a minimum uncertainty in the physical length is unavoidable non-perturbatively as well as perturbatively, except perhaps for DSR3 theories that admit a non-perturbative quantization in which $`X`$ can be represented as an explicitly time-independent observable.
## VI First Order Corrections in the Perturbative Case
In this section we will study the physical length uncertainty that arises in the perturbative quantization when the operator $`\widehat{L}_\eta `$ is approximated up to first order corrections in the energy. To obtain this approximation, we expand the functions $`f`$ and $`g`$ (which we suppose smooth) in the variables $`H_0`$ and $`\pi `$ around their minimum values. Motivated by the case of free particles in special relativity, we assume that the minimum magnitude of the pseudo momentum is zero, whereas the minimum of the pseudo energy $`\mu `$ will be just non-negative PG . We then denote $`_0:=H_0\mu `$ and keep only up to quadratic terms in $`_0`$ and $`\pi `$ in the expansions of the two functions; this truncation will suffice for our purposes. In addition, we suppose that $`\mu `$ is small compared with the invariant energy/momentum scale of the DSR theory, so that the leading terms in the region of expansion are $`f(H_0,\pi )\pi `$ and $`g(H_0,\pi )H_0`$ (because the map $`U`$ determined by $`f`$ and $`g`$ must approach the identity in the low energy-momentum sector).
From Eq. (3.6), one then gets
$`M(H_0,\pi )`$ $``$ $`(_{H_0}_\pi g)|_0_0+(_\pi ^2g)|_0\pi ,`$
$`N(H_0,\pi )`$ $``$ $`1(_{H_0}_\pi f)|_0_0(_\pi ^2f)|_0\pi ,`$ (6.1)
where the symbol $`|_0`$ represents evaluation at $`_0=\pi =0`$. Substituting these results and the expression $`H_0(\pi )`$ of the background Hamiltonian in Eqs. (3.9) and (3.10) \[and recalling definitions (3.6)\], we deduce the first order approximation for the operators $`\widehat{Y}`$ and $`\widehat{Z}_\eta `$. An extrapolation of the situation found in special relativity PG leads us to consider the following cases.
1) *Massive* case: $`\mu 0`$, with $`H_0^{}|_{\pi =0}=0`$.
We obtain $`H_0(\pi )\mu +b\pi ^2`$, where $`2b:=H_0^{\prime \prime }|_{\pi =0}`$. Assuming that $`b>0`$, we have that $`\pi \sqrt{_0/b}`$. Thus, we can neglect terms proportional to $`_0`$ with respect to those linear in $`\pi `$. In this way, one finds
$`\widehat{Y}`$ $``$ $`\left[2b+(_\pi ^2g)|_0\right]\widehat{\pi },`$ (6.2)
$`\widehat{Z}_\eta `$ $``$ $`\eta (_\pi ^2f)|_0\widehat{s}_0,`$ (6.3)
where we have employed that $`s_0=\pi _jq^j|_{T=0}`$ is of the same order as $`\pi `$.
The function $`𝒴`$, defined in Sec. IV, is given in this approximation by the classical analog of Eq. (6.2) with $`\pi =\sqrt{_0/b}`$. The resulting function is strictly monotonic in $`H_0`$ if the constant coefficient $`2b+(_\pi ^2g)|_0`$ does not vanish, as it must happen if our truncation provides indeed the first order approximation. Therefore, the first set of hypotheses considered in Sec. IV is applicable in this case, leading us to the conclusion that it is impossible to achieve an infinite resolution in the physical length.
2) *Massless* case: $`\mu =0`$, with $`H_0^{}|_{\pi =0}=k0`$.
Now $`_0=H_0k\pi `$, so that corrections proportional to either $`H_0`$ or $`\pi `$ are of the same order. We then arrive at
$`\widehat{Y}`$ $``$ $`k+[{\displaystyle \frac{2b}{k}}(_\pi ^2f)|_0k(_{H_0}_\pi f)|_0`$ (6.4)
$`+{\displaystyle \frac{(_\pi ^2g)|_0}{k}}+(_{H_0}_\pi g)|_0]\widehat{H}_0,`$
$`\widehat{Z}_\eta `$ $``$ $`\eta \left[k(_{H_0}_\pi f)|_0+(_\pi ^2f)|_0\right]\widehat{s}_0.`$
The constant $`b`$ is defined as in the *massive* case. The next-to-leading order approximation to the function $`𝒴`$ is thus given by the classical counterpart of Eq. (6.4). Again, provided that the constant coefficient of the first order correction in $`H_0`$ differs from zero, the function $`𝒴`$ is strictly monotonic. The physical length uncertainty is hence greater than zero in this approximation.
## VII First Order Corrections: Behavior at Large Times
In this section, we will analyze in more detail the physical length uncertainty in the perturbative quantization for the *massless* case adopting the next-to-leading order approximation for low energies. We will pay a special attention to the behavior displayed at large values of the background time. We will show that this behavior is of the kind that was first discussed by Salecker and Wigner SalWig . Since a similar study was not considered in Ref. PG for the physical time uncertainty, we will carry out our analysis in a way that is also valid for it.
From the results of Ref. PG and our comments above, the physical time lapse $`t_\eta `$ is affected in the perturbative quantization by the uncertainty:
$$(\mathrm{\Delta }t_\eta )^2=[\mathrm{\Delta }(V\overline{T}+W_\eta )]^2+\widehat{V}^2(\mathrm{\Delta }T)^2+(\mathrm{\Delta }T\mathrm{\Delta }V)^2,$$
(7.1)
where the operators $`\widehat{V}`$ and $`\widehat{W}_\eta `$ have these expressions in the considered approximation for the *massless* case:
$`\widehat{V}`$ $``$ $`1+[k(_{H_0}^2f)|_0+(_{H_0}_\pi f)|_0`$ (7.2)
$`(_{H_0}^2g)|_0{\displaystyle \frac{(_{H_0}_\pi g)|_0}{k}}]\widehat{H}_0,`$
$`\widehat{W}_\eta `$ $``$ $`\eta \left[(_{H_0}_\pi f)|_0+k(_{H_0}^2f)|_0\right]\widehat{s}_0.`$
We then introduce the notation $`\{L_{\alpha ,\eta }\}:=\{t_\eta ,L_\eta \}`$, $`\{Y_\alpha \}:=\{V,Y\}`$, and $`\{Z_{\alpha ,\eta }\}:=\{W_\eta ,Z_\eta \}`$ to describe simultaneously the formulas for the physical time and length uncertainties. Let us emphasize that $`\alpha =0,1`$ is just an abstract subscript notation.
After a trivial elaboration, we can rewrite Eqs. (3.11) and (7.1) as
$`(\mathrm{\Delta }L_{\alpha ,\eta })^2`$ $`=`$ $`\overline{T}^2(\mathrm{\Delta }Y_\alpha )^2+(\mathrm{\Delta }Z_{\alpha ,\eta })^2+\overline{T}\mathrm{cov}(\widehat{Y}_\alpha ,\widehat{Z}_{\alpha ,\eta })`$ (7.3)
$`+\widehat{Y}_\alpha ^2(\mathrm{\Delta }T)^2+(\mathrm{\Delta }T\mathrm{\Delta }Y_\alpha )^2.`$
No sum over $`\alpha `$ is implied and
$$\mathrm{cov}(\widehat{Y}_\alpha ,\widehat{Z}_{\alpha ,\eta }):=\widehat{Y}_\alpha \widehat{Z}_{\alpha ,\eta }+\widehat{Z}_{\alpha ,\eta }\widehat{Y}_\alpha 2\widehat{Y}_\alpha \widehat{Z}_{\alpha ,\eta }.$$
(7.4)
In addition, in the studied approximation for the *massless* case, we can write the operators $`\widehat{Y}_\alpha `$ and $`\widehat{Z}_{\alpha ,\eta }`$ in the form $`\widehat{Y}_\alpha =\kappa _\alpha +\lambda _\alpha \widehat{H}_0/E_P`$ and $`\widehat{Z}_{\alpha ,\eta }=\eta \delta _\alpha \widehat{s}_0/E_P`$ \[see Eqs. (6.4) and (7.2)\], where $`E_P`$ is the Planck energy ($`E_P=1/\sqrt{G}`$ in our units, $`G`$ being Newton constant), $`\lambda _\alpha `$ and $`\delta _\alpha `$ are appropriate constant coefficients that differ from zero, $`\kappa _0:=1`$, and $`\kappa _1:=k=H_0^{}|_{\pi =0}`$.
The last term in Eq. (7.3) is then
$$(\mathrm{\Delta }T\mathrm{\Delta }Y_\alpha )^2=\frac{\lambda _\alpha ^2(\mathrm{\Delta }T\mathrm{\Delta }H_0)^2}{E_P^2}\frac{\lambda _\alpha ^2l_P^2}{4}.$$
(7.5)
In the last step, we have used the fourth Heisenberg relation for the background time and energy, and introduced the Planck length $`l_P=1/E_P`$ (in our units). Recalling that the other contributions to the physical uncertainty are positive, we conclude that $`\mathrm{\Delta }L_{\alpha ,\eta }|\lambda _\alpha |l_P/2`$. Therefore, we see that the uncertainty in both the physical time lapse and the physical length is bounded from below by a contribution of quantum gravitational origin that is of the order of the Planck length Luis ; Pad ; unc .
From the rest of contributions to the physical uncertainty (7.3), one gets in a similar way the bound
$`(\mathrm{\Delta }L_{\alpha ,\eta })^2`$ $`>`$ $`\lambda _\alpha ^2\overline{T}^2{\displaystyle \frac{(\mathrm{\Delta }H_0)^2}{E_P^2}}+{\displaystyle \frac{\widehat{Y}_\alpha ^2}{4(\mathrm{\Delta }H_0)^2}}+(\mathrm{\Delta }Z_{\alpha ,\eta })^2`$ (7.6)
$`+\overline{T}\mathrm{cov}(\widehat{Y}_\alpha ,\widehat{Z}_{\alpha ,\eta }).`$
The r.h.s. of this inequality can be regarded as a function of the uncertainty in the background energy $`\mathrm{\Delta }H_0`$, once the next-to-leading order expressions for the operators $`\widehat{Y}_\alpha `$ and $`\widehat{Z}_{\alpha ,\eta }`$ have been substituted. Hence, for uncertainties $`\mathrm{\Delta }H_0`$ in a certain interval, one can deduce a more general bound for $`\mathrm{\Delta }L_{\alpha ,\eta }`$ by minimizing that function. The extrema can be deduced by imposing the vanishing of the first derivative with respect to $`\mathrm{\Delta }H_0`$:
$`0`$ $`=`$ $`2\lambda _\alpha ^2\overline{T}^2{\displaystyle \frac{(\mathrm{\Delta }H_0)^4}{E_P^2}}{\displaystyle \frac{\widehat{Y}_\alpha ^2}{2}}+(\mathrm{\Delta }H_0)^3_\mathrm{\Delta }(\mathrm{\Delta }Z_{\alpha ,\eta })^2`$ (7.7)
$`+`$ $`{\displaystyle \frac{\mathrm{\Delta }H_0_\mathrm{\Delta }(\widehat{Y}_\alpha ^2)}{4}}+\overline{T}(\mathrm{\Delta }H_0)^3_\mathrm{\Delta }\mathrm{cov}(\widehat{Y}_\alpha ,\widehat{Z}_{\alpha ,\eta }).`$
Here, we have introduced the notation $`_\mathrm{\Delta }`$ to denote the derivative with respect to $`\mathrm{\Delta }H_0`$.
Provided that $`\widehat{Y}_\alpha `$ can be considered independent of both $`\mathrm{\Delta }H_0`$ and the (mean value of the) background time $`\overline{T}`$, the first two terms in the r.h.s. of Eq. (7.6) are in fact the kind of contributions that lead to the emergence of a minimum uncertainty of the Salecker and Wigner type (see Appendix A for details) SalWig ; BA . Namely, we get a contribution that is linear in $`(\mathrm{\Delta }H_0)^2`$ and another one that is proportional to its inverse. If these two terms were the only ones that appeared in our equations, an analysis similar to the standard one for Salecker-Wigner devices would prove that the bound for $`\mathrm{\Delta }L_{\alpha ,\eta }`$ reaches its minimum at a value of $`\mathrm{\Delta }H_0`$ that scales with the background time like $`\mathrm{\Delta }H_0^{min}1/\sqrt{\overline{T}}`$, whereas the lower bound obtained for the physical uncertainty at $`\mathrm{\Delta }H_0^{min}`$ increases in time like $`\sqrt{\overline{T}}`$.
Motivated by these remarks, we will now show that, at least in the region of small $`\mathrm{\Delta }H_0`$ and for large values of the background time $`\overline{T}`$, the terms in Eqs. (7.6) and (7.7) other than the first two ones do not invalidate the above conclusions about the existence of a (local) minimum and its associated bound. The restriction to small values of $`\mathrm{\Delta }H_0`$ is natural in the context of the low-energy approximation that we are discussing. Moreover, for unboundedly large times $`\overline{T}`$, the sector of vanishingly small values of $`\mathrm{\Delta }H_0`$ contains the relevant region for the analysis of the Salecker-Wigner bound on the uncertainty, i.e. the region around the minimum $`\mathrm{\Delta }H_0^{min}1/\sqrt{\overline{T}}`$.
In this sector of background energy uncertainties and time, one can demonstrate that a set of sufficient conditions to deduce a Salecker-Wigner behavior are:
a) $`\underset{\mathrm{\Delta }H_00}{lim}`$ $`\widehat{Y}_\alpha ^2=c_\alpha ^{(1)},`$
b) $`\underset{\mathrm{\Delta }H_00}{lim}`$ $`(\mathrm{\Delta }H_0)^2(\mathrm{\Delta }Z_{\alpha ,\eta })^2=c_\alpha ^{(2)},`$
c) $`\underset{\mathrm{\Delta }H_00}{lim}`$ $`(\mathrm{\Delta }H_0)^3_\mathrm{\Delta }(\mathrm{\Delta }Z_{\alpha ,\eta })^2=c_\alpha ^{(3)},`$
d) $`\underset{\mathrm{\Delta }H_00}{lim}`$ $`\mathrm{\Delta }H_0_\mathrm{\Delta }\widehat{Y}_\alpha ^2=0,`$
e) $`\underset{\mathrm{\Delta }H_00}{lim}`$ $`\mathrm{cov}(\widehat{Y}_\alpha ,\widehat{Z}_{\alpha ,\eta })=0,`$
f) $`\underset{\mathrm{\Delta }H_00}{lim}`$ $`\mathrm{\Delta }H_0_\mathrm{\Delta }\mathrm{cov}(\widehat{Y}_\alpha ,\widehat{Z}_{\alpha ,\eta })=0,`$ (7.8)
where $`c_\alpha ^{(n)},\text{ n=1,2,3}`$, are constants (with $`c_\alpha ^{(1)}2c_\alpha ^{(3)}0`$ and $`c_\alpha ^{(1)}+2c_\alpha ^{(2)}c_\alpha ^{(3)}0`$). Conditions a), b), and c) allow one to absorb the third term in the r.h.s. of Eqs. (7.6) and (7.7) just as a modification to $`\widehat{Y}_\alpha ^2`$ and treat this (square) expectation value as a constant when calculating the value of our function around its extrema in the region $`\mathrm{\Delta }H_01`$. In such a calculation and for sufficiently large background times, conditions d), e), and f) guarantee that all but the first three terms in Eqs. (7.6) and (7.7) can be neglected.
Taking into account that $`\widehat{Z}_{\alpha ,\eta }`$ vanishes when $`\eta =0`$, the only non-trivial requirements in that case are conditions a) and d). Regardless of the value of $`\eta `$, we prove in Appendix B that all the above conditions are satisfied at least for quantum states that are described by gaussian wave packets nota3 . Since we are assuming the feasibility of a (perturbative) quantization with canonical variables given by the background flat spatial coordinates and the pseudo momentum, and in addition we have focused our discussion on free systems, it seems reasonable to suppose that such states exist and provide the analog of classical particles in our quantum theory. Besides, the limitation to wave packets is already present in the deduction of the Salecker-Wigner bound for the spacetime uncertainty (in order to justify the assumption that the position and momentum operators have vanishing covariance) BA . So, it is natural to incorporate the same restriction to our analysis.
Substituting the values of the constants $`c_n`$ computed in Appendix B (under the simplifying assumption of only one spatial dimension), one obtains the following bounds for large background times from the corresponding minima in the region $`\mathrm{\Delta }H_01`$:
$$(\mathrm{\Delta }L_{\alpha ,\eta })^2>d_{\alpha ,\eta }l_p\overline{T},$$
(7.9)
where
$$d_{\alpha ,\eta }=\lambda _\alpha \left[\eta k^2\frac{\delta _\alpha ^2}{E_P^2}\nu ^2+\left(\kappa _\alpha +k\frac{\lambda _\alpha }{E_P}|\nu |\right)^2\right]^{\frac{1}{2}}.$$
(7.10)
Here, $`\nu `$ denotes the expectation value of the pseudo momentum.
In conclusion, in the perturbative quantization of free *massless* systems in DSR theories and within the low-energy approximation, we have seen that the physical time and length uncertainties are always bounded from below by a quantum gravitational contribution of the order of the Planck length, while for large values of the background time the uncertainties increase like $`\sqrt{l_P\overline{T}}`$ (at least for wave packets), just like in Salecker-Wigner devices.
## VIII Conclusion
In this work, we have analyzed the emergence of a minimum non-vanishing length uncertainty in the framework of a gravity’s rainbow formalism, derived from a dual realization of DSR theories in spacetime. This realization leads to a set of spacetime coordinates that are canonically conjugate to the physical energy and momentum. Therefore, the transformation from the background energy-momentum and spacetime coordinates (also called pseudo variables) to those that we consider as physical is provided by a canonical transformation. In particular, the physical spacetime variables are linear in the background ones, but in general depend nonlinearly on the pseudo energy and momentum of the particle.
We have specialized our analysis to systems that admit a Hamiltonian formulation, with the energy determined by the value of the Hamiltonian, and concentrated our attention on the case of a free dynamics, motivated by the consideration of DSR theories as (effective) descriptions of free particles in special relativity modified by gravity. In these free systems, the background Hamiltonian is a function of only the (magnitude of the) pseudo momentum. We have studied the behavior of the physical position, understanding as such the scalar obtained by projecting the physical position vector in the momentum direction. More specifically, we have investigated the quantum uncertainty that affects the physical length, defined by the difference between this physical position and the initial value of the position, either in the background or in the physical variables of the system. This study has been carried out in two possible quantization schemes, referred as perturbative and non-perturbative quantizations.
The perturbative approach corresponds to a quantization in which the evolution is generated by the background Hamiltonian, so that the background time $`T`$ plays the role of evolution parameter. We have assumed that a quantum theory of this kind is feasible. In this quantization, the physical time and length are represented by genuine operators that depend explicitly on the time parameter. We have been able to generalize the analysis of Ref. PG for the physical time uncertainty, and prove that the uncertainty in the physical length is also strictly positive in this approach.
Rigorously speaking, we have demonstrated this positivity under two different sets of generic assumptions. Both sets contain the more than reasonable hypothesis that the considered quantum state has a finite expectation value of the background energy, $`\widehat{H}_0<\mathrm{}`$. Besides, the two sets include an assumption about the functional dependence of the background energy on the pseudo momentum, namely, that the function $`H_0=H_0(\pi )`$ be invertible. The rest of hypotheses concern the detailed form of the DSR theory, and more concretely the properties of the function $`𝒴(H_0):=Y[H_0,\pi (H_0)]`$ introduced in Sec. IV.
One set of assumptions requires this function to be strictly monotonic, i.e. $`𝒴^{}(H_0)0`$ for all values of $`H_0`$.
The other set involves several requirements. The most important ones are: i) the positivity of $`𝒴`$, $`𝒴0`$; and ii) a linear or faster increase of $`𝒴`$ with $`H_0`$ at infinity, $`lim_{H_0\mathrm{}}(𝒴/H_0)>r`$ for a certain constant $`r>0`$. In addition, it is demanded that: iiia) the kernel of $`𝒴`$ be empty, or either iiib) this kernel consist of a single point $`\overline{H}_0`$ where the derivative of $`𝒴`$ does not vanish, $`(d𝒴/dH_0)|_{\overline{H}_0}0`$.
In the non-perturbative approach, on the other hand, the evolution is generated by the physical Hamiltonian, and the physical time $`t`$ is identified with the evolution parameter. Starting with the perturbative quantization that we have assumed to exist, it is in general possible to construct a non-perturbative quantum theory of this kind, in which the physical length is represented by an operator that depends explicitly on the time parameter $`t`$. We have proved that the quantum uncertainty in this operator is strictly positive under similar sets of assumptions to those discussed for the case of the perturbative quantization. Therefore, it is again impossible to reach an infinite resolution in the physical length.
It might also happen that the system admits a different non-perturbative quantization in which the evolution is indeed generated by the physical Hamiltonian, but the physical position variable gets promoted to an operator that is explicitly independent of time and canonically conjugate to the operator which represents the magnitude of the physical momentum. In general, the existence of such a quantum theory is not granted from the sole assumption of the viability of the perturbative quantization. Supposing besides that the quantum spectrum of the physical momentum is contained in its classical domain, Heisenberg principle implies that the uncertainty in the physical position can be made to vanish only if the physical momentum is not bounded from above. The same result holds for the physical length if it is determined by the difference of two uncorrelated physical positions.
The existence of an upper bound for the physical momentum, with the consequent limit in the spatial resolution, occurs only in the DSR1 and DSR2 families, but not in DSR3 theories. Remarkably, for such theories the physical time uncertainty is always bounded away from zero in the non-perturbative quantum theory PG . As a result, it is never possible to reach an infinite resolution, both in the physical time and position, in the non-perturbative quantization of Hamiltonian free systems within the context of DSR theories.
Finally, we have also analyzed the uncertainty in the perturbative quantization when the operator corresponding to the physical length is approximated up to first order corrections in the energy. The study has lend support to the conclusion that this uncertainty is generically greater than zero. Special attention has been paid to the *massless* case, in which the background energy is proportional to the magnitude of the pseudo momentum in the considered approximation. We have proved that, in that case, the uncertainty is always bounded by a quantity of the order of the Planck length. This bound can be interpreted as a contribution of quantum gravitational origin. In addition we have proved that, in the low-energy regime and for large values of the background time, the uncertainties in the physical time and length admit lower bounds that increase with the square root of time. This is precisely the kind of behavior that was suggested by Salecker and Wigner for spacetime measurements made with quantum devices.
###### Acknowledgements.
The authors want to thank B. Barceló and J. Cortés for helpful conversations. P.G. acknowledges CSIC and the European Social Fund for the financial support provided by an I3P grant. This work was supported by funds provided by the Spanish MEC-MCYT projects BFM2002-04031-C02-02 and FIS2004-01912.
## Appendix A: Salecker-Wigner Devices
In this appendix we will briefly summarize the rationale of Salecker and Wigner about the quantum uncertainty that is inherent to the measurement of spacetime distances SalWig ; BA . The analysis starts with the consideration of a measurement device, regarded as a free system with mass $`m`$ and uncertainties in its initial position and momentum $`\mathrm{\Delta }q`$ and $`\mathrm{\Delta }\pi `$. The (square) uncertainty in its position at a later instant of time $`t`$ is
$`[\mathrm{\Delta }q(t)]^2`$ $`=`$ $`\left[\mathrm{\Delta }(q+{\displaystyle \frac{t}{m}}\pi )\right]^2`$
$`=`$ $`(\mathrm{\Delta }q)^2+{\displaystyle \frac{t^2}{m^2}}(\mathrm{\Delta }\pi )^2+{\displaystyle \frac{t}{m}}\mathrm{cov}(\widehat{q},\widehat{\pi }),`$
where $`\mathrm{cov}(\widehat{q},\widehat{\pi }):=\widehat{q}\widehat{\pi }+\widehat{\pi }\widehat{q}2\widehat{q}\widehat{\pi }.`$ This expression gets simplified when the (initial) position and momentum observables are not correlated. This occurs, for instance, if the states of the system are plane waves modulated by a Gaussian. In that case $`\mathrm{cov}(\widehat{q},\widehat{\pi })=0`$. Making use of the fourth Heisenberg relation, one then obtains the inequality
$$[\mathrm{\Delta }q(t)]^2\frac{t^2}{m^2}(\mathrm{\Delta }\pi )^2+\frac{1}{4(\mathrm{\Delta }\pi )^2}.$$
(A.1)
The r.h.s. of this equation can be viewed as a function of $`\mathrm{\Delta }\pi `$. Its extrema can be determined by imposing the vanishing of the first derivative:
$$0=\frac{4t^2}{m^2}(\mathrm{\Delta }\pi )^41.$$
The minimum value of the uncertainty is hence reached at $`\mathrm{\Delta }\pi ^{min}=\sqrt{m/(2t)}`$. Substituting this value in (A.1) one gets a lower bound for the position uncertainty at the instant $`t`$:
$$\mathrm{\Delta }q(t)\sqrt{\frac{t}{m}}.$$
Therefore, the arguments of Salecker and Wigner imply that the uncertainty increases with the square root of time.
## Appendix B: Calculations for Wave Packets
This appendix contains the calculation of the mean values, uncertainties and covariance of the operators $`\widehat{Y}_\alpha `$ and $`\widehat{Z}_{\alpha ,\eta }`$ introduced in Sec. VII, adopting the next-to-leading order approximation for low energies and restricting the quantum states to be gaussian wave packets (in the free quantum theory with elementary variables given by the background spatial coordinates and momenta). Moreover, in order to simplify our calculations, we will carry out our analysis not in three, but just in one spatial dimension. We do not expect this reduction to qualitatively affect our results.
Explicitly, we will adopt a standard momentum representation in one dimension, with wave packets given by the following wave functions nota3 :
$$\mathrm{\Psi }(\pi _1)=\frac{1}{(2\mathrm{\Pi }\sigma ^2)^{1/4}}e^{(\pi _1\nu )^2/(4\sigma ^2)}e^{i\mu \pi _1}.$$
Here, $`\nu :=\widehat{\pi _1}`$, $`\sigma :=\mathrm{\Delta }\pi _1`$, and $`\mu :=\widehat{q^1}`$, with $`q^1`$ being the initial background position (we obviate its subscript $`0`$ to simplify the notation). The number $`\mathrm{\Pi }`$ is denoted in this appendix with a capital Greek letter in order to distinguish it from the magnitude of the pseudo momentum $`\pi `$. Besides, note that in one dimension $`\pi =|\pi _1|`$.
From the functional form of the wave packets, it is clear that the quantities that we want to compute will depend on the parameters $`\mu `$, $`\nu `$, and $`\sigma `$. So, in order to calculate the limiting values (VII), we need to express the limit $`\mathrm{\Delta }H_00`$ in terms of those parameters. In the studied approximation, $`H_0=k\pi `$ for the *massless* case, and a trivial calculation shows that the uncertainty $`\mathrm{\Delta }H_0`$ for the wave packets is given by
$`(\mathrm{\Delta }H_0)^2`$ $`=`$ $`k^2(\mathrm{\Delta }\pi )^2=k^2\left(\sigma ^2+\nu ^2\widehat{\pi }^2\right)`$ (B.1)
$`:=`$ $`G^2(\sigma ,\nu ),`$
$`\widehat{\pi }`$ $`=`$ $`|\nu |\mathrm{erf}\left({\displaystyle \frac{|\nu |}{\sqrt{2}\sigma }}\right)+\sqrt{{\displaystyle \frac{2}{\mathrm{\Pi }}}}\sigma e^{\nu ^2/(2\sigma ^2)}.`$ (B.2)
It is worth emphasizing that $`\widehat{\pi }`$, the expectation value of the magnitude of the pseudo momentum, differs in general from $`\nu `$. We have introduced the error function
$$\mathrm{erf}(x)=\frac{2}{\sqrt{\mathrm{\Pi }}}_0^x𝑑ye^{y^2},\text{ with}\underset{x\mathrm{}}{lim}\mathrm{erf}(x)=1.$$
From the above equations, we see that $`\widehat{\pi }|\nu |`$ and $`\mathrm{\Delta }H_0k\sigma `$ for small uncertainties $`\mathrm{\Delta }H_0`$. Via the implicit function theorem, we can then use the relation $`\mathrm{\Delta }H_0=G(\sigma ,\nu )`$ ($`G`$ being the square root of $`G^2`$) to define $`\sigma `$ as a function of $`\mathrm{\Delta }H_0`$ in a neighborhood of the origin of these quantities, provided that $`_\sigma G`$ does not vanish there. Actually, one has that $`lim_{\sigma 0}_\sigma G=k0`$. Therefore, one is allowed to replace the limit $`\mathrm{\Delta }H_00`$ with $`\sigma 0`$. In addition, one can substitute the partial derivative with respect to $`\mathrm{\Delta }H_0`$ (i.e., $`_\mathrm{\Delta }`$) by $`_\mathrm{\Delta }\sigma _\sigma `$, where $`lim_{\sigma 0}_\mathrm{\Delta }\sigma =1/k`$. These considerations lead to the results given in the rest of this appendix, where we analyze simultaneously the cases of the physical time and length uncertainties.
In the first order approximation for the *massless* case, the operators $`\widehat{Y}_\alpha `$ and $`\widehat{Z}_{\alpha ,\eta }`$ adopt expressions of the form \[see Eqs. (6.4) and (7.2)\]:
$`\widehat{Y}_\alpha `$ $`=`$ $`\kappa _\alpha +k{\displaystyle \frac{\lambda _\alpha }{E_P}}\widehat{\pi },`$
$`\widehat{Z}_{\alpha ,\eta }`$ $`=`$ $`\eta {\displaystyle \frac{\delta _\alpha }{E_P}}\widehat{s}_0=\eta {\displaystyle \frac{\delta _\alpha }{2E_P}}\left(\widehat{\pi _1}\widehat{q^1}+\widehat{q^1}\widehat{\pi _1}\right),`$
where $`\lambda _\alpha `$ and $`\delta _\alpha `$ are certain non-vanishing constants, $`\eta `$ can take the values 0 or 1, $`\kappa _0=1`$, and $`\kappa _1=k`$. We have employed that in this approximation $`\widehat{H}_0=k\widehat{\pi }`$.
A straightforward calculation along the lines explained above shows that for wave packets
$$\underset{\mathrm{\Delta }H_00}{lim}\widehat{Y}_\alpha ^2=\underset{\sigma 0}{lim}\widehat{Y}_\alpha ^2=\left(\kappa _\alpha +k\frac{\lambda _\alpha }{E_P}|\nu |\right)^2:=c_\alpha ^{(1)}.$$
In the same way, one finds
$`\mathrm{\Delta }H_0_\mathrm{\Delta }\widehat{Y}_\alpha ^2`$ $`=`$ $`2k{\displaystyle \frac{\lambda _\alpha }{E_P}}(\kappa _\alpha +k{\displaystyle \frac{\lambda _\alpha }{E_P}}\widehat{\pi })_\sigma \widehat{\pi }\mathrm{\Delta }H_0_\mathrm{\Delta }\sigma ,`$
$`\mathrm{\Delta }H_0_\mathrm{\Delta }\sigma `$ $`=`$ $`{\displaystyle \frac{\sigma ^2+\nu ^2\widehat{\pi }^2}{\sigma \widehat{\pi }_\sigma \widehat{\pi }}}.`$ (B.3)
From Eq. (B.2) one can check that $`_\sigma \widehat{\pi }`$ tends fast enough to zero when $`\sigma 0`$ ($`\mathrm{\Delta }H_00`$) as to guarantee that
$$\underset{\mathrm{\Delta }H_00}{lim}\mathrm{\Delta }H_0_\mathrm{\Delta }\widehat{Y}_\alpha ^2=0.$$
On the other hand, a similar computation leads to the following uncertainty for the operator $`\widehat{Z}_{\alpha ,\eta }`$
$`(\mathrm{\Delta }Z_{\alpha ,\eta })^2`$ $`=`$ $`\eta {\displaystyle \frac{\delta _\alpha ^2}{E_P^2}}\left(\widehat{s}_0^{\mathrm{\hspace{0.33em}2}}\widehat{s}_0^2\right)`$
$`=`$ $`\eta {\displaystyle \frac{\delta _\alpha ^2}{E_P^2}}\left({\displaystyle \frac{\nu ^2}{4\sigma ^2}}+\mu ^2\sigma ^2+{\displaystyle \frac{1}{2}}\right).`$
From this and Eqs. (B.1) and (B.3), it is not difficult to prove that
$`\underset{\mathrm{\Delta }H_00}{lim}(\mathrm{\Delta }H_0)^2(\mathrm{\Delta }Z_{\alpha ,\eta })^2`$ $`=`$ $`\eta k^2{\displaystyle \frac{\delta _\alpha ^2}{4E_P^2}}\nu ^2:=c_\alpha ^{(2)},`$
$`\underset{\mathrm{\Delta }H_00}{lim}(\mathrm{\Delta }H_0)^3_\mathrm{\Delta }(\mathrm{\Delta }Z_{\alpha ,\eta })^2`$ $`=`$ $`\eta k^2{\displaystyle \frac{\delta _\alpha ^2}{2E_P^2}}\nu ^2:=c_\alpha ^{(3)}.`$
Finally, the covariance of $`\widehat{Y}_\alpha `$ and $`\widehat{Z}_{\alpha ,\eta }`$ is given by
$$\mathrm{cov}(\widehat{Y}_\alpha ,\widehat{Z}_{\alpha ,\eta })=\eta k\frac{\lambda _\alpha \delta _\alpha }{E_P^2}\left(\widehat{\pi }\widehat{s}_0+\widehat{s}_0\widehat{\pi }2\widehat{\pi }\widehat{s}_0\right)$$
which for wave packets gives
$$\mathrm{cov}(\widehat{Y}_\alpha ,\widehat{Z}_{\alpha ,\eta })=2\eta k\frac{\lambda _\alpha \delta _\alpha }{E_P^2}\mu \sigma ^2\mathrm{sign}(\nu )\mathrm{erf}\left(\frac{|\nu |}{\sqrt{2}\sigma }\right).$$
Therefore, one can check that
$`\underset{\mathrm{\Delta }H_00}{lim}\mathrm{cov}(\widehat{Y}_\alpha ,\widehat{Z}_{\alpha ,\eta })=0,`$
$`\underset{\mathrm{\Delta }H_00}{lim}\mathrm{\Delta }H_0_\mathrm{\Delta }\mathrm{cov}(\widehat{Y}_\alpha ,\widehat{Z}_{\alpha ,\eta })`$
$`=`$ $`\underset{\sigma 0}{lim}\mathrm{\Delta }H_0_\mathrm{\Delta }\sigma _\sigma \mathrm{cov}(\widehat{Y}_\alpha ,\widehat{Z}_{\alpha ,\eta })=0.`$
In conclusion, we see that conditions (VII) are satisfied. |
warning/0507/physics0507128.html | ar5iv | text | # Low energy electron scattering from DNA and RNA bases: shape resonances and radiation damage
## I Introduction
Reactions induced by electrons drive nearly all the important chemical processes in radiation chemistry, plasma etching in semiconductors, stability of waste repositories, and are also fundamental in the dynamics of the atmosphere and interstellar clouds, with processes such as dissociative recombination and electron attachment.
In recent years, an increasing importance has been recognized to these processes in biological environments, especially in relation to radiation damage to nucleic acids (DNA and RNA). These processes consist in the interaction of ionizing radiation (like $`\alpha `$, $`\beta `$ and $`\gamma `$-rays) with living tissue, generating possibly mutagenic and carcinogenic byproducts, through a wide variety of ionization, excitation and energy transfer processes, that can interest many molecular species in the complex cell environment.
The important work of Sanche and coworkers Boudaiffa et al. (2000); Levesque et al. (2003); Sanche (2002) has shown that damage to nucleic acids from ionizing radiation Mozejko and Sanche (2003) (single and double strand breaks in particular) can be generated through a mechanism involving low energy electron attachment to the nucleic acid and subsequent bond breaking due to energy transfer to a vibrational mode of the temporary anion formed in the electron capture step. These low-energy secondary electrons are generated by electron-impact ionization caused by high energy electrons, originally produced directly by the ionizing radiation. In the electron-impact ionization process, the scattered electron loses part of its kinetic energy, while another electron is ejected, with energy much lower than the first one.
In the past few years many studies have been devoted to understanding the mechanism for the action of the low-energy electrons and their capability to cause strand breaks. Barrios et al. (2002); Gu et al. (2005); Scheer et al. (2004); Martin et al. (2004); Aflatooni et al. (1998) A first general feature on which there is wide agreement is that the electron capture is mainly due to the DNA and RNA bases. These molecules have extended aromatic systems, therefore there is a wide range of low-lying unoccupied $`\pi ^{}`$ orbitals where an electron can be captured, giving rise to a shape resonance, a temporary anion, in the range of energies between 0 and 15 eV, where the experiments have found signatures of electron-induced damage to nucleic acids.
The simplest of these bases are thymine, cytosine, uracil (pyrimidines, monocyclic) and adenine and guanine (purines, bicyclic and generally larger than pyrimidines). Their structures are shown in Fig. 1. In this paper we will present theoretical predictions of cross sections for elastic electron scattering from these large molecules. Determination of the location, width, and electronic structure of resonances for a single target molecule is an important step towards understanding and possibly modeling the complex dynamics of DNA, which consists of multiple components. Specifically, besides the bases, there are also the sugar backbone, the phosphates, and also the solvation water,Blackburn and Gait (1996) that probably plays a major role in stabilizing the temporary anions. Sommerfeld (2004); Grandi et al. (2004) No previous theoretical or experimental study of low-energy electron scattering from all DNA bases is available for comparison (although a study at intermediate energy has been carried out recently Mozejko and Sanche (2003)), but our method has proven its reliability in the study of small molecules Tonzani and Greene (2005a) and more extended systems like C<sub>60</sub>, SF<sub>6</sub> and XeF<sub>6</sub>. Tonzani and Greene (2005b)
Much experimental work also has been carried out on dissociative electron attachment from DNA bases,Abdoul-Carime et al. (2000); Gohlke et al. (2003); Martin et al. (2004) to understand what fragments are generated. We will discuss the possible connections with measured dissociation branching ratios that we can infer from the examination of the spatial shape and nodal surfaces of the resonant wavefunctions.
## II Theory
A detailed description of our method is available in Ref. Tonzani and Greene, 2005a. For this reason we will only review here the main points and the changes we have implemented since that work, notably a new polarization-correlation potential Lee et al. (1988) which permits enhanced predictive capabilities.
The interaction between the $`N`$ electrons in a molecule and the scattered electron is a many body problem that can be recast, Morrison and Collins (1978) using the so called static exchange approximation, into a one body problem for the continuum electron with a nonlocal potential: Tonzani and Greene (2005a)
$$(\frac{1}{2}^2+V_sE)\varphi _0(\stackrel{}{r})=\underset{i,j}{}c_ic_j\underset{k=1}{\overset{N}{}}\varphi _{ki}(\stackrel{}{r})𝑑\stackrel{}{r^{}}\frac{\varphi _{kj}^{}(\stackrel{}{r^{}})\varphi _0(\stackrel{}{r^{}})}{\stackrel{}{r}\stackrel{}{r^{}}}$$
(1)
here $`V_s`$ is the electrostatic potential, $`\varphi _0`$ is the scattered electron orbital, while the other orbitals refer to target electrons, and the $`c_i`$ are configuration interaction (CI) coefficients. In this approximation only one state for the target is considered (the ground state), whereby it is only suitable to describe electronically-elastic processes. The nonlocal interaction consists of three main pieces: the direct electrostatic contribution, the exchange term and the polarization-correlation potential. Of these three, only the first is a local potential.
To describe electron scattering from a general, possibly very complicated molecule, we use the R-matrix method, Greene et al. (1996) which partitions space into two regions: an internal region within which all the short-range physics is confined and an outer region where only long-range interactions (like Coulomb or dipole potentials) are important. Our calculation begins from a variational principle Greene et al. (1996) for the logarithmic derivative of the wave function at the boundary between the two regions
$$b\frac{\mathrm{log}(r\mathrm{\Psi })}{r}=2\frac{_V\mathrm{\Psi }^{}(E\widehat{H}\widehat{L})\mathrm{\Psi }𝑑V}{_V\mathrm{\Psi }^{}\delta (rr_0)\mathrm{\Psi }𝑑V}$$
(2)
where $`\widehat{L}`$ is the Bloch operator, Tonzani and Greene (2005a) and $`r_o`$ is the boundary between the internal and external regions. It is possible, after expanding the internal region wavefunction in a suitable basis set, to recast the solution of Eq. 2 as an eigenvalue problem:
$$\underset{¯}{\mathrm{\Gamma }}\stackrel{}{C}=(E\underset{¯}{H}\underset{¯}{L})\stackrel{}{C}=\underset{¯}{\mathrm{\Lambda }}\stackrel{}{C}b.$$
(3)
The external region is then treated matching the solution of Eq. 3 to the exact wavefunctions for the long-range tail of the molecular potential. We show in Sec. II.3 how the contribution from a long range dipole field can be included in our method. The basis set we use for the internal region of the $`R`$-matrix is a product of finite element cubic polynomials in all three dimensions, using a grid in spherical coordinates.
### II.1 Local Density Approximation (LDA)
To simplify further the description of our system we have to deal with the nonlocality inherently present in the potential. To do this we use a local density approximation for the exchange potential, which reduces it to a functional only of the local density:
$$V_{ex}(\stackrel{}{r})=\frac{2}{\pi }k_FF(k_F,E),$$
(4)
where $`k_F`$ is the local Fermi momentum:
$$k_F(\stackrel{}{r})=(3\pi ^2\rho (\stackrel{}{r}))^{1/3}$$
(5)
and $`F`$ is a functional of the energy and the local density $`\rho (\stackrel{}{r})`$ (through the local Fermi momentum). The functional form we use for $`F`$ is called the Hara exchange. Hara (1969) It has been extensively employed in continuum state calculations, and it is energy-dependent. The local exchange approximation, widely used also in density functional calculations, Parr and Yang (1989) has proven itself to give qualitatively correct results, Tonzani and Greene (2005a); Morrison and Collins (1978) while being sufficiently simple to implement computationally that it permits an exploration of complex molecular species.
### II.2 DFT Polarization potential
We have recently added to our computer code the capability to use a parameter free correlation-polarization potential,Gianturco and Rodriguez-Ruiz (1993); Colle and Salvetti (1983) based on density functional theory (DFT) ideas. As shown in Ref. Lane, 1980 the polarization-correlation contribution is physically related to the distortion-relaxation effect on the molecule generated by the incoming electron. This is extremely important for an accurate description of the scattering process. The long range part of this potential is a simple multipole expansion, of which we retain only the induced dipole polarization terms:
$$V_{pol}=\frac{1}{2r^4}(\alpha _0+\alpha _2P_2(\mathrm{cos}\theta ))$$
(6)
where $`\alpha _0`$ and $`\alpha _2`$ are the totally symmetric and nontotally symmetric components of the polarizability tensor, and are calculated $`ab`$ $`initio`$ using electronic structure codes.
In the volume where the electronic density of the target is not negligible, this potential is nonlocal. The interaction can be approximated again as a local potential, different forms of which have been suggested in the literature. The one we use is based on DFT (specifically on the LYP potential of Ref. Lee et al., 1988) and it has yielded reliable results in the work of Gianturco and coworkers. Lucchese et al. (1999) This form makes use of the electron density, its gradient and laplacian, which have to be calculated for each target molecule. The short and long range potentials are matched unambiguously at the innermost crossing point, whose radius is dependent on the angles.
### II.3 Dipole physics
Since the molecules we considered in this work have large electric dipole moments, there is a need to consider the long range effect of the dipole field on the scattered electron. Two possible options might be considered, either extending the boundary of the R-matrix box far out to a region where the dipole potential is very small, which would be extremely time-consuming for our calculations, or matching to outer region functions adapted to the dipole interaction. We choose this second route and, following the example of Clark, Clark (1979) we define the matrix of the operator:
$$(l^22D\mathrm{cos}\theta )\mathrm{\Omega }_N=N(N+1)\mathrm{\Omega }_N$$
(7)
where $`l`$ is the angular momentum operator, $`\theta `$ is the angle between the incoming electron and the dipole direction, $`D`$ is the dipole moment, $`N(N+1)`$ and $`\mathrm{\Omega }_N`$ are eigenvalues and eigenfunctions. We expand $`\mathrm{\Omega }_N`$ in a basis of spherical harmonics to diagonalize the system in Eq. 7. The order of the spherical Bessel functions that are matched in the outer region will be now $`N`$ (not an integer, in general) rather than the usual orbital angular momentum quantum number $`l`$. Since the dipole moments of the molecules in question are very large, the dipole plus centrifugal potential becomes attractive for the first few channels which may thus have a complex $`N`$. In such cases we define $`N=1/2+i\mu `$ and the matching functions will become:Clark (1979)
$$\overline{j}_N(kr)=\sqrt{\frac{\pi }{2r}}\frac{1}{\mathrm{sinh}\frac{1}{2}\pi \mu }Im(J_{i\mu }(kr))$$
(8)
$$\overline{n}_N(kr)=\sqrt{\frac{\pi }{2r}}\frac{1}{\mathrm{cosh}\frac{1}{2}\pi \mu }Re(J_{i\mu }(kr))$$
(9)
where $`J_{i\mu }`$ is a cylindrical Bessel function. This allows us to keep the functions in Eq. 8 always real, and therefore have real $`K`$-matrices.
It should be mentioned, however, that at extremely low energies these functions oscillate rapidly in energy as $`\mathrm{sin}(\mu \mathrm{log}kr)`$, giving rise to $`K`$-matrices that are not smooth functions of energy. Defining the base pair as in Ref. Greene et al., 1979 solves the problem, but since we are not interested in energies below about 0.5 eV in this study, the functions in Eqs. 8-9 will be sufficient.
The dipole plus centrifugal potential is attractive if the value of the dipole moment is larger than a critical value ($`D_c`$=1.625 Debye for a nonrotating dipole). In this case the dipole interaction can bind the electron all by itself. In general, when rotation is included, the critical value of the dipole moment to have a bound state is aroundJordan and Wang (2003) 2-2.5 D and the number of dipole-bound states is finite. In the case of uracil, Dolgounitcheva et al. (1999); Sommerfeld (2004); Scheer et al. (2004) such a dipole-bound state exists at roughly 0.1 eV below the neutral ground state energy, at the equilibrium geometry of the target molecule.
### II.4 Calculation details
To calculate the target properties we have used the GAUSSIAN 98 program suite, at the Hartree-Fock (HF) level of theory. We have noticed in the past Tonzani and Greene (2005a) that CI calculations are in this case much more expensive and make comparatively little difference in the final cross sections. All of the molecules we treat here are spin singlets in their ground state. For the scattering calculations we have used an IBM Power 4 supercomputer, each calculation taking about 6 wall clock hours on 16 processors working in parallel. The size of the matrices generated is about 180000 by 180000 and the direct solution of the linear system requires about 10 minutes per energy point when distributed over 16 processors. The matrices are very sparse (fewer than 0.5% nonzero elements), and we use a direct sparse factorization method to solve the linear system. The convergence of the calculation is such that incrementing the number of sectors by 30% lowers the energy of the resonances by a further 0.1 eV in uracil; since it is already quite cumbersome to carry out these calculations we deemed this level of convergence as sufficient for the purposes of this study. The geometry of the molecules is chosen to be the equilibrium target geometry, optimized at the HF level with a 6-31G\* basis set.
## III Results
To our knowledge there are no available experimental data or calculations of low energy electron scattering from the complete set of DNA bases. A study of electron attachment has been presented in Ref. Aflatooni et al., 1998 and the resonance positions are clearly marked. Compared to these results, our calculations show resonances shifted typically by about 2 eV higher in energy, but the energy spacing of the resonances is comparable to what is observed in the experiment.Moreover the relative values of the widths of successive resonances resemble the measured widths. There is also a theoretical study at intermediate energies, Mozejko and Sanche (2003) and calculations for scattering from uracil; Grandi et al. (2004); Gianturco and Lucchese (2004) in the following we compare these results to ours.
We have already mentioned that the heterocyclic DNA bases have many low-lying unoccupied orbitals, so it is not surprising that their elastic cross sections for electron scattering exhibit many shape resonances. These may be viewed as a capture of the scattered electron into one of these antibonding orbitals to form a short lived negative ion state.Loomba et al. (1981); Dehmer (1984)
Since all these molecules have, in their equilibrium configuration, only one symmetry element - reflection through the molecular plane - we will characterize the resonances as being of $`\sigma `$ type (no node in the plane) or $`\pi `$ type (when they have instead a node in the plane) rather than using the $`A^{}`$ and $`A^{\prime \prime }`$ labels as is customary for the $`C_s`$ group.
### III.1 Positions and widths of resonances
A general comparison of partial elastic cross sections for all five of these molecules is shown in Fig. 2, while in the following we give a more detailed description and compare with information available in the literature. Also, a plot of total time-delays (see also Sec. III.2 for details) is provided in Fig. 3 to show the resonances is more detail.
Since we are dealing with polar molecules, applying the fixed-nuclei approximation as it stands makes the partial wave expansion of the forward scattering amplitude divergent. Due to the long-range nature of the dipole interaction, in fact, all partial waves would contribute to the scattering process, causing an infinite scattering in the forward direction and therefore infinite integral cross sections. There is a method, extensively discussed in the literature, Gianturco et al. (1998); Clark (1977) to deal with this problem by means of a Born closure formula, which yields a finite integral cross section once molecular rotations are included. We will not pursue this further, since existing experiments are not likely to deal with such detailed rotational structures. Therefore our cross sections and time-delays include only up to $`l_{max}=10`$ and omit all higher partial waves. The correction would be proportional to the dipole moment and inversely proportional to the smallest rotational spacing. For the DNA bases the dipole moment is large, while the rotational spacing is small. Therefore the correction can be quite large especially at very low energy. The correction would thus tend to mask the resonant structures, which are the most interesting observables and which have been measured in experiments. All of the calculated cross sections grow rapidly when the incident electron energy decreases below 1 eV, which is a signature of the role played by the dipole field in pulling in the electron and which is very common in electron scattering from polar molecules. Gorfinkiel et al. (2002); Aflatooni et al. (1998)
A comparison o four resonance patterns with the electron transmission spectroscopy (ETS) data of Burrow and coworkers Aflatooni et al. (1998); Scheer et al. (2004) can be found in Figs. 10-12, where the time-delay data from our calculations has been rescaled by an overall constant and shifted down in energy to facilitate the comparison of energies, widths and spacing. All of the resonances obtained in our calculations are listed in Table 1.
#### III.1.1 Uracil
In the cross section of uracil we find 3 resonances, at 2.16 eV (of width 0.2 eV), at 5.16 (0.6 eV wide) and a very broad resonance at 7.8 eV. The resonance at 2.16 eV is dominated by the $`l=3`$ partial wave (50%) and has contributions from $`l=1`$ (35%) and $`l=2`$ (11%). at 5.16 eV the main partial wave is d (66%), at 7.8 eV f-wave is the dominant contribution (64%).
In the work of Gianturco $`et`$ $`al.`$ (see Ref. Gianturco and Lucchese, 2004) three $`\pi ^{}`$ resonances are found at energies of 2.2, 3.5 and 6.5 eV. The second and third $`\pi ^{}`$ resonances from that work fall at lower energies than ours, a somewhat surprising discrepancy since the theoretical models are very similar.
The contribution of the dipole field at distances larger than 12 Bohr is neglected in Ref. Gianturco and Lucchese, 2004, but we have noticed that this influences only the overall magnitude of the cross sections (roughly an increase of 20% at very low energy, that is reduced to about 5% around 10eV), the dipole physics only weakly affects the resonance positions and widths.
Resonances are measured Ref. Aflatooni et al., 1998 to occur at 0.3, 1.5 and 3.8 eV. They are all assigned as $`\pi `$ resonances, Scheer et al. (2004) so our results should be shifted by about 2 eV down, whereas the spacing between the resonances is larger than experiment. The relative resonance widths are similar to Ref. Scheer et al., 2004, in that the first resonance is very narrow, the second broader and the third very broad, a comparison is shown in Fig. 10, where an integration of the experimental data has been performed to show more clearly the resonance positions and widths.
#### III.1.2 Cytosine
For cytosine we find 3 main resonances, a very sharp one at 1.7 eV (width 0.5 eV), then at 4.3 eV (width 0.7 eV) and a third at 8.1 eV (width 0.8 eV). The dominant angular momentum character of the resonances is the same as for the three corresponding resonances of uracil. Comparing the resonance positions with the data of Ref. Aflatooni et al., 1998 we see the same general trend already observed with uracil, of an overall shift higher than experiment of all resonances by about 2 eV. Interestingly, the first two resonances are measured to occur at an energy lower than in uracil, a trend that we verify in our calculations.
#### III.1.3 Thymine
The scattering cross section for thymine is closely similar to uracil, which is not surprising in view of their close structural similarities, this applies to both the magnitude and the position of the resonances, which are slightly shifted to higher energies. Specifically, we find resonances at 2.4 eV (width 0.2 eV) at 5.5 eV (width 0.6 eV) and at 7.9 eV (width 1 eV).
#### III.1.4 Adenine
The electron scattering spectrum for adenine presents many resonances, due to the complexity of the target structure, as expected. Also very interesting is the fact that the cross section drops sharply at energies below 2 eV, a behavior opposite to that found for the other molecules, if we do not consider the dipole physics outside the R-matrix box, whereas a zero-energy peak appears in the full calculation, a possible sign of a dipole bound state right below threshold.
The first resonance occurs at 2.4 eV (width 0.2 eV), the second at 3.2 eV (sharp, width 0.2 eV), then another centered at 4.4 eV (0.3 eV wide), while at 9 eV we have a broader resonance of width 0.5 eV. The dominant partial wave of the first two resonances is $`l=2`$ (65% and 62% respectively). The third resonance is $`l=3`$ at 51% and $`l=4`$ at 33% The resonance at 9 eV is dominantly $`l=5`$ (53%) with an $`l=3`$ contribution (22%).
Compared to experiment we have a shift of all resonances roughly 1.5 eV higher, as in guanine, in this case the spacings are correct (about 1 eV between the first three resonances, while the fourth falls too high in energy and it is not measured in experiment). Also the experimental widths of the first three resonances are very similar, as in our data. A comparison with the data of Ref. Aflatooni et al., 1998 is shown in Fig. 11.
#### III.1.5 Guanine
For guanine we find 4 resonances: at 2.4 eV (width 0.2 eV), at 3.8 eV (width 0.25eV), a third at 4.8 eV (width 0.35 eV), then at 8.9 eV (width 0.6 eV) and a broad resonance around 12 eV. Each of the first three resonances has strong contributions from d, f, and g-waves. At 2.4 eV the contributions are 46% for $`l=2`$ and 37% for $`l=3`$, for the second resonance $`l=2`$ is 44% while $`l=4`$ is 32%, the third is 38% of $`l=4`$ character and 35% of $`l=3`$. The resonance at 8.9 eV is 33% h-wave, 28% f-wave and 20% g-wave. At 12 eV the composition is: $`l=4`$ and $`l=5`$ equally at 23%, while $`l=6`$ contributes a further 13%.
Comparison to experimental data (see Fig. 12) shows again a shift of 1.5 eV overall, while the resonance spacing is well reproduced, and the second and third resonances fall at higher energies with respect to adenine, as in our calculations. Also the widths seem close to experiment.
### III.2 Resonance molecular structures
From the shapes and nodal structures of the resonant states it is possible to attempt a discussion of the dissociation patterns observed experimentally, if we consider the resonant states as being precursors for dissociative states. Caution must be used, though, in drawing conclusions from this analysis, because this involves a certain degree of speculation. In fact, to establish once and for all the dissociation patterns of these complicated systems, scattering calculations at many different geometries would have to be carried out, and the nuclear dynamics should be included, At present this is computationally too expensive to contemplate. The first two resonances observed for uracil are shown in Figs. 5 and 5. The quantity plotted is a projection on the molecular plane of the eigenvector corresponding to the maximum eigenvalue of the time-delay matrix Greene et al. (1996)
$$Q=iS\frac{dS}{dE},$$
(10)
where $`S`$ is the scattering matrix. At the energy where the time delay of the resonance is a maximum, this eigenvector constitutes the dominant contribution to the resonant structure, since it corresponds to the partial wave that experiences the maximum time delay in the scattering process. For sufficiently narrow resonances one eigenvalue is always dominant, making the resonance analysis much easier. The eigenvectors of the time-delay matrix are complex, so we adopt a phase factor such that the highest peak of the wavefunction is a purely real number. We then find that the resonance wavefunction is real everywhere, to a good approximation (the imaginary part is about $`10^6`$ smaller than the value of the real part), and we plot only the real part. We analyze in detail only the cases of uracil and adenine, the other pyrimidines being very similar to the former and guanine to the latter.
The nodal patterns for uracil are very similar to the ones showed in Ref. Gianturco and Lucchese, 2004, which is not surprising since the approximations made in that work are similar to ours, as already discussed, therefore we show only the first two resonances. Incidentally, we notice the close resemblance of these resonant wavefunctions to the first virtual orbitals of uracil from a HF calculation performed with a small basis set (6-31G\*), see for example Fig. 6. In Ref. Gianturco and Lucchese, 2004 also a very low $`\sigma ^{}`$ resonance is found at 0.012 eV. Our cross section grows substantially at low energy. If we plot the eigenstate corresponding to the largest eigenvalue of the time-delay matrix (as described in Sec. III.2), as in Fig. 5, corresponding to this low energy range, it looks similar to Fig. 5 in Ref. Gianturco and Lucchese, 2004, with the main differences being that in our wavefunction the N<sub>3</sub>-H bond has a node, there is a large excess charge on N<sub>3</sub> and on the oxygen attached to C<sub>4</sub>, while another nodal surface cuts diagonally from C<sub>2</sub> to C<sub>5</sub>. This resonance anyway does not appear to be so relevant in the experimental data,Scheer et al. (2004) where mainly the $`\pi ^{}`$ resonances are detected.
We can also see that there is accumulation of electronic density (the peaks of the wavefunction) on the ring structure, and that many of the ring bonds have nodal surfaces cutting through them, so capture in these resonant states can be reasonably thought as leading to a fragmentation of the molecule in which the aromatic ring is broken. Experimental dissociation patterns are illustrated for Br-uracil in Ref. Abdoul-Carime et al., 2000, where evidence for breaking of the ring structure lies in the peaks at 1.6 and 3.5 eV produced by (OCN)<sup>-</sup> and other fragments. These fragments can be generated by capture into shape resonances, appearing in our calculations at 2.2, 5.2 and 7.8 eV respectively. In particular there is a nodal surface in the 5.2 eV resonance that encloses the C<sub>4</sub>-N<sub>3</sub> bond, which could generate a CN<sup>-</sup> fragment.
Since our calculations do not take into account core-excited states or vibrations, our results do not include any Feshbach resonant structures. These appear to cause at least some of the patterns observed in experiment, as in the case of uracil, Scheer et al. (2004) and they will presumably constitute the dominant trapping pathways for energies higher than 7 eV, where the number of electronic Feshbach resonances starts to become very large. Tonzani and Greene (2005b)
In the case of uracil, we looked carefully for a $`\sigma ^{}`$ resonance that might be similar to the state shown in Fig. 3 of Ref. Scheer et al., 2004 around 3 eV at equilibrium geometry, a dissociative state most likely responsible for N<sub>1</sub>-H bond cleavage. The Ref. Scheer et al., 2004 calculation was performed by scaling Hartree-Fock continuum orbital energies, so no information about the width was provided. Such a state, taking into account an expected shift of 2-3 eV upward in our calculations, should have appeared at around 6-7 eV, and it was not found. This is probably due to the fact that this resonance is extremely broad, since it is also not seen even in experiment. Scheer et al. (2004) Moreover in calculations carried out using complex absorbing potentials, in connection with Green’s function methods, Santra and Cederbaum (2002); Feuerbacher et al. (2003) for similar systems (like benzene Feuerbacher and Santra (2005)), analogous $`\sigma `$ resonances were extremely hard to detect. They became narrower (around 1eV width at equilibrium) only when the relevant hydrogen was substituted with a heavy atom like chlorine; this was also demonstrated experimentally in the case of Cl-uracil in Ref. Scheer et al., 2004.
For adenine there is less experimental information available to compare. In Ref. Gohlke et al., 2003 it is stated that the dominant breakup channel for low energy (0-4eV) electron attachment leads to hydrogen atom loss, and very prominent resonant structures are present in the range 1 to 3 eV. If we look at the resonance wavefunction maps in Fig. 8-9 we can see that there is no significant buildup of electronic density on any of the hydrogens, consistently with the fact that the negative charge stays on the molecular frame, and therefore there is no H<sup>-</sup> formation.
The first few unoccupied molecular orbitals that can be obtained from a Hartree-Fock calculation for adenine, as in the case of uracil, are extremely similar in their nodal structures to our resonant wavefunctions, so these shape resonances can be viewed quite reasonably as the trapping of the scattered electron in a virtual orbital.
Most of the C-C and C-N bonds have nodal surfaces passing through them, This might suggest that other channels that involve the breakup of C-C and C-N bonds could also be available at these energies, although probably they are less important than the hydrogen loss products.
## IV Conclusions
We have presented results for electron scattering from DNA and RNA bases. In showing some of their resonant wavefunctions, we have attempted to link the molecular breakup patterns and products to the structure of the nodal surfaces of these wavefunctions. The results for cross sections and resonances show an overall shift of about 1.5-2 eV higher for all the resonances in our calculations, compared to experiment, we believe that this shift is due to the approximate nature of our model. Apart from this, though, we seem to reproduce trends observed in experiment, with respect to the resonance spacing, their widths and also in relative positions for different molecules, which gives us guarded confidence in our results. We have presented the first calculations of cross sections for all the main DNA bases, and discussed the relationship of our results to experimental data. Ample room exists for improvement in our model, and it will be desirable to eventually calculate resonant surfaces, and not just equilibrium values, in order to characterize the dissociation processes. Work is presently underway to meet some of these challenges.
## Acknowledgments
This work was supported in part by the Department of Energy, Office of Science, by an allocation of NERSC supercomputing resources, and by the Keck Foundation through computational resources. We would like to thank P. Burrow, R. Santra and S. Feuerbacher for useful discussions and for sharing their unpublished data. |
warning/0507/hep-th0507186.html | ar5iv | text | # IS IT PHYSICALLY SOUND TO ADD A TOPOLOGICALLY MASSIVE TERM TO THREE-DIMENSIONAL MASSIVE ELECTROMAGNETIC OR GRAVITATIONAL MODELS ?
## 1 Introduction
In the last two decades much attention has been devoted to the study of the remarkable properties of gauge theories in (2 + 1) dimensions. Certainly, it would not be an exaggeration to claim that by now these properties are not only well-appreciated but also well-understood. Therefore, it should be natural, at least from a naive point of view, to expect that the addition of a Chern-Simons term to massive electromagnetic or gravitational models would produce systems endowed with properties that, in principle, should be as exciting as those concerning the well-known theories of Maxwell-Chern-Simons or Einstein-Chern-Simons. Our aim here is to analyze these massive, topologically massive, models. Since, currently, there are two distinct non-topological mass-generating mechanisms for gauge fields: adding the well-known Proca/Fierz-Pauli, or the more sophisticated higher-derivative electromagnetic/higher-derivative gravitational, terms, our analysis will comprise topologically massive Proca electromagnetism (TMPE), topologically massive higher-derivative electromagnetism (TMHDE), which is also known as Podolsky-Chern-Simons planar electromagnetism, topologically massive Fierz-Pauli gravity (TMFPG), and topologically massive higher-derivative gravity (TMHDG). These systems will be examined for both possible sign choices of the Maxwell/Einstein Lagrangian, as well as in its absence, which implies that they are the most general such models. Both TMPE and TMFPG are not gauge-invariant due the presence of an explicit mass term, but the three-term models with higher-derivatives are gauge-invariant.
We are particularly interested in two issues that are somewhat correlated:
* The compatibility—from the point of view of the unitarity—between massive electromagnetic or gravitational models and topologically massive terms.
* The exciting physics resulting from the utilization of the gauge-invariant three-term systems as effective field models. We remark that it was recently shown that boson-boson bound states do exist in the framework of three-dimensional higher-derivative electromagnetism augmented by a topological Chern-Simons term.
To probe the unitarity of the massive, topologically massive, models, we will make use of an uncomplicated and easily handling algorithm that converts the task of checking the unitarity, which in general demands much work, into a straightforward algebraic exercise. The prescription consists basically in saturating the propagator with external conserved currents, compatible with the symmetries of the system, and in examining afterwards the residues of the saturated propagator ($`SP`$) at each of their simple poles. We use natural units throughout.
## 2 Massive, Topologically Massive, Electromagnetic Models
The Lagrangian for TMPE is the sum of Maxwell, standard Proca mass, Chern-Simons, terms, namely,
$$_{\mathrm{TMPE}}=a\frac{F_{\mu \nu }F^{\mu \nu }}{4}+\frac{1}{2}m^2A^\mu A_\mu +\frac{s}{2}\epsilon _{\mu \nu \rho }A^\mu ^\nu A^\rho ,$$
(1)
while the Lagrangian for TMHDE is the sum of Maxwell, higher-derivative, gauge-fixing (Lorentz-gauge), and Chern-Simons, terms, i.e.,
$$_{\mathrm{TMHDE}}=a\frac{F_{\mu \nu }F^{\mu \nu }}{4}+\frac{l^2}{2}_\nu F^{\mu \nu }^\lambda F_{\mu \lambda }\frac{1}{2\lambda }(_\nu A^\nu )^2+\frac{s}{2}\epsilon _{\mu \nu \rho }A^\mu ^\nu A^\rho .$$
(2)
Here, $`F_{\mu \nu }=_\nu A_\mu _\mu A_\nu `$ is the usual electromagnetic tensor field, $`l`$ is a cutoff, $`s>0`$ is the topological mass, and $`a`$ is a convenient parameter that can take the values $`+1`$ (Maxwell’s term with the standard sign), $`1`$ (Maxwell’s term with the “wrong sign”), or $`0`$ (absence of the Maxwell’s term). The corresponding propagators are given by
$$𝒫_{\mathrm{TMPE}}=\frac{ak^2m^2}{(ak^2m^2)^2s^2k^2}\theta +\frac{1}{m^2}\omega \frac{s}{(ak^2m^2)^2s^2k^2}S,$$
(3)
$$𝒫_{\mathrm{TMHDE}}=\frac{l^2k^4ak^2}{(l^2k^4ak^2)^2s^2k^2}\theta \frac{\lambda }{k^2}\omega \frac{s}{(l^2k^4ak^2)^2s^2k^2}S,$$
(4)
where $`\theta _{\mu \nu }\eta _{\mu \nu }\frac{_\mu _\nu }{\mathrm{}}`$ and $`\omega _{\mu \nu }=\frac{_\mu _\nu }{\mathrm{}}`$ are, respectively, the usual transverse and longitudinal vector projector operators, $`S_{\mu \nu }\epsilon _{\mu \rho \nu }^\rho `$ is the operator associated with the topological term, and $`\eta _{\mu \nu }`$ is the Minkowski metric. Our signature conventions are $`(+,,)`$, $`\epsilon ^{012}=+1=\epsilon _{012}`$. Now, the algorithm from Ref. 4 says that all we have to do in order to check the unitarity of a massive, topologically massive, electromagnetic model is to verify whether the residues at each simple pole of the $`\theta `$-component of the propagator in the basis $`\{\theta ,\omega ,S\}`$ which, for short, we will designate as $`f_\theta `$, are $`0`$. We use this recipe in the following to test the unitarity of TMPE and TMHDE, in this order.
### 2.1 Checking the unitarity of TMPE
We start our analysis by setting the parameter $`a`$ in Eq. (1) equal to $`+1`$ because it must be positive both in the Proca ($`s=0`$) and Maxwell-Chern-Simons ($`m=0`$) limits. Now, the $`\theta `$-component of the propagator in the basis $`\{\theta ,\omega ,S\}`$ is, according to Eq. (3), equal to $`f_\theta =\frac{m^2k^2}{(k^2m_+^2)(k^2m_{}^2)}`$, where $`m_\pm =\frac{1}{2}\left[\sqrt{4m^2+s^2}\pm s\right]`$. Therefore, the model has two degrees of freedom with masses $`m_+`$ and $`m_{}`$, which is precisely what Deser and Tekin have found using a rather different approach. Our result is also in agreement with another works existing in the literature. On the other hand, it is trivial to show that both $`\mathrm{Res}f_\theta |_{k^2=m_+^2}`$ and $`\mathrm{Res}f_\theta |_{k^2=m_{}^2}`$ are less than zero. Thence, TMPE with the Maxwell’s term with the usual sign is unitary. Choosing $`a=1`$, we see that if $`s^2>4m^2`$, then $`f_\theta =\frac{m^2+k^2}{(k^2m_+^2)(k^2m_{}^2)}`$, with $`m_\pm =\frac{1}{2}\left[s\pm \sqrt{s^24m^2}\right]`$. A straightforward calculation allows us to conclude that $`\mathrm{Res}f_\theta |_{k^2=m_+^2}>0`$, and $`\mathrm{Res}f_\theta |_{k^2=m_{}^2}<0`$, implying that, if $`s^2>4m^2`$, TMPE with Maxwell’s term with the wrong sign is non-unitary . It is worth mentioning that this system, despite having acceptable values for the masses, faces ghost problems. Of course, if $`s^2<4m^2`$ the two roots of $`x^2+x(2m^2s^2)+m^4=0`$, where $`k^2x`$, are imaginary; note also that for $`s^2=4m^2`$ the two hitherto complex roots coalesce and the masses are simply $`m_+=m_{}=\frac{s}{2}`$: these models are never viable. We focus, at last, on the case $`a=0`$ (absence of the Maxwell’s term). This model was analyzed long ago by Deser and Jackiw, who came to the conclusion that setting $`a=0`$ yields just another version of the Maxwell-Chern-Simons theory and so it is equivalent to choosing $`m^2=0`$.
#### 2.1.1 Discussion
Note that Proca electromagnetism $`(s=0)`$ with $`a=1`$ contains tachyons; however, TMPE with $`a=1`$ and $`s^2>4m^2`$ is plagued by ghosts but not by tachyons: the particle content of the model is one non-tachyonic spin-1 ghost of mass $`m_+`$ and one massive spin-1 normal particle of mass $`m_{}`$. Thence, a field theory built from this model would not be satisfactory from the point of view of their fundamentals. It could regarded, perhaps, as an effective field theory, i.e., a low-energy approximation to a more fundamental theory. Nonetheless, the condition $`s>2m`$ is greatly discouraging as far as the possibility of applying this kind of model, for instance, to some condensed matter systems where one deals, in general, with low-energy excitations. Interesting enough, only the model with $`a=+1`$ may be viewed as physically sound. Why is this so? Because the aforementioned system has a Lagrangian that reproduces the Lagrangians of well-behaved physical models when the appropriate limits are taken. Indeed, if $`s=0`$, we recover Proca electromagnetism; on the other hand, setting $`m=0`$, we obtain Chern-Simons electromagnetism. It is remarkable that we also arrive at a nice physical model by removing the Maxwell’s term: the system with $`a=0`$ and the “self-dual” model of Ref. 8 are equivalent.
### 2.2 Checking the unitarity of TMHDE
Based on the above informations, we have every reason to begin the unitarity analysis of TMHDE by setting $`a=+1`$ in Eq. (2). The calculations are now more complicated because, unlike the previous model, this represents in general three massive excitations rather than two massive ones. Since the $`\theta `$-component of the propagator concerning TMHDE with $`a=+1`$ can be written as $`f_\theta =\frac{M^2(xM^2)}{x^32M^2x^2+M^4xM^4s^2}`$, where $`M\frac{1}{l}`$, we have to analyze the nature, as well as the signs, of the roots of the cubic equation $`x^3+a_2x^2+a_1x+a_0=0`$, where $`a_22M^2,a_1M^4,`$ and $`a_0M^4s^2`$. Taking into account that we are only interested in those roots that are both real and unequal, we require $`D<0`$, where $`DQ^3+R^2`$, with $`Q`$ and $`R`$ being, in this order, equal to $`\frac{3a_1a_2^2}{9}`$ and $`\frac{9a_1a_227a_02a_2^3}{54}`$, is the polynomial discriminant. Performing the computations we get $`D=M^8s^2\left[\frac{s^2}{4}\frac{M^2}{27}\right]`$, implying that only and if only $`s^2<\frac{4M^2}{27}`$ will the roots be real and unequal. Our next step is to verify whether or not these roots are positive. This can be accomplished by building the Routh-Hurwitz array, namely,
Noting that there are three signs changes in the first column of the array above, we conclude that all the three roots are positive. In summary, if $`s^2<\frac{4m^2}{27}`$, TMHDE with $`a=+1`$ is a model with acceptable values for the masses. Denoting these roots as $`x_1,x_2`$, and $`x_3`$, and assuming without any loss of generality that $`x_1>x_2>x_3`$, we get
$`f_\theta `$ $`=`$ $`{\displaystyle \frac{M^2(x_1M^2)}{(x_1x_2)(x_1x_3)}}{\displaystyle \frac{1}{xx_1}}+{\displaystyle \frac{M^2(x_2M^2)}{(x_2x_1)(x_2x_3)}}{\displaystyle \frac{1}{xx_2}}`$
$`+{\displaystyle \frac{M^2(x_3M^2)}{(x_3x_1)(x_3x_2)}}{\displaystyle \frac{1}{xx_3}}.`$
Hence, TMHDE with $`a=+1`$ will be unitary if the conditions $`x_1M^2<0,x_2M^2>0,`$ and $`x_3M^2<0`$ hold simultaneously. Obviously, this will never occur, which allows us to conclude that TMHDE with the Maxwell’s term with the standard sign is non-unitary. If $`a=1`$, $`f_\theta =\frac{M^2(x+M^2)}{x^3+2M^2x^2+M^4xM^4s^2}`$. Since the polynomial discriminant, $`D=M^8s^2\left[\frac{s^2}{4}+\frac{M^2}{27}\right]`$, for the cubic equation $`x^3+2M^2x^2+M^4xM^4s^2=0`$ is greater than zero, one of the roots of the equation is real and the other two are complex conjugates, which means that the system with $`a=1`$ is forbidden. To finish our analysis we set $`a=0`$ in Eq. (2). In this case $`f_\theta =\frac{xM^2}{x^3M^4s^2}`$, and the polynomial discriminant related to $`x^3M^4s^2=0`$ is greater than zero. This model, as the previous one, is also forbidden.
#### 2.2.1 Discussion
Should we expect intuitively that TMHDE with $`a=+1`$ faced unitary problems? The answer is affirmative. In fact, setting $`s=0`$, for instance, in its Lagrangian, we recover the Lagrangian for the usual Podolsky electromagnetism which is non-unitary. Nonetheless, Podolsky-Chern-Simons (PCS) planar electromagnetism with $`a=+1`$ and $`s^2<\frac{4M^2}{27}`$, despite being haunted by ghosts, has normal massive modes. Since the existence of these well-behaved excitations is subordinated to the condition $`s<\frac{2M}{\sqrt{27}}`$, we are really encouraged to regard this system as an effective field model. It is quite remarkable that the coupling of PCS planar electrodynamics with a charged scalar field, produces an attractive interaction between equal charge bosons. To see this we need to know beforehand the expression of the effective non-relativistic potential for the interaction of two charged-bosons in the center-of-mass frame. A somewhat involved computation, yields
$`V(r)`$ $`=`$ $`{\displaystyle \frac{sQ^2}{\pi ml^4}}\left[{\displaystyle \frac{l^4}{s^2}}{\displaystyle \frac{1}{r^2}}+{\displaystyle \frac{1}{r}}{\displaystyle \underset{j}{}}B_j\sqrt{|x_j|}K_1(\sqrt{|x_j|}r)\right]𝐋`$ (5)
$`+{\displaystyle \frac{Q^2}{2\pi l^4}}\left[{\displaystyle \underset{j}{}}A_jK_0(\sqrt{|x_j|}r)\right],`$
where $`A_1\frac{1+l^2x_1}{(x_1x_2)(x_1x_3)},A_2\frac{1+l^2x_2}{(x_2x_1)(x_2x_3)},A_3\frac{1+l^2x_3}{(x_3x_1)(x_3x_2)},B_1\frac{(1+l^2x_1)^2}{s^2(x_1x_2)(x_1x_3)},B_2\frac{(1+l^2x_2)^2}{s^2(x_2x_1)(x_2x_3)},`$ and $`B_3\frac{(1+l^2x_3)^2}{s^2(x_3x_1)(x_3x_2)},`$ $`x_1,x_2,`$ and $`x_3`$ are the roots of the equation $`x^3+\frac{2x^2}{l^2}+\frac{x}{l^4}+\frac{s^2}{l^4}=0,`$ L is the angular momentum, $`K`$ is the modified Bessel function, and $`Q`$ and $`m`$ are, respectively, the charge and the mass of the scalar boson. On the other hand, the radial Schrödinger equation associated with this potential is given by
$`\left[{\displaystyle \frac{d^2}{dr^2}}+{\displaystyle \frac{1}{r}}{\displaystyle \frac{d}{dr}}\right]_{n\overline{l}}+m\left[E_{n\overline{l}}V_{\overline{l}}^{eff}\right]_{n\overline{l}}=0,`$ (6)
where
$`V_{\overline{l}}^{eff}(r)`$ $`=`$ $`{\displaystyle \frac{sQ^2}{\pi ml^4}}\left[{\displaystyle \frac{l^4}{s^2}}{\displaystyle \frac{1}{r^2}}+{\displaystyle \frac{1}{r}}{\displaystyle \underset{j}{}}B_j\sqrt{|x_j|}K_1(\sqrt{|x_j|}r)\right]\overline{l}`$
$`+{\displaystyle \frac{Q^2}{2\pi l^4}}\left[{\displaystyle \underset{j}{}}A_jK_0(\sqrt{|x_j|}r)\right]+{\displaystyle \frac{\overline{l}^2}{mr^2}}.`$
Here $`\overline{l}`$ denotes the eigenvalues of the operator $`𝐋`$. In terms of the dimensionless parameters $`ysr,\alpha \frac{Q^2}{\pi s},b_j\frac{s^2}{l^4}B_j,X_j\frac{|x_j|}{s},\beta \frac{m}{s},a_j\frac{A_j}{l^4},`$ and $`\stackrel{~}{E}_{n\overline{l}}\frac{mE_{n\overline{l}}}{s^2}`$, Eq. (6) reads
$`\left[{\displaystyle \frac{d^2}{dy^2}}+{\displaystyle \frac{1}{y}}{\displaystyle \frac{d}{dy}}\right]_{n\overline{l}}+\left[\stackrel{~}{E}_{n\overline{l}}\stackrel{~}{V}_{\overline{l}}^{eff}\right]_{n\overline{l}}=0,`$
with
$`\stackrel{~}{V}_{\overline{l}}^{eff}{\displaystyle \frac{\overline{l}(\alpha \overline{l})}{y^2}}+{\displaystyle \frac{\alpha \beta }{2}}{\displaystyle \underset{j}{}}a_jK_0(X_jy){\displaystyle \frac{\alpha \overline{l}}{y}}{\displaystyle \underset{j}{}}b_jX_jK_1(X_jy).`$
We call attention to the fact that $`\stackrel{~}{V}_{\overline{l}}^{eff}`$ behaves as $`\frac{\overline{l}^2}{y^2}`$ at the origin and as $`\frac{\overline{l}(\overline{l}\alpha )}{y^2}`$ asymptotically. Now, in four dimensions, the anomalous factor of $`\frac{4}{3}`$ in the Abraham-Lorentz model for the electron does not show up if $`l>\frac{1}{2}r_e`$, where $`r_e`$ denotes the classical radius of the electron. Therefore, we assume $`l1`$. In this limit the derivative of the potential with respect to $`y`$ reduces to
$`{\displaystyle \frac{d}{dy}}\stackrel{~}{V}_{\overline{l}}^{eff}{\displaystyle \frac{2\overline{l}(\alpha \overline{l})}{y^3}}\left[{\displaystyle \frac{2\alpha \overline{l}}{y^2}}+{\displaystyle \frac{\alpha \beta }{2}}\right]K_1(y){\displaystyle \frac{\alpha \overline{l}}{y}}K_0(y).`$
Supposing $`\overline{l}>0`$, without any loss of generality, we promptly see that, if $`\overline{l}>\alpha `$, the potential is strictly decreasing. The remaining possibility is $`\overline{l}<\alpha `$. In this interval $`\stackrel{~}{V}_{\overline{l}}^{eff}`$ approaches $`+\mathrm{}`$ at the origin and $`0^{}`$ for $`y+\mathrm{}`$, which is indicative of a local minimum. Consequently, the existence of the attractive potential is subordinated to the conditions $`a1`$ and $`0<\overline{l}<\alpha `$. One can show that the effective potential with $`l1`$ and $`0<\overline{l}<\alpha `$ can bind a pair of identical charged-scalar bosons. Accordingly, the addition of the topologically massive term to Podolsky’s electromagnetism with $`a=+1`$—an admittedly non-unitary model—did not cure its non-unitary problem; nonetheless, the condition for the resulting three-term model to be free of tachyons gives rise to a constraint between the topological and Podolsky masses which is responsible for a scalar attractive interaction.
## 3 Massive, Topologically Massive, Gravitational Models (MTMG)
TMFPG is defined by the Lagrangian
$`_{\mathrm{TMFPG}}`$ $`=`$ $`a{\displaystyle \frac{2}{\overline{\kappa }^2}}\sqrt{g}R{\displaystyle \frac{m^2}{2}}\left(h_{\mu \nu }^{}{}_{}{}^{2}h^2\right)`$ (7)
$`+{\displaystyle \frac{1}{\mu }}ϵ^{\lambda \mu \nu }\mathrm{\Gamma }_{}^{\rho }{}_{\lambda \sigma }{}^{}\left(_\mu \mathrm{\Gamma }_{}^{\sigma }{}_{\rho \nu }{}^{}+{\displaystyle \frac{2}{3}}\mathrm{\Gamma }_{}^{\sigma }{}_{\mu \beta }{}^{}\mathrm{\Gamma }_{}^{\beta }{}_{\nu \rho }{}^{}\right),`$
at quadratic order in $`\overline{\kappa }`$, where $`\overline{\kappa }^2`$ is a suitable constant that in four dimensions is equal to $`24\pi G`$, with $`G`$ being Newton’s constant. Here $`g_{\mu \nu }\eta _{\mu \nu }+\overline{\kappa }h_{\mu \nu }`$, $`\mu >0`$ is a dimensionless parameter, and $`h\eta _{\mu \nu }h^{\mu \nu }`$. Indices are raised (lowered) with $`\eta ^{\mu \nu }`$ ($`\eta _{\mu \nu }`$). The Lagrangian related to TMHDG, in turn, is given by
$`_{\mathrm{TMHDG}}`$ $`=`$ $`\sqrt{g}\left(a{\displaystyle \frac{2R}{\kappa ^2}}+{\displaystyle \frac{\alpha }{2}}R^2+{\displaystyle \frac{\beta }{2}}R_{\mu \nu }^2\right)`$ (8)
$`+{\displaystyle \frac{1}{\mu }}ϵ^{\lambda \mu \nu }\mathrm{\Gamma }_{}^{\rho }{}_{\lambda \sigma }{}^{}\left(_\mu \mathrm{\Gamma }_{}^{\sigma }{}_{\rho \nu }{}^{}+{\displaystyle \frac{2}{3}}\mathrm{\Gamma }_{}^{\sigma }{}_{\mu \beta }{}^{}\mathrm{\Gamma }_{}^{\beta }{}_{\nu \rho }{}^{}\right),`$
where $`\alpha `$ and $`\beta `$ are suitable constants with dimension $`L`$. For the sake of simplicity, the gauge-fixing term was omitted. Note that the parameter $`a`$ appearing in Eqs. (7) and (8) allows for choosing the Einstein sign’s term or even removing it.
The propagator related to TMFPG is
$`𝒫_{\mathrm{TMFPG}}`$ $`=`$ $`{\displaystyle \frac{1}{m^2}}P^1{\displaystyle \frac{M^2\left(m^2+a\mathrm{}\right)}{\mathrm{}^3+M^2a^2\mathrm{}^2+2am^2M^2\mathrm{}+M^2m^4}}P^2`$ (9)
$`{\displaystyle \frac{M}{\mathrm{}^3+M^2a^2\mathrm{}^2+2am^2M^2\mathrm{}+M^2m^4}}P`$
$`{\displaystyle \frac{m^2+a\mathrm{}}{2m^4}}\overline{P}^{\mathrm{\hspace{0.33em}0}}+{\displaystyle \frac{1}{2m^2}}\overline{\overline{P}}^{\mathrm{\hspace{0.33em}0}},`$
where $`M\mu /\overline{\kappa }^2`$. On the other hand, linearizing Eq. (8) and adding to the result the gauge-fixing term $`_{\mathrm{g}f}=\frac{1}{2\lambda }(h_{\mu \nu }^{}{}_{}{}^{,\nu }\frac{1}{2}h_{,\mu })^2`$ (de Donder gauge), we find that the propagator concerning TMHDG takes the form
$`𝒫_{\mathrm{TMHDG}}`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{}[a+b(\frac{3}{2}+4c)\mathrm{}]}}\overline{\overline{P}}^{\mathrm{\hspace{0.33em}0}}+{\displaystyle \frac{2\lambda }{k^2}}P^1+{\displaystyle \frac{1}{\mathrm{}[a+b(\frac{3}{2}+4c)\mathrm{}]}}P^0`$ (10)
$`+{\displaystyle \frac{4M}{\mathrm{}[M^2b^2\mathrm{}^24(abM^21)\mathrm{}+4M^2a^2]}}P`$
$`{\displaystyle \frac{2M^2(2ab\mathrm{})}{\mathrm{}[M^2b^2\mathrm{}^24(abM^21)\mathrm{}+4M^2a^2]}}P^2`$
$`+\left[{\displaystyle \frac{4\lambda }{\mathrm{}}}+{\displaystyle \frac{2}{\mathrm{}[a+b(\frac{3}{2}+4c)\mathrm{}]}}\right]\overline{P}^{\mathrm{\hspace{0.33em}0}},`$
where $`b\frac{\beta \kappa ^2}{2}`$, $`c\frac{\alpha }{\beta }`$, and $`M\frac{\mu }{\kappa ^2}`$. Our conventions are $`R_{}^{\alpha }{}_{\beta \gamma \delta }{}^{}=_\delta \mathrm{\Gamma }_{}^{\alpha }{}_{\beta \gamma }{}^{}+\mathrm{},R_{\mu \nu }=R_{}^{\alpha }{}_{\mu \nu \alpha }{}^{},R=g^{\mu \nu }R_{\mu \nu }`$, where $`g_{\mu \nu }`$ is the metric tensor, and signature $`(+,,)`$. The propagators were calculated using the basis
$$\{P^1,P^2,P^0,\overline{P}^{\mathrm{\hspace{0.33em}0}},\overline{\overline{P}}^{\mathrm{\hspace{0.33em}0}},P\},$$
where $`P^1,P^2,P^0,\overline{P}^{\mathrm{\hspace{0.33em}0}},`$ and $`\overline{\overline{P}}^{\mathrm{\hspace{0.33em}0}}`$, are the usual three-dimensional Barnes-Rivers operators, namely,
$`P_{\mu \nu ,\rho \sigma }^1={\displaystyle \frac{1}{2}}\left(\theta _{\mu \rho }\omega _{\nu \sigma }+\theta _{\mu \sigma }\omega _{\nu \rho }+\theta _{\nu \rho }\omega _{\mu \sigma }+\theta _{\nu \sigma }\omega _{\mu \rho }\right),`$
$`P_{\mu \nu ,\rho \sigma }^2={\displaystyle \frac{1}{2}}\left(\theta _{\mu \rho }\theta _{\nu \sigma }+\theta _{\mu \sigma }\theta _{\nu \rho }\theta _{\mu \nu }\theta _{\rho \sigma }\right),`$
$`P_{\mu \nu ,\rho \sigma }^0={\displaystyle \frac{1}{2}}\theta _{\mu \nu }\theta _{\rho \sigma },\overline{P}_{\mu \nu ,\rho \sigma }^{\mathrm{\hspace{0.33em}0}}=\omega _{\mu \nu }\omega _{\rho \sigma },`$
$`\overline{\overline{P}}_{\mu \nu ,\rho \sigma }^{\mathrm{\hspace{0.33em}0}}=\theta _{\mu \nu }\omega _{\rho \sigma }+\omega _{\mu \nu }\theta _{\rho \sigma },`$
and $`P`$ is the operator associated with the linearized Chern-Simons term, i.e.,
$`P_{\mu \nu ,\rho \sigma }{\displaystyle \frac{\mathrm{}^\lambda }{4}}[ϵ_{\mu \lambda \rho }\theta _{\nu \sigma }+ϵ_{\mu \lambda \sigma }\theta _{\nu \rho }+ϵ_{\nu \lambda \rho }\theta _{\mu \sigma }+ϵ_{\nu \lambda \sigma }\theta _{\mu \rho }].`$
According to Ref. 4, the saturated propagator concerning the MTMG, is given by
$`SP_{\mathrm{MTMG}}=\left[T^{\mu \nu }T_{\mu \nu }{\displaystyle \frac{1}{2}}T^2\right]f_{P^2}+{\displaystyle \frac{1}{2}}T^2f_{P^0},`$ (11)
where $`T^{\mu \nu }`$ is the external conserved current that, obviously, is symmetric in the indices $`\mu `$ and $`\nu `$, $`T\eta _{\mu \nu }T^{\mu \nu }`$, and $`f_{P^2}`$ and $`f_{P^0}`$ are, respectively, the components $`P^2`$ and $`P^0`$ of the propagator in the basis $`\{P^1,P^2,P^0,\overline{P}^{\mathrm{\hspace{0.33em}0}},\overline{\overline{P}}^{\mathrm{\hspace{0.33em}0}},P\}`$. Therefore, to find out whether or not the gravitational model is unitary, we must compute $`SP_{\mathrm{MTMG}}`$ using Eq. (11) and determine afterwards the the residue at each simple pole of $`SP_{\mathrm{MTMG}}`$: If all the residues are $`0`$, the model is unitary; however, if at least one of them is negative, the system is non-unitary. The unitarity analysis is greatly facilitated if we take into account that $`\left[T^{\mu \nu }T_{\mu \nu }\frac{1}{2}T^2\right]_{k^2=m^2}>0`$ and $`\left[T^{\mu \nu }T_{\mu \nu }T^2\right]_{k^2=0}=0`$, where $`m0`$ is the mass of a generic physical particle associated with the MTMG, and $`k`$ is the corresponding momentum exchanged. Using this prescription, we check in the following the unitarity of TMFPG and TMHDE, in this order.
### 3.1 Checking the unitarity of TMFPG
To begin with, we set $`a=1`$ in Eq. (8) because we want to recover the Einstein-Chern-Simons Lagrangian in the $`m=0`$ limit—topologically massive gravity is a theory that requires $`a=1`$ to be ghost-free. The corresponding saturated propagator is given by
$`SP_{\mathrm{TMFPG}}=\left[T^{\mu \nu }T_{\mu \nu }{\displaystyle \frac{1}{2}}T^2\right]{\displaystyle \frac{M^2(m^2+k^2)}{k^6M^2k^42m^2M^2k^2M^2m^4}}.`$
Our next step is to study the roots of the cubic equation $`x^3M^2x^22m^2M^2x^2M^2m^4=0`$. Since the discriminant, $`D=M^4m^6\left[\frac{M^2}{27}+\frac{m^2}{4}\right]`$, related to this equation is greater than zero, the model at hand is unphysical and must be rejected. Consequently, we turn our attention to the system with $`a=+1`$. Now, we have to consider the roots of the equation $`x^3M^2x^2+2m^2M^2x^2M^2m^4=0`$, whose polynomial discriminant can be written as $`D=M^4m^6\left[\frac{m^2}{4}\frac{M^2}{27}\right]`$. Therefore, if $`\frac{m^2}{M^2}<\frac{4}{27}`$, our equation has three distinct real roots. The corresponding Routh-Hurwitz array is
Accordingly, the system with $`a=+1`$ and $`\frac{m^2}{M^2}<\frac{4}{27}`$ has acceptable values for the masses. Proceeding just as we have done for TMHDE with $`a=+1`$ and $`s^2<\frac{4M^2}{27}`$, we promptly obtain
$`SP_{\mathrm{TMFPG}}`$ $`=`$ $`{\displaystyle \frac{F(k)(m^2x_1)}{(x_1x_2)(x_1x_3)}}{\displaystyle \frac{1}{k^2x_1}}+{\displaystyle \frac{F(k)(m^2x_2)}{(x_2x_1)(x_2x_3)}}{\displaystyle \frac{1}{k^2x_2}}`$
$`+{\displaystyle \frac{F(k)(m^2x_3)}{(x_3x_1)(x_3x_2)}}{\displaystyle \frac{1}{k^2x_3}},`$
where $`F(k)\{T^{\mu \nu }(k)T_{\mu \nu }(k)\frac{1}{2}\left[T(k)\right]^2\}M^2`$. From the above, we clearly see that this model will be unitary if $`m^2>x_1,m^2<x_2,`$ and $`m^2>x_3`$. We thus come to the conclusion that TMFPG with $`a=+1`$ and $`\frac{m^2}{M^2}<\frac{4}{27}`$ is non-unitary, which means that the topological term is responsible for breaking down the unitarity of the harmless Fierz-Pauli gravity. If $`a=0`$, the discriminant associated with the equation $`x^3M^2m^4=0`$ is greater than zero, which implies that this model is physically unsound. We remark that our conclusions are in complete agreement with those of Ref. 6 where a quite different approach to the unitarity problem was employed.
#### 3.1.1 Discussion
The above results points to an important and at the same time interesting question: Why can unitary massive electromagnetic models coexist in peace with topologically massive terms, whereas unitary massive gravitational ones cannot? The root of the problem lies in the rather odd way Einstein-Chern-Simons theory is constructed: The presence of the ghosts in the dynamical field is avoided by choosing the Einstein’s term with the wrong sign. This is trivial to show. Indeed, writing the Einstein-Chern-Simons Lagrangian as
$`=a\sqrt{g}{\displaystyle \frac{2R}{\kappa ^2}}+{\displaystyle \frac{1}{\mu }}ϵ^{\lambda \mu \nu }\mathrm{\Gamma }_{}^{\rho }{}_{\lambda \sigma }{}^{}\left(_\mu \mathrm{\Gamma }_{}^{\sigma }{}_{\rho \nu }{}^{}+{\displaystyle \frac{2}{3}}\mathrm{\Gamma }_{}^{\sigma }{}_{\mu \beta }{}^{}\mathrm{\Gamma }_{}^{\beta }{}_{\nu \rho }{}^{}\right),`$ (12)
with $`a=\pm 1`$, we promptly see that the corresponding saturated propagator is given by
$$SP=\frac{1}{a}\left(T^{\mu \mu }T_{\mu \nu }\frac{1}{2}T^2\right)\frac{1}{k^2M^2}\frac{1}{a}\left(T^{\mu \mu }T_{\mu \nu }T^2\right)\frac{1}{k^2}.$$
(13)
Thus, to render the theory unitary we are obliged to set $`a=1`$ in Eq. (13). Note that as far as these three-term systems are concerned, we are always in a dilemma: Which value should we assign to $`a`$, $`1`$ or $`+1`$? If we single out $`a=1`$, for instance, we recover Einstein-Chern-Simons theory when the non-topological massive term is removed; however, in the absence of the topological term ,we do not get a nice physical theory because now the Einstein’s term has the wrong sign. On the other hand, if we pick out $`a=+1`$, we do not recover Einstein-Chern-Simons theory when the non-topological massive term is removed. In other words, due to the unusual Einstein sign’s term in the Lagrangian concerning Einstein-Chern-Simons theory, the augmented systems do not reduce to well-behaved physical models in the suitable limits. Note that these idiosyncrasies do not occur in the framework of massive, topologically massive, electromagnetic models because the Maxwell sign’s term concerning Maxwell-Chern-Simons theory is the same as that of the usual Maxwell’s theory.
### 3.2 Checking the unitarity of TMHDG
Assuming $`a0`$, the $`SP`$ concerning TMHDG can be written as
$`SP_{\mathrm{TMHDG}}`$ $`=`$ $`{\displaystyle \frac{M^2b}{2}}{\displaystyle \frac{(T^{\mu \nu }T_{\mu \nu }\frac{1}{2}T^2)}{k^2M_1^2}}{\displaystyle \frac{1+\sqrt{12abM^2}}{\sqrt{12abM^2}\left[1abM^2\sqrt{12abM^2}\right]}}`$ (14)
$`+{\displaystyle \frac{M^2b}{2}}{\displaystyle \frac{(T^{\mu \nu }T_{\mu \nu }\frac{1}{2}T^2)}{k^2M_2^2}}{\displaystyle \frac{1+\sqrt{12abM^2}}{\sqrt{12abM^2}\left[1abM^2+\sqrt{12abM^2}\right]}}`$
$`+{\displaystyle \frac{T^{\mu \nu }T_{\mu \nu }T^2}{ak^2}}+{\displaystyle \frac{\frac{1}{2}T^2}{a(k^2m^2)}},`$
where
$`M_1^2`$ $``$ $`\left({\displaystyle \frac{2}{b^2M^2}}\right)[1abM^2\sqrt{12abM^2}],`$
$`M_2^2`$ $``$ $`\left({\displaystyle \frac{2}{b^2M^2}}\right)[1abM^2+\sqrt{12abM^2}],`$
$`m^2`$ $``$ $`{\displaystyle \frac{a}{b(3/2+4c)}}.`$
It is interesting to note that $`M_1^2M^2`$, and $`M_2^2+\mathrm{}`$, as $`b0`$, which implies that Eq. (14) reduces to Eq. (13) when $`\alpha ,\beta 0`$, as expected. We are now ready to analyze the excitations and mass counts concerning TMHDG for both allowed signs of $`a`$. To avoid needless repetitions, we restrict ourselves to presenting a summary of the main results in Table 1. The systems that do not appear in this table are tachyonic, i.e., unphysical.
In conclusion, we consider TMHDG with $`a=0`$. In this case,
$`SP_{\mathrm{TMDHDG}}={\displaystyle \frac{M^2b}{2}}\left[{\displaystyle \frac{T^{\mu \nu }T_{\mu \nu }\frac{1}{2}T^2}{k^2}}+{\displaystyle \frac{T^{\mu \nu }T_{\mu \nu }\frac{1}{2}T^2}{k^2\frac{4}{M^2b^2}}}\right]+{\displaystyle \frac{1}{b(\frac{3}{2}+4c)}}{\displaystyle \frac{\frac{1}{2}T^2}{k^4}}.`$
The pole of order two at $`k^2=0`$ indicates that these models are unphysical.
#### 3.2.1 Discussion
As intuitively expected, TMHDG is non-unitary for $`a=\pm 1`$; nonetheless, these models are in general non-tachyonic which means that under certain circumstances they may be viewed as effective field models. Our aim here is to investigate, in passing, the novel features of one of these non-unitary gauge-invariant three-term effective field models. To be more specific, we fix our attention on the first model of Table 1, i.e., TMHDG with $`a=1`$, $`b>0`$, and $`\frac{3}{2}+4c<0`$ Ref. 13. We have chosen the $`a=1`$ system because it reduces, in the absence of the topologically massive term, to higher-derivative gravity with $`a=1`$—an effectively multimass model of the fourth-derivative order with interesting properties of its own. Now, it can be shown that the effective non-relativistic potential for the interaction of two scalars bosons in the framework of TMHDG with $`a=1`$, $`b>0`$, and $`\frac{3}{2}+4c<0`$, is given by
$`V(r)=2\overline{m}^2\overline{G}\left[K_0(rm){\displaystyle \frac{K_0(rM_+)}{1+\frac{bM_+^2}{2}}}{\displaystyle \frac{K_0(rM_{})}{1+\frac{bM_{}^2}{2}}}\right],`$ (15)
where $`\overline{m}`$ is the mass of one of the neutral bosons , $`\overline{G}\frac{\kappa ^2}{32\pi }`$, and
$$M_\pm =\frac{1}{bM}\left[\sqrt{1+2bM^2}\pm 1\right].$$
Note that $`V(r)`$ behaves as $`2\overline{m}^2\overline{G}\mathrm{ln}\left(\frac{M_+^{1+\frac{bM_+^2}{2}}M_{}^{1+\frac{bM_{}^2}{2}}}{m}\right)`$ at the origin and as
$$2\overline{m}^2\overline{G}\left[\sqrt{\frac{\pi }{2rm}}e^{rm}\frac{1}{1+\frac{bM_+^2}{2}}\sqrt{\frac{\pi }{2rM_+}}e^{rM_+}\frac{1}{1+\frac{bM_{}^2}{2}}\sqrt{\frac{\pi }{2rM_{}}}e^{rM_{}}\right].$$
asymptotically. Accordingly, $`V(r)`$ is finite at the origin and zero at infinity. The derivative of this potential with respect to $`r`$ is in turn given by
$`{\displaystyle \frac{dV}{dr}}=2\overline{m}^2\overline{G}\left[mK_1(rm)+{\displaystyle \frac{M_+}{1+\frac{bM_+^2}{2}}}K_1(rM_+)+{\displaystyle \frac{M_{}}{1+\frac{bM_{}^2}{2}}}K_1(rM_{})\right]`$ (16)
On the other hand, it was shown recently that in four dimensions the propagation of photons in the context of higher-derivative gravity (HDG) is dispersive. In other words, gravitational rainbows and semiclassical HDG can coexist without conflict. On the basis of the fact that the rainbow effect is currently undetectable, it is possible to show that $`|\beta |10^{60}`$ Ref. 17. How reliable is this result? The aforementioned constraint is of the same order as that obtained by testing the gravitational inverse-square law in the submillimeter regime. Thence, we assume $`b1`$. As a consequence, Eq. (16) reduces to
$$\frac{dV}{dr}2\overline{m}^2\overline{G}\left[mK_1(rm)+\sqrt{\frac{2}{b}}K_1\left(r\sqrt{\frac{2}{b}}\right)\right],$$
implying that the potential $`V(r)`$, which in this approximation may be expressed as
$`V(r)2\overline{m}^2\overline{G}\left[K_0(rm)K_0\left(r\sqrt{{\displaystyle \frac{2}{b}}}\right)\right],`$ (17)
is everywhere attractive if $`\sqrt{\frac{2}{b}}>m`$, is repulsive if $`m>\sqrt{\frac{2}{b}}`$, and vanishes if $`m=\sqrt{\frac{2}{b}}`$. If we appeal to the usual tools of Einstein’s geometrical theory, we arrive at the same conclusions. In fact, in the weak field approximation the gravitational acceleration , $`\gamma ^l=\frac{dv^l}{dt}`$, of a slowly moving test particle is given by $`\gamma ^l=\kappa \left[\frac{}{t}h_0^l\frac{1}{2}\frac{}{l}h_{00}\right]`$, which for time-independent fields reduces to $`\gamma ^l=\frac{\kappa }{2}\frac{}{l}h_{00}`$. Now, taking into account that $`h_{00}=\frac{2V}{\overline{m}\kappa }`$, we obtain
$$\gamma ^l=2\overline{m}\overline{G}\frac{x^l}{r}\left[mK_1(rm)+\sqrt{\frac{2}{b}}K_1\left(r\sqrt{\frac{2}{b}}\right)\right].$$
Therefore, the gravitational force exerted on the particle,
$$F^l=2\overline{m}^2\overline{G}\frac{x^l}{r}\left[mK_1(rm)+\sqrt{\frac{2}{b}}K_1\left(r\sqrt{\frac{2}{b}}\right)\right],$$
is everywhere attractive if $`\sqrt{\frac{2}{b}}>m`$, is repulsive if $`m>\sqrt{\frac{2}{b}}`$ (antigravity), and vanishes if $`m=\sqrt{\frac{2}{b}}`$ (gravitational shielding). It is remarkable that this force does not exist in general relativity. It is peculiar to topologically massive higher-derivative gravity.
## 4 Final remarks
We have shown that topologically massive terms cannot be used as a panacea for curing the non-unitarity of massive electromagnetic/gravitational models. In truth, the addition of a Chern-Simons term to a massive electromagnetic/gravitational model is physically sound only and if only the resulting three-term system reduces to well-behaved physical models in the suitable limits. A direct consequence of this fact is that we will never be able to construct an unitary, massive, topologically massive, gravitational model. Indeed, the fancy way Einstein-Chern-Simons theory is built, i.e., with the Einstein’s term with the opposite sign, precludes the existence of ghost-free, massive, topologically massive, gravitational models. Therefore, from a conceptual point of view, the addition of a topologically massive term to a massive gravitational model is a complete nonsense: On the one hand, it does not cure the non-unitarity of the original model; on the other hand, it spoils the unitarity of admittedly unitary models. An interesting and elucidatory example is furnished by $`R+R^2`$ gravity, which is defined by the Lagrangian $`=\left[a\frac{2R}{\kappa ^2}+\frac{\alpha }{2}R^2\right]\sqrt{g}`$, with $`a=\pm 1`$. If $`\alpha >0`$, this theory is non-tachyonic regardless of the sign of $`a`$; in addition, it is unitary if $`a=+1`$, and non-unitary if $`a=1`$. Incidentally, $`R+R^2`$ gravity with $`a=+1`$ is the only known gravity theory with higher-derivatives that is unitary. However, topologically massive $`R+R^2`$ gravity is non-unitary for both possible sign choices of $`a`$ Ref. 19. Yet, a new and surprising physics emerges when we analyze the three-term effective field models that are both gauge-invariant and non-unitary. In the framework of the electromagnetic models, an attractive interaction between equal charge particles can be produced that leads to an unusual planar dynamics: scalar pairs can condense into bound states. In the framework of the gravity systems, in turn, unlike what occurs in the context of the insipid and odorless three-dimensional Einstein’s general relativity, we have a gravitational interaction that can be both attractive and repulsive. We can also have a null gravitational interaction, such as in three-dimensional gravity that is trivial outside the sources. Certainly, these effective field models deserve to be both much better known and further investigated.
## Acknowledgments
We are very grateful to Prof. S. Deser for calling our attention to Ref. 6. A. Accioly thanks CNPq-Brazil for partial support while M. Dias is indebted to CAPES-Brazil for full support. |
warning/0507/quant-ph0507056.html | ar5iv | text | # Lorentz–covariant reduced spin density matrix and Einstein–Podolsky–Rosen–Bohm correlations
## I Introduction
Relativistic aspects of quantum mechanics have recently attracted much attention, especially in the context of the theory of quantum information. One of the important questions in this context is how to define the reduced spin density matrix. Such a matrix should enable one to make statistical predictions for the outcomes of ideal spin measurements which are not influenced by the particle momentum. We consider this problem in detail in the case of massive particles. The reduced spin density matrix is usually defined by the following formula Peres et al. (2002):
$$\tau _{\sigma \lambda }=𝑑\mu (k)k,m,s,\sigma \left|\widehat{\rho }\right|k,m,s\lambda ,$$
(1)
where $`\widehat{\rho }`$ denotes the complete density matrix of a single particle with mass $`m`$, $`d\mu (k)=\frac{d^3k}{2k^0}`$ is the Lorentz–invariant measure on the mass shell and four-momentum eigenvectors $`|k,m,s,\lambda `$ (i.e., $`P_\mu |k,m,s,\lambda =k_\mu |k,m,s,\lambda `$) span the space of the irreducible representation of the Poincaré group. They are normalized as follows
$$p,m,s,\sigma |k,m,s,\lambda =2k^0\delta ^3(𝐤𝐩)\delta _{\sigma \lambda }.$$
(2)
The action of the Lorentz transformation $`\mathrm{\Lambda }`$ on the vector $`|k,m,s,\lambda `$ is of the form
$$U(\mathrm{\Lambda })|k,m,s,\lambda =𝒟_{\sigma \lambda }^s(R(\mathrm{\Lambda },k))|\mathrm{\Lambda }k,m,s,\sigma ,$$
(3)
where $`𝒟^s`$ is the matrix spin $`s`$ representation of the $`SO(3)`$ group, $`R(\mathrm{\Lambda },k)=L_{\mathrm{\Lambda }k}^1\mathrm{\Lambda }L_k`$ is the Wigner rotation and $`L_k`$ designates the standard Lorentz boost defined by the relations $`L_k\stackrel{~}{k}=k`$, $`L_{\stackrel{~}{k}}=I`$, $`\stackrel{~}{k}=(m,\mathrm{𝟎})`$.
The key question is whether the reduced density matrix is covariant. In Peres et al. (2002) it was stressed that the matrix (1) is not covariant under Lorentz boosts. It means that when we calculate the complete density matrix as seen by the boosted observer
$$\widehat{\rho }^{}=U(\mathrm{\Lambda })\widehat{\rho }U^{}(\mathrm{\Lambda })$$
(4)
and then the reduced spin density matrix $`\tau _{\sigma \lambda }^{}`$ (using Eq. (1) with $`\widehat{\rho }`$ replaced with $`\widehat{\rho }^{}`$) we find that we cannot express $`\tau ^{}`$ only in terms of $`\tau `$ and $`\mathrm{\Lambda }`$. The reason is quite obvious — the Wigner rotation in the transformation law (3) is momentum dependent, except of the case $`\mathrm{\Lambda }O(3)`$. From the group theoretical point of view it is related to the fact that the Lorentz group and the rotation group are not homomorphic. Notice that in the nonrelativistic quantum mechanics it is possible to define the Galilean–covariant reduced density matrix by the formula analogous to Eq. (1) Caban et al. (2003) because such a homomorphism exists.
## II Covariant reduced density matrix
As was pointed out in Czachor (2005) matrix (1) is not always relevant to the discussion of relativistic aspects of polarization experiments (see, however, Peres et al. (2005)). For this reason we propose here another definition of the reduced density matrix. This definition relies on the analogy with the polarization tensors formalism used in quantum field theory. As a result we obtain the finite–dimensional matrix which contains not only the information about the polarization of the particle but also the information about average values of the kinematical degrees of freedom. Moreover, such a matrix transforms covariantly under the Lorentz group action.
To begin with we introduce vectors $`|\alpha ,k`$ such that
$$|\alpha ,k=v_{\alpha \sigma }(k)|k,m,s,\sigma $$
(5)
which are assumed to transform under Lorentz transformation $`\mathrm{\Lambda }`$ due to the following, manifestly covariant, rule
$$U(\mathrm{\Lambda })|\alpha ,k=\text{D}(\mathrm{\Lambda }^1)_{\alpha \beta }|\beta ,\mathrm{\Lambda }k,$$
(6)
where $`\text{D}(\mathrm{\Lambda })`$ is a given finite–dimensional Lorentz group representation. Consistency of the rules (3, 5, 6) leads to the Weinberg–like condition Weinberg (1964a, b) which has to be fulfilled:
$$\text{D}(\mathrm{\Lambda })v(k)𝒟_{}^{s}{}_{}{}^{T}(R(\mathrm{\Lambda },k))=v(\mathrm{\Lambda }k),$$
(7)
where $`v(k)`$ denotes matrix $`[v_{\alpha \sigma }(k)]`$. Thus to calculate $`v(k)`$ it is enough to determine $`v(\stackrel{~}{k})`$ and use the formula $`v(k)=\text{D}(L_k)\stackrel{~}{k}`$ which is a consequence of Eq. (7).
Assuming that the condition (7) can be solved we can define the following (unnormalized) covariant reduced density matrix:
$$\theta _{\alpha \beta }=𝑑\mu (k)\beta ,k\left|\widehat{\rho }\right|\alpha ,k.$$
(8)
We can easily check that this matrix is manifestly covariant under the transformation (4), namely we have
$$\theta ^{}=\text{D}(\mathrm{\Lambda })\theta \text{D}^{}(\mathrm{\Lambda }),$$
(9)
where $`\theta =[\theta _{\alpha \beta }]`$.
One can also easily verify that the matrix (8) is Hermitian and positive semidefinite (similarly as (1)). Transformation (9) preserves Hermicity and positive semidefiniteness of $`\theta `$ but changes its trace.
It is clear that we can define also normalized density matrix
$$\stackrel{~}{\theta }=\frac{\theta }{\text{Tr}\theta }.$$
(10)
Such a matrix transforms according to the rule
$$\stackrel{~}{\theta }^{}=\frac{\text{D}(\mathrm{\Lambda })\stackrel{~}{\theta }\text{D}^{}(\mathrm{\Lambda }}{\text{Tr}(\stackrel{~}{\theta }\text{D}^{}(\mathrm{\Lambda })\text{D}(\mathrm{\Lambda }))}.$$
(11)
One can check immediately that Eq. (11) gives a nonlinear realization of the Lorentz group connected with the quotient space $`SO(1,3)_0/SO(3)`$. Therefore this realization is linear on the rotation group. However, to extract information about polarization of the particle it does not matter which matrix we use, $`\theta `$ or $`\stackrel{~}{\theta }`$. Moreover, when we consider representations of the full Lorentz group (i.e., including inversions) the most convenient choice is to consider the matrix
$$\mathrm{\Omega }=\theta \mathrm{\Gamma },$$
(12)
where $`\mathrm{\Gamma }`$ fulfills the condition
$$\text{D}^{}\mathrm{\Gamma }=\mathrm{\Gamma }\text{D}^1,$$
(13)
which means that in this representation $`\mathrm{\Gamma }`$ represents space inversions. Thus the matrix $`\mathrm{\Omega }`$ transforms under the Lorentz group action in the following way:
$$\mathrm{\Omega }^{}=\text{D}(\mathrm{\Lambda })\mathrm{\Omega }\text{D}^1(\mathrm{\Lambda }).$$
(14)
We see that transformation (14) does not change the trace of $`\mathrm{\Omega }`$. Of course, having $`\mathrm{\Omega }`$ we can easily determine $`\theta `$ and normalized density matrix $`\stackrel{~}{\theta }`$.
Hereafter we restrict ourselves to the case of a spin-1/2 particle; generalization to the higher spin is immediate. In this case the Weinberg condition (7) can be easily solved. We want to consider representations of the full Lorentz group thus we choose as the representation D the bispinor representation $`D^{(\frac{1}{2},0)}D^{(0,\frac{1}{2})}`$, so $`\mathrm{\Gamma }=\gamma ^0`$ in this case. Explicitly, if $`ASL(2,)`$ and $`\mathrm{\Lambda }(A)`$ is an image of $`A`$ in the canonical homomorphism of the $`SL(2,)`$ group onto the Lorentz group, we take the chiral form of $`D^{(\frac{1}{2},0)}D^{(0,\frac{1}{2})}`$, namely
$$\text{D}(\mathrm{\Lambda }(A))=\left(\begin{array}{cc}A& 0\\ 0& (A^{})^1\end{array}\right).$$
(15)
The canonical homomorphism between the group $`SL(2,)`$ (universal covering of the proper ortochronous Lorentz group $`L_+^{}`$) and the Lorentz group $`L_+^{}SO(1,3)_0`$ Barut and Ra̧czka (1977) is defined as follows: With every four-vector $`k^\mu `$ we associate a two-dimensional hermitian matrix k such that
$$\text{k}=k^\mu \sigma _\mu ,$$
(16)
where $`\sigma _i`$, $`i=1,2,3`$, are the standard Pauli matrices and $`\sigma _0=I`$. In the space of two-dimensional hermitian matrices (16) the Lorentz group action is given by $`𝗄^{}=A\text{k}A^{}`$, where $`A`$ denotes the element of the $`SL(2,)`$ group corresponding to the Lorentz transformation $`\mathrm{\Lambda }(A)`$ which converts the four-vector $`k`$ to $`k^{}`$ (i.e., $`k_{}^{}{}_{}{}^{\mu }=\mathrm{\Lambda }_\nu ^\mu k^\nu `$) and $`\text{k}^{}=k_{}^{}{}_{}{}^{\mu }\sigma _\mu `$.
Now, the explicit solution of the Weinberg condition (7) under our choice of D (Eq. (15)) is given by
$$v(k)=\frac{1}{2\sqrt{1+\frac{k^0}{m}}}\left(\begin{array}{c}(I+\frac{1}{m}\text{k})\sigma _2\\ (I+\frac{1}{m}\text{k}^P)\sigma _2\end{array}\right),$$
(17)
where k is given by Eq. (16) and $`\text{k}^P=(k^P)^\mu \sigma _\mu `$ with $`k^P=(k^0,𝐤)`$. As is well known, the intertwining matrix $`v(k)`$ fulfills the Dirac equation
$$(k\gamma mI)v(k)=0,$$
(18)
where $`\gamma ^\mu `$ are Dirac matrices, $`k\gamma =k_\mu \gamma ^\mu `$. The explicit representation of Dirac matrices used in the present paper is summarized in Appendix A.
Now we discuss the general structure of the reduced density matrix (8) for $`s=\frac{1}{2}`$. We show that this matrix contains information about both average polarization as well as kinematical degrees of freedom. Recall that the polarization of the relativistic particle is determined by the Pauli–Lubanski four-vector
$$W^\mu =\frac{1}{2}\epsilon ^{\mu \nu \sigma \lambda }P_\nu J_{\sigma \lambda },$$
(19)
where $`P_\nu `$ is a four-momentum operator and $`J_{\sigma \lambda }`$ denotes generators of the Lorentz group, i.e., $`U(\mathrm{\Lambda })=\mathrm{exp}i\omega ^{\mu \nu }J_{\mu \nu }`$. We will also use the spin tensor $`S_{\mu \nu }`$ defined by the formula Anderson (1967)
$$S_{\mu \nu }=\frac{1}{m^2}\epsilon _{\mu \nu \sigma \tau }P^\sigma W^\tau .$$
(20)
Now, the $`4\times 4`$ reduced spin density matrix $`\theta `$ can be written as the following combination
$$\begin{array}{c}\theta =\frac{1}{4}(a\gamma ^0+bi\gamma ^5\gamma ^0+u_\mu \gamma ^\mu \gamma ^0+\frac{2w_\mu }{m}\gamma ^5\gamma ^\mu \gamma ^0\hfill \\ \hfill +2s_{\mu \nu }\frac{i}{4}[\gamma ^\mu ,\gamma ^\nu ]\gamma ^0).\end{array}$$
(21)
Real coefficients $`a`$, $`b`$, $`u_\mu `$, $`w_\mu `$, $`s_{\mu \nu }`$ can be determined by calculating corresponding traces. Thus, after some algebra, using Eqs. (8, 12, 1720) and (5758) we get
$`a`$ $`=\text{Tr}(\mathrm{\Omega })=1,`$ (22)
$`b`$ $`=\text{Tr}(i\mathrm{\Omega }\gamma _5)=0`$ (23)
$`u_\mu `$ $`=\text{Tr}(\mathrm{\Omega }\gamma _\mu )=\frac{1}{m}P_\mu _{\widehat{\rho }},`$ (24)
$`w_\mu `$ $`={\displaystyle \frac{m}{2}}\text{Tr}(\mathrm{\Omega }\gamma _\mu \gamma _5)=W_\mu _{\widehat{\rho }},`$ (25)
$`s_{\mu \nu }`$ $`=\text{Tr}(\mathrm{\Omega }\frac{i}{4}[\gamma _\mu ,\gamma _\nu ])=S_{\mu \nu }_{\widehat{\rho }},`$ (26)
where $`A_{\widehat{\rho }}`$ denotes the mean value of the observable $`A`$ in the state described by the complete density matrix $`\widehat{\rho }`$, $`A_{\widehat{\rho }}=\text{Tr}(\widehat{\rho }A)`$. Notice that the above relations are not accidental, since $`\gamma ^0\gamma ^\mu `$ is a canonical four-velocity operator for the Dirac particle and $`\frac{i}{4}[\gamma ^\mu ,\gamma ^\nu ]`$ are Lorentz group generators in the bispinor representation. Thus, finally, the matrix $`\mathrm{\Omega }=\theta \gamma ^0`$ has the following form:
$$\begin{array}{c}\mathrm{\Omega }=\frac{1}{4}I+\frac{1}{4m}P_\mu _{\widehat{\rho }}\gamma ^\mu +\frac{1}{2m}W_\mu _{\widehat{\rho }}\gamma ^5\gamma ^\mu \hfill \\ \hfill +\frac{1}{2}S_{\mu \nu }_{\widehat{\rho }}\frac{i}{4}[\gamma ^\mu ,\gamma ^\nu ].\end{array}$$
(27)
It can be also checked that in the nonrelativistic limit we have
$$\frac{1}{m}P^\mu _{\widehat{\rho }}\delta _0^\mu ,$$
(28a)
$$W_0_{\widehat{\rho }}0,$$
(28b)
$$S_{0\mu }_{\widehat{\rho }}0,$$
(28c)
$$S_{ij}_{\widehat{\rho }}\epsilon _{ijk}\frac{1}{m}W^k_{\widehat{\rho }}.$$
(28d)
The formalism we have introduced above can be straightforward generalized to the multiparticle case. As an example we shall discuss briefly the reduced spin density matrix for two massive particles. Two–particle Hilbert space is spanned by vectors $`|\alpha ,k|\beta ,p`$, where $`|\alpha ,k`$ is defined by Eq. (5). Therefore we define the two-particle unnormalized reduced density matrix as follows:
$$\theta _{\alpha ^{}\beta ^{},\alpha \beta }=𝑑\mu (k)𝑑\mu (p)\alpha ,k|\beta ,p\left|\widehat{\rho }\right|\alpha ^{},k|\beta ^{},p,$$
(29)
where $`\widehat{\rho }`$ denotes the complete two–particle density matrix. It is obvious that the matrix (29) is Hermitian, positive–semidefinite and can be easily normalized similarly like in the one–particle case. Moreover, in the case of two spin $`1/2`$ particles we define
$$\mathrm{\Omega }=\theta (\gamma ^0\gamma ^0).$$
(30)
## III Particle with a sharp momentum
Now let us discuss the case of the particle with a sharp momentum, say $`𝐪`$, and polarization determined by the Bloch vector $`𝝃`$, $`|𝝃|1`$, i.e., we assume that the complete density matrix has the following matrix elements
$$\begin{array}{c}k,m,s,\tau \left|\widehat{\rho }\right|p,m,s,\lambda \hfill \\ \hfill =\frac{2q^0}{\delta ^3(\mathrm{𝟎})}\delta ^3(𝐤𝐪)\delta ^3(𝐩𝐪)\frac{1}{2}(I𝝃𝝈)_{\tau \lambda }.\end{array}$$
(31)
Of course the normalization factor $`\frac{1}{\delta ^3(\mathrm{𝟎})}`$ should be understood as the result of the proper regularization procedure. Now, using Eqs. (8) and (58) we can find the corresponding matrix $`\mathrm{\Omega }`$. We have
$$\mathrm{\Omega }=\frac{1}{4}\left(\frac{q\gamma }{m}+I\right)\left(I+2\gamma ^5\frac{w\gamma }{m}\right),$$
(32)
where the four-vector $`w^\mu =W^\mu _{\widehat{\rho }}`$ is given in this case by
$$w^0=\frac{𝐪𝝃}{2},𝐰=\frac{1}{2}\left(m𝝃+\frac{𝐪(𝐪𝝃)}{q^0+m}\right),$$
(33)
i.e., $`w`$ is obtained from $`(0,m\frac{𝝃}{2})`$ by applying the Lorentz boost $`L_q`$. It should also be noted that $`w^\mu q_\mu =0`$. The matrix (32) is known in the literature as the spin density matrix for Dirac particle Berestetzki et al. (1968).
Now, to connect the density matrix introduced above with some macroscopic experiments like the Stern–Gerlach one let us consider a charged particle with sharp momentum moving in the external electromagnetic field. We assume that the giromagnetic ratio $`g=2`$. The momentum and polarization of such a particle vary in time, thus they can be regarded as functions of its proper time $`\tau `$:
$$q=q(\tau ),𝝃=𝝃(\tau ),$$
(34)
The expectation value of the operators representing the spin and the momentum will necessarily follow the same time dependence as one would obtain from the classical equations of motion Bargmann et al. (1959); Corben (1961); Anderson (1967); Costella and McKellar (1994):
$`{\displaystyle \frac{dq^\mu }{d\tau }}=`$ $`{\displaystyle \frac{e}{m}}F_\beta ^\mu q^\beta +{\displaystyle \frac{\zeta }{m^2}}q^\beta w^\nu _\nu \stackrel{~}{F}_\beta ^\mu `$
$`+{\displaystyle \frac{\zeta ^2}{m}}\stackrel{~}{F}_\beta ^\mu \left[F_\nu ^\beta w^\nu +{\displaystyle \frac{1}{m^2}}q^\beta (w^\sigma F_{\sigma \alpha }q^\alpha )\right],`$ (35)
$`{\displaystyle \frac{dw^\mu }{d\tau }}=`$ $`\zeta \left[F_\nu ^\mu w^\nu +{\displaystyle \frac{1}{m^2}}q^\mu (w^\sigma F_{\sigma \alpha }q^\alpha )\right],`$ (36)
where $`e`$ denotes the charge of the particle, $`m`$ its mass, $`\zeta `$ is the proportionality constant between the magnetic moment of the particle $`\mu ^\alpha `$ and $`w^\alpha `$ i.e. $`\mu ^\alpha =\frac{\zeta }{m}w^\alpha `$<sup>1</sup><sup>1</sup>1For the particle with charge $`e`$ and giromagnetic ratio $`g`$ $`\zeta =\frac{ge}{2m}`$Costella and McKellar (1994). and $`F_{\mu \nu }`$ is the tensor of the external electromagnetic field, $`\stackrel{~}{F}_{\alpha \beta }=\frac{1}{2}\epsilon _{\alpha \beta }^{\mu \nu }F_{\mu \nu }`$. Eq. (36) describes Thomas precession of the spin vector in the electromagnetic field Bargmann et al. (1959) while Eq. (35) allows one to determine the trajectory of the spinning particle moving in the electromagnetic field $`F_{\mu \nu }`$. The slow motion limit of the above equations takes the well–known form Costella and McKellar (1994)
$$\frac{d𝐪}{dt}=\frac{e}{m}𝐪\times 𝐁+\frac{\zeta }{2}𝝃𝐁,$$
(37)
$$\frac{d𝝃}{dt}=\zeta 𝝃\times 𝐁,$$
(38)
where we assumed that the electric component of the electromagnetic field is equal to zero. Eqs. (3738) describe forces acting on the particle in the Stern–Gerlach experiment, therefore we can really identify $`𝝃`$ with the polarization of the particle.
In this simple case of the monochromatic particle we can also calculate explicitly the von Nuemann entropy of the reduced density matrix. The matrix $`\mathrm{\Omega }`$ in the rest frame of the particle can be written as
$$\mathrm{\Omega }_0=\frac{1}{2}\left(\begin{array}{cc}1& 1\\ 1& 1\end{array}\right)\frac{1}{2}(I+𝝃𝝈).$$
(39)
To calculate entropy we have to use the normalized density matrix $`\stackrel{~}{\theta }_0`$, but in this particular case $`\stackrel{~}{\theta }_0=\mathrm{\Omega }_0`$. Thus the von Neumann entropy of the state (39) is equal to
$$S_{\stackrel{~}{\theta }_0}=\frac{1}{2}\left((1+|𝝃|)\mathrm{ln}\frac{1+|𝝃|}{2}+(1|𝝃|)\mathrm{ln}\frac{1|𝝃|}{2}\right).$$
(40)
Now, to find the entropy in the arbitrary Lorentz frame we apply to the matrix $`\stackrel{~}{\theta }_0`$ the Lorentz transformation (11) with $`\text{D}(\mathrm{\Lambda })`$ given by (15) and we find that entropy of the corresponding reduced density matrix $`\stackrel{~}{\theta }^{}`$ is given by (40) too, i.e., $`S_{\stackrel{~}{\theta }_0}=S_{\stackrel{~}{\theta }^{}}`$. Therefore for a particle with the sharp momentum the entropy of the reduced density matrix does not change under Lorentz transformations. However, in the case of an arbitrary momentum distribution, the entropy of the reduced density matrix $`\stackrel{~}{\theta }`$ is not in general Lorentz–invariant.
## IV Spin operator
In the next section we will use our formalism to calculate the Einstein–Podolsky–Rosen–Bohm (EPR–Bohm) correlation function. Thus we have to introduce the spin operator for a relativistic massive particle. The choice is not obvious since in the discussion of relativistic EPR–Bohm experiments various spin operators have been used Ahn et al. (2003); Czachor (1997); Rembieliński and Smoliński (2002); Lee and Chang-Young (2004); Li and Du (2003); Terashima and Ueda (2003a, b). However our previous considerations (Eqs. (33)–(38)) as well as the classical definition of the relativistic spin Anderson (1967) suggest that the best candidate for the spin operator is
$$\widehat{𝐒}=\frac{1}{m}\left(\widehat{𝐖}\widehat{W}^0\frac{\widehat{𝐏}}{\widehat{P}^0+m}\right),$$
(41)
which corresponds to the classical polarization vector $`𝝃`$ (precisely to $`𝝃/2`$) in Eq. (33). This operator is also used in quantum field theory Bogolubov et al. (1975). It fulfills the following standard commutation relations:
$$[\widehat{J}^i,\widehat{S}^j]=i\epsilon _{ijk}\widehat{S}^k,$$
(42a)
$$[\widehat{S}^i,\widehat{S}^j]=i\epsilon _{ijk}\widehat{S}^k,$$
(42b)
$$[\widehat{P}^\mu ,\widehat{S}^j]=0,$$
(42c)
which should be satisfied for the spin operator. Here $`\widehat{J}^i=\frac{1}{2}\epsilon _{ijk}\widehat{J}^{jk}`$ and one can show that it is the only operator which is a linear function of $`\widehat{W}^\mu `$ and fulfills relations (42) Bogolubov et al. (1975).
Therefore the operator corresponding to the spin projection along arbitrary direction $`𝐧`$ ($`𝐧^2=1`$) in the representation of gamma matrices (56) reads explicitly
$$\begin{array}{c}𝐧\widehat{𝐒}=\frac{1}{2m}\{\widehat{P}^0\left(\begin{array}{cc}𝐧𝝈& 0\\ 0& 𝐧𝝈\end{array}\right)\hfill \\ \hfill i\left(\begin{array}{cc}(𝐧\times \widehat{𝐏})𝝈& 0\\ 0& (𝐧\times \widehat{𝐏})𝝈\end{array}\right)\\ \hfill \frac{𝐧\widehat{𝐏}}{\widehat{P}^0+m}\left(\begin{array}{cc}\widehat{𝐏}𝝈& 0\\ 0& \widehat{𝐏}𝝈\end{array}\right)\},\end{array}$$
(43)
where we have used Eqs. (62).
Eq. (42b) implies that eigenvalues of the operator $`𝐧\widehat{𝐒}`$ are integers or half–integers. As one can easily check by direct calculation the eigenvalues of the operator (43) are equal to $`\pm \frac{1}{2}`$. This observation supports our choice of the operator $`\widehat{𝐒}`$ as the spin operator.
Now we want to express the average of the spin operator (41) in terms of the reduced matrix $`\mathrm{\Omega }`$. One can check that in an arbitrary state $`\widehat{\rho }`$
$$(\widehat{P}^0+m)\widehat{𝐒}_{\widehat{\rho }}=\frac{m}{2}\text{Tr}\left(\mathrm{\Omega }𝜸\gamma ^5(I+\gamma ^0)\right).$$
(44)
Thus a reasonable choice for the normalized average of the spin component is
$$𝚺=\frac{(\widehat{P}^0+m)\widehat{𝐒}_{\widehat{\rho }}}{(\widehat{P}^0+m)_{\widehat{\rho }}}=\frac{\text{Tr}\left(\mathrm{\Omega }𝜸\gamma ^5(I+\gamma ^0)\right)}{2\text{Tr}\left(\mathrm{\Omega }(I+\gamma ^0)\right)}.$$
(45)
When $`\widehat{\rho }(k)`$ describes a particle with a sharp momentum $`k`$ the normalized average is simply the average of $`\widehat{S}`$, i.e., inserting reduced density matrix $`\mathrm{\Omega }`$ (32) into (45) we get
$$𝚺=\widehat{𝐒}_{\widehat{\rho }(k)}=\frac{𝝃}{2}.$$
(46)
It should also be noted that in the nonrelativistic limit we recover the result (46) for an arbitrary state $`\widehat{\rho }`$
$$𝚺=\widehat{𝐒}_\rho =\frac{𝝃}{2}.$$
(47)
## V Quantum correlations
Using the formalism introduced above, we now calculate the correlation between measurements of spin components performed by two observers, A and B, along two arbitrary directions, $`𝐚`$ and $`𝐛`$, respectively. We consider the simplest situation in which both observers are at rest with respect to a certain inertial frame of reference $`𝒪`$. We assume also that both measurements are performed simultaneously in the frame $`𝒪`$.
We calculate the EPR–Bohm correlation function in the pure state of two particles with sharp momenta
$$|\psi =\underset{\alpha \beta }{}c_{\alpha \beta }|\alpha ,k|\beta ,p.$$
(48)
The corresponding reduced density matrix (30) has the following form:
$$\begin{array}{c}\mathrm{\Omega }_{\alpha \beta ,\alpha ^{}\beta ^{}}^\psi =\frac{4k^0p^0(\delta ^3(\mathrm{𝟎}))^2}{\psi |\psi }\hfill \\ \hfill \left[v(k)\overline{v}(k)\gamma ^0C^{}\gamma ^0\left(v(p)\overline{v}(p)\right)^T\right]_{\alpha \beta }\\ \hfill \left[\left(v(k)\overline{v}(k)\right)^TCv(p)\overline{v}(p)\right]_{\alpha ^{}\beta ^{}},\end{array}$$
(49)
where the matrix $`C=(c_{\alpha \beta })`$ determines the state (48) while $`v(k)\overline{v}(k)`$ and $`v(p)\overline{v}(p)`$ are given by (57b).
Observers A and B use observables $`2𝐚\widehat{𝐒}I`$ and $`I2𝐛\widehat{𝐒}`$, respectively ($`𝐚^2=𝐛^2=1`$). Thus the correlation function has the form (see Eqs. (44) and (64))
$`𝒞(𝐚,𝐛)`$ $`=4{\displaystyle \frac{(\widehat{P}^0+m)𝐚\widehat{𝐒}(\widehat{P}^0+m)𝐛\widehat{𝐒}_\psi }{(\widehat{P}^0+m)(\widehat{P}^0+m)_\psi }}`$
$`={\displaystyle \frac{\text{Tr}\left[\mathrm{\Omega }^\psi \left((𝐚𝜸\gamma ^5(I+\gamma ^0))(𝐛𝜸\gamma ^5(I+\gamma ^0))\right)\right]}{\text{Tr}\left[\mathrm{\Omega }^\psi \left((\gamma ^0+I)(\gamma ^0+I)\right)\right]}}`$
$`=4{\displaystyle \frac{\psi \left|𝐚\widehat{𝐒}𝐛\widehat{𝐒}\right|\psi }{\psi |\psi }}.`$ (50)
After some algebra we find that
$$\begin{array}{c}𝒞(𝐚,𝐛)=\hfill \\ \hfill \frac{\text{Tr}\left\{\left(𝐛𝐒(p)v(p)\overline{v}(p)\gamma ^0\right)C^{}\left(𝐚𝐒(k)v(k)\overline{v}(k)\gamma ^0\right)^TC\right\}}{\text{Tr}\left\{\left(v(p)\overline{v}(p)\gamma ^0\right)C^{}\left(v(k)\overline{v}(k)\gamma ^0\right)^TC\right\}},\end{array}$$
(51)
where matrices $`𝐚𝐒(k)`$ and $`𝐛𝐒(p)`$ have the same form as (43) with $`𝐧`$, $`\widehat{P}`$ equal to $`𝐚`$, $`k`$ and $`𝐛`$, $`p`$ respectively. Now, for the sake of simplicity, we specify the state $`|\psi `$. We choose
$$C=a\left(\begin{array}{cc}\sigma _2& 0\\ 0& \sigma _2\end{array}\right),$$
(52)
where $`a`$ is a normalization constant. This choice is rather natural because the state described by Eqs. (48), (52) has the same form for all inertial observers, namely
$$U(\mathrm{\Lambda })U(\mathrm{\Lambda })|\psi =\underset{\alpha \beta }{}c_{\alpha \beta }|\alpha ,\mathrm{\Lambda }k|\beta ,\mathrm{\Lambda }p,$$
(53)
where we used Eq. (15). Moreover in the center of mass frame it is an ordinary singlet state.
Now, provided that $`C`$ is given by Eq. (52), after straightforward calculation we arrive at
$$\begin{array}{c}𝒞(𝐚,𝐛)=𝐚𝐛+\frac{(𝐤\times 𝐩)}{m^2+kp}((𝐚\times 𝐛)\hfill \\ \hfill +\frac{(𝐚𝐤)(𝐛\times 𝐩)(𝐛𝐩)(𝐚\times 𝐤)}{(k^0+m)(p^0+m)}).\end{array}$$
(54)
We see that the correction to the nonrelativistic correlation function $`\mathrm{\Delta }𝒞=𝒞(𝐚,𝐛)𝒞_{\text{nonrel}}=𝒞(𝐚,𝐛)+\mathrm{𝐚𝐛}`$ is of order $`\beta ^2`$, where $`\beta =\frac{v}{c}`$, $`v`$ denotes the velocity of the particle, and $`c`$ the velocity of light. Let us note first that when momenta of both particles are parallel or antiparallel the correlation function has the same form as in the nonrelativistic case. This result differs from Czachor’s results<sup>2</sup><sup>2</sup>2The Czachor’s result can be obtained in our framework by calculating the appriopriately normalized average of $`𝐚\widehat{𝐖}𝐛\widehat{𝐖}`$. Czachor (1997). The reason is that we use a different, and in our opinion more adequate, spin operator.
Now let us consider the configuration in which the nonrelativistic correlation vanishes ($`𝐚𝐛`$). For simplicity let us assume also that $`|𝐤|=|𝐩|`$ and $`𝐚𝐤`$, $`𝐛𝐩`$ or $`𝐚𝐤`$, $`𝐛𝐩`$. In such fixed configurations the correlation function has the very simple form
$$𝒞(𝐚,𝐛)=\mathrm{\Delta }𝒞=\frac{p_0^2m^2}{p_0^2+m^2}=\frac{\beta ^2}{2\beta ^2}.$$
(55)
Dependence of the above correlation on $`\beta `$ is depicted in Fig. 1. Notice that (55) was also obtained by Czachor Czachor (1997) but for a different configuration.
## VI Conclusions
To conclude, we have constructed a Lorentz–covariant reduced spin density matrix for a single massive particle. It contains not only information about average polarization of the particle but also information about its average kinematical state. We have also showed that this matrix has the proper nonrelativistic limit.
Our results shows that we can define a Lorentz–covariant finite–dimensional matrix describing polarization of a massive particle. However in the relativistic case (contrary to the nonrelativistic one) we cannot completely separate kinematical degrees of freedom if we want to construct a finite-dimensional covariant description of the polarization degrees of freedom.
With help of our covariant formalism we have also calculated the correlation function in the EPR–Bohm type experiment with massive relativistic particles. We have showed that relativistic correction $`\mathrm{\Delta }𝒞`$ to the nonrelativistic correlation function $`𝒞_{\text{nonrel}}=𝐚𝐛`$ vanishes when momenta of both particles are parallel or antiparallel, i.e., in the standard configuration of EPR–Bohm type experiments. We have found also the configurations in which the nonrelativistic correlation vanishes while the relativistic correction $`\mathrm{\Delta }𝒞`$ survives and is of order $`\beta ^2`$ (Eq. (55)).
###### Acknowledgements.
The authors thank Marek Czachor for interesting discussions. This paper has been partially supported by the Polish Ministry of Scientific Research and Information Technology under Grant No. PBZ-MIN-008/P03/2003 and partially by the University of Lodz grant.
## Appendix A Dirac matrices
In this paper we use the following conventions. Dirac matrices fulfills the condition $`\gamma ^\mu \gamma ^\nu +\gamma ^\nu \gamma ^\mu =2g^{\mu \nu }`$ where $`g^{\mu \nu }=\text{diag}(1,1,1,1)`$ denotes Minkowski metric tensor; moreover we adopt the convention $`\epsilon ^{0123}=1`$. We use the following explicit representation of gamma matrices:
$$\gamma ^0=\left(\begin{array}{cc}0& I\\ I& 0\end{array}\right),𝜸=\left(\begin{array}{cc}0& 𝝈\\ 𝝈& 0\end{array}\right),\gamma ^5=\left(\begin{array}{cc}I& 0\\ 0& I\end{array}\right),$$
(56)
where $`𝝈=(\sigma _1,\sigma _2,\sigma _3)`$ and $`\sigma _i`$ are standard Pauli matrices.
## Appendix B Useful formulas
The matrix (17) is normalized as follows
$$\overline{v}(k)v(k)=I,$$
(57a)
$$v(k)\overline{v}(k)=\frac{1}{2m}(k\gamma +mI),$$
(57b)
where $`\overline{v}(k)=v^{}(k)\gamma ^0`$. Moreover it can be verified that it fulfills the following relation
$$\overline{v}(k)\gamma ^\mu v(k)=\frac{k^\mu }{m}I.$$
(58)
Vectors $`|\alpha ,k`$ fulfill the orthogonality relation:
$$\alpha ,k|\beta ,p=2k^0\delta ^3(𝐤𝐩)\left(v(k)v^{}(k)\right)_{\beta \alpha },$$
(59)
and one can check that
$$I=\underset{\alpha \beta }{}𝑑\mu (k)\gamma _{\beta \alpha }^0|\alpha ,k\beta ,k|,$$
(60)
$$\left(v(k)\overline{v}(k)\right)_{\alpha \beta }|\beta ,k=|\alpha ,k.$$
(61)
In the representation of gamma matrices (56) we have
$$\widehat{W}^0=\frac{1}{2}\left(\begin{array}{cc}\widehat{𝐏}𝝈& 0\\ 0& \widehat{𝐏}𝝈\end{array}\right),$$
(62a)
$$\widehat{𝐖}=\frac{1}{2}\widehat{P}^0\left(\begin{array}{cc}𝝈& 0\\ 0& 𝝈\end{array}\right)\frac{i}{2}\left(\begin{array}{cc}\widehat{𝐏}\times 𝝈& 0\\ 0& \widehat{𝐏}\times 𝝈\end{array}\right).$$
(62b)
It can be also checked, that when
$$\widehat{F}_\rho =\text{Tr}(\mathrm{\Omega }f),$$
(63a)
$$\widehat{G}_\rho =\text{Tr}(\mathrm{\Omega }g),$$
(63b)
we have
$$\widehat{F}\widehat{G}_\rho =\text{Tr}(\mathrm{\Omega }(fg)),$$
(64)
where in Eqs. (63) and (64) $`\rho `$ and $`\mathrm{\Omega }`$ are complete and reduced density matrices for one and two particles, respectively. |
warning/0507/hep-ph0507248.html | ar5iv | text | # References
Photon Irradiated Enhancement as a Tool of Investigating Fundamental Physics beyond Standard Model
M. Yoshimura
Department of Physics, Okayama University
Tsushima-naka 3-1-1 Okayama Japan 700-8530
ABSTRACT
We clarify how intense laser irradiation leads to an enhancement of rare processes that may occur within atoms. Non-perturbative calculation using a coherent laser beam gives an exact, time dependent formula of the enhancement factor in the large power limit. A high quality laser may provide a new tool of experimentally investigating physics far beyond the standard model of particle physics such as lepton and baryon nonconservation.
In our recent paper it was pointed out that intense laser beam, when irradiated to appropriate heavy atoms, can enhance rare processes related to atomic electron capture by nucleus, otherwise difficult to detect. In the present work we extend the formalism such that non-perturbative effects of laser-atom interaction are fully taken into account by directly solving dynamics of laser irradiation onto target atoms. In the dipole approximation of the two level problem of atoms the result suggests an onset of a repeated process of compression and expansion of the atomic electron cloud. Rare processes that may subsequently occur via the overlap of atomic electrons with nucleus may thus be enhanced. We propose to call this compression mechanism as PHIRAC to abbreviate PHoton IRrAdiated Compression of the electron cloud. For magnetic transitions such as those between hyperfine split levels this interpretation of electron cloud compression is not appropriate, but one can obtain a large enhancement of rare processes as well.
Our discussion in the present work is independent of any particular rare process $`X`$ subsequent to photon absorption. It may thus provide a new method of exploring physics far beyond the standard model of particle physics. In this way one can experimentally explore lepton number nonconservation (LENNON) of the type $`e^{}e^+`$, baryon number nonconserving (BARNNON) process of the kind $`e^{}+N`$ many $`\pi ^{}`$s, and hopefully many other fundamental processes that face physics beyond the standard model.
Consider laser irradiation which induces transitions between two electronic levels of a target atom, denoted by $`|e`$ and $`|g`$. We imagine that the target is irradiated continuously. For instance putting the target into a resonator may be useful. The system of laser and atom is then described by a Hamiltonian of the form ,
$$H=\frac{\omega _{eg}}{2}\sigma _3+\omega a^{}a+\stackrel{~}{s}(\sigma _+a+\sigma _{}a^{}).$$
(1)
Here the Pauli matrices $`\sigma _i`$ act on two levels of the ground $`|g`$ and the excited $`|e`$ state, and $`\omega _{eg}=\omega _e\omega _g`$ is the energy level difference. We made what is called the rotating wave approximation in the literature , by neglecting $`\sigma _+a^{}`$ and its conjugate, which should be justified for mostly energy-conserving processes. The coupling strength $`s`$ is given by
$$e|\stackrel{~}{s}\sigma _+|gs=e|\stackrel{}{d}|g\stackrel{}{e},\stackrel{}{e}=i\frac{\omega }{\sqrt{2\omega V}}\stackrel{}{ϵ}_k$$
(2)
for the dipole transtion, with $`\stackrel{}{d}`$ the dipole operator of atomic electron and $`\stackrel{}{e}`$ the electric field of a single photon beam with polarization $`\stackrel{}{ϵ}_k`$. When multiplied by a photon number $`N`$, this strength is expressed as $`Ns^2=\pi \gamma _dP/\omega ^3`$ where $`P`$ is the laser power in the unit of energy/ (time $`times`$ area).
The Hamiltonian is block-diagonal, hence the effect of laser irradiation can be solved by decomposing the infinite dimensional Fock space and using a mixture of two states, $`|e,n=(a^{})^n/\sqrt{n!}|e`$ and $`|g,n=(a^{})^n/\sqrt{n!}|g`$; $`|\psi (t)=_n(c_{gn+1}(t)|g,n+1+c_{en}(t)|e,n).`$ The Schroedinger equation to be solved is
$`i{\displaystyle \frac{d}{dt}}\left(\begin{array}{c}c_{en}\\ c_{gn+1}\end{array}\right)=\left(\begin{array}{c}𝒸_𝓃\\ 𝒸_{𝓃+\mathcal{1}}\end{array}\right),=\left(\begin{array}{cc}𝓃\omega +\frac{\omega _{}}{\mathcal{2}}& 𝓈\sqrt{𝓃+\mathcal{1}}\\ 𝓈\sqrt{𝓃+\mathcal{1}}& (𝓃+\mathcal{1})\omega \frac{\omega _{}}{\mathcal{2}}\end{array}\right).`$ (9)
The solution may be written in terms of what is called dressed states denoted by $`|\pm ,n`$; assuming a spacially constant (valid in the long wavelength approximation) laser field, the Hamiltonian diagonalization is possible with
$`|+,n=\mathrm{cos}{\displaystyle \frac{\phi _n}{2}}|e,n+\mathrm{sin}{\displaystyle \frac{\phi _n}{2}}|g,n+1,`$ (10)
$`|,n=\mathrm{sin}{\displaystyle \frac{\phi _n}{2}}|e,n+\mathrm{cos}{\displaystyle \frac{\phi _n}{2}}|g,n+1,`$ (11)
where
$`\mathrm{tan}\phi _n={\displaystyle \frac{2s\sqrt{n+1}}{\omega \omega _{eg}}},\omega _\pm =(n+{\displaystyle \frac{1}{2}})\omega \pm {\displaystyle \frac{\mathrm{\Omega }_n}{2}},`$ (12)
with $`\mathrm{\Omega }_n=\sqrt{(\omega \omega _{eg})^2+4s^2(n+1)}`$ the Rabi frequency . Unless the photon energy $`\omega `$ is very far from the resonance energy $`\omega _{eg}`$, the mixing is nearly maximal; $`\mathrm{sin}^2\phi _n1.`$ We thus assume the maximal mixing for simplicity; $`\phi _n=\pi /2`$.
It is reasonable under the continuous laser irradiation to assume that the target is initially in a superposed state of two levels, $`|\frac{1}{2}\frac{1}{\sqrt{2}}(e^{i\delta }|g+|e)`$ with $`\delta =\pi /2`$. The result is insensitive to the choice of this phase and the initial condition as a whole.
Time evolution is then given by
$`g,n+1|{\displaystyle \frac{1}{2}},n;t=ie^{i(n+1/2)\omega t}\mathrm{cos}({\displaystyle \frac{\mathrm{\Omega }_nt}{2}}+{\displaystyle \frac{\pi }{4}})c_n^{(\gamma )}(0).`$ (13)
In the rest of this paper we take the coherent state of laser, thus $`c_n^{(\gamma )}(0)=e^{N/2}N^n/\sqrt{n!},`$ where the average photon number $`n=N`$ and the dispersion $`(\mathrm{\Delta }n)^2=N`$.
It is conceptually important to distinguish from the Schroedinger (S) picture and use the interaction (I) picture, since the subsequent rare process is treated in the perturbation theory. Fortunately, since the laser-atom system is exactly solvable, one may straightforwardly use the identity $`{}_{S}{}^{}a|b_{S}^{}=_Ia|b_I,`$ to simplify computation.
We imagine a circumstance under which rare processes occur via the ground state $`|g`$ of zero angular momentum, a $`ms`$ state. It is assumed that at a time $`t`$ the electron in the excited $`|e`$ goes to the $`ms`$ $`|g`$ state due to the stimulated emmision, and is subsequently captured by nucleus, with a probability proportional to the wave function factor $`|\psi _{ms}(0)|^2`$. An example of subsequent rare processes of this sort is LENNON electron capture of the kind $`e^{}e^+`$, as discussed in . The probability amplitude of laser irradiated rare process is then given by
$`\stackrel{~}{}_{1/2}(t)=e^{(\gamma _e+\gamma _g)t/4}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}_X(t){\displaystyle \frac{s\sqrt{n+1}}{\omega \omega _{eg}+i\gamma /2}}g,n+1|{\displaystyle \frac{1}{2}},n;t`$ (14)
where $`_X(t)`$ is the amplitude for the rare $`X`$ process from the $`ms`$ state.
Computation of the discrete $`n`$ sum (14) is well approximated in the large $`N`$ limit by a continuous $`n`$ integral. Using the large $`n`$ limit formula of $`n!`$, and the coherent state expression for $`c_n^{(\gamma )}(0)`$, the integrand is found to change violently, and one may estimate the integral by a gaussian approximation around the stationary phase, or the saddle point. The saddle point $`n_0`$ of the integrand is determined by the minimal variation of the exponent of the integrand and is given by taking the $`n`$derivative of the exponent to vanish. A complex saddle is thus obtained; $`n_0Ne^{2i\omega t}.`$ Making the gaussian approximation around this saddle gives the rate formula,
$`|\stackrel{~}{}_{1/2}(t)|^2({\displaystyle \frac{\pi }{2}})^{1/2}N^{3/2}{\displaystyle \frac{s^2|_X(t)|^2}{(\omega \omega _{eg})^2+\gamma ^2/4}}`$
$`\mathrm{exp}[N(1\mathrm{cos}2\omega t)(\gamma _e+\gamma _g)t/2][\mathrm{sinh}^2(\sqrt{N}st\mathrm{sin}\omega t)+\mathrm{cos}^2(\sqrt{N}st\mathrm{cos}\omega t)].`$
As a funtion of time, the rate becomes very large of order $`N^{3/2},`$ periodically at $`\mathrm{cos}2\omega t=1`$. Outside regions of a time range $`1/(\omega \sqrt{N})`$ the rate is very small, of order $`N^{3/2}e^N`$. Thus, the rate is sizable only around infinitely many discrete times of $`t=k\pi /\omega `$, with $`k`$ any positive integer. In other words, the rate has a spiky time profile with a period $`\pi /\omega `$. A formula valid for $`N1`$
$`\mathrm{exp}[N(1\mathrm{cos}2\omega t)]\sqrt{{\displaystyle \frac{\pi }{2}}}{\displaystyle \frac{1}{\omega \sqrt{N}}}{\displaystyle \underset{k}{}}\delta (t{\displaystyle \frac{k\pi }{\omega }}),`$ (16)
may then be used. For $`\omega t1`$, it is reasonable to take a time average over $`\mathrm{\Delta }t\pi /\omega \times `$ a few, which gives a time variant averaged rate,
$`\stackrel{~}{}_{1/2}(t)={\displaystyle \frac{N}{2}}{\displaystyle \frac{s^2e^{(\gamma _e+\gamma _g)t/2}}{(\omega \omega _{eg})^2+\gamma ^2/4}}_X(t),`$ (17)
with $`_X(t)=d|_X(t)|^2/dt.`$
The last factor $`_X(t)`$ may differ in rare processes in which one is interested. For exmaple, the LENNON conversion of the type $`e^{}e^+`$ has $`_{e^{}e^+}(t)=|\psi _{ms}(0)|^2\sigma _{e^{}e^+},`$ where $`\sigma _{e^{}e^+}`$ is the cross section of free electron capture, a virtual process considered for our gedanken experiment. This quantity is computed using the perturbation theory.
One may summarize arguments so far by defining a quality factor $`Q(\omega )`$ which signifies the rate enhancement,
$`Qr{\displaystyle \frac{\pi }{2}}{\displaystyle \frac{\gamma _d}{\omega _0^2}}{\displaystyle 𝑑\omega \frac{𝒫(\omega )}{\omega [(\omega \omega _0)^2+\gamma ^2/4]}},`$ (18)
where $`𝑑\omega 𝒫(\omega )=P`$ is the total laser power in the unit of energy /(time $`\times `$ area). The factor $`r`$ is the wave function ratio sqaured, for instance, for LENNON
$`r={\displaystyle \frac{|\psi _{ms}(0)|^2}{|\psi _{ns}(0)|^2}}=({\displaystyle \frac{r_{ns}}{r_{ms}}})^3O[({\displaystyle \frac{n}{m}})^6],n=1.`$ (19)
For a laser beam of the energy resolution $`\mathrm{\Delta }E\gamma `$ we may replace $`\mathrm{\hspace{0.25em}1}/[(\omega \omega _0)^2+\gamma ^2/4]`$ by $`\frac{2\pi }{\gamma }\delta (\omega \omega _0)`$ for a laser beam of Lorentzian energy distribution. When the laser tuning is perfect, $`𝒫(\omega _0)P/\mathrm{\Delta }E.`$ It is thus found that
$`Qr{\displaystyle \frac{\pi ^2P}{\omega _0^4}}{\displaystyle \frac{\omega _0}{\mathrm{\Delta }E}}1.6\times 10^6r{\displaystyle \frac{P}{Wmm^2}}({\displaystyle \frac{\omega _0}{eV}})^4(10^9{\displaystyle \frac{\omega _0}{\mathrm{\Delta }E}}).`$ (20)
Note a strong dependence on the photon energy $`\omega _0^4`$, which should be important to get a large quality factor for BARRNON. This formula for the quality (20) agrees with the result of after a minor correction . The spiky time profile is however a result of the present nonperturbative formalism.
In order to advance an intuitive understanding of the enhancement mechanism, it might be instructive to compute the electron displacement squared
$`(\delta 𝒟)^2(t)=d^2{\displaystyle \underset{n}{}}{\displaystyle \frac{1}{2}},n;t|\sigma _{}\sigma _+a^{}a|{\displaystyle \frac{1}{2}},n;t`$
$`{\displaystyle \frac{d^2}{2}}{\displaystyle \underset{n}{}}(1\mathrm{sin}\mathrm{\Omega }_nt)|c_n^{(\gamma )}(0)|^2.`$ (21)
The leading term of $`O[d^2N/2]`$ simply shows that quantum mechanics gives a result consistent with the classical Lorentz oscillator model.
The behavior of the next leading term in the large $`N`$ limit of $`(\delta 𝒟)^2(t)`$ and the linear dipole $`𝒟(t)`$ is more complicated. They exhibit a spiky time profile in much the same way as the rate formula, but with an important difference of time scale; these quantities that have classical analogues vary in time much more slowly, the spike interval being given by $`4\pi /\mathrm{\Omega }_R`$ with $`\mathrm{\Omega }_R=2s\sqrt{N}`$ the Rabi resonance frequency. We may only say that the spiky time profile observed in the rate enhancement is a signal of the onset of the oscillatory behavior of the electron displacement. The link is however indirect.
Baryon nonconservation may experimentally be investigated by searching for the atomic electron capture of the type , $`e^{}+N\pi +\pi .`$ Enhancement factor $`P`$ may compensate the small nuclear overlap factor of order $`(a_Bm_\pi )^310^{15}`$ of atomic electrons that otherwise disfavors this process. A rough estimate of the rate gives the enhancement factor of order, $`Q(r_AA^{1/3}m_\pi )^3,`$ where the pion mass $`m_\pi `$ times $`A^{1/3}`$ gives a measure of the inverse nuclear size, and $`r_A`$ is the atomic size. It is important to go to a low frequency range for a large enhancement of order $`10^{25}`$. Thus, the use of the Zeeman split hyperfine levels is promissing, which involves the microwave region.
In summary, the prospect of search for new physics far beyond the standard model appears bright if one uses an appropriate high quality laser.
I would like to thank Y. Kuno, I. Nakano, and N. Sasao for stimulating discussions. |
warning/0507/hep-th0507117.html | ar5iv | text | # 𝑁=4 Super Yang-Mills NMHV Loop Amplitude in Superspace
## I 1.Introduction
N=4 SuperYang-Mills one loop amplitudes have the special property of being cut constructible 1 2 , that is they are uniquely determined by their unitary cuts. It was shown in 1 that loop amplitudes can be written as a combination of scalar box integrals (for definition of the scalar box integral and scalar box function see appendix for 1 ) with rational coefficients. Therefore the calculation of one loop amplitude is reduced to determining the coefficients in front of these box integrals, which is done by analyzing the cuts of the amplitude and matching them with the cuts of the scalar box integrals. Unfortunately complication arises from the fact that some of the cuts are shared by more than one box integral. This was dramatically simplified in 3 by using generalized unitarity (quadruple cuts) cuts 4 to analyze the leading singularities which turns out to be unique in the box integrals.
Recently 5 it was shown by using Supersymmetric Ward Identity (SWI)6 one can derive N $`=4`$ SYM NMHV 6 point tree and loop amplitudes with gluinos or scalars from their pure gluonic partners. Since SWI corresponds to a transformation in superspace, one would guess this implies the existence of a full superspace amplitude while amplitudes with different external particle species are considered as different component of the superspace expansion. The simplest superspace amplitude was the MHV and $`\overline{MHV}`$ tree amplitude written down by Nair 7 . Since the work by Witten 8 which showed that perturbative N=4 SYM is dual to a particular string theory with super-twistor space as its target space, various new techniques have been developed to calculate N=4 SYM amplitudes more efficiently9 10 . The MHV vertex construction 10 , which uses MHV vertices as the basic building block of the scattering amplitudes, provides a convenient method to construct the amplitudes in superspace form. This was done for the NMHV tree amplitude in 11 . At loop level the valedictory of MHV vertex approach was proven to give the same result as that in field theory in 12 for MHV loop amplitudes, and 13 reproduces the relationship between the color leading amplitudes and sub-leading amplitudes. At this point it is natural to continue with MHV vertices to compute NMHV loop which would require three MHV vertices connected by three propagators, and this should give the full superspace form of the NMHV loop. At this point it is not clear how the correct scalar box functions should arise in this formalism. One of the complication is for more than two fermionic delta functions (there is one for each MHV vertex), after the expansion in superspace there will be multiple spinor products that contain the off shell continuation spinor of the propagator, which takes different form with different external particle specie. Since these spinor products should be integrated over, the integrand for the gluonic amplitudes will be dramatically different from the ones with gluinos, implying one can only derive the box functions from the superspace expansion one term at a time and not in the original superspace full form.
In 5 the SWI identities were not used directly upon the coefficients in front of the box integrals for the gluonic amplitude, but rather the coefficients in front of a particular combination of box integrals, which originated from the three different cut channels 2 . To realize the superspace amplitude all one needs to observe is that for the six point amplitude the three channels from which the cuts were computed, the tree graphs on either side of the cuts always come in MHV and $`\overline{MHV}`$ pair. Since MHV and $`\overline{MHV}`$ tree can be written straight forwardly in superspace form, one naturally derives the six point one loop NMHV amplitude for all helicity configuration and external species as one superspace amplitude by fusing the two tree amplitude. In the following we present the amplitude in its full superspace form and confirm our result by explicitly expanding out the terms that give the correct amplitudes with two gluino obtained in 5 . We will also give a brief demonstration of how one could obtain the field theory result for the loop amplitude from the MHV vertex approach (CSW) 10 .
## II 2. The Construction
The n point MHV and $`\overline{MHV}`$ tree level amplitudes have a remarkable simple form. For MHV tree 7 :
$$A(\mathrm{}j^{}\mathrm{}..i^{}\mathrm{})_{tree}=\frac{\delta ^8(_{i=1}^n\lambda _i\eta _i^A)}{\mathrm{\Pi }_{i=1}^n<ii+1>}$$
(1)
where
$$\delta ^8(\underset{i=1}{\overset{n}{}}\lambda _i\eta _i^A)=\frac{1}{2}\underset{A=1}{\overset{4}{}}(\underset{i=1}{\overset{n}{}}\lambda _i^\alpha \eta _i^A)(\underset{i=1}{\overset{6}{}}\lambda _{i\alpha }\eta _i^A)$$
(2)
as for $`\overline{MHV}`$ tree:
$$A(\mathrm{}j^+\mathrm{}..i^+\mathrm{})=\frac{\delta ^8(_{i=1}^6\stackrel{~}{\lambda _i}\stackrel{~}{\eta _i^A)}}{\mathrm{\Pi }_{i=1}^6[ii+1]}$$
(3)
Here we’ve omitted the energy momentum conserving delta function and the group theory factor. After expansion in the fermionic parameters $`\eta _i^A`$, one can obtain MHV amplitudes with different helicity ordering ($`++,++`$…etc) and different particle content.
We proceed to construct the full N=4 SYM NMHV 1-loop six point amplitudes by following the original gluonic calculation 2 , where the amplitude was computed from the cuts of the three channels $`t_{123}`$ $`t_{234}`$ $`t_{345}`$ $`(t_{ijk}=(k_i+k_j+k_l)^2)`$, except now the tree amplitudes across the cuts are written in supersymmetric form. We find that the propagator momentum integrals from which the various scalar box functions arise are the same for different external particles. Thus with the gluon amplitude already computed all we need to do is extract away the part of the gluon coefficient that came from the expansion of the two fermionic delta function, the remaining pre factor will be universal and has its origin from the denominator of eq.(1) and (3). The N=4 SYM 6 point NMHV loop amplitude for the gluonic case was given 2 as
$$A(\mathrm{}j^{}\mathrm{}..i^{}\mathrm{})_{loop}=c_\mathrm{\Gamma }[B_1W_6^{(1)}+B_2W_6^{(2)}+B_3W_6^{(3)}]$$
(4)
where $`W_6^{(i)}`$ contains particular combination of the two-mass-hard and one-mass box functions 1 . The full 6 point NMHV loop amplitude for any given set of external particle and helicity ordering are then given with the following coefficients :
$`B_1={\displaystyle \frac{\delta ^8(_{i=1}^3\stackrel{~}{\lambda _i}\stackrel{~}{\eta _i}\stackrel{~}{l_1}\stackrel{~}{\eta _1}+\stackrel{~}{l_2}\stackrel{~}{\eta _2})\delta ^8(_{i=4}^6\lambda _i\eta _il_2\eta _2+l_1\eta _1)}{t_{123}}}B_0`$ (5)
$`+{\displaystyle \frac{\delta ^8(_{i=1}^3\lambda _i\eta _il_1\eta _1+l_2\eta _2)\delta ^8(_{i=4}^6\stackrel{~}{\lambda _i}\stackrel{~}{\eta _i}\stackrel{~}{l_2}\stackrel{~}{\eta _2}+\stackrel{~}{l_1}\stackrel{~}{\eta _1})}{t_{123}}}B_0^{}`$
$`B_2={\displaystyle \frac{\delta ^8(_{i=2}^4\stackrel{~}{\lambda _i}\stackrel{~}{\eta _i}\stackrel{~}{l_1}\stackrel{~}{\eta _1}+\stackrel{~}{l_2}\stackrel{~}{\eta _2})\delta ^8(_{i=5}^1\lambda _i\eta _il_2\eta _2+l_1\eta _1)}{t_{234}}}B_+`$ (6)
$`+{\displaystyle \frac{\delta ^8(_{i=2}^4\lambda _i\eta _il_1\eta _1+l_2\eta _2)\delta ^8(_{i=5}^1\stackrel{~}{\lambda _i}\stackrel{~}{\eta _i}\stackrel{~}{l_2}\stackrel{~}{\eta _2}+\stackrel{~}{l_1}\stackrel{~}{\eta _1})}{t_{234}}}B_+^{}`$
$`B_3={\displaystyle \frac{\delta ^8(_{i=3}^5\stackrel{~}{\lambda _i}\stackrel{~}{\eta _i}\stackrel{~}{l_1}\stackrel{~}{\eta _1}+\stackrel{~}{l_2}\stackrel{~}{\eta _2})\delta ^8(_{i=6}^2\lambda _i\eta _il_2\eta _2+l_1\eta _1)}{t_{345}}}B_{}`$ (7)
$`+{\displaystyle \frac{\delta ^8(_{i=3}^5\lambda _i\eta _il_1\eta _1+l_2\eta _2)\delta ^8(_{i=6}^2\stackrel{~}{\lambda _i}\stackrel{~}{\eta _i}\stackrel{~}{l_2}\stackrel{~}{\eta _2}+\stackrel{~}{l_1}\stackrel{~}{\eta _1})}{t_{345}}}B_{}^{}`$
where we define :
$$B_0=i\frac{1}{[12][23]<45><56><1|K_{123}|4><3|K_{123}|6>}$$
(8)
and
$$B_+=B_0|_{jj+1}B_{}=B_0|_{jj1}$$
(9)
with $`<A|K_{ijk}|B>=[Ai]iB+[Aj]jB+[Ak]kB`$. Each coefficient is expressed in two terms, this corresponds to the assignment of helicity for the propagators $`l_1`$ and $`l_2`$ which for specific assignments will reverse the MHV and $`\overline{MHV}`$ nature of the two tree amplitude across the cut (fig-1). The presence of the loop momenta seems perplexing at this point since all loop momenta should have been integrated out to give the box functions. As we will see on a case by case basis this comes as a blessing. The actual expansion for a particular set of helicity ordering and external particles contains multiple terms, the presence of loop momentum forces one to regroup the terms such that the loop momentum forms kinematic invariants, it is after this regrouping that one obtains previous known results.
The amplitudes for different external particles are computed as an expansion in the $`SU(4)_R`$ anti-commuting fermionic variables $`\eta `$. Choosing particular combinations following 14
$`g_i^{}=\eta _i^1\eta _i^2\eta _i^3\eta _i^4,\varphi _i^{AB}=\eta _i^A\eta _i^B,\mathrm{\Lambda }_i^1=\eta _i^2\eta _i^3\eta _i^4,\mathrm{\Lambda }_i^2=\eta _i^1\eta _i^3\eta _i^4`$ (10)
$`\mathrm{\Lambda }_i^3=\eta _i^1\eta _i^2\eta _i^4,\mathrm{\Lambda }_i^4=\eta _i^1\eta _i^2\eta _i^3,\mathrm{\Lambda }_i^{A+}=\eta _i^A,g_i^+=1`$
The superscript represents which flavor the particle carries, in the N=4 multiplet there are four gluinos and six scalars. Corresponding combination in the $`\stackrel{~}{\eta }`$ follows:
$`g_i^+=\stackrel{~}{\eta _i^1}\stackrel{~}{\eta _i^2}\stackrel{~}{\eta _i^3}\stackrel{~}{\eta _i^4},\varphi _i^{AB}=\stackrel{~}{\eta _i^C}\stackrel{~}{\eta _i^C},\mathrm{\Lambda }_i^{1+}=\stackrel{~}{\eta _i^2}\stackrel{~}{\eta _i^3}\stackrel{~}{\eta _i^4},\mathrm{\Lambda }_i^{2+}=\stackrel{~}{\eta _i^1}\stackrel{~}{\eta _i^3}\stackrel{~}{\eta _i^4}`$ (11)
$`\mathrm{\Lambda }_i^{3+}=\stackrel{~}{\eta _i^1}\stackrel{~}{\eta _i^2}\stackrel{~}{\eta _i^4},\mathrm{\Lambda }_i^{4+}=\stackrel{~}{\eta _i^1}\stackrel{~}{\eta _i^2}\stackrel{~}{\eta _i^3},\mathrm{\Lambda }_i^A=\stackrel{~}{\eta _i^A},g_i^{}=1`$
Thus a particular term in the expansion corresponds to a particular assignment of the fermionic variables to the external particle and results in an amplitude with a particular set of external particle specie and helicity ordering. In the next two section we show by expanding eq.(5),(6),(7) and following the above dictionary one can recover the amplitudes containing two same color gluino with different helicity ordering computed in 5 .
### II.1 2.1 $`B_1`$ Coefficient $``$ $`t_{123}`$ cut
First we look at the $`t_{123}`$ cut which correspond to the $`B_1`$ coefficient. For the purely gluonic amplitude $`A(g_1^{}g_2^{}g_3^{}|g_4^+g_5^+g_6^+)`$ (we use a bar to indicate the cut ), we have only one particle assignment for the loop propagators:
$$l_1=g^+,l_2=g^+$$
(12)
Here the assignment of helicity is labelled with respect to the $`\overline{MHV}`$ vertex. Therefore we get only contribution from the first term in eq.(5), the expansion from the delta function gives $`ł_1l_2^4[l_1l_2]^4=(l_1l_2)^8=t_{123}^4`$ and therefore $`B_1=t_{123}^3B_0`$ which matches eq.(5.4) in 5 .
For the two gluino amplitudes first we look at $`A(\mathrm{\Lambda }_1^{}g_2^{}g_3^{}|\mathrm{\Lambda }_4^+g_5^+g_6^+)`$ from the delta function expansion the we have helicity assignments :
$$l_1=\mathrm{\Lambda }^+l_2=g^+,+(exchagebetweenl_1andl_2)$$
(13)
Again only the first term in eq.(5) gives contribution :
$$l_1l_2^3[l_1l_2]^3([1l_1]l_14[1l_2]l_24)=t_{123}^31|K_{123}|4$$
(14)
Note that only when the external gluino carry the same flavor will this term contribute. Since in 5 the two gluino amplitude was derived using N=1 SWI, the two gluinos carry the same flavor. Thus we have
$$B_1(\mathrm{\Lambda }_1^{}g_2^{}g_3^{}|\mathrm{\Lambda }_4^+g_5^+g_6^+)=i\frac{t_{123}^21|K_{123}|4}{[12][23]<45><56><1|K_{123}|4><3|K_{123}|6>}$$
(15)
This is exactly the result of 5 . Other non-cyclic permutations of two gluino amplitude calculated in 5 at this cut do not change the assignment of the propagators thus the amplitude remains the same form apart from the labelling of the position of the two gluinos.
### II.2 2.2 $`B_2`$ Coefficient $``$$`t_{234}`$ cut
For this cut with different helicity assignment of the propagators, contribution can arise from both terms. Propagators with the same helicities (here we mean they are both plus or minus regardless of the specie) get its contribution from one term while the rest from the other, this is why $`B_2`$ was split in two terms in the original computation of the gluon amplitude 2 . We deal with the same helicity first since there is only one way of assigning propagators.
$`B_2(\mathrm{\Lambda }_1^{}|g_2^{}g_3^{}\mathrm{\Lambda }_4^+|g_5^+g_6^+)_{samehelicity}=0`$ since there is no way of assigning same helicity particles to the propagators.
For $`B_2(g_1^{}|\mathrm{\Lambda }_2^{}g_3^{}\mathrm{\Lambda }_4^+|g_5^+g_6^+)_{samehelicity}`$ we have
$$l_1=g^{}l_2=g^{}$$
(16)
this receives contribution from the second term in eq.(6) which is $`23^343[56]^4`$ thus giving
$$B_2(g_1^{}|\mathrm{\Lambda }_2^{}g_3^{}\mathrm{\Lambda }_4^+|g_5^+g_6^+)_{samehelicity}=(\frac{23^343[56]^4}{t_{234}})B_+^{}$$
(17)
For $`B_2(g_1^{}|\mathrm{\Lambda }_2^{}g_3^{}g_4^+|\mathrm{\Lambda }_5^+g_6^+)_{samehelicity}`$ we have
$$l_1=g^{}l_2=\mathrm{\Lambda }^{},+(exchangebetweenl_1,l_2)$$
(18)
This gives contribution $`23^3[56]^3(3l_1[l_16]3l_2[l_26])=23^3[56]^33|K_{234}|6`$ giving
$$B_2(g_1^{}|\mathrm{\Lambda }_2^{}g_3^{}g_4^+|\mathrm{\Lambda }_5^+g_6^+)_{samehelicity}=(\frac{23^3[56]^33|K_{234}|6}{t_{234}})B_+^{}$$
(19)
Now we move to configurations with different helicity. For $`B_2(\mathrm{\Lambda }_1^{}|g_2^{}g_3^{}\mathrm{\Lambda }_4^+|g_5^+g_6^+)_{DiffHelicity}`$ we have :
$$l_1=\mathrm{\Lambda }^+l_2=g^{},l_1=\mathrm{\Lambda }^{}l_2=\varphi ,+(exchangebetweenl_1,l_2)$$
(20)
For fix flavored $`\mathrm{\Lambda }_4^+`$ and $`\mathrm{\Lambda }_1^{}`$ we have to sum up all possible flavors for the internal gluino. This gives a contribution of
$`[1l_1]^3[l_1l_2]4l_2^3l_1l_23[l_1l_2][l_24][l_14]^2l_1l_2l_21l_11^2`$ (21)
$`+3[l_1l_2][l_14][l_24]^2l_1l_2l_11l_21^2`$
$`[1l_2]^3[l_1l_2]4l_2^3l_1l_2=t_{123}(1|l_1l_2|4)^3=t_{123}1|K_{123}|4^3`$
Therefore
$$B_2(\mathrm{\Lambda }_1^{}|g_2^{}g_3^{}\mathrm{\Lambda }_4^+|g_5^+g_6^+)_{DiffHelicity}=\frac{1|K_{123}|4^3}{t_{123}^3}B_+$$
(22)
For $`B_2(g_1^{}|\mathrm{\Lambda }_2^{}g_3^{}\mathrm{\Lambda }_4^+|g_5^+g_6^+)_{DiffHelicity}`$ we have:
$$l_1=g^{}l_2=g^+,l_1=\mathrm{\Lambda }^+l_2=\mathrm{\Lambda }^{},l_1=\varphi l_2=\varphi ,+(exchangebetweenl_1,l_2)$$
(23)
Here whether or not $`\mathrm{\Lambda }_4^+`$ and $`\mathrm{\Lambda }_1^{}`$ carry the same flavor will effect the number of ways one can assign flavor to the internal gluino and scalar. For the same flavor we have
$$([4l_1]l_11[4l_2]l_21)^3([2l_1]l_11[2l_2]l_21)=(4|K_{234}|1)^3(2|K_{234}|1)$$
(24)
Thus
$$B_2(g_1^{}|\mathrm{\Lambda }_2^{}g_3^{}\mathrm{\Lambda }_4^+|g_5^+g_6^+)_{DiffHelicity}=(\frac{(4|K_{234}|1)^3(2|K_{234}|1)}{t_{234}})B_+$$
(25)
For $`B_2(g_1^{}|\mathrm{\Lambda }_2^{}g_3^{}g_4^+|\mathrm{\Lambda }_4^+g_6^+)_{Diffhelicity}`$ we have :
$$l_1=g^{}l_2=\mathrm{\Lambda }^+,l_1=\mathrm{\Lambda }^{}l_2=\varphi ,+(exchangebetweenl_1,l_2)$$
(26)
This gives contribution :
$`1l_1^315[l_14]^3[42]+31l_1^2[l_14]^21l_2[l_24][42]15`$ (27)
$`31l_2^2[l_24]^21l_1[l_14][42]15+1l_1^315[l_14]^3[42]`$
$`=(4|K_{234}|1)^3[42]15`$
Thus
$$B_2(g_1^{}|\mathrm{\Lambda }_2^{}g_3^{}g_4^+|\mathrm{\Lambda }_4^+g_6^+)_{Diffhelicity}=\frac{(4|K_{234}|1)^3[42]15}{t_{234}}B_+$$
(28)
Adding eq.(17),(19),(22),(25) and (28) together gives the $`B_2`$ coefficient of the gluino anti-gluino pair amplitudes computed in 5 . Coefficients for the next cut can be calculated in similar way, we’ve checked it gives the same result as that derived in 5 .
It is straight forward to compute amplitudes that involve more than one pair of gluino or scalar. The new amplitudes are :
$`A(g^{}g^+\mathrm{\Lambda }^+\mathrm{\Lambda }^{}\mathrm{\Lambda }^{}\mathrm{\Lambda }^+),A(g^{}g^+\varphi \varphi \varphi \varphi ),A(\varphi \varphi \varphi \varphi \varphi \varphi )A(\mathrm{\Lambda }^{}\mathrm{\Lambda }^{}\mathrm{\Lambda }^{}\mathrm{\Lambda }^+\mathrm{\Lambda }^+\mathrm{\Lambda }^+),`$ (29)
$`A(\mathrm{\Lambda }^{}\mathrm{\Lambda }^+\varphi \varphi \varphi \varphi ),A(\mathrm{\Lambda }^{}\mathrm{\Lambda }^{}\mathrm{\Lambda }^+\mathrm{\Lambda }^+\varphi \varphi ),A(\mathrm{\Lambda }^{}\mathrm{\Lambda }^+\varphi g^{}g^+g^+)`$
Complication arises for these amplitudes because non-gluon particles carry less superspace variables and increase the amount of spinor combination. Luckily with the specification of the flavor for the external particles, the propagators are restricted to take certain species. This is discussed in detail in the next section where we calculate the all gluino and all scalar amplitude.
## III 3. Amplitudes with all gluinos and all scalars
Here we present N=4 SYM NMHV loop amplitudes with all gluino and all scalars. These amplitudes were derived from explicit expansion of eq.(5)-(7). Since scalars and gluinos carry less fermionic parameters as seen in eq.(10)(11), the spinor product that arises from the fermionic delta function becomes complicated. The final coefficient should not contain the off shell propagator spinor, thus one can use this as a guideline to group the spinor products to form kinematic invariant terms. With specific flavors this also restrict the possible species for propagators.
### III.1 3.1 $`A(\mathrm{\Lambda }_1^{\text{1}+}\mathrm{\Lambda }_2^{\text{2}+}\mathrm{\Lambda }_3^{\text{3}+}\mathrm{\Lambda }_4^\text{1}\mathrm{\Lambda }_5^\text{2}\mathrm{\Lambda }_6^\text{3})`$
For the six gluino amplitude we look at amplitudes with all three positive helicity gluino carrying different flavor. The negative helicities also carry different flavor and is the same set as the positive. For $`t_{123}`$ the flavors of the internal particles are uniquely determined.
$$l_1=\mathrm{\Lambda }^{}l_2=g^+,l_1=\mathrm{\Lambda }^+l_2=\varphi ,+exchange$$
(30)
This gives
$`B_1(\mathrm{\Lambda }_1^{\text{1}+}\mathrm{\Lambda }_2^{\text{2}+}\mathrm{\Lambda }_3^{\text{3}+}|\mathrm{\Lambda }_4^\text{1}\mathrm{\Lambda }_5^\text{2}\mathrm{\Lambda }_6^\text{3})=(1|K_{123}|52|K_{123}|63|K_{123}|4`$ (31)
$`+1|K_{123}|42|K_{123}|53|K_{123}|6+1|K_{123}|62|K_{123}|43|K_{123}|5)B_0^{}`$
Next we look at $`t_{345}`$ cut. The propagator assignment with same helicity (the definition of same or different helicity again follows that of the previous paragraph ) would be :
$$l_1=g^{}l_2=\mathrm{\Lambda }^{},+exchangepropagator$$
(32)
this gives
$`B_3(\mathrm{\Lambda }_1^{\text{1}+}\mathrm{\Lambda }_2^{\text{2}+}|\mathrm{\Lambda }_3^{\text{3}+}\mathrm{\Lambda }_4^\text{1}\mathrm{\Lambda }_5^\text{2}|\mathrm{\Lambda }_6^\text{3})_{SameHelicity}=45^2[12]^2\{34[61]5|K_{345}|2`$ (33)
$`+34[62]5|K_{345}|1+35[61]4|K_{345}|2+35[62]4|K_{345}|1\}B_{}^{}`$
There are two ways of assigning different helicity propagators
$$l_1=g^{}l_2=\mathrm{\Lambda }^+,orl_1=\mathrm{\Lambda }^{}l_2=\varphi ,+exchange$$
(34)
Note however for the present set of flavors, there is no consistent way of assigning flavors when the propagators are a gluon and a gluino. Thus we are left with the gluino scalar possibility with it’s flavor uniquely determined.
$`B_3(\mathrm{\Lambda }_1^{\text{1}+}\mathrm{\Lambda }_2^{\text{2}+}|\mathrm{\Lambda }_3^{\text{3}+}\mathrm{\Lambda }_4^\text{1}\mathrm{\Lambda }_5^\text{2}|\mathrm{\Lambda }_6^\text{3})_{DiffHelicity}=1662[43][35]6|K_{345}|3t_{345}B_{}`$ (35)
Luckily there is no need to compute $`B_2`$ coefficients since it is related to $`B_3`$ by symmetry.
### III.2 3.2 $`A(\varphi _1\varphi _2\varphi _3\varphi _4\varphi _5\varphi _6)`$
The power of deriving amplitudes from a superspace expansion is that one can rule out certain amplitudes just by inspection. Amplitudes with more than two scalars carrying the same color vanishes since there is no way of assigning the correct fermionic variables. Here we look at six scalar amplitude all carrying different flavor. This should be the simplest amplitude since the flavor carried by the internal particle is uniquely determined. We give the result for cut $`t_{123}`$ while the other cuts are related by symmetry.
$`B_1(\varphi _1\varphi _2\varphi _3\varphi _4\varphi _5\varphi _6)=\{(12[56]3|K_{123}|4+12[64]3|K_{123}|5+12[45]3|K_{123}|6`$ (36)
$`+31[56]2|K_{123}|4+31[64]2|K_{123}|5+31[45]2|K_{123}|6`$
$`+23[56]1|K_{123}|4+23[64]1|K_{123}|5+23[45]1|K_{123}|6)^2\}B_0`$
$`+complexconjugate.`$
## IV 4. A brief discussion on the MHV vertex approach
As discussed in the introduction, the straight forward way to compute amplitudes in superspace is the generalization of the MHV vertex 10 approach. It is also of conceptual interest to see if this approach actually works for the NMHV loop amplitude. Here we give a brief discussion of the extension.
The MHV vertex approach was shown to be successful 12 in constructing the n point MHV loop amplitude. This is partly due to the similarity between the cut diagrams 1 originally used to compute the amplitude and the MHV vertex diagram, so that one can use a dispersion type integral to reconstruct the box functions from it’s discontinuity across the branch cut. For the NMHV loop amplitude, one requires three propagator for the three MHV vertex one-particle-irreducible(1PI) diagram and two propagators for the one-particle-reducible(1PR) diagram (fig-2)13 .
We would then encounter the following integration:
$`{\displaystyle \frac{1}{_{i=1}^nii+1}}{\displaystyle \frac{d^4L_1}{L_1^2}\frac{d^4L_2}{L_2^2}\frac{d^4L_3}{L_3^2}\delta (P_\alpha +L_2L_3)\delta (P_\beta +L_3L_1)\delta (P_\gamma +L_1L_2)}`$ (37)
$`{\displaystyle d^8\eta _{l_1}d^8\eta _{l_2}d^8\eta _{l_3}\frac{\delta ^8(\mathrm{\Theta }_1)\delta ^8(\mathrm{\Theta }_2)\delta ^8(\mathrm{\Theta }_3)m_2m_2+1m_1m_1+1m_3m_3+1}{l_2l_1l_3l_2l_1l_3l_1m_2+1m_2l_1l_2m_3+1m_3l_2l_3m_1+1m_1l_3}}`$
$`+{\displaystyle \frac{\delta (L_1P_\gamma )}{_{i=1}^nii+1}}{\displaystyle \frac{d^4L_2}{L_2^2}\frac{d^4L_3}{L_3^2}\delta (P_\alpha +L_3L_2)\delta (P_\beta +L_2+L_1L_3)d^8\eta _{l_1}d^8\eta _{l_2}d^8\eta _{l_3}}`$
$`\times {\displaystyle \frac{\delta ^8(\mathrm{\Theta }_1)\delta ^8(\mathrm{\Theta }_2)\delta ^8(\mathrm{\Theta }_3)m_2m_2+1m_1m_1+1m_3m_3+1m_4m_4+1}{L_1^2l_3l_2^2m_1l_2l_2m_{1+1}l_3m_2+1m_2l_3l_1m_3+1m_3l_1l_1m_4+1m_4l_1}}`$
where for the first term
$`\mathrm{\Theta }_1={\displaystyle \underset{i=\alpha }{}}\eta _i\lambda _i+l_2\eta _{l_2}l_3\eta _{l_3}`$ (38)
$`\mathrm{\Theta }_2={\displaystyle \underset{i=\beta }{}}\eta _i\lambda _i+l_3\eta _{l_3}l_1\eta _{l_1}`$
$`\mathrm{\Theta }_3={\displaystyle \underset{i=\gamma }{}}\eta _i\lambda _i+l_1\eta _{l_1}l_2\eta _{l_2}`$
for the second term
$`\mathrm{\Theta }_1={\displaystyle \underset{i=\alpha }{}}\eta _i\lambda _il_2\eta _{l_2}+l_3\eta _{l_3}`$ (39)
$`\mathrm{\Theta }_2={\displaystyle \underset{i=\beta }{}}\eta _i\lambda _i+l_1\eta _{l_1}+l_2\eta _{l_2}l_3\eta _{l_3}`$
$`\mathrm{\Theta }_3={\displaystyle \underset{i=\gamma }{}}\eta _i\lambda _il_1\eta _{l_1}`$
$`\alpha `$$`\beta `$$`\gamma `$ labels the external momenta assigned to the three MHV vertex and the $`l_i`$s are the off shell continuation spinor following the CSW prescription 10 . We can reorganize the delta functions to reproduce the overall momentum conservation. For the first term in eq.(37) we have
$$\delta (P_{\alpha +\beta +\gamma })\delta (P_{\beta +\gamma }+L_3L_2)\delta (P_\gamma +L_1L_2)$$
(40)
For the second term
$$\delta (P_{\alpha +\beta +\gamma })\delta (P_\alpha +L_2L_3)$$
If we integrate the last delta function away in the first term and combine with the 1PR graphs, it is equivalent to using two MHV vertex to construct NMHV tree amplitude on one side of the two remaining propagator, namely this combines vertex $`\gamma `$ and $`\beta `$ through propagator $`L_1`$. To see this note that the momentum conserving delta function forces $`L_1`$ propagator to carry the correct momentum as it would for the CSW method and the $`\frac{1}{P_{L_1}^2}`$ is present in the integral measure in the first place. This would obviously affect the off shell spinor in the following way.
$$l_1=L_1\stackrel{~}{\eta }(L_2P_\gamma )\stackrel{~}{\eta }$$
(41)
This simply fixes the off shell spinor to be computed from the correct momentum as the CSW method. Thus we have come to a two propagator integral with two tree level amplitude on both side constructed from the CSW method. This is exactly the picture one would have if one apply the standard cut, except the propagators are off shell instead of on shell. For higher number of MHV vertices this can be applied straight forward, by integrating the momentum conserving propagator one at a time one can reduce the number of propagators until one arrive at the standard cut picture. As shown in 12 one can then proceed to recast the two propagator integral into a dispersion integral which computes the discontinuity across the cut of the integrand, by using the cut constructibility of N=4 SYM loop amplitudes, one can reconstruct the box function and it’s coefficient. However there is one subtlety. In the original standard cut one has to analyze every cut channel, and then disentangle the information since more than one box integral share the same cuts. If we follow the CSW prescription we can always reduce the loop diagrams down to two propagator one loop diagrams with a MHV vertex on one side of the two propagators. Thus this implies if the CSW approach is valid at one loop, then the full loop amplitude should be able to be reconstructed from the cuts of a subgroup of two propagator diagrams which always have a MHV vertex on one side of the cut.
This construction makes the connection between MHV vertex and $`\overline{MHV}`$ loop amplitude more transparent. $`\overline{MHV}`$ loop are just the parity transformation of the MHV loop, where one simply take the complex conjugate of the MHV loop:
$$A(\overline{MHV})_{loop}=A(\overline{MHV})_{tree}\underset{i=1}{\overset{n}{}}\underset{r=1}{\overset{[n/2]1}{}}(1\frac{1}{2}\delta _{\frac{n}{2}},r)F_{n:r;i}^{2me}$$
(42)
It’s derivation from MHV vertex is as follows. In 15 it was shown that using MHV vertices one can reconstruct the $`\overline{MHV}`$ tree amplitude in it’s complex conjugate spinor form. Since by integrating out one loop propagator corresponds to using MHV vertex to construct NMHV tree amplitude, one can proceed in a specific manner to reduce the number of loop propagators down to two with two $`\overline{MHV}`$ tree on both side. Since from 15 the two $`\overline{MHV}`$ tree amplitudes on both side is expressed in complex conjugate form, following exactly the same lines in 12 one can reproduce eq.(43).
## V 5. Conclusion
In this paper we constructed the 6 point NMHV loop amplitude for N=4 SuperYang-Mills in a compact form using its cut constructible nature. The expansion with respect to the fermionic parameter gives amplitudes with different particle content and helicity ordering. To extend further to higher point NMHV loops one may have to resolve to the MHV vertex approach since the tree level amplitudes on both side of the cut in general will not be in simple MHV and $`\overline{MHV}`$ combination. We also give a general discussion on how to proceed with the MHV vertex construction for higher than MHV loop(more than two negative helicities). The fact that it reproduces the two propagator picture for any one loop diagram combined with earlier results that have reproduced the MHV loop12 and the relationship between the leading order and sub leading order amplitudes13 , gives a strong support for the CSW approach beyond tree level.
## VI 6. Acknowledgements
It is a great pleasure to thank Warren Siegel for suggesting the question, Gabriele Travaglini, L. J. Dixon and Freddy Cachazo for their useful comments. |
warning/0507/hep-ph0507335.html | ar5iv | text | # Neutrino Factory Based on Linear Collider
## 1 Introduction
The project of Linear Collider (LC) contains one essential element that is not present in other colliders. Here each electron (or positron or photon) bunch will be used only once, and physical collision leave two very dense and strongly collimated beams of high energy electrons or/and photons with precisely known time structure. We consider, for definiteness, electron beam parameters of the TESLA project
$$\begin{array}{c}particleenergyE_e=250GeV,\\ numberofelectronspersecondN_e=2.710^{14}/s,\\ meanbeampowerP_b11MWt,\\ transversesizeandangularspreadnegligible.\end{array}$$
(1)
The problem, how to deal with this powerful beam dump, is under intensive discussion.
Main discussed variant is to destruct these used beams with minimal radioactive pollution (see e. g. ). It looks natural also to use these once–used beams for fixed target experiments with unprecedented precision.
Recently we suggested to utilize these used beams to initiate work of subcritical fission reactor and to construct neutrino factory \[LCWS05\]. Here we present estimates for one of these options. Real choice and optimization of parameters should be the subject of detail subsequent studies.
$``$ The study of neutrino oscillations is one of the most important problems in modern particle physics. In this problem the neutrino factories promise most detailed and important results. The existing projects of neutrino factories (see e.g. ) are very expensive and their physical potential is limited by expected neutrino energy and productivity of neutrino source.
The neutrino factory based on LC is much less expensive than those discussed nowadays . The combination of a high number of particles in the beam and high particle energy (1) provides very favorable properties of neutrino factory. The initial beam will be prepared in LC irrelevantly to the neutrino factory construction. The construction demands no special electronics except for that for detectors. The initial beam is very well collimated so that the additional efforts for beam cooling are not necessary. The use of the Ice-cub in Antarctic as a far distance detector (FDD) allows to see possible oscillations $`\nu _\mu sterile\nu `$ via measurement of deficit of $`\nu _\mu N\mu X`$ events.
The neutrino beam will have very well known discrete time structure that repeats the same structure in the LC. This fact allows to separate cosmic and similar backgrounds with high precision during operations. Very simple structure of neutrino generator allows to calculate the energy spectrum and content of the main neutrino beam with high accuracy. It must be verified with high precision in nearby detector (NBD).
In this project neutrino beam will contain mainly muon neutrino’s and antineutrino’s with small admixture $`\nu _e`$ and $`\overline{\nu }_e`$ and tiny dope of $`\nu _\tau `$ and $`\overline{\nu }_\tau `$ (the latter can be calculated with low precision). The neutrino energies are spread up to about 80 GeV with mean energy about 30 GeV, providing reliable observation of $`\tau `$, produced by $`\nu _\tau `$ from $`\nu _\mu \nu _\tau `$ oscillations. In the physical program of discussed $`\nu `$ factory we consider only problem of oscillations $`\nu _\mu \nu _\tau `$ and/or $`\nu _\mu sterile\nu `$. The potential of this $`\nu `$ factory in other problems of $`\nu `$ physics should be studied after detailed consideration of the project.
## 2 Elements of neutrino factory
### 2.1 Scheme
The proposed scheme deals with the electron beam used in LC and contains the following parts (see Fig. 1).
$``$ Pion producer (PP), $``$ Neutrino transformer (NT), $``$ Nearby detector (NBD),
$``$ Far distance detector (FDD), $``$ Beam turning magnet (BM) before PP.
### 2.2 Beam turning magnet
The system should start with the magnetic system which turns the used beam at an angle necessary to reach FDD with sacrifice of monochromaticity but without growth of angular spread. The vertical component of turning angle $`\alpha _V`$ is determined by Earth curvature. Let us denote the distance from LC to FDD by $`L_F`$. To reach FDD the initial beam (and therefore NT) should be turned before PP at the angle $`\alpha _V=\mathrm{arcsin}[L_F/(2R_E)]`$ below horizon (here $`R_E`$ is Earth radius).
The horizontal component of turning angle can be minimized by suitable choice of the proper LC orientation (orientation of incident beam near the LC collision point).
### 2.3 Pion producer (PP)
The next stage is pion production in the PP in the form of a 20 cm long water cylinder (20 cm is one radiation length). The water in cylinder should rotate for cooling. In this PP almost each electron will produce bremsstrahlung photon with energy $`E_\gamma =100200`$ GeV. The angular spread of these photons can be estimated as angular spread of initial beam (about 0.1 mrad). The bremsstrahlung photons have additional angular spread of about $`1/\gamma 210^6`$. These two spreads are negligible for our problem.
Then these photons collide with nuclei and produce pions,
$$\gamma NN^{}+\pi ^{}s,\sigma 110\mu b.$$
(2)
This process gives about $`10^3`$ $`\gamma N`$ collisions per 1 electron, which is about $`310^{11}`$ $`\gamma N`$ collisions per second. On average, each of this collisions produces one pion with high energy $`E_\pi >E_\gamma /2`$ (for estimates $`E_\pi ^h=70`$ GeV) and at least 2-3 pions with lower energy (for estimates, $`E_\pi ^{\mathrm{}}20`$ GeV).
Mean transverse momentum of these pions is 350-500 MeV. The angular spread of high energy pions with energy $`E_\pi ^h`$ is within 7 mrad. The increase of angular spread of pions with decrease of energy is compensated by growth of the number of produced pions. Therefore, for estimates we accept that the pion flux within angular interval 7 mrad contains $`310^{11}`$ pions with $`E_\pi =E_\pi ^h`$ and the same number of pions with $`E_\pi =E_\pi ^{\mathrm{}}`$ per second. Let us denote the energy distribution of pions near forward direction by $`f(E)`$.
Certainly, more refined calculations should also consider production and decay of $`K`$ mesons, etc. Reaction mentioned in ref.
$$\gamma ND_s^\pm X\nu _\tau \overline{\tau }X.$$
(3a)
plays the most essential role for our estimates. Its cross section rapidly increases with energy growth and
$$\sigma 210^{33}cm^2atE_\gamma 50\text{ GeV}.$$
(3b)
### 2.4 Neutrino transformer (NT). Neutrino beams
For the neutrino transformer (NT) we suggest a low vacuum pipe of length $`L_{NT}1`$ km and radius $`r_{NT}2`$ m. Here muon neutrino $`\nu _\mu `$ and $`\overline{\nu }_\mu `$ are created from $`\pi \mu \nu `$ decay. This length $`L_{NT}`$ allows more than one quarter of pions with $`E_\pi E_\pi ^h`$ to decay. The pipe with radius $`r_{NT}`$ gives an angular coverage of 2 mrad, which cuts out 1/12 part of total flux of low and medium energy neutrinos. With the growth of pion energy two factors act in opposite ways. First, with this growth initial angular spread of pions decreases, therefore the fraction of flux cut out by the pipe increases. Second, with this growth the number of pion decays within relatively short pipe decreases. These two tendencies compensate each other in the resulting flux.
The energy distribution of neutrino’s obtained from $`\pi `$ decay with energy $`E`$ is uniform in the interval $`(aE,\mathrm{\hspace{0.17em}0})`$ with $`a=1(m_\mu /m_\pi )^2`$. Therefore, the energy distribution in neutrino energy $`\epsilon `$ is obtained from energy distribution of pions near forward direction $`f(E)`$ as (note that $`f(E)=0`$ at $`E>E_e`$)
$$F(\epsilon )=\underset{\epsilon /a}{\overset{\mathrm{}}{}}f(E)𝑑E/(aE),a=1\frac{m_\mu ^2}{m_\pi ^2}0.43.$$
(4)
The increase of angular spread in the decay is negligible in the rough approximation. Finally, at the end of NT we expect to have the neutrino flux within the angle 2 mrad
$$\begin{array}{c}0.610^{10}\nu /swithE_\nu =E_\nu ^h30GeV,\\ \text{ and }\mathrm{\hspace{0.33em}\hspace{0.33em}0.6}10^{10}\nu /swithE_\nu =E_\nu ^{\mathrm{}}9GeV.\end{array}$$
(5)
We denote below neutrino’s with $`E_\nu =30`$ GeV and $`9`$ GeV as high energy neutrino’s and low energy neutrino’s respectively.
$``$ The background $`\nu _\tau `$ beam.
$``$ The background $`\nu _\tau `$ beam.
The $`\tau `$ neutrino are produced in PP. Two mechanisms were discussed in this respect, the Bethe-Heitler process $`\gamma N\tau \overline{\tau }+X`$ and process (3) which is is dominant . The cross section (3) is 5 orders less than $`\sigma (\gamma NX)`$. Mean transverse momentum of $`\nu _\tau `$ is given by $`m_\tau `$, which is more than 3 times higher than that for $`\nu _\mu `$. Along with e.g. $`\overline{\nu }_\tau `$ produced in this process, in NT each $`\tau `$ decays to $`\nu _\tau `$ plus other particles. Therefore, each such reaction is a source of a $`\nu _\tau +\overline{\nu }_\tau `$ pair. Finally, for flux density we have
$$N_{\nu _\tau }10^4\nu _\tau /(smrad^2)810^6N_{\nu _\mu }.$$
(6)
The $`\nu _\tau `$ (or $`\overline{\nu }_\tau `$) energy is typically higher than that of $`\nu _\mu `$ by factor $`2÷2.5`$.
Besides, $`\nu _\tau `$ will be produced by non-decayed pions within the protecting wall behind NP in the process like $`\pi ND_sX\tau \nu _\tau X`$. The cross section of this process increases rapidly with energy growth and equals $`0.13\mu `$b at $`E_\pi =200`$ GeV . Rough estimate shows that the number of additional $`\nu _\tau `$ propagating in the same angular interval is close to the estimates given in (6). In the numerical estimates below we consider for definiteness first contribution only. The measurement of $`\nu _\tau `$ flux in the NBD is a necessary component for the study of neutrino oscillations in FDD.
$``$ Other sources of $`\nu _\mu `$ and $`\nu _e`$ change these numbers only weakly.
### 2.5 Nearby detector (NBD)
$``$ Main goal of nearby detector (NBD) is to measure the energy and angular distribution of neutrino’s within the beam as well as $`N_{\nu _e}/N_{\nu _\mu }`$ and $`N_{\nu _\tau }/N_{\nu _\mu }`$.
$``$ We suggest to position the NBD behind NT and a concrete wall (to eliminate pions and other particles from initial beam). For estimates, we consider the NBD in a form of water cylinder of radius about 2-3 m (roughly the same as NT) and length $`\mathrm{}_{NBD}100`$ m.
For $`E_\nu =30`$ GeV the cross section for $`\nu `$ absorbtion is
$$\begin{array}{c}\sigma (\overline{\nu }N\mu ^+h)=0.1\frac{m_pE_{\overline{\nu }}}{\pi v^4}10^{37}cm^2,\\ \sigma (\nu N\mu ^{}h)=0.22\frac{m_pE_\nu }{\pi v^4}210^{37}cm^2.\end{array}$$
(7)
At these numbers the free path length in water is $`\lambda _{\overline{\nu }}=10^{13}`$ cm and $`\lambda _\nu =0.4510^{13}`$ cm. That gives
$$\begin{array}{c}(3÷6)10^7\mu /\mathrm{𝐲𝐞𝐚𝐫}(withE_\mu 30GeV);\\ 400÷800\tau /\mathrm{𝐲𝐞𝐚𝐫}(withE_\tau 50GeV)\end{array}$$
(8)
(here 1 year =$`10^7`$ s). These numbers look sufficient for detailed measurements of muon neutrino spectra and verification of calculations of $`\nu _\tau `$ backgrounds.
### 2.6 Far Distance Detector (FDD)
$``$ Main goal of FDDstudy of neutrino oscillations. We consider $`\nu _\mu \nu _\tau `$ oscillations and oscillations with sterile neutrino. We consider separately the potentials of two possible positions of FDD, assuming the length of oscillations to be
$$L_{osc}E_\nu /(50GeV)10^5km.$$
(9)
$``$ FDD I
The first opportunity is to place FDD at the distance $`L_F=200`$ km. In this case the initial beam should be turned at 16 mrad angle. This angle can be reduced by 3 mrad (one half of angular spread of initial pion beam).
We consider this FDD in the form of water channel of length 1 km with radius $`R_F40`$ m. The transverse size is limited by water transparency.
The fraction of neutrino’s reaching this FDD is given by ratio $`k=(R_F/L_F)^2/[(r_{NT}/L_{NT})^2]`$. In our case $`k=0.01`$. Main effect under interest here is $`\nu _\mu \nu _\tau `$ oscillation. They add $`(L_F/L_{osc})^2N_{\nu _\mu }`$ to initial $`N_{\nu _\tau }`$.
In FDD of chosen sizes we expect the counting rate to be just 10 times lower than that in NBD (8) for $`\nu N\mu X`$ reactions with high energy neutrino. We also expect the rate of $`\nu _\tau N\tau X`$ events to be another $`10^5`$ times lower (that is about 10 times higher than the background given by initial $`\nu _\tau `$ flux,
$$\begin{array}{c}N(\nu _\mu N\mu X)(3÷6)10^6/year,\\ N(\nu _\tau N\tau X)(30÷60)/year\end{array}inFDDI.$$
(10)
For neutrino of lower energies effect increases. Indeed, $`\sigma (\nu N\tau X)E_\nu `$ while $`L_{osc}E_\nu `$. Therefore, observed number of $`\tau `$ from oscillations increases $`1/E_\nu `$ at $`E_\nu 10`$ GeV. The additional counting rate for $`\nu _\tau N\tau X`$ reaction with low energy neutrino (with $`E_\nu =9`$ GeV) cannot be estimated so simply, but rough estimates give numbers similar to (10).
These numbers look sufficient for separation of $`\nu _\mu \nu _\tau `$ oscillations and rough measurement of $`s_{23}`$.
Note that at given FDD size the counting rate of $`\nu _\tau N\tau X`$ reaction is independent on FDD distance from LC, $`L_F`$. The growth of $`L_F`$ improves the signal to background ratio for oscillations. The value of signal naturally increases with growth of volume of FDD.
$``$ FDD II
The second opportunity is to use for FDD well known Ice-cub detector in Antarctica with volume 1 km<sup>3</sup>. The distance to FDD in this case is $`L_F10^4`$ km. This opportunity requires relatively expensive excavation work for NT and NBD at the angle about $`60\mathrm{deg}`$ under horizon.
At this $`L_F`$ for $`\nu `$ with energy about 30 GeV we expect the conversion of $`(L_F/L_{osc})^21/36`$ for $`\nu _\mu \nu _\tau `$ or $`\nu _\mu sterile\nu `$.
In this FDD the number of expected events $`\nu _\mu \mu X`$ with high energy neutrino will be about 0.01 of that in NBD,
$$\begin{array}{c}N(\nu _\mu N\mu X)(3÷6)10^5/year,\\ N(\nu _\tau N\tau X)10^4/year\end{array}inFDDII.$$
(11)
The contribution of low energy neutrino increases both these counting rates.
Therefore, one can hope that a few years of experimentation with reasonable $`\tau `$ detection efficiency will allow to measure $`s_{23}`$ with percent accuracy, and similar period of observation of $`\mu `$ production will allow to observe the loss of $`\nu _\mu `$ due to transition this neutrino to sterile $`\nu `$.
## 3 Discussion
$``$ All technical details of proposed scheme including sizes of all elements, construction, and materials of detectors can be modified in the forthcoming simulations and optimization of parameters. The numbers obtained above represent first rough estimates only. In particular, we did not discuss here methods of $`\mu `$ and $`\tau `$ registration and their efficiency. Next, large fraction of residual electrons, photons and pions leaving the PP will reach the walls of the NT pipe. The heat sink and radiation protection of this pipe must be taken into account.
$``$ More detailed physical program of this neutrino factory will be similar to the one discussed in other projects .
## 4 Other possible applications of LC used beam
$``$ Pion producer of neutrino factory in the fixed target experiment. The proposed PP can be used also as $`\gamma N`$ collider with luminosity $`310^{39}`$ cm<sup>-2</sup>s<sup>-1</sup>. Therefore, the PP with additional standard detector behind PP can be used for precise experiments in the fixed target regime for the $`\gamma N`$ collider with huge luminosity. Here one can study rare processes in $`\gamma N`$ collisions, $`B`$ physics, etc.
$``$ Additional opportunity for using NBD of neutrino factory. High rate of $`\nu _\mu N\mu X`$ processes expected in NBD allows to study new problems of high energy physics. The simplest example is the opportunity to study charged and axial current induced diffraction ($`\nu N\mu N^{}\rho ^\pm N`$, $`\nu N\mu N^{}b_1^\pm N`$,…) with high precision. Measurements of charged current induced structure functions present the second example.
$``$ Material sciences. The interaction of beam having exceptional energy density (1) with different materials will be of great interest for material physics (for example, to understand what happens at collision of micro-meteorite with spacecraft).
I am thankful to D. Naumov, L. Okun, V. Saveliev, A. Sessler, A. Skrinsky, V. Telnov, M. Vysotsky, M. Zolotarev for comments and new information. This work is supported by grants RFBR 05-02-16211, NSh-2339.2003.2. |
warning/0507/hep-th0507197.html | ar5iv | text | # Dirac equation for the supermembrane in a background with fluxes from a component description of the D=11 supergravity–supermembrane interacting system
## 1 Introduction
There has been a growing interest in the fermionic equations for superbranes in a supergravity background with fluxes (see and for earlier papers) as these are needed to study nonperturbative effects in string theory. To find such equations, one takes the super–$`p`$–brane action in curved superspace and expands it in powers of the fermionic coordinate function $`\widehat{\theta }^{\stackrel{ˇ}{\beta }}(\xi )`$ or proceeds in the same manner directly from the ‘superfield’ fermionic equation itself . To make such a decomposition one uses the Wess–Zumino (WZ) gauge for a superfield supergravity background plus the superspace supergravity constraints. For $`D=10,11`$ these superspace constraints imply the ‘free’ supergravity equations of motion, without any contribution from the superbrane source. Hence the ‘Dirac’ equation for $`\widehat{\theta }(\xi )`$, as derived in from superembedding approach and in from the $`𝒪(\widehat{\theta }(\xi )^2)`$ actions for M2 and various super–Dp-branes, apparently involves a gravitational and gauge field background that satisfies the ‘free’ supergravity equations without any contribution from the superbrane. Thus the consistency of the results obtained in this approach, although widely believed, is not manifest. A way to check this consistency would be to re-derive the same equations from a system of equations providing a fully dynamical description of the supergravity–super-$`p`$-brane interacting system by using a well defined approximation.
As, by definition, a super-$`p`$-brane is a $`p`$-brane moving in superspace, a complete system of equations, including those for the supergravity fields with contributions from the super-$`p`$-brane, could be derived from the sum of the superbrane action and a superfield action for supergravity. Such an action can be studied in lower dimensions (see ), but a superfield action for $`D=10,11`$ supergravity is unknown. This made difficult the study the $`D=11`$ supergravity–M-brane and $`D=10`$ type II supergravity–D$`p`$-brane interacting systems, the most interesting ones in an M–theoretical perspective.
This difficulty may be overcome by using the gauge–fixed description provided by the sum of the spacetime (component) action for supergravity (without auxiliary fields) and that for a bosonic brane, as given by the purely bosonic limit of the superbrane action. From the point of view of the superfield formulation of the interacting system (hypothetical for $`D=10,11`$) the gauge is provided by the conditions of the WZ gauge for the supergravity superfields plus the condition $`\widehat{\theta }^\alpha (\xi )=0`$ for the superbrane coordinate functions. The resulting gauge–fixed description is complete in the sense that it contains gauge–fixed versions of all the dynamical equations of motion of the interacting system, including a ‘fermionic equation for the bosonic brane’ . This equation, which formally coincides with the leading component of the superfield fermionic equation for the superbrane in a superspace supergravity background, appears in this component scheme as a selfconsistency condition for the bulk gravitino equations. Note that the ‘fermionic equation for bosonic brane’ is actually a non-dynamical ‘boundary’ condition for the gravitino on the brane worldvolume $`W^{p+1}`$. However, we will see how this algebraic equation allows us to obtain the superbrane fermionic field dynamics (the ‘Dirac equation’), which is hidden in it.
The above gauge–fixed action for the supergravity–superbrane interacting system can be derived from the superfield description in the dimensions where a superfield supergravity action exists (see for $`D=4`$, $`N=1`$ interacting systems). In the general case its form may be also deduced if one assumes the existence of a superfield supergravity action and exploits its defining properties . Then one concludes that, whether a superfield description of a supergravity–superbrane interacting system exists or not, the description of this system by means of the sum of the spacetime component supergravity action without auxiliary fields and the action of the ‘limiting’ bosonic brane (obtained by taking the purely bosonic, $`\widehat{\theta }^\alpha (\xi )=0`$, ‘limit’ of the supermembrane) does exist. Such an action preserves one–half of the local supersymmetry characteristic of the pure supergravity action. This one-half of the local supersymmetry reflects the $`\kappa `$–symmetry of the original superbrane action while the existence of a non–preserved one-half reflects the spontaneous breaking of the local supersymmetry by the superbrane.
In this paper we show that this complete but gauge–fixed description of the supergravity–superbrane interacting system may be used, despite it corresponds to the $`\widehat{\theta }^\alpha (\xi )=0`$ gauge for the superbrane variables, to reproduce the superbrane fermionic equation i.e., the dynamical Dirac equation for the fermionic field $`\widehat{\theta }^\alpha (\xi )`$, in a supergravity background with fluxes. This is related to the known fact that the superbrane fermionic coordinate functions are the Goldstone fields for the supersymmetry spontaneously broken by the superbrane. We also discuss the possibility of using the Goldstone nature of the $`\widehat{\theta }^\alpha (\xi )`$’s to find their lower order ($`𝒪(\widehat{\theta }^k)`$) contributions to the bulk supergravity equations, i.e. to search for a lower–order approximation in $`\widehat{\theta }^\alpha (\xi )`$ to the system of interacting equations that would possess full local supersymmetry (not just the one half preserved by the gauge–fixed description of ) in the same (actually, $`𝒪(\widehat{\theta }^{(k1)})`$) approximation.
For definiteness we consider here the case of the $`D=11`$ supergravity–supermembrane (SG–M2) interaction, although the method could be applied to other systems like the SG–M5 one involving the M5–brane, or the SG-Dp system, with $`D=10`$ type II Dirichlet $`p`$–branes.
## 2 Supergravity interacting with a bosonic membrane as a gauge–fixed description of the D=11 supergravity–supermembrane (SG–M2) interacting system and its properties
### 2.1 D=11 Supermembrane in the on-shell superfield supergravity background
The supermembrane action in a supergravity background is
$`S_{_{M2}}[\widehat{E}^a,\widehat{A}_3]={\displaystyle \underset{W^3}{}}\left(d^3\xi {\displaystyle \frac{1}{2}}\sqrt{|g(\xi )|}\widehat{A}_3\right)={\displaystyle \underset{W^3}{}}\left({\displaystyle \frac{1}{3!}}\widehat{E}^a\widehat{}\widehat{E}^aA_3(\widehat{Z})\right),`$ (1)
where the pull–backs to the supermembrane worldvolume $`W^3`$
$`\widehat{E}^a:=d\widehat{Z}^M(\xi )E_M{}_{}{}^{a}(\widehat{Z}(\xi ))=:d\xi ^m\widehat{E}_m^a`$
and
$$\widehat{A}_3=\frac{1}{3!}d\widehat{Z}^{M_3}d\widehat{Z}^{M_2}d\widehat{Z}^{M_1}A_{M_1M_2M_3}(\widehat{Z}(\xi ))$$
(2)
of the bosonic supervielbein $`E^a=dZ^ME_M{}_{}{}^{a}(Z)`$ and the 3–superform $`A_3=\frac{1}{3!}dZ^{M_3}dZ^{M_2}dZ^{M_1}A_{M_1M_2M_3}(Z)`$ of the superspace formulation of $`D=11`$ supergravity are denoted by a caret. These are obtained by replacing the superspace coordinate $`Z^M=(x^\mu ,\theta ^{\stackrel{ˇ}{\alpha }})`$ by the coordinate functions $`\widehat{Z}^M(\xi )=(\widehat{x}^\mu (\xi ),\widehat{\theta }^{\stackrel{ˇ}{\alpha }}(\xi ))`$ that ‘locate’ the worldvolume $`W^3`$ as a hypersurface in superspace, $`Z^M=\widehat{Z}^M(\xi )`$. The worldvolume Hodge star operator $`\widehat{}`$ in (1) is defined by
$`\widehat{}\widehat{E}^a={\displaystyle \frac{1}{2}}d\xi ^nd\xi ^mϵ_{mnk}\sqrt{|g(\xi )|}g^{kl}_l\widehat{Z}^ME_M{}_{}{}^{a}(\widehat{Z}),`$ (3)
where $`g_{mn}(\xi )`$ is the induced metric on $`W^3`$,
$`g_{mn}(\xi ):=\widehat{E}_m^a\widehat{E}_{na}:=_m\widehat{Z}^M_n\widehat{Z}^NE_N{}_{}{}^{a}(\widehat{Z})E_{Ma}(\widehat{Z}),|g(\xi )|:=|det(g_{mn}(\xi ))|.`$ (4)
The supermembrane equations of motion
$`\widehat{D}(\widehat{}\widehat{E}_a)={\displaystyle \frac{1}{3}}\widehat{E}^d\widehat{E}^c\widehat{E}^bF_{abcd}(\widehat{Z}){\displaystyle \frac{1}{2}}\widehat{E}^d\widehat{E}^\alpha \widehat{E}^\beta \mathrm{\Gamma }_{ab\alpha \beta },`$ (5)
$`\widehat{}\widehat{E}^a\widehat{E}^\beta \left(\mathrm{\Gamma }_a(1\overline{\overline{\gamma }})\right)_{\beta \alpha }=0\widehat{E}_m^\beta \left(\widehat{\mathrm{\Gamma }}^m(1\overline{\overline{\gamma }})\right)_{\beta \alpha }=0,`$ (6)
where $`D`$ is the standard covariant derivative involving the spin connection, $`\widehat{E}_m^A:=_m\widehat{Z}^ME_M^A(\widehat{Z})`$, $`\widehat{\mathrm{\Gamma }}^m:=g^{mn}(\xi )\widehat{E}_n^a\mathrm{\Gamma }_a`$ and
$`\overline{\overline{\gamma }}=`$ $`{\displaystyle \frac{i}{3!\sqrt{|g(\xi )|}}}ϵ^{mnk}_k\widehat{Z}^K_n\widehat{Z}^N_m\widehat{Z}^ME_M^a(\widehat{Z})E_N^b(\widehat{Z})E_K^c(\widehat{Z})`$ $`\mathrm{\Gamma }_{abc},`$ (7)
are obtained making use of the superspace supergravity constraints
$`T^a`$ $`:=`$ $`DE^a:=dE^aE^b\omega _b{}_{}{}^{a}=iE^\alpha E^\beta \mathrm{\Gamma }_{\alpha \beta }^a,`$ (8)
$`dA_3`$ $`=`$ $`{\displaystyle \frac{1}{4}}E^\alpha E^\beta E^bE^a\mathrm{\Gamma }_{ab\alpha \beta }+F_4(Z),F_4(Z):={\displaystyle \frac{1}{4!}}E^{a_4}\mathrm{}E^{a_1}F_{a_1\mathrm{}a_4}(Z).`$ (9)
These are known to be on–shell constraints i.e., they include the equations of motion for the physical spacetime or ‘component’ fields $`e_\mu ^a(x)`$, $`\psi _\mu ^\alpha (x)`$, $`A_{\mu \nu \rho }(x)`$,
$`e_\mu ^a(x)`$ $`:`$ $`e^a(x):=dx^\mu e_\mu {}_{}{}^{a}(x)=E^a(Z)|_{\theta =0,d\theta =0}`$ (10)
$`\psi _\mu ^\alpha (x)`$ $`:`$ $`\psi ^\alpha (x):=dx^\mu \psi _\mu ^\alpha (x)=E^\alpha (Z)|_{\theta =0d\theta =0}`$ (11)
$`A_{\mu \nu \rho }(x)`$ $`:`$ $`A_3(x):=1/3!dx^\mu dx^\nu dx^\rho A_{\mu \nu \rho }(x)=A_3(Z)|_{\theta =0,d\theta =0},`$ (12)
among their consequences. For our present discussion it is important to note that these are ‘free’ supergravity equations in the sense that they do not contain any source contribution from the supermembrane.
The on–shell supergravity constraints (8), (9) are also necessary conditions for the $`\kappa `$–symmetry of the supermembrane in a curved superspace background . This local fermionic symmetry, first found in a superparticle context in , manifests itself by the presence of the projector $`(1\overline{\overline{\gamma }})`$ in the fermionic equations (6). Its explicit form is given by
$$\delta _\kappa \widehat{Z}^M(\xi )=(1+\overline{\overline{\gamma }})^\alpha {}_{\beta }{}^{}\kappa _{}^{\beta }(\xi )E_\alpha {}_{}{}^{M}(\widehat{Z}).$$
(13)
When the supergravity background superfields obey the Wess–Zumino (WZ) gauge conditions
$`i_\theta E^a`$ $`:=`$ $`\theta ^{\stackrel{ˇ}{\beta }}E_{\stackrel{ˇ}{\beta }}{}_{}{}^{a}(x,\theta )=0,`$
$`i_\theta E^\alpha `$ $`:=`$ $`\theta ^{\stackrel{ˇ}{\beta }}E_{\stackrel{ˇ}{\beta }}{}_{}{}^{\alpha }(x,\theta )=\theta ^{\stackrel{ˇ}{\beta }}\delta _{\stackrel{ˇ}{\beta }}{}_{}{}^{\alpha }=:\theta ^\alpha ,`$
$`i_\theta A_3(Z)=0,i_\theta \omega ^{ab}(Z)=0`$
(see and refs. therein; see also )<sup>2</sup><sup>2</sup>2$`i_\theta `$ is a shorthand for the contraction with the vector field $`\theta ^{\stackrel{ˇ}{\alpha }}\frac{}{\theta ^{\stackrel{ˇ}{\alpha }}}`$. the index of the Grassmann coordinate $`\theta `$ and of the coordinate function $`\widehat{\theta }(\xi )`$ is identified with the spinor index and, furthermore, one can extract from (13) the $`\kappa `$–symmetry transformations for the $`\widehat{\theta }^\alpha `$ variable,
$$\delta _\kappa \widehat{\theta }^\alpha (\xi )=(1+\overline{\overline{\gamma }})^\alpha {}_{\beta }{}^{}\kappa _{}^{\beta }(\xi )+𝒪(\widehat{\theta })=(1+\overline{\gamma })^\alpha {}_{\beta }{}^{}\kappa _{}^{\beta }(\xi )+𝒪(\widehat{\theta }),$$
(15)
where in Eq. (15) $`\overline{\gamma }=\overline{\overline{\gamma }}|_{\widehat{\theta }=0}`$.
The name Dirac equation for the superbrane is usually given for (an approximation to) the equation of motion for the superbrane fermion $`\widehat{\theta }(\xi )`$ in a spacetime supergravity background. We will also call it below superbrane fermionic equation. In the standard approach to derive this equation one considers the action (1) or the superfield fermionic equation (6) for the on–shell supergravity background taken in the WZ gauge and expands it in powers of $`\widehat{\theta }(\xi )`$ keeping the lower orders in $`\widehat{\theta }(\xi )`$; the first order is usually considered to be sufficient. Then one uses the $`\kappa `$–symmetry (13) to gauge away half ($`16`$ out of $`32`$) of the $`\widehat{\theta }(\xi )`$ components to retain only the physical supermembrane fermions.
The fact that both the very derivation of the superfield fermionic equation (6) and its decomposition in powers of $`\widehat{\theta }(\xi )`$ makes an essential use of the on–shell superspace supergravity constraints, which cannot incorporate any supermembrane source contribution, makes the consistency of the standard background superfield approach not obvious. The check of its consistency is one of the motivations of this paper.
### 2.2 On the properties of a (hypothetical) superfield Lagrangian description of the D=11 supergravity–superbrane interaction
A complete, supersymmetric description of the SG–M2 interaction would be provided by the sum
$`𝐒_{SGM2}=𝐒_{SG}[E^a,E^\alpha ,A_3(Z)]+S_{M2}[\widehat{E}^a,A_3(\widehat{Z})]`$ (16)
of the supermembrane action (1) and the hypothetical superfield action for $`D=11`$ supergravity $`𝐒_{SG}[E^a,E^\alpha ,A_3(Z)]`$. This action is not known and it is not even clear whether it exists. Nevertheless, if exists, such a supergravity action would possess certain properties. In particular, it would be invariant under arbitrary changes of the superspace coordinates, i.e. superdiffeomorphisms $`\delta _{sdiff}`$. The same is true of the full interacting action (16) provided that the transformations of the coordinate functions of superbrane, $`\widehat{Z}^M(\xi )=(\widehat{x}^\mu (\xi ),\widehat{\theta }^{\stackrel{ˇ}{\alpha }}(\xi ))`$ are given by the pull–backs $`\widehat{b}^M=b^M(\widehat{Z}(\xi ))`$ to $`W^3`$ of the superspace diffeomorphism parameters $`b^M(Z)`$, i.e.
$`\delta _{sdiff}Z^M=b^M(Z),\delta _{sdiff}\widehat{Z}^M(\xi )=b^M(\widehat{Z}(\xi )).`$ (17)
Eq. (17) implies, in particular,
$`\delta _{sdiff}\theta ^{\stackrel{ˇ}{\alpha }}=b^{\stackrel{ˇ}{\alpha }}(Z),\delta _{sdiff}\widehat{\theta }^{\stackrel{ˇ}{\alpha }}(\xi )=b^{\stackrel{ˇ}{\alpha }}(\widehat{Z}(\xi )).`$ (18)
Clearly, the transformations $`\delta _{sdiff}Z^M=b^M(Z)`$ cannot be used to set the fermionic coordinates $`\theta ^\alpha `$ equal to zero since such a transformation would have a vanishing superdeterminant and, hence, would not be a superdiffeomorphism. However, in contrast, the transformations (18) can be used to make the fermionic coordinate functions $`\widehat{\theta }^\alpha (\xi )`$ vanishing, i.e. one can fix the gauge
$`\widehat{\theta }^\alpha (\xi )=0,`$ (19)
which might be considered the analogue to the ‘unitary gauge’ of the Higgs model.
Another expected property of the hypothetical superfield interacting action (16) is that, in addition to the superspace diffeomorphism gauge symmetry (Eqs. (17), (18)), it would possess a local $`16`$–parametric fermionic $`\kappa `$–symmetry $`\delta _\kappa `$ acting on the supermembrane variables $`\widehat{Z}^M(\xi )`$ only. It is also plausible to assume that such a $`\kappa `$–symmetry would be characterized by Eq. (13) with some superfield projector $`1/2(1+\overline{\overline{\gamma }})`$, $`\overline{\overline{\gamma }}\overline{\overline{\gamma }}(Z)`$. Thus the set of fermionic gauge symmetries of the action would contain $`\delta _{gauge}=\delta _{sdiff}+\delta _\kappa `$. These transformations act on $`\widehat{\theta }^\alpha (\xi )`$ as
$`\delta _{gauge}\widehat{\theta }^\alpha (\xi )=b^\alpha (\widehat{Z}(\xi ))+\delta _\kappa \widehat{\theta }^\alpha (\xi )=\epsilon ^\alpha (\widehat{x})+𝒪(\widehat{\theta }(\xi ))+\delta _\kappa \widehat{\theta }^\alpha (\xi ),`$ (20)
where the leading component of the superfield superdiffeomorphism parameter has been denoted by $`\epsilon ^\alpha (\widehat{x})`$,
$`b^\alpha (\widehat{Z}(\xi ))=\epsilon ^\alpha (\widehat{x})+𝒪(\widehat{\theta }(\xi ))`$ (21)
to identify $`\epsilon ^\alpha (\widehat{x})`$ with the spacetime local supersymmetry parameter.
Irrespective of the details of the superspace formulation of supergravity, the WZ gauge (2.1) can be fixed on the supergravity superfields (see e.g. , and refs. therein) by using superdiffeomorphism symmetry (17) and the superspace structure group symmetry, $`SO(1,10)`$ in the present case. The WZ gauge is then preserved by a certain combination of the superdiffeomorphism and the superspace local Lorentz group transformations expressed in terms of a number of independent parameters, $`ϵ^\alpha (x)`$ of the spacetime local supersymmetry, $`b^\mu (x)`$ of spacetime diffeomorphisms and $`L^{ab}(x)`$ of spacetime local Lorentz transformations. In the WZ gauge the transformations of the fermionic coordinate function of the superbrane, $`\widehat{\theta }^\alpha (\xi )`$, read
$`\delta _\kappa \widehat{\theta }^\alpha (\xi )=(1+\overline{\overline{\gamma }}(\widehat{Z}))^\alpha {}_{\beta }{}^{}\kappa _{}^{\beta }(\xi )=(1+\overline{\gamma }(\widehat{x}))^\alpha {}_{\beta }{}^{}\kappa _{}^{\beta }(\xi )+𝒪(\widehat{\theta }(\xi )),`$ (22)
where $`\overline{\gamma }\overline{\gamma }(x)`$ is the leading component of $`\overline{\overline{\gamma }}\overline{\overline{\gamma }}(\widehat{Z})`$ in the $`\kappa `$–symmetry projector, $`\overline{\overline{\gamma }}(\widehat{Z})|_{\widehat{\theta }=0}=\overline{\gamma }(x)`$. Substituting (22) for $`\delta _\kappa \widehat{\theta }^\alpha (\xi )`$ in (20) one finds
$`\delta _{gauge}\widehat{\theta }^\alpha (\xi )=\epsilon ^\alpha (\widehat{x})+(1+\overline{\gamma })^\alpha {}_{\beta }{}^{}\kappa _{}^{\beta }(\xi )+𝒪(\widehat{\theta }(\xi )),`$ (23)
Eq. (23) exhibits, first of all, the Goldstone nature of the superbrane fermionic coordinate functions $`\widehat{\theta }^\alpha (\xi )`$: $`\widehat{\theta }^\alpha (\xi )`$ are the Goldstone fermions corresponding to the supersymmetry spontaneously broken by the superbrane (see , and refs. therein). In the supergravity–superbrane interacting system this supersymmetry is the spacetime local gauge symmetry which can be used to remove the Goldstone field by fixing the gauge (19). Secondly, Eq. (23) makes transparent that the spontaneous breaking of the local supersymmetry by superbrane is partial. Indeed, the simple observation
$`\widehat{\theta }^\alpha (\xi )=0`$ $``$ $`0=\delta _{gauge}\widehat{\theta }^\alpha (\xi )|_{\widehat{\theta }(\xi )=0}=\epsilon ^\alpha (\widehat{x})+(1+\overline{\gamma })^\alpha {}_{\beta }{}^{}\kappa _{}^{\beta }(\xi )`$ (24)
implies that the gauge (19) is preserved by a local supersymmetry of parameter $`\epsilon (x)`$ whose pull–back to the brane is restricted by being of the form
$`\widehat{\epsilon }^\alpha (\xi ):=\epsilon ^\alpha (\widehat{x})=(1+\overline{\gamma }(\widehat{x}))_\beta \kappa ^\beta (\xi ).`$ (25)
### 2.3 Gauge–fixed description of the SG–M2 interacting system and its properties
Hence, as shown in and discussed above, in a hypothetical superfield description (16) of the supergravity–superbrane interacting system the gauge
$`\widehat{\theta }^\alpha (\xi )=0`$ (26)
(Eq. (19)) and the WZ gauge (2.1) may be fixed simultaneously. In the WZ gauge the integration over the Grassmann superspace coordinates $`\theta ^\alpha `$ in such a superfield action $`𝐒_{SG}[E^a,E^\alpha ,A_3(Z)]`$ would produce a component spacetime supergravity action involving a (hypothetical) set of auxiliary fields. By definition, these auxiliary fields would satisfy algebraic equations which, used in the supergravity action, would lead to the standard supergravity action (in our case that of ) involving only the physical fields of the supergravity multiplet. This action is invariant under the local supersymmetry the algebra of which closes on–shell.
Notice that the auxiliary fields would be contained in the higher order components of $`E_M^A(Z)`$, $`A_{MNK}(Z)`$ (and, perhaps, in some additional auxiliary superfields). The leading ($`\theta =0`$) components of the $`E_M^A(Z)`$ and $`A_{MNK}(Z)`$ superfields in the WZ gauge are either zero, unity, or, in the case of $`E_\mu ^A(Z)`$ and $`A_{\mu \nu \rho }(Z)`$, determine the physical fields $`e_\mu ^a(x)`$, $`\psi ^\alpha (x)`$ and $`A_{\mu \nu \rho }(Z)`$ of the Cremmer–Julia–Scherk (CJS) supergravity multiplet (Eqs. (10), (11) and (12)). As a result, in the gauge defined by $`\widehat{\theta }(\xi )=0`$ (Eq. (26)) plus the WZ gauge (Eqs. (2.1)), the supermembrane action (1) reduces to the action $`S_{M2}^0`$ of a purely bosonic membrane coupled to the physical bosonic fields of the supergravity multiplet only; neither the gravitino nor the auxiliary fields enter the membrane part of the gauge–fixed interacting action.
Hence, in the supergravity part of such a gauge–fixed action for the SG–M2 interacting system one may remove the auxiliary fields through their algebraic equations in the same manner that one would do for the (also hypothetical) pure supergravity action with auxiliary fields. As a result one would arrive at the following gauge–fixed action for the SG–M2 interacting system (see also )
$`S_{SGM2}^0=S_{SG}[e^a,\psi ^\alpha ,A_3]+S_{M2}^0={\displaystyle _{M^{11}}}_{11}[e^a,\psi ^\alpha ,A_3]+{\displaystyle _{W^3}}\left({\displaystyle \frac{1}{3!}}\widehat{e}^a\widehat{}\widehat{e}_aA_3(\widehat{x})\right),`$ (27)
where $`S_{SG}=S_{SG}[e^a,\psi ^\alpha ,A_3]`$ is the standard CJS action for D=11 supergravity and the second term is the action for a purely bosonic brane where the relative coefficient between its two terms is fixed (for a given supergravity action $`S_{SG}[e^a,\psi ^\alpha ,A_3]`$ invariant under definite supersymmetry transformations) since $`S_{M2}^0=S_{M2}[e^a(\widehat{x}),A_3(\widehat{x})]`$ is the bosonic limit of the M2–superbrane action $`S_{M2}[E^a(\widehat{Z}),A_3(\widehat{Z})]`$ of Eq. (1).
The following properties of the gauge–fixed action (27) will be important
* 1) The gauge–fixed description (27) of the supergravity–superbrane interacting system (16) is complete in the sense that it produces a gauge–fixed version of all the dynamical equations that would be obtained from a possible superfield action, including the ‘fermionic equation for bosonic brane’ , which is given by an algebraic condition on the pull–back $`\widehat{\psi }^\alpha :=d\xi ^m\widehat{\psi }_m^\alpha (\xi )`$ of the gravitino to $`W^3`$. It states that a projection of a gamma–trace of $`\widehat{\psi }_m^\alpha `$ vanishes, i.e. that
$`\widehat{\psi }_m\widehat{\mathrm{\Gamma }}^m(1\overline{\gamma })=0`$ (28)
(see Sec. 2.3.2), where
$`\widehat{\psi }_m^\beta :=\widehat{e}_m^a(\xi )\psi _a^\beta (\widehat{x}(\xi ))=_m\widehat{x}^\mu \psi _\mu ^\alpha (\widehat{x}),\widehat{e}_m^a:=_m\widehat{x}^\mu (\xi )e_\mu {}_{}{}^{a}(\widehat{x}(\xi )),`$
$`\widehat{\mathrm{\Gamma }}_{\alpha \beta }^n:=g^{nm}(\xi )\widehat{e}_m^a(\xi )\mathrm{\Gamma }_{a\alpha \beta },`$ (29)
$`g^{nm}(\xi )`$ is inverse of the induced metric $`g_{mn}(\xi )=\widehat{e}_m{}_{}{}^{a}\widehat{e}_{na}^{}`$ (Eq. (4) with $`\widehat{\theta }(\xi )=0`$) and
$`\overline{\gamma }:={\displaystyle \frac{i}{3!\sqrt{|g(\xi )|}}}ϵ^{mnk}\widehat{e}_m{}_{}{}^{a}\widehat{e}_{n}^{}{}_{}{}^{b}\widehat{e}_{k}^{}{}_{}{}^{c}\mathrm{\Gamma }_{abc}^{}`$ (30)
(cf. Eq. (7) for $`\widehat{\theta }(\xi )=0`$) has the properties $`\overline{\gamma }^2=I`$, $`tr(\overline{\gamma })=0`$.
* 2) The equations of motion for the bosonic supergravity fields get (or may get) a source term contribution from the superbrane, while the gauge–fixed equations for the bulk fermionic fields are sourceless (see Sec.2.3.1).
* 3) The action (27) possesses half of the local supersymmetries of the pure supergravity action $`S_{SG}[e^a,\psi ^\alpha ,A_3]`$. This is characterized by the standard transformation rules for the supergravity fields (see Eqs. (47), (48), (49) below) and by the following conditions (see Sec. 2.3.3) restricting the local supersymmetry parameter on the worldvolume $`W^3`$,
$`\widehat{\epsilon }^\alpha :=\widehat{\epsilon }^\alpha (\widehat{x}(\xi ))=(1+\overline{\gamma })^\alpha {}_{\beta }{}^{}\kappa _{}^{\beta }(\xi ).`$ (31)
* 4) The local supersymmetry algebra closes on-shell in exactly the same manner as it does for the case of free supergravity (Sec. 2.3.4).
Property 1) might seem strange since no worldvolume fermionic degrees of freedom are seen directly in the gauge–fixed interacting action (27) involving the bosonic brane action. But this ‘fermionic equation for the bosonic brane’ can be derived from the selfconsistency condition $`𝒟\mathrm{\Psi }_{10\alpha }=0`$ for the gravitino equation $`\mathrm{\Psi }_{10\alpha }=0`$ (Eq. (35) below) which, according to 2), remains sourceless in the presence of the bosonic brane . Thus, it is convenient to discuss property 2) first.
#### 2.3.1 Field equations for the SG–M2 system (property 2)
Varying the CJS action with respect to differentail forms $`\delta S_{SG}=2i\mathrm{\Psi }_{10\alpha }\delta \psi ^\alpha +𝒢_8\delta A_3+M_{10a}\delta E^a`$ one can write the ‘free’ supergravity equations in differential form notation (see and, e.g. ). The same can be done for the (gauge–fixed) field equations of the SG–M2 interacting system, $`\delta (S_{SG}+S_{M2}^0)=2i\mathrm{\Psi }_{10\alpha }\delta \psi ^\alpha +(𝒢_8J_8)\delta A_3+(M_{10a}J_{10a})\delta E^a`$. The variation of the bosonic membrane part $`S_{M2}^0`$ in the action is written as an integral over spacetime $`M^{11}`$ with the use of the currents (see ; $`\widehat{e}^b:=d\xi ^n\widehat{e}_n{}_{}{}^{b}(\xi )`$, $`\widehat{e}_n{}_{}{}^{b}:=_n\widehat{x}^\mu (\xi )e_\mu {}_{}{}^{b}(\widehat{x})`$)
$`J_{10a}={\displaystyle \frac{1}{2e(x)}}e_b^{\mathrm{\hspace{0.17em}10}}{\displaystyle \underset{W^3}{}}\widehat{}\widehat{e}_a\widehat{e}^b\delta ^{11}(x\widehat{x}(\xi ))`$ $`=e_b^{\mathrm{\hspace{0.17em}10}}{\displaystyle \underset{W^3}{}}d^3\xi {\displaystyle \frac{\sqrt{|g|}g^{mn}\widehat{e}_{na}\widehat{e}_m^b}{2|det(e_\nu {}_{}{}^{c}(\widehat{x}))|}}\delta ^{11}(x\widehat{x}(\xi )),`$ (32)
$`J_8={\displaystyle \frac{1}{e(x)}}e_{abc}^{\mathrm{\hspace{0.17em}8}}{\displaystyle \underset{W^3}{}}\widehat{e}^a\widehat{e}^b\widehat{e}^c\delta ^{11}(x\widehat{x}(\xi ))`$ $`=e_{abc}^{\mathrm{\hspace{0.17em}8}}{\displaystyle \underset{W^3}{}}d^3\xi {\displaystyle \frac{ϵ^{mnk}\widehat{e}_m^a\widehat{e}_n^b\widehat{e}_k^c}{|det(e_\nu {}_{}{}^{c}(\widehat{x}))|}}\delta ^{11}(x\widehat{x}(\xi )),`$ (33)
which describe the brane source terms in the Einstein and gauge field equations. The Einstein and the Rarita–Schwinger equations of the interacting system are written in terms of the ten–forms
$`M_{10a}`$ $`:={\displaystyle \frac{1}{4}}R^{bc}e_{abc}^8+{\displaystyle \frac{1}{2}}(i_aF_4F_4+F_4i_a(F_4))+𝒪(\psi ^2)+𝒪(\psi ^4)=J_{10a},`$ (34)
$`\mathrm{\Psi }_{10\alpha }`$ $`:=𝒟\psi ^\beta \overline{\mathrm{\Gamma }}_{\beta \alpha }^{(8)}=0,`$ (35)
while the eight–form expression of the three–form gauge field equation reads
$`𝒢_8`$ $`:=d(F_4+b_7A_3dA_3)=J_8,b_7:={\displaystyle \frac{i}{2}}\psi ^\alpha \psi ^\beta \overline{\mathrm{\Gamma }}_{\alpha \beta }^{(5)}.`$ (36)
In Eqs. (34), (35) and (36) the eight–forms $`e_{abc}^8`$, $`\overline{\mathrm{\Gamma }}_{\beta \alpha }^{(8)}`$ and the five form $`\overline{\mathrm{\Gamma }}_{\beta \alpha }^{(5)}`$ are defined by
$`e_{a_1\mathrm{}a_q}^{(11q)}:={\displaystyle \frac{1}{(11q)!}}ϵ_{a__1\mathrm{}a_qb__1\mathrm{}b_{_{11q}}}e^{b__1}\mathrm{}e^{b_{_{11q}}},\overline{\mathrm{\Gamma }}{}_{}{}^{(q)}:={\displaystyle \frac{1}{q!}}e^{a_q}\mathrm{}e^{a__1}\mathrm{\Gamma }_{a__1\mathrm{}a_q},`$ (37)
the four–form $`F_4`$ is the ‘supersymmetric’ field strength of the three–form gauge field $`A_3`$,
$`F_4:={\displaystyle \frac{1}{4!}}e^{a_4}\mathrm{}e^{a_1}F_{a_1\mathrm{}a_4}(x)=dA_3{\displaystyle \frac{1}{2}}\psi ^\alpha \psi ^\beta \overline{\mathrm{\Gamma }}_{\alpha \beta }^{(2)},`$ (38)
and $`F_4:=e_{abcd}^7F^{abcd}`$. The spin connection $`\omega ^{ab}=\omega ^{ba}`$ are expressed through the graviton and gravitino by the solution of the torsion constraint
$`T^a(x):=De^a=de^ae^b\omega _b{}_{}{}^{a}=i\psi \mathrm{\Gamma }^a\psi ,`$ (39)
which formally coincides with the leading component of the on–shell superspace constraint (8). The explicit expressions for the two–fermionic and four–fermionic contributions, $`𝒪(\psi ^2)`$ and $`𝒪(\psi ^4)`$, to the Einstein equations (34) will not be needed in this paper. The generalized covariant derivative $`𝒟\psi ^\alpha `$ in (35) is defined by
$`𝒟\psi ^\alpha =d\psi ^\alpha \psi ^\beta \omega _\beta {}_{}{}^{\alpha },`$ $`\omega _\alpha {}_{}{}^{\beta }={\displaystyle \frac{1}{4}}\omega ^{ab}\mathrm{\Gamma }_{ab\alpha }{}_{}{}^{\beta }+e^at_{a\alpha }^\beta `$ (40)
and contains, in addition to the spin connection $`\frac{1}{4}\omega ^{ab}\mathrm{\Gamma }_{ab}{}_{\alpha }{}^{}^\beta `$, the covariant contribution $`e^at_{a\alpha }^\beta `$,
$`t_{a\beta }{}_{}{}^{\alpha }:={\displaystyle \frac{i}{18}}(F_{ac_1c_2c_3}\mathrm{\Gamma }^{c_1c_2c_3}{}_{\beta }{}^{\alpha }{\displaystyle \frac{1}{8}}F^{c_1c_2c_3c_4}\mathrm{\Gamma }_{ac_1c_2c_3c_4}{}_{\beta }{}^{\alpha }),`$ (41)
expressed through the ‘supersymmetric’ field strength $`F_{abcd}(x)`$ of $`A_3`$, Eq. (38). This covariant part of the generalized connection thus describes the coupling of the bulk gravitino field to the fluxes of the three–form gauge field $`A_3`$.
The reason for the absence of source in the fermionic equation (35) obtained by varying the gauge–fixed action (27) with respect to the gravitino field is, clearly, that the bosonic brane action $`S_{M2}^0`$ in (27) does not include the gravitino $`\widehat{\psi }^\alpha `$; this, in turn, follows from the absence of the fermionic supervielbein ($`E^\alpha (\widehat{Z})`$) in the supermembrane action $`𝐒_{M2}`$ of Eq. (1). Nevertheless, the absence of an explicit source term in (35) does not imply that the gravitino is decoupled from the brane source since Eq. (35) includes the vielbein $`e_\mu ^a(x)`$ (entering also through the composite spin connection $`\omega ^{ab}`$) and the field strength of the three–form gauge field $`A_3`$ that do obey the sourceful Eqs. (34) and (36).
Notice that the system of interacting equations, including Eqs. (34), (36), (35), (28) as well as the bosonic equation for the brane, admits particular solutions with $`\psi ^\alpha (x)=0`$. Inserting $`\psi ^\alpha (x)=0`$ back into the equations one arrives at the well–known system of purely bosonic supergravity equations in .
#### 2.3.2 ‘Fermionic equations’ for the bosonic brane interacting with supergravity (property 1)
To understand how the ‘fermionic equation for the bosonic brane’ results from the consistency conditions of the gravitino equation one can use the identity (see e.g. )
$$𝒟\mathrm{\Psi }_{10\alpha }=i\psi ^\beta \left(M_{10a}\mathrm{\Gamma }_{\beta \alpha }^a+\frac{i}{2}𝒢_8\overline{\mathrm{\Gamma }}_{\beta \alpha }^{(2)}\right)$$
(42)
that expresses the generalized covariant derivative of the l.h.s. of the fermionic equation (35) in terms of the left hand sides of the Einstein and the $`A_3`$ gauge field equations, $`M_{10a}`$ and $`𝒢_8`$ in Eqs. (34) and (36), respectively. For ‘free’ $`D=11`$ supergravity the equations of motion are $`\mathrm{\Psi }_\alpha =0`$, $`M_{10a}=0`$ and $`𝒢_8=0`$, and the eleven–form identity (42) shows their interdependence. This Noether identity reflects a local gauge symmetry of the CJS supergravity action $`S_{SG}[e^a,\psi ^\alpha ,A_3]`$, the local supersymmetry of $`D=11`$ supergravity (Eqs. (47)–(49) below).
When supergravity interacts with a bosonic membrane, $`S_{SG}S_{SG}+S_{M2}^0`$ like in the gauge fixed description, the bosonic field equations acquire source terms and read $`M_{10a}=J_{10a}`$ \[Eqs. (34), (32)\] and $`𝒢_8=J_8`$ \[Eqs. (36), (33)\]; the fermionic equation, however, remains sourceless, $`\mathrm{\Psi }_{10\alpha }=0`$ \[Eq. (35)\]. Hence, Eq. (42) produces the following equation for the M2–brane currents $`J_{10a}`$ and $`J_8`$ (see Eqs. (32) and (33))
$$i\psi ^\beta (J_{10a}\mathrm{\Gamma }_{\beta \alpha }^a+\frac{i}{2}J_8\overline{\mathrm{\Gamma }}{}_{\beta \alpha }{}^{(2)})=0.$$
(43)
Due to the currents, this eleven–form equation has support on the M2–brane worldvolume $`W^3`$ and so it can be written as a three–form equation on $`W^3`$ in terms of the pull–backs of the graviton and gravitino . When $`S_{SG}+S_{M2}^0`$ provides the gauge fixed description of the supergravity–supermembrane interaction the currents $`J_{10a}`$ and $`J_8`$ are defined by Eqs. (32) and (33) and the equivalent form of Eq. (43) reads
$`\widehat{\psi }^\beta \left(i\widehat{}\widehat{e}^a\mathrm{\Gamma }_{a\beta \alpha }+{\displaystyle \frac{1}{2}}\widehat{e}^b\widehat{e}^c\mathrm{\Gamma }_{bc\beta \alpha }\right)=0.`$ (44)
A simple algebra allows us to present Eq. (44) in the form of
$`\widehat{\mathrm{\Xi }}_{3\alpha }`$ $`:=`$ $`\widehat{}\widehat{e}^a\widehat{\psi }^\beta (\mathrm{\Gamma }_a(1\overline{\gamma }))_{\beta \alpha }=0`$ (45)
where the action of $`\widehat{}`$ is defined by $`\widehat{}\widehat{e}^a=\frac{1}{2}d\xi ^nd\xi ^mϵ_{mnk}\sqrt{|g(\xi )|}g^{kl}_l\widehat{x}^\mu e_\mu {}_{}{}^{a}(\widehat{x})`$ (Eq. (3) with $`\widehat{\theta }=0`$, i.e. for $`\widehat{}=\widehat{}|_{\widehat{\theta }=0}`$). Eq. (45) is an equivalent form of Eq. (28) written in a conventional differential form notation,
$`\widehat{\mathrm{\Xi }}_{3\alpha }`$ $``$ $`d^3\xi \widehat{\psi }_m\widehat{\mathrm{\Gamma }}^m(1\overline{\gamma }).`$ (46)
#### 2.3.3 Supersymmetry of the gauge–fixed action (property 3)
Eq. (42) is the Noether identity for the local supersymmetry of the pure supergravity action $`S_{SG}[e^a,\psi ^\alpha ,A_3]`$,
$`\delta _\epsilon e^a`$ $`=`$ $`2i\psi \mathrm{\Gamma }^a\epsilon :=2i\psi ^\alpha \mathrm{\Gamma }_{\alpha \beta }^a\epsilon ^\beta ,`$ (47)
$`\delta _\epsilon \psi ^\alpha `$ $`=`$ $`𝒟\epsilon ^\alpha (x)=D\epsilon ^\alpha (x)\epsilon ^\beta (x)e^at_{a\beta }{}_{}{}^{\alpha }(x),`$ (48)
$`\delta _\epsilon A_3`$ $`=`$ $`\psi ^\alpha \overline{\mathrm{\Gamma }}_{\alpha \beta }^{(2)}\epsilon ^\beta ,`$ (49)
where $`𝒟`$ is the generalized covariant derivative of Eq. (40), $`D=d\omega `$ is the standard covariant derivative, and the tensorial part of the generalized connection $`t_{a\beta }{}_{}{}^{\alpha }(x)`$ is defined in Eq. (41). The fact that Eq. (42) is not identically satisfied in the presence of a bosonic brane, i.e. when $`S_{SG}S_{SG}+S_{M2}^0`$, reflects the fact that the bosonic brane action $`S_{M2}^0`$ breaks the local supersymmetry (47)–(49).
When the bosonic brane is the purely bosonic ($`\widehat{\theta }=0`$) ‘limit’ of a superbrane, the sum of the supergravity action and the action of bosonic brane provides a gauge–fixed description of the supergravity—superbrane (SG–M2) interacting system and preserves one–half of the local supersymmetry of $`S_{SG}`$. This half of the local supersymmetry is defined by the restriction (31) on the pull–back of the supersymmetry parameter to the membrane worldvolume $`W^3`$, $`\widehat{\epsilon }^\alpha :=\widehat{\epsilon }^\alpha (\widehat{x}(\xi ))=(1+\overline{\gamma })^\alpha {}_{\beta }{}^{}\kappa _{}^{\beta }(\xi )`$. Its preservation can be shown in two ways, either explicitly (see also Sec. 3.3 below) or using the fact that the action (27) provides a gauge–fixed version of the hypothetical superfield description of the supergravity–superbrane interaction as discussed also in Sec. 2.2.
#### 2.3.4 On–shell closure of the local supersymmetry algebra in the spacetime gauge–fixed description of SG–M2 system (property 4)
As known from the pioneering paper , the local supersymmetry transformations (47)–(49) that leave invariant the supergravity action $`S_{SG}[e^a,\psi ^\alpha ,A_3(x)]`$ form an algebra which is closed on shell, i.e. using the ‘free’ supergravity equations. The structure of this algebra is schematically
$`[\delta _{ϵ_1},\delta _{ϵ_1}][fields]=\left(\delta _{b^\mu (ϵ_1,ϵ_2)}+\delta _{L^{ab}(ϵ_1,ϵ_2)}+\delta _{ϵ_3(ϵ_1,ϵ_2)}+𝒦_{(ϵ_1,ϵ_2)}\right)[fields],`$ (50)
where $`\delta _{b^\mu }`$ determines the general coordinate transformations, $`\delta _{L^{ab}}`$ the local Lorentz transformations, $`\delta _ϵ`$ the local supersymmetry transformations (47)–(49) and $`𝒦_{(ϵ_1,ϵ_2)}[fields]`$ denotes terms that express the non–closure of the algebra and that become zero on shell.
Let us now consider the SG–M2 interacting system. The form of the supersymmetry transformations leaving invariant the coupled supergravity–bosonic brane action (27) (i.e. preserving the gauge $`\widehat{\theta }(\xi )=0`$ for the interacting supergravity–superbrane system ) is exactly the same as that of the supersymmetry of ‘free’ supergravity <sup>3</sup><sup>3</sup>3Notice that this is not the case for the supersymmetric brane world models in . There, the brane actions also contain the pull–back of the gravitino field. Probably these two facts are related and prevent or hamper a superfield formulation of the brane actions of . In all other respects the models of are similar to the dynamical systems of supergravity interacting with standard superbranes as they are presented in the gauge–fixed description of and Sec. 2 of this paper. The breaking of 1/2 of the supersymmetry in the gauge–fixed description corresponds to imposing a kind of boundary conditions on the supersymmetry parameter in .. However, in principle, the last term in (50) might spoil the closure of the local supersymmetry algebra of supergravity–bosonic membrane system that provides the gauge–fixed description of the SG–M2 interacting system). This is not so, however. On the bosonic fields of the supergravity multiplet, $`e_\mu {}_{}{}^{a}(x)`$ and $`A_{\mu \nu \rho }(x)`$, the algebra (50) is closed off shell , i.e. without any use of the equations of motion. This means that
$`𝒦_{(ϵ_1,ϵ_2)}[e^a(x)]0,𝒦_{(ϵ_1,ϵ_2)}[A_3(x)]0.`$ (51)
Hence the on–shell character of the supersymmetry algebra comes from the fermionic fields since only $`𝒦_{(ϵ_1,ϵ_2)}[\psi ]0`$ off–shell <sup>4</sup><sup>4</sup>4The statement that in the absence of the auxiliary fields $`𝒦_{(ϵ_1,ϵ_2)}[fermionicfields]0`$ off–shell while $`𝒦_{(ϵ_1,ϵ_2)}[bosonicfields]=0`$ seems to be quite general, i.e. valid for many supersymmetric theories in various dimensions, see e.g. . To our knowledge, the only exception is provided by the supersymmetry transformations that preserve the equations of motion for supermultiplets that include self–dual gauge fields, where the selfduality condition for the bosonic gauge field is also needed to close the supersymmetry algebra.. Moreover, and this is the key point, only the fermionic equations are necessary to close supersymmetry algebra on the fermionic fields; schematically,
$`𝒦_{(ϵ_1,ϵ_2)}[\psi ^\alpha (x)]\mathrm{\Psi }_{10\alpha }.`$ (52)
But as noticed above (following ), the fermionic equation for the interacting system in the gauge–fixed description given by the sum of supergravity action and the action for bosonic brane preserving a half of the local supersymmetry remains formally the same (i.e., sourceless) as that for ‘free’ supergravity, $`\mathrm{\Psi }_{10\alpha }=0`$. Hence,
$`𝒦_{(ϵ_1,ϵ_2)}[\psi ^\alpha (x)]|_{_{_{interactingsystem}^{onshellforthe}}}=0.`$ (53)
Further, the local supersymmetry transformations act only on the fields of supergravity multiplet, Eqs. (47)–(49), since the only supermembrane field in the gauge–fixed description, the bosonic $`\widehat{x}(\xi )`$, is inert under the local spacetime supersymmetry. Thus the on-shell closure of the local supersymmetry algebra of the gauge–fixed description of the supergravity–supermembrane interacting system follows from that of the pure $`D=11`$ supergravity theory.
## 3 Goldstone nature of the supermembrane fermionic fields and Dirac equation for the supermembrane in a D=11 supergravity background with fluxes
The Goldstone nature of the superbrane coordinate functions, in particular of the fermionic functions $`\widehat{\theta }^\alpha (\xi )`$, has been known for a long time . For a superbrane interacting with dynamical supergravity the $`\widehat{\theta }^\alpha (\xi )`$ are Goldstone (or compensator) fields for the local supersymmetry, a fact that explains the possibility of taking the gauge (26), $`\widehat{\theta }(\xi )=0`$, by using this local supersymmetry (see Sec. 2.2).
In this gauge the Lagrangian description of the system is provided by the sum (27) of the spacetime supergravity action without auxiliary fields and of the bosonic M2–brane action . The full set of equations of motion is given by the supergravity field equations (34), (36), (35), the bosonic brane equations (cf. Eq. (5))
$`D(\widehat{}\widehat{e}_a)=2i_aF_4{\displaystyle \frac{1}{2}}\widehat{e}^b\widehat{\psi }\mathrm{\Gamma }_{ab}\widehat{\psi },i_aF_4:={\displaystyle \frac{1}{3!}}\widehat{e}^d\widehat{e}^c\widehat{e}^bF_{abcd}(\widehat{x}),`$ (1)
and the ‘fermionic equation for the bosonic brane’ , Eq. (45), (cf. (6))
$`\widehat{\mathrm{\Xi }}_{3\alpha }=\widehat{e}^a\widehat{\psi }^\beta \left(\mathrm{\Gamma }_a(1\overline{\gamma })\right)_{\beta \alpha }=0.`$ (2)
In the $`\widehat{\theta }(\xi )=0`$ gauge, the fermionic degrees of freedom of the superbrane, usually associated with $`\widehat{\theta }(\xi )`$, are contained in the pull–back $`\widehat{\psi }^\beta `$ of the bulk gravitino to the worldvolume $`W^3`$ as zero modes corresponding to the supersymmetry broken by the brane <sup>5</sup><sup>5</sup>5See for a discussion of the brane degrees of freedom as zero modes, but starting from certain brane solutions of the supergravity equations.. Namely, the fact that half of the supersymmetry is spontaneously broken by the presence of the supermembrane is reflected by the explicit breaking of half of the local supersymmetry by the bosonic brane in the gauge–fixed description (27). Hence, in the presence of the supermembrane, the remaining local supersymmetry does not produce the same number of gauge–fixing conditions on the gravitino field as the full local supersymmetry of ‘free’ supergravity. Nevertheless, as shown in , this super–Higgs effect in the presence of a superbrane does not make the gravitino massive, because the ‘fermionic equation for the bosonic brane’, Eq. (45), takes the rôle of the lost gauge–fixing conditions and keeps the number of polarizations of the gravitino equal to those in ‘free’ supergravity. However, the fermionic zero modes corresponding to the supersymmetry broken by the membrane remain in the pull–back $`\widehat{\psi }`$ of the bulk gravitino $`\psi `$ to $`W^3`$ <sup>6</sup><sup>6</sup>6To illustrate the above statement one can consider the weak field case, making a linear approximation in the fields to find the general solution of the equations; schematically, $`\varphi =\varphi _0+(b(p)e^{+ipx}+b^{}(p)e^{ipx})`$. Then the number of polarizations (which distinguishes between the massive and massless cases) is determined by the oscillating contributions, and depends on the number of conditions imposed on the creation and annihilation operators $`b^{}(p)`$ and $`b(p)`$, while the zero modes are associated with the non–oscillating contribution $`\varphi _0`$. The same occurs in the spacetime Higgs effect in general relativity interacting with branes . A bulk graviton does not get mass and keeps the same number of polarizations in the presence of a $`p`$–brane because the bosonic equations of the brane replace the gauge fixing conditions that were lost due to the spontaneous breaking of the diffeomorphism symmetry by the $`p`$–brane . However, the corresponding gauge–fixing conditions for the case of free gravity allow for a residual gauge symmetry, which is absent when the bosonic equations of the brane take the rôle of the gauge fixing conditions. This set of residual gauge symmetries broken by the brane is the origin of the zero modes of the graviton on $`W^{p+1}`$ that describe the brane degrees of freedom in the ‘static’ gauge (the statement in that a $`p`$–brane does not carry any local degrees of freedom in the presence of dynamical gravity refers to the oscillating degrees of freedom -polarizations- of the graviton and gravitino, as discussed above, not to the zero modes). We thank W. Siegel for an illuminating discussion on this point.. Precisely these zero modes represent the $`16`$ fermionic degrees of freedom of supermembrane in the gauge–fixed description of (27).
To summarize, in the gauge–fixed description of the supergravity–supermembrane interaction provided by the set of equations (34), (36), (35), (1), (2), the bulk gravitino carries both the supergravity and the superbrane fermionic degrees of freedom as determined by the solution of field equation (35) with the boundary conditions (2) on the 3–dimensional ‘defect’, the brane worldvolume $`W^3`$. This description is convenient in studying the cases where both the effects from the bulk and from the worldvolume fermions are equally important and there is no need to separate their contribution.
However, in some cases (interesting e.g. for M-theory–based ‘realistic’ model building, see ) it may happen that the effects from the worldvolume fermions, and in particular the explicit form of their interaction with the flux, constitute the main interest. Then, when starting from our gauge–fixed description, one faces the problem of visualizing the fermionic degrees of freedom of the superbrane, i.e. the supermembrane coordinate functions $`\widehat{\theta }(\xi )`$. This will be the main subject of the study below.
In the light of Goldstone nature of $`\widehat{\theta }(\xi )`$, the general answer should not be too surprising: the recovery of the $`\widehat{\theta }(\xi )`$ contributions to the action and equations of motion can be done by making (consistently) a local supersymmetry transformation the parameter of which is identified with the Goldstone fermion field $`\widehat{\theta }(\xi )`$. We begin by showing how the supermembrane fermionic equations in a supergravity background with fluxes, this is to say with nonvanishing $`F_{abcd}`$, can be obtained on this way.
### 3.1 Dirac equation for the supermembrane in a supergravity background with fluxes from the gauge–fixed approach
When supergravity is treated as a background, one concentrates on the supermembrane equations. In our gauge–fixed description these are given by the bosonic equation (1) and the fermionic Eq. (2) which is more a condition on the pull–back of the gravitino than a dynamical equation. To separate the contribution form the bulk fermions and from the supermembrane fermions one makes, following the above prescription, the local supersymmetry transformations (47)–(49) of the supergravity fields in (45) and identifies the (pull–back of the) parameter of these transformations with the supermembrane fermionic field, $`ϵ(\widehat{x}(\xi ))=\widehat{\theta }(\xi )`$. The result at first order in $`\widehat{\theta }(\xi )`$ is given by
$`\widehat{e}^a\widehat{\psi }^\beta \left(\mathrm{\Gamma }_a(1\overline{\gamma })\right)_{\beta \alpha }+\delta _{\widehat{ϵ}=\widehat{\theta }}(\widehat{e}^a\widehat{\psi }^\beta \left(\mathrm{\Gamma }_a(1\overline{\gamma })\right)_{\beta \alpha })=0,`$ (3)
or, in more detail ($`\widehat{e}_m^a=_m\widehat{x}^\mu e_\mu ^a(\widehat{x})`$, $`\widehat{\psi }_m^\alpha =_m\widehat{x}^\mu \psi _\mu ^\alpha (\widehat{x})`$, $`\widehat{\mathrm{\Gamma }}^k=g^{km}(\xi )\widehat{e}_m^a\mathrm{\Gamma }_a`$, $`g_{mn}(\xi )=\widehat{e}_m^a\widehat{e}_{na}`$),
$`\widehat{}\widehat{e}^a`$ $`\left(\widehat{\psi }\mathrm{\Gamma }_a(1\overline{\gamma })+𝒟\widehat{\theta }\mathrm{\Gamma }_a(1\overline{\gamma })+2i\widehat{\psi }_k\widehat{\mathrm{\Gamma }}^k\widehat{\theta }\widehat{\psi }\mathrm{\Gamma }_a+{\displaystyle \frac{ϵ^{mnk}\widehat{e}_m^{b_1}\widehat{e}_n^{b_2}}{\sqrt{|g(\xi )|}}}(\widehat{\psi }_k\mathrm{\Gamma }^{b_3}\widehat{\theta })\widehat{\psi }\mathrm{\Gamma }_a\mathrm{\Gamma }_{b_1b_2b_3}\right)_\alpha +`$ (4)
$`+2i\widehat{}\widehat{e}^b(\widehat{\psi }\mathrm{\Gamma }_a(1\overline{\gamma }))_\alpha (\widehat{\psi }^l\mathrm{\Gamma }_b\widehat{\theta }+\psi _b(\widehat{x})\widehat{\mathrm{\Gamma }}^l\widehat{\theta })\widehat{e}_l{}_{}{}^{a}2i\widehat{}\widehat{\psi }\mathrm{\Gamma }^a\widehat{\theta }(\widehat{\psi }\mathrm{\Gamma }_a(1\overline{\gamma }))_\alpha =0,`$
where again the generalized covariant derivative $`𝒟`$ is given by Eqs. (40), (41) and, thus, includes a contribution from the fluxes $`F_{abcd}`$. We have checked explicitly that Eq. (4) formally coincides with the first order equation that can be obtained within the standard ‘background superfield’ approach (without setting $`\widehat{\psi }=0`$ as in ). By ‘formally’ we mean that in the equations obtained in the standard framework the graviton, the gravitino and the gauge field strength are, strictly speaking, solutions of the ‘free’ supergravity equations, while in our case such a restriction is absent and one can use, e.g., solutions of the interacting system of equations.
Eq. (4) is rather complicated. A simpler one results when in (4) the gravitino field is set equal to zero. This gives
$`\widehat{\psi }^\alpha =0:\widehat{}\widehat{e}_a𝒟\widehat{\theta }^\beta \left(\mathrm{\Gamma }^a(1\overline{\gamma })\right)_{\beta \alpha }=0,`$ (5)
or, equivalently,
$$𝒟\widehat{\theta }^\beta (i\widehat{}\widehat{e}_a\mathrm{\Gamma }_{\beta \alpha }^a+\widehat{\overline{\mathrm{\Gamma }}}{}_{\beta \alpha }{}^{(2)})=0.$$
Eq. (5) formally coincides with the M2–brane Dirac equation which is obtained in within the on–shell background superfield approach, namely by expanding Eq. (6) in $`\widehat{\theta }`$ for $`\psi =0`$. To see this explicitly, one may use the expression (40), (41) for the generalized covariant derivative $`𝒟`$ in (5) and the worldvolume tensor notation ($`\widehat{\mathrm{\Gamma }}_n:=\widehat{e}_n^a\mathrm{\Gamma }_a`$, $`\widehat{\mathrm{\Gamma }}^n:=g^{nm}(\xi )\widehat{e}_m^a\mathrm{\Gamma }_a`$ etc.) to arrive at
$`\left(D_n\widehat{\theta }+{\displaystyle \frac{i}{18}}\widehat{e}_n^a\left(F_{ab_1b_2b_3}\widehat{\theta }\mathrm{\Gamma }^{b_1b_2b_3}{\displaystyle \frac{1}{8}}F^{b_1b_2b_3b_4}\widehat{\theta }\mathrm{\Gamma }_{ab_1b_2b_3b_4}\right)\right)^\beta \left(\widehat{\mathrm{\Gamma }}^n(1\overline{\gamma })\right)_{\beta \alpha }=0,`$ (6)
$`D_n\widehat{\theta }^\alpha :=_n\widehat{\theta }^\alpha (\xi ){\displaystyle \frac{1}{4}}\widehat{e}_n^c\omega _c^{ab}\widehat{\theta }^\beta (\xi )\mathrm{\Gamma }_{ab}{}_{\beta }{}^{}{}_{}{}^{\alpha }.`$
In this form the interaction of the supermembrane fermionic field with the $`A_3`$ ‘fluxes’, this is to say with the field strength $`F_{abcd}`$, is manifest.
Linearizing Eq. (4) in all the fermions, i.e. ignoring $`𝒪(\widehat{\theta }\widehat{\psi })`$, $`𝒪(\widehat{\psi }^2)`$ together with the $`𝒪(\widehat{\theta }^2)`$ contributions, we find the equation
$`\widehat{}\widehat{e}^a`$ $`(𝒟\widehat{\theta }+\widehat{\psi })\mathrm{\Gamma }_a(1\overline{\gamma })=0`$ (7)
which includes the pull–back of the gravitino $`\widehat{\psi }`$ and the Goldstone fermion $`\widehat{\theta }`$ in the combination $`(𝒟\widehat{\theta }+\widehat{\psi })`$ only, which is invariant under the linearized supersymmetry. This observation supports the discussed fact (see footnote 6 and above) that, in the gauge $`\widehat{\theta }=0`$, the zero modes describing the brane fermionic degrees of freedom appear in the pull–back $`\widehat{\psi }`$ of the bulk gravitino to $`W^3`$.
### 3.2 On the contribution of the supermembrane fermionic field to the full set of interacting equations
From the point of view of the interacting system, the setting $`\widehat{\psi }=0`$ above (and in ) or, taking into account the previous supersymmetry transformations that make manifest the Goldstone degrees of freedom, $`\widehat{\psi }^\alpha :=\psi ^\alpha (\widehat{x}(\xi ))=𝒟\widehat{\theta }^\alpha (\xi )`$ is a kind of ansatz, or boundary condition, for the gravitino field on $`W^3`$. As such, its consistency with the supergravity equations should be checked. This is a convenient point to begin discussing the contribution of the supermembrane fermionic fields to the complete system of interacting equations, which includes the field equations whose gauge–fixed form is given by Eqs. (34), (36) and (35).
It is natural to consider the above relation $`\widehat{\psi }^\alpha =𝒟\widehat{\theta }^\alpha (\xi )`$ on $`W^3`$ as produced by the ansatz
$`\psi ^\alpha (x)=𝒟\stackrel{~}{\theta }^\alpha (x)`$ (8)
for the bulk gravitino, where the tilde denotes function on spacetime. Here the defining property of the Volkov–Akulov Goldstone fermion (see ) $`\stackrel{~}{\theta }^\alpha (x)`$ is that its pull–back on $`W^3`$ coincides with the supermembrane fermionic field,
$`\stackrel{~}{\theta }^\alpha (\widehat{x}(\xi ))=\widehat{\theta }^\alpha (\xi ).`$ (9)
The irrelevance of the properties of $`\stackrel{~}{\theta }^\alpha (x)`$ outside the brane worldvolume $`W^3`$ is just the statement of the local supersymmetry of the ‘free’ supergravity action.
However, a direct substitution of the ansatz (8), (9) into the gravitino equations (35) would produce a problem. After some algebra (e.g. using identities from and (42)) one finds that such a Volkov–Akulov Goldstone fermion $`\stackrel{~}{\theta }^\alpha (x)`$ would obey $`i\stackrel{~}{\theta }^\beta (J_{10a}\mathrm{\Gamma }_{\beta \alpha }^a+\frac{i}{2}J_8\overline{\mathrm{\Gamma }}{}_{\beta \alpha }{}^{(2)})=0`$. This is equivalent (cf. Sec. 2.3.2) to the condition $`\widehat{\theta }^\beta (\xi )(i\widehat{}\widehat{e}^a\mathrm{\Gamma }_{a\beta \alpha }+\frac{1}{2}\widehat{e}^b\widehat{e}^c\mathrm{\Gamma }_{bc\beta \alpha })=0`$ or, equivalently, $`\widehat{\theta }^\beta (\xi )(\widehat{\mathrm{\Gamma }}_n(1\overline{\gamma }))_{\beta \alpha }=0`$ which implies the effective vanishing of the supermembrane fermionic field (actually $`(1\overline{\gamma })\widehat{\theta }=0`$, but this in turn implies $`\widehat{\theta }=0`$, since the $`(1+\overline{\gamma })\widehat{\theta }`$ part can be removed by the preserved supersymmetry gauge transformations which correspond to the $`\kappa `$–symmetry of the superbrane).
The reason for this apparent problem lies in the fact that the correct prescription to recover the supermembrane fermionic fields is to make the supersymmetry transformations of the gauge–fixed equations rather than using an ansatz like (8), (9) in them. Despite that the r.h.s. of (8) coincides with the gravitino supersymmetry transformations, its substitution into (35) does not automatically give the supersymmetry transformations of this equations. The point is that in a gauge–fixed equation where some Goldstone fields are set equal to zero, e.g. $`\widehat{\theta }^\alpha (\xi )=0`$, a zero in the r.h.s. of this equation may come from a term proportional to $`\widehat{\theta }^\alpha (\xi )`$. As it is suggested by the study of the superfield description of the $`D=4`$, $`N=1`$ supergravity–superparticle and supergravity–superstring systems , this is exactly the case for the gauge fixed form of the gravitino equation (35). Namely, the fully supersymmetric (not gauge–fixed) counterpart of this equation contains a r.h.s. proportional to $`\widehat{\theta }^\alpha (\xi )`$. Schematically,
$`\mathrm{\Psi }_{10\alpha }`$ $`:=𝒟\psi ^\beta \overline{\mathrm{\Gamma }}_{\beta \alpha }^{(8)}=𝒪(\widehat{\theta }^\alpha (\xi )).`$ (10)
In other words, $`\mathrm{\Psi }_{10\alpha }\widehat{\theta }^\alpha `$ rather than zero like in Eq. (35) which comes from (10) in the gauge $`\widehat{\theta }=0`$. Then, taking into account the presence of a right hand side proportional to $`\widehat{\theta }^\alpha `$ in a fully supersymmetric (not gauge–fixed) counterpart of Eq. (35), one can use a local supersymmetry transformation to find an approximate expression for this r.h.s. ($`𝒪(\widehat{\theta })`$ in (10)) up to the first order in $`\widehat{\theta }^\alpha `$. This suggests a way of deriving the contributions of the supermembrane Goldstone fermion $`\widehat{\theta }^\alpha (\xi )`$ to the supergravity equations from the local supersymmetry transformations of the gauge–fixed system of interacting equations (34), (35), (36) or of the gauge–fixed interacting action (27), which should work at least in low orders in $`\widehat{\theta }^\alpha (\xi )`$.
As a first step in this direction let us derive the gravitino vertex operator of and, thus, find the contribution proportional to $`\widehat{\theta }^\alpha `$ in the right hand side of the fermionic field equation (10).
### 3.3 Gravitino vertex operator, a simple derivation of the ‘fermionic equation for bosonic brane’ and the Dirac action for the supermembrane fermionic field
The supersymmetry variation of the supergravity fields in the bosonic membrane action gives
$`\delta _\epsilon S_{M2}^0`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{W^3}}\widehat{}\widehat{e}_a\delta _\epsilon e^a{\displaystyle _{W^3}}\delta _\epsilon \widehat{A}_3=i{\displaystyle _{W^3}}\widehat{}\widehat{e}_a\widehat{\psi }^\beta \left(\mathrm{\Gamma }^a(1\overline{\gamma })\right)_{\beta \alpha }\epsilon ^\alpha (\widehat{x}),`$ (11)
Notice that the requirement that this variation is zero for an arbitrary value of the parameter $`\epsilon ^\alpha (\widehat{x}(\xi ))`$, $`\delta _\epsilon S_{M2}=0`$ results in
$`\delta _\epsilon S_{M2}`$ $`=`$ $`0\epsilon ^\alpha (\widehat{x})\widehat{}\widehat{e}_a\widehat{\psi }^\beta \left(\mathrm{\Gamma }^a(1\overline{\gamma })\right)_{\beta \alpha }=0,`$ (12)
which is exactly the ‘fermionic equation for the bosonic brane’, Eq. (45). This easy way to derive Eq. (45) is, actually, equivalent to a more involved derivation through the consistency conditions for the bulk field equations (see Sec. 2.3.2 and ). As a byproduct one also easily sees that a supersymmetry transformation with parameter restricted by (31), $`\widehat{\epsilon }^\alpha :=\widehat{\epsilon }^\alpha (\widehat{x}(\xi ))=(1+\overline{\gamma })^\alpha {}_{\beta }{}^{}\kappa _{}^{\beta }(\xi )`$, leaves the action $`S_{M2}^0`$ and hence $`S_{SG}+S_{M2}^0`$ invariant,
$`\delta _\epsilon S_{M2}^0|_{\widehat{\epsilon }=(1+\overline{\gamma })\kappa (\xi )}=0\delta _\epsilon (S_{SG}[e^a,\psi ^\alpha ,A_3]+S_{M2}[\widehat{e}^a,\widehat{A}_3])|_{\widehat{\epsilon }=(1+\overline{\gamma })\kappa (\xi )}=0,`$ (13)
which follows from the fact that $`(1\overline{\gamma })(1+\overline{\gamma })=0`$. This preserved supersymmetry, coming from the $`\kappa `$–symmetry of the superbrane, allows one to extend the identification of broken supersymmetry with physical fermionic degrees of freedom of supermembrane, $`(1\overline{\gamma })\widehat{\epsilon }=(1\overline{\gamma })\widehat{\theta }`$ (as in ), to a full identification of the pull–back to $`W^3`$ of supersymmetry parameter with the fermionic field,
$`\widehat{\epsilon }^\alpha =\epsilon ^\alpha (\widehat{x}(\xi ))=\widehat{\theta }^\alpha (\xi ).`$ (14)
With such an identification the expression for the interacting action becomes
$`S_{SGM2}=S_{SG}[e^a,\psi ^\alpha ,A_3]+S_{M2}=S_{SG}[e^a,\psi ^\alpha ,A_3]+S_{M2}^0[\widehat{e}^a,\widehat{A}_3]i{\displaystyle _{W^3}}\mathrm{\Xi }_{3\alpha }\widehat{\theta }^\alpha +𝒪(\widehat{\theta }^2),`$ (15)
where the first two terms in the l.h.s. describe the gauge–fixed action (27), while the third term (cf. (45)),
$`{\displaystyle _{W^3}}\mathrm{\Xi }_{1\alpha }(\xi )\widehat{\theta }^\alpha :=i\delta _{\widehat{\epsilon }=\widehat{\theta }}S_{M2},\mathrm{\Xi }_{3\alpha }(\xi ):=\widehat{}\widehat{e}_a\widehat{\psi }^\beta \left(\mathrm{\Gamma }^a(1\overline{\gamma })\right)_{\beta \alpha },`$ (16)
is given by the supersymmetry variation of $`S_{M2}^0`$ by substituting $`\widehat{\theta }`$ for $`\widehat{\epsilon }`$. This first order contribution determines the supermembrane fermionic vertex operator $`𝒱`$ as defined in ,
$`i\delta _{\widehat{\epsilon }=\widehat{\theta }}S_{M2}={\displaystyle _{W^3}}\mathrm{\Xi }_{3\alpha }(\xi )\widehat{\theta }^\alpha ={\displaystyle _{W^3}}d^3\xi \psi _\mu {}_{}{}^{\beta }(\widehat{x})𝒱_\beta {}_{}{}^{\mu }(\xi ),`$ (17)
$`d^3\xi 𝒱_\beta {}_{}{}^{\mu }(\xi )=d\widehat{x}^\mu (\xi )\widehat{}\widehat{e}_a\left(\mathrm{\Gamma }^a(1\overline{\gamma })\widehat{\theta }(\xi )\right)_\beta =d\widehat{x}^\mu (\widehat{}\widehat{e}_a(\mathrm{\Gamma }^a\widehat{\theta })_\beta i(\widehat{\overline{\mathrm{\Gamma }}}{}_{}{}^{(2)}\widehat{\theta })_\beta ).`$ (18)
Thus, starting from a full but gauge–fixed description of the supergravity–superbrane interaction of , we reproduce the supermembrane vertex operator from .
Notice that calculations as those above may apply equally well to the action (27) of the interacting system and to the action $`S_{M2}`$ of a bosonic membrane in a spacetime supergravity background. Of course in the latter case the local supersymmetry is not a gauge transformation of the action but rather a transformation of the background fields.
By construction, the action (15) is invariant under full local supersymmetry (not just one–half as the gauge–fixed action (27)) up to contributions proportional to $`\widehat{\theta }`$. Indeed, the Goldstone nature of $`\widehat{\theta }`$ implies $`\delta _\epsilon \widehat{\theta }(\xi )=\epsilon (\widehat{x})+𝒪(\widehat{\theta })`$ which, in the light of (17) and of the supersymmetry invariance of $`S_{SG}`$, gives $`\delta _{\widehat{\epsilon }}(S_{SG}+S_{M2})=\delta _{\widehat{\epsilon }}(S_{SG}+S_{M2}^0i\mathrm{\Xi }_3\widehat{\theta })=𝒪(\widehat{\theta })`$ for the action (15). To reach the supersymmetry invariance up to the first order in $`\widehat{\theta }`$ one needs to recover the $`𝒪(\widehat{\theta }^2)`$ components in the action. In our approach this can be done by adding $`\frac{1}{2}\delta _{\widehat{\epsilon }=\widehat{\theta }}(S_{SG}+S_{M2}^0i\mathrm{\Xi }_3\widehat{\theta })`$ to the action (15). This is just the term that should produce the supermembrane fermionic equation (4). One easily checks this for $`\psi =0`$. Indeed,
$`{\displaystyle \frac{1}{2}}\delta _{\widehat{\epsilon }=\widehat{\theta }}\left(S_{M2}^0i{\displaystyle \mathrm{\Xi }_3\widehat{\theta }}\right)_{|\widehat{\psi }=0}={\displaystyle \frac{i}{2}}{\displaystyle \delta _{\widehat{\epsilon }=\widehat{\theta }}\left(\mathrm{\Xi }_3\right)_{|\widehat{\psi }=0}\widehat{\theta }}={\displaystyle \frac{i}{2}}{\displaystyle \widehat{}\widehat{e}_a}𝒟\widehat{\theta }\mathrm{\Gamma }^a(1\overline{\gamma })\widehat{\theta }`$ (19)
produces the Dirac equation (5).
The action (15), linear in $`\widehat{\theta }`$, allows us to derive a supersymmetric set of interacting equations for the supergravity–supermembrane system with the same accuracy. For instance, the supersymmetric gravitino equation reads (cf. (35); notice that $`\mathrm{\Gamma }^a(1\overline{\gamma })\widehat{\theta }=\widehat{\theta }(1+\overline{\gamma })\mathrm{\Gamma }^a=\widehat{\theta }\mathrm{\Gamma }^a(1\overline{\gamma })`$)
$`\mathrm{\Psi }_{10\alpha }`$ $`:=`$ $`𝒟\psi ^\beta \overline{\mathrm{\Gamma }}_{\beta \alpha }^{(8)}=J_{10\alpha }[\widehat{\theta }]+𝒪(\widehat{\theta }^2),`$ (20)
$`J_{10\alpha }[\widehat{\theta }]`$ $`={\displaystyle \frac{i}{2e(x)}}e_b^{10}{\displaystyle _{W^3}}\widehat{}\widehat{e}_a\widehat{e}^b\left(\widehat{\theta }(\xi )\mathrm{\Gamma }^a(1\overline{\gamma })\right)_\alpha \delta ^{11}(x\widehat{x}(\xi )).`$ (21)
Now, removing the bulk fermion by inserting the ansatz (8) for $`\psi (x)`$ in (20) and ignoring higher order terms in $`\widehat{\theta }`$, one finds the relation between the bosonic currents (32), (33) and the fermionic current (21),
$`\stackrel{~}{\theta }(x)(iJ_{10a}\mathrm{\Gamma }_{\beta \alpha }^a{\displaystyle \frac{1}{2}}J_8\overline{\mathrm{\Gamma }}{}_{\beta \alpha }{}^{(2)})=J_{10\alpha }[\widehat{\theta }],`$ (22)
which is satisfied identically for a Goldstone fermion obeying (9). This shows that it is consistent to use the ansatz (8) to study particular solutions for the interacting system of supergravity and superbrane. Although this consistency is widely believed, the above is, to our knowledge, its first explicit check within the fully interacting system.
The study of the first order contribution in $`\widehat{\theta }`$ to the full system of interacting equations for the $`D`$=11 supergravity–supermembrane system, as well as for systems including M5–brane and $`D`$=10 Dirichlet superbranes is a problem for further study. Another interesting question is whether one can extend the present approach to include contributions of higher order in $`\widehat{\theta }`$ by using a counterpart of Noether method (see ) or, better still, the gauge completion procedure (see ) but applied to the action as a whole rather than to the construction of the supervielbein and other separate superfields.
## 4 Conclusions and discussion
In this paper we have shown how the Dirac equation (5) \[(4) for $`\psi 0`$\] for the fermionic coordinate field $`\widehat{\theta }(\xi )`$ of the supermembrane (see ) can be reproduced from a complete but gauge–fixed Lagrangian description of the D=11 supergravity–supermembrane interacting system . This component spacetime Lagrangian description is provided by the sum of the Cremmer–Julia–Scherk supergravity action and a bosonic brane action given by the purely bosonic ($`\widehat{\theta }=0`$) ‘limit’ of the supermembrane action . It preserves half of the local supersymmetry reflecting the $`\kappa `$–symmetry of the superbrane action. From the point of view of the hypothetical superfield action for the supergravity–supermembrane interacting system the above spacetime description appears as a result of fixing the superdiffeomorphism and superspace Lorentz symmetry by choosing the Wess–Zumino gauge for the supergravity superfields and of fixing (half of) the local supersymmetry by the $`\widehat{\theta }^\alpha (\xi )=0`$ gauge for the superbrane.
Formulated as a general prescription, our way of deriving the superbrane equations of motion consists in performing a spacetime local supersymmetry transformation \[$`\delta _\epsilon `$ of Eqs. (47)–(49)\] on the component fields that appear in the ‘fermionic equation for bosonic brane’ \[$`\widehat{\mathrm{\Xi }}_{3\alpha }=0`$, Eq. (45)\], and then identifying the (pull–back of the) parameter of this transformation with the superbrane fermionic field $`\widehat{\theta }(\xi )`$ \[thus $`\widehat{\mathrm{\Xi }}_{3\alpha }+\delta _{\widehat{\epsilon }=\widehat{\theta }}\widehat{\mathrm{\Xi }}_{3\alpha }=0`$\]. The identification of the $`\widehat{\theta }(\xi )`$ with the parameter of the supersymmetry ($`\widehat{\epsilon }=\widehat{\theta }`$) is made possible by the Goldstone nature of this superbrane fermionic field: its (non–pure gauge with respect to the $`\kappa `$–symmetry) components are the Goldstone fermions for the supersymmetries spontaneously broken by the superbrane .
The original ‘fermionic equation for the bosonic brane’ ( $`\widehat{\mathrm{\Xi }}_{3\alpha }=0`$, Eq. (45)) is obtained as a consistency condition for the bosonic and fermionic field equations of the gauge fixed description of the supergravity–superbrane interacting system which does not involve the superbrane fermionic $`\widehat{\theta }^\alpha (\xi )`$ variable explicitly<sup>7</sup><sup>7</sup>7Our approach makes particularly clear why the Dirac equation for the superbrane in a supergravity background with $`\widehat{\psi }=0`$ contains the same generalized covariant derivative ($`𝒟=Dt=d\omega t`$) involved in the gravitino supersymmetry transformation rules, a point also emphasized in Sec. 3 of a recent paper , where our $`𝒟\widehat{\theta }`$ is denoted by $`\delta \psi \theta `$. In the standard on-shell superfield approach such a coincidence can be traced to the fact that the component supersymmetry transformations may be deduced from the on–shell superspace constraints of .. Here, in Sec. 3.1, we have also shown how this ‘fermionic equation for the bosonic brane’ (45) can be obtained in an equivalent but very simple way, using as above the local supersymmetry transformation with $`\widehat{\epsilon }=\widehat{\theta }`$, but for the bosonic brane action. In this way one also recovers the gravitino vertex operator of . One may also notice that the ‘fermionic equation for bosonic brane’ formally coincides with the result of setting $`\widehat{\theta }^\alpha (\xi )=0`$ in the most general form of the superfield fermionic equations for superbranes in an on–shell superfield supergravity background, Eqs. (6). Namely the leading component of (6) gives (45) but with the graviton and the gravitino satisfying the ‘free’ supergravity equations of motion, which is not the case for Eq. (45) derived from the complete spacetime Lagrangian description. Moreover, this situation holds at least at first order in $`\widehat{\theta }`$ for the fermionic equations of motion and at second order in $`\widehat{\theta }`$ for the action, namely our equation for the supermembrane Goldstone fermion $`\widehat{\theta }(\xi )`$ also coincides (formally) with the equations derived in .
This shows, as widely believed, that the linearized equation for $`\widehat{\theta }(\xi )`$ derived from the standard on–shell superfield approach to the supergravity background is still valid for the case of background fields that are not restricted by the ‘free’ supergravity equations, in spite of the fact that the on–shell constraints implying these ‘free’ supergravity equations were an essential ingredient in the derivation of the Dirac equation within the usual background on-shell superfield approach. Notice that our results also fit with those of where it was found that, although the complete $`\kappa `$–symmetry of the supermembrane action (1) in curved superspace requires that the supervielbein $`E_M^A(Z)`$ and the super-3-form $`A_3(Z)`$ obey the on–shell supergravity constraints, the requirement of $`\kappa `$–symmetry up to the first order in $`\widehat{\theta }`$ for the action written up to the second order in $`\widehat{\theta }`$ does not impose any restrictions on the component background fields. Namely , if the on-shell supergravity constraints are used to decompose the action (1) in powers of $`\widehat{\theta }`$ neglecting $`𝒪(\widehat{\theta }^3)`$ terms and, then, the $`\kappa `$–symmetry is checked neglecting $`𝒪(\widehat{\theta }^2)`$ terms, the result is that, surprisingly, such a weakened $`\kappa `$–symmetry requirement does not restrict the background fields of the supergravity multiplet by any equations of motion. An important question is whether this is also the case for the decomposition of the standard supermembrane action including higher order $`𝒪(\widehat{\theta }^3)`$ terms in $`\widehat{\theta }`$, and, if so, whether such a decomposition would coincide with the action obtained by a development of the approach of the present paper.
Within the on–shell background superfield approach such calculations, also technically involved, are possible using the recent results of . To obtain equations of motion with higher order $`\widehat{\theta }(\xi )`$ terms in present approach one has to perform a ‘non–infinitesimal’ supersymmetry transformation up to some power in the parameter; the finite supersymmetry transformation, if found, might produce the fully supersymmetric (not gauge–fixed) action, if exists. For the existence of such finite transformation it is important that the local supersymmetry of the component gauge fixed description of the supergravity–supermembrane system is closed at least on shell. We have shown in Sec. 2.2.4 that this is indeed the case and that this follows from the closure of the local supersymmetry of free supergravity<sup>8</sup><sup>8</sup>8It would be interesting to study the algebra of the spacetime local supersymmetry of the $`D=11`$ supergravity interacting with M5–brane and of the $`D=10`$ supergravity interacting with higher Dirichlet branes. The (spacetime, gauge–fixed) Lagrangian description of such interactions implies the use of the duality–invariant formulations of supergravity (see for $`D=11`$, for $`D=10`$ type IIA and for $`D=10`$ type IIB) where the commutator of two supersymmetry transformations leaving invariant the supergravity action would involve the PST (Pasti–Sorokin–Tonin) gauge transformations.. A practical way to pursue the above proposed procedure method to find the action up to the terms of higher order in $`\widehat{\theta }(\xi )`$ is to use a counterpart of the gauge completion method (see ), but applied to the action itself. Namely, one makes an ‘infinitesimal’ supersymmetry transformation in the action written up to $`𝒪(\widehat{\theta }^k)`$ and recovers the next order in $`\widehat{\theta }(\xi )`$, $`𝒪(\widehat{\theta }^{(k+1)})`$, by identifying $`\widehat{\epsilon }=\widehat{\theta }(\xi )`$; then one tries making such an action supersymmetric up to order $`𝒪(\widehat{\theta }^k)`$ by modifying the supersymmetry transformation rules of the $`\widehat{x}^\mu (\xi )`$ and $`\widehat{\theta }^\alpha (\xi )`$ <sup>9</sup><sup>9</sup>9To lowest order $`\delta _\epsilon \widehat{x}^\mu (\xi )=i\widehat{\theta }\mathrm{\Gamma }^a\epsilon (\widehat{x})e_a{}_{}{}^{\mu }(\widehat{x})+𝒪(\widehat{\theta })`$, $`\delta _\epsilon \widehat{\theta }^\alpha (\xi )=\epsilon ^\alpha (\widehat{x}(\xi ))+𝒪(\widehat{\theta })`$.. Such a procedure would also answer the question of whether a fully supersymmetric (not gauge–fixed) interacting action $`S_{SG}(e^a(x),\psi ^\alpha (x),A_3(x))+S_{M2}(e^a(\widehat{x}),A_3(\widehat{x});\widehat{\theta }(\xi ),\psi ^\alpha (\widehat{x}))`$, with $`S_{M2}(\widehat{e}^a,A_3(\widehat{x});0,\psi ^\alpha (\widehat{x}))=S_{M2}^0(\widehat{e}^a,A_3(\widehat{x}))`$, exists formulated only in terms of the physical fields of the supergravity multiplet and the superbrane Goldstonions $`\widehat{x}(\xi )`$ and $`\widehat{\theta }(\xi )`$. As we discussed in this paper (and may gathered from the results of ), the answer to this question is affirmative up to second order in $`\widehat{\theta }(\xi )`$. Notice that, if an obstruction were found at some higher order in $`\widehat{\theta }`$, it would pose an interesting dilemma: whether such an obstruction is the result of a non-Lagrangian nature of the equations of motion for the physical fields of the supergravity multiplet in the interacting system, or whether it is the application of the above procedure to the equations of motion for the physical fields of the supergravity multiplet that fails. The second alternative would imply the impossibility of finding a fully supersymmetric system of equations for the physical fields of the supergravity multiplet and the superbrane Goldstone fields. Although at first glance this would look discouraging, it might also point towards some hidden ingredients of M-theory.
## Acknowledgments
The authors thank Dima Sorokin for several valuable discussions. We also wish to thank Eric Berghoeff, Sergio Ferrara, Toine Van Proeyen, G. Moore and W. Siegel for useful conversations at different stages of this work. This paper has been partially supported by the research grants BFM2002-03681 from the Ministerio de Educación y Ciencia and EU FEDER funds, N 383 of the Ukrainian State Fund for Fundamental Research, from Generalitat Valenciana and by the EU network MRTN–CT–2004–005104 ‘Constituents, Fundamental Forces and Symmetries of the Universe’. |
warning/0507/math-ph0507026.html | ar5iv | text | # On the geometrical thermodynamics of chemical reactions
## 1. Introduction
In the traditional approach to present the basic structure of homogeneous thermodynamics, it is customary to fix a set of variables describing the state of a system, the processes going on in the system and interactions with the outside world. Such a set usually includes the internal energy U of the thermodynamic system. A set of values of these functions form the extended state space of the system (referred to as the energy-phase space in $`[8]`$, and thermodynamic phase space in $`[11]`$). All but one of these variables, named thermodynamic potential in this representation, are set into couples ($`y`$,$`x`$) of intensive and extensive variables in such a way that to each extensive variable $`x`$ there corresponds an intensive variable, $`y`$, and the infinitesimal work (change of energy U or a chosen thermodynamic potential) related to the change in the extensive variable $`x`$ is,
$$dW=ydx$$
These couples are often collected into larger groups, corresponding to the tensorial type of the process they describe or to the process in whose description they participate. Collecting all such pairs, the first law of thermodynamics in its geometrical form postulates that during the process, the change in the internal energy, $`U`$, is given by integration over the trajectory in the state space of the one form,
$$dU=dQ+\underset{i}{}y_idx^i$$
Here $`dQ`$ is the heat change one form. The second law of thermodynamics in the formulation of C. Caratheodory states that the form,
$$dQ=dU\underset{i}{}y_idx^i$$
has an integrating factor (see also ). After some effort , this integrating factor is determined to be $`\frac{1}{T}`$ ($`T`$ is the absolute temperature). Thus,
$$dQ=TdS$$
where the new state variable S is the entropy.
The couple of variables ($`S`$,$`T`$), (entropy/temperature), plays a special role in this formulation. If one lists all the extensive variables $`x_i`$ (including entropy) and the corresponding intensive variables $`y_i`$, the infinitesimal change of the internal energy of the system is given by,
$$dU=\underset{i=1}{\overset{k}{}}y_idx^i$$
Thus, if the thermodynamic phase space in variables $`(U;x^i,y_i)`$ is denoted as $`P`$, there is a one-form $`\theta `$ defined by the choice of the process and the variables related to it, namely,
$$\theta =dU\underset{i=1}{\overset{k}{}}y_idx^i$$
(1.1)
Processes that might occur in the system should be such that, along the curve $`tr(t)=(U(t),x^i(t),y_i(t))`$, $`\theta (r^{}(t))=0`$. Thus, in this geometrical situation the 1-form $`\theta `$ defines the contact structure on $`P`$ and all the physically admissible processes should be integrable curves of the contact distribution $`D_p=ker(\theta _p)`$, with $`pP`$, of this structure .
### 1.1. Geometrical thermodynamics
Geometrical interpretations of equilibrium thermodynamics have proved that state space is endowed with a canonical contact structure that underlines the first law of thermodynamics. Different representations of this structure in a canonical D’Arbois-chart are related to different forms of the law of conservation of energy which can be expressed through the internal energy, entropy, Helmholtz free energy, or other extensive variables . Hermann and later Mrugala *argued that extended phase space of a homogeneous thermodynamic system, endowed with the contact structure, is the natural geometric space for descriptions of equilibrium thermodynamics*. Until now, explicit application of this geometric analysis to gain insight into, among other things, the critical behavior of real chemical systems has not been presented. Applications of this geometric approach to the analysis of simple thermodynamic systems is the focus of the present study.
To begin, R.Hermann and later R.Mrugala defined the extended state space of a homogeneous thermodynamic system as a (2k+1)-dimensional manifold $`P`$ endowed with the contact structure given by a differential $`1`$-form $`\theta `$ such that
$$\theta (d\theta )^k0$$
(1.2)
where $`\theta `$ is called the contact form. This condition is equivalent to the property of the smooth subbundle $`DT(P)`$ being as far from integrable as possible .
Locally, the association $`pker(\theta _p)`$ defines a 2k-dimensional distribution
$$pD_pT_p(P)$$
independent of the choice of $`\theta `$.
Moreover, it is possible to show that replacement of $`\theta `$ by $`f(p)\theta `$, with some function $`f(p)`$ positive at all points of domain of $`\theta `$, does not violate this condition. Any contact form $`\theta `$, in an appropriate local canonical chart $`𝒞`$ of variables $`(x^0,x^i,y_i,i=1,\mathrm{}k)`$, called the D’Arbois chart, is expressed by ,
$$\theta =dx^0\underset{i=1}{\overset{k}{}}y_idx^i.$$
(1.3)
In such a canonical chart, the distribution D at the point $`p`$ is generated by the vector fields,
$$D_p=<\frac{}{y_i},y_j\frac{}{x^0}+\frac{}{x^j}>$$
(1.4)
and the differential of the form $`\theta `$, $`d\theta `$, defines on each hyperplane, $`D_p`$, the symplectic structure $`(D_p,d\theta )`$. In local coordinates $`d\theta `$ has the canonical form,
$$d\theta _p=\underset{i=1}{\overset{k}{}}dx^idy_i$$
Replacement of $`\theta `$ by $`f(p)\theta `$ leads to the replacement of $`d\theta `$ by $`df\theta +fd\theta `$ which, after restriction to a hyperplane $`D_p`$, becomes $`fd\theta `$. Thus, contact structure alone defines only conformally symplectic structure on the distribution D.
On the manifold $`P`$ there exists a unique smooth vector field $`Y`$ called the Reeb vector field such that,
$$\theta (Y)=1\iota _Y(d\theta )=0$$
(1.5)
In particular one gets the canonical splitting,
$$T_p(P)=D_p\mathrm{ker}d\theta _p$$
(1.6)
of the tangent bundle of $`P`$ into the direct sum of two subbundles, the first being the subbundle of horizontal vectors of the distribution D, while the second being the characteristic subbundle of the form $`d\theta `$ . Correspondingly, the cotangent bundle $`T^{}(P)`$ splits as well.
In the canonical D’Arbois chart, the Reeb vector field is just,
$$Y=\frac{}{x^0}$$
Couples of variables $`(x^i,y_i)`$ denote pairs of independent parameters ($`x^i`$) and corresponding conjugate variables ($`y_i`$) with respect to the chosen thermodynamic potential $`x^0`$. Examples of $`(x^i,y_i)`$ pairs are: (1) temperature and entropy $`(T,S)`$; (2) pressure and volume $`(p,V)`$; (3) mole number of $`i`$-th component and corresponding chemical potential $`(N_i,\mu _i)`$; extent of reaction and corresponding affinity $`(\xi ,A)`$.
### 1.2. Thermodynamic Equilibrium
Thermodynamic equilibrium is a key notion in thermodynamics. In particular,
in all systems there is a tendency to evolve toward states in which the properties are determined by intrinsic factors and not by previously applied external influences. Such simple terminal states are, by definition, time independent. They are called equilibrium states .
In a state of thermodynamic equilibrium, intensive variables are functions of the extensive variables, namely
$$y_i=y_i(x^j)$$
(1.7)
Then, choosing the internal energy $`U`$ as the potential, the form $`\theta `$ becomes
$$\theta =dU\underset{i=1}{\overset{k}{}}y_i(x^j)dx^i$$
(1.8)
Along the contact distribution $`D=ker(\theta )`$,
$$dU=\underset{i}{}y_i(x^j)dx^i$$
(1.9)
Thus, the relation $`y_i(x^j)=\frac{U}{x^i}`$ exists just on the maximal integral submanifold of the contact manifold $`P`$.
Denoting $`\mathrm{\Phi }`$ as a generic thermodynamic potential, it is known that equilibrium states belong to a maximal integrable surface of contact form $`\theta `$ in the space P determined by the choice of k *independent variables* $`x^i`$ and by the thermodynamic potential $`\mathrm{\Phi }(x^i)`$ as a function of these variables (constitutive relations). Another choice of independent variables together with some other specification of $`\mathrm{\Phi }(x^i)`$ leads to another equilibrium surface corresponding in general to another constitutive relation. The core of the present study focuses on Legendre submanifolds defined to be maximal integral k-dimensional submanifolds of $`P`$ on which the Pfaff equation $`\theta =0`$ holds . The standard approach to locally defining such a submanifold, $`𝒮_\mathrm{\Phi }`$, in terms of a generating function, $`\mathrm{\Phi }`$, is given by the following theorem:
###### Theorem 1.
(V.Arnold, ).
For any partition $`IJ`$ of the set of indices (1,…,k) into two disjoint subsets I,J and for a function $`\mathrm{\Phi }(y_I,x^J)`$ of k variables $`y_i`$ with $`iI`$ and $`x^j`$ with $`jJ`$, the following equations,
$$x^0=\mathrm{\Phi }y_i\frac{\mathrm{\Phi }}{y_i},x^i=\frac{\mathrm{\Phi }}{y_i},y_j=\frac{\mathrm{\Phi }}{x^j}$$
(1.10)
define a Legendre submanifold $`𝒮_\mathrm{\Phi }`$ of a contact manifold $`(P^{2k+1},D)`$.
Conversely, every Legendre submanifold of $`(P^{2k+1},D)`$, in a neighborhood of any point, is defined by these equations for at least one of $`2^k`$ possible choices of the subset I.
In the special case in which $`\mathrm{\Phi }`$ is a function of only the independent variables ($`x^1`$,…,$`x^k`$), the Legendre submanifold (submanifold of equilibria states) $`𝒮_\mathrm{\Phi }`$ is given by,
$$𝒮_\mathrm{\Phi }=\{(\mathrm{\Phi },x^i,y_j,i=1,\mathrm{}k)P|\mathrm{\Phi }=\mathrm{\Phi }(x^i),y_i=\frac{\mathrm{\Phi }}{x^i},iI)$$
(1.11)
On the integral submanifold $`𝒮_\mathrm{\Phi }`$, the function $`\mathrm{\Phi }(x^i)`$ can be defined in terms of the variables $`x^i,y_i`$ as,
$$\mathrm{\Phi }(x^1,\mathrm{},x^k)=\underset{i}{}x^iy_i(x^1,\mathrm{},x^k)$$
(1.12)
In most cases, $`(1.12)`$ is homogeneous of order one and satisfies the Euler equation, i.e.
$$\mathrm{\Phi }(\lambda x^1,\mathrm{},\lambda x^k)=\lambda \mathrm{\Phi }(x^1,\mathrm{}x^k)$$
(1.13)
(see ). The expression in $`(1.12)`$ leads directly to the Gibbs-Duhem relation between variables along the integral submanifold $`𝒮_\mathrm{\Phi }`$. Indeed, taking the differential of $`\mathrm{\Phi }`$, we obtain
$$0=\theta |_{𝒮_\mathrm{\Phi }}=d\mathrm{\Phi }\underset{i}{}y_idx^i=\underset{i}{}x^idy_i$$
Thus, the basic contact condition $`\theta =0`$ is equivalent to the Gibbs-Duhem constitutive relation,i.e.
$$\underset{i}{}x^idy_i=0$$
In order to apply the contact formalism to equilibrium thermodynamics, all variables $`x^0`$, $`y_i`$, $`x^i`$, $`i=1,\mathrm{},k`$, must be identified with the thermodynamic parameters such that $`\theta =dx^0y_idx^i`$ satisfies the first law of thermodynamics. These parameters *only have physical meaning on the k-dimensional Legendre submanifold* defined by the Pfaff equation $`\theta =0`$. In this context, $`x^0=\mathrm{\Phi }`$ is a thermodynamic potential function of the independent variables $`x^1`$,…$`x^k`$; and $`y_1`$,…,$`y_k`$ are the corresponding conjugate parameters with respect to the potential. On the Legendre submanifold, the parameters $`y_i=\frac{\mathrm{\Phi }}{x^i}`$ .
One of the goals of the present study is, in the context of geometrical thermodynamics, to show how the choice of extensive variables and thermodynamic potential is crucial for a reasonable physical interpretation of geometrical objects.
In the present development, two sets of variables are chosen. In the case of a single component system, $`k=3`$ and $`x^0`$ is identified with the internal energy, while, $`x^1`$, $`x^2`$, $`x^3`$ are equated to the independent variables S, V, N, respectively. Likewise $`p_1`$, $`p_2`$, $`p_3`$ are the corresponding conjugate variables T, -p, $`\mu `$, respectively, and the contact form becomes $`\theta =dUTdS+pdV\mu dN`$. This is an extension of previous work with the addition of a more detailed exposition on the relation between geometry and thermodynamics.
In the case of an $`r`$-component system, $`k=r+2`$ and $`x^0`$ is identified with the Gibbs free energy; $`x^1`$, $`x^2`$, $`x^3`$,…,$`x^{r+2}`$ with the independent variables T, p, $`N_1`$,…,$`N_r`$; and $`p_1`$, $`p_2`$, $`p_3`$,…,$`p_{r+2}`$ with the conjugate variables -S, V, $`\mu _1`$,…,$`\mu _r`$, respectively. Then, the contact form becomes $`\theta =dG+SdTVdp\mu _1dN_1\mathrm{}\mu _rdN_r`$.
### 1.3. Thermodynamic Metric
A thermodynamic metric $`\eta _\mathrm{\Phi }`$ defined by the constitutive relation $`\mathrm{\Phi }=\mathrm{\Phi }(x^i)`$ on the Legendre submanifold $`𝒮_\mathrm{\Phi }`$ of the contact structure $`\theta `$ has the form(),
$$\eta _\mathrm{\Phi }=\underset{ij}{}\frac{^2\mathrm{\Phi }}{x^ix^j}dx^idx^j.$$
(1.14)
For the case where $`\mathrm{\Phi }`$ is the internal energy, $`U`$, the metric, $`\eta _U`$, is called the Weinhold metric(). When $`\mathrm{\Phi }`$ is the entropy, $`S`$, the metric, $`\eta _S`$, is called the Ruppeiner metric(). Here a new metric, $`\eta _G`$, is introduced for the case where $`\mathrm{\Phi }`$ is the Gibbs free energy, G.
###### Remark 1.
A thermodynamic metric, $`\eta _\mathrm{\Phi }`$, of the form $`(1.14)`$ is induced on $`𝒮_\mathrm{\Phi }`$ by the following symmetrical tensor(),
$$\stackrel{~}{\eta }=\frac{1}{2}\underset{i=1}{\overset{k}{}}(dy_idx^i+dx^idy^i).$$
(1.15)
Up to a conformal factor, this tensor is the symmetrical tensor in $`P`$ annihilating the Reeb vector field, $`Y`$, of structure $`\theta `$, ($`Y=\frac{}{x^0}`$), and is invariant under substitution of indices, $`i`$. $`\stackrel{~}{\eta }`$ is obtained as the sum of symmetrical tensors in the 2-dimensional subspaces $`D_i^{}`$ of $`D_x^{}`$ spanned by pairs of covectors $`(dx^i,dy_i)`$ of thermodynamic conjugate variables .
Moreover, thermodynamic metrics are generally degenerate and non-definite. In particular the Legendre submanifold or (as it is referred to hereafter) the thermodynamic state space, is the union of domains where these metrics have different signatures separated by the submanifold (generically of codimension one) of states where these metrics are degenerate .
### 1.4. Scalar Curvature of the Thermodynamic Metric
Now consider a thermodynamic potential, $`\mathrm{\Phi }=\mathrm{\Phi }(x^i)`$, a function of the independent variables, $`x^i`$, and calculate the scalar curvature of the corresponding thermodynamic metric defined on the Legendre submanifold $`𝒮_\mathrm{\Phi }`$. According to expression $`(1.14)`$,
$$\eta _{\mathrm{\Phi }ij}=\frac{^2\mathrm{\Phi }}{x^ix^j}.$$
(1.16)
Christoffel symbols for this metric are given by (see also ),
$$\mathrm{\Gamma }_{ij}^k=\frac{1}{2}\underset{m}{}\eta _{ij,m}\eta ^{km},$$
(1.17)
where $`\eta _{ij,m}=\frac{\eta _{ij}}{x^m}.`$
It can be shown that the curvature tensor of the metric $`\eta _\mathrm{\Phi }`$ is given by,
$$_{ijk}^l=\mathrm{\Gamma }_{ki,j}^l\mathrm{\Gamma }_{ji,k}^l+\mathrm{\Gamma }_{jp}^l\mathrm{\Gamma }_{ki}^p\mathrm{\Gamma }_{kp}^l\mathrm{\Gamma }_{ji}^p=\frac{1}{4}(\eta _{ij,m}\eta _{sn,k}\eta _{sn,j}\eta _{ki,m})\eta ^{mn}\eta ^{ls}$$
(1.18)
Therefore, the Ricci Tensor of metric $`\eta _\mathrm{\Phi }`$ is given by,
$$_{ik}=_{ijk}^j=\frac{1}{4}(\eta _{ij,m}\eta _{sn,k}\eta _{sn,j}\eta _{ki,m})\eta ^{mn}\eta ^{js},$$
(1.19)
and the scalar curvature $`_{\eta _\mathrm{\Phi }}`$ by (see also ),
$$_{\eta _\mathrm{\Phi }}=_{ik}\eta ^{ik}=\frac{1}{4}(\eta _{ij,m}\eta _{sn,k}\eta _{sn,j}\eta _{ki,m})\eta ^{mn}\eta ^{js}\eta ^{ik}.$$
(1.20)
###### Remark 2.
Consider the scalar curvature of the metric defined on two-dimensional integral surfaces $`𝒮_\mathrm{\Phi }`$. Components of the Ricci tensor are given by(),
$$_{11}=\frac{1}{4}(((\eta _{21,_1})^2\eta _{11,_1}\eta _{21,_2})\eta ^{11}\eta ^{22}+(\eta _{21,_1}\eta _{21,_2}\eta _{11,_1}\eta _{22,_2})\eta ^{12}\eta ^{22}+$$
$$((\eta _{21,_2})^2\eta _{11,_2}\eta _{22,_2})(\eta ^{22})^2)$$
(1.21)
$$_{12}=_{21}=\frac{1}{4}((\eta _{11,_1}\eta _{21,_2}(\eta _{21,_1})^2)\eta ^{11}\eta ^{12}+(\eta _{11,_1}\eta _{22,_2}\eta _{21,_2}\eta _{11,_2})(\eta ^{12})^2+$$
$$(\eta _{11,_2}\eta _{22,_2}(\eta _{21,_2})^2)\eta ^{12}\eta ^{22})$$
(1.22)
$$_{22}=\frac{1}{4}((\eta _{21,_1})^2\eta _{22,_1}\eta _{11,_1})(\eta ^{11})^2+(\eta _{21,_2}\eta _{11,_2}\eta _{11,_1}\eta _{22,_2})\eta ^{11}\eta ^{12}+$$
$$((\eta _{21,_2})^2\eta _{12,_1}\eta _{22,_2})\eta ^{11}\eta ^{22})$$
(1.23)
It is straight forward to show that,
$$_{11}\eta ^{11}=_{22}\eta ^{22}$$
(1.24)
Thus, the scalar curvature is given by(),
$$=2(_{11}\eta ^{11}+_{12}\eta ^{12})$$
(1.25)
These general expressions for curvature are employed to characterize the thermodynamics of several systems.
### 1.5. Closed and Open Thermodynamic Systems
Consider the thermodynamic state of a system as a function of a certain number of independent variables such as the entropy, S, volume, V, and number of moles, $`N_1`$,…,$`N_r`$, of r components. Any function expressed in terms of these variables is a state function of the system. The First Law of Thermodynamics or Conservation of Energy postulates the existence of a state function, called the energy function, such that the change in internal energy of the universe, given as the sum of the energies of our system and of the surroundings, is always constant, namely(),
$$0=dU_{univ.}=dU_{sys.}+dU_{surr.}$$
(1.26)
This expression states that,
$$dU_{sys.}=dU_{surr.}$$
(1.27)
The minus sign indicates a loss of energy by the surroundings. Denoting $`d_EU=dU_{surr.}`$ as the change in energy supplied to the system by the surroundings, Conservation of Energy dictates,
$$dU_{sys.}=dU=d_EUord_IU=0$$
(1.28)
Subscript $`I`$ indicates the energy change of the system. Note, $`d_IU=0`$, is equivalent to $`dU_{univ.}=0`$.
The Second Law of Thermodynamics or the principle of Entropy Production postulates the existence of a state function, called the entropy function, which possesses the following properties: the entropy is an extensive variable, and the change in entropy $`dS`$ can be separated into the flow of entropy, $`d_ES`$, due to interactions with surroundings and a term, $`d_IS`$, corresponding to entropy changes in the system (). That is,
$$dS=d_ES+d_IS$$
(1.29)
where $`d_IS`$ is denoted as the entropy production. $`d_IS`$ is always non-negative, zero for reversible processes and positive for irreversible ones.
For closed systems, conservation of energy in $`(1.28)`$ can be expressed as,
$$dU=dQ+dW_M=dQpdV$$
(1.30)
with pressure, $`p`$, normal to the surface. For open systems,
$$dU=d\mathrm{\Psi }+dW_M=d\mathrm{\Psi }pdV$$
(1.31)
where $`d\mathrm{\Psi }`$ is the infinitesimal rate of change of heat transfer and exchange of matter of the system with the external environment . These expressions are employed to examine the geometries of three different thermodynamic systems. First, work done on single component thermodynamic systems is reviewed while stressing the importance of choosing the right thermodynamic variables suitable for a particular situation. For a one component thermodynamic system physical interpretations are deduced from geometrical objects such as the degeneracy and scalar curvature of the Weinhold metric. Choosing the Gibbs free energy as the preferable thermodynamic potential certain physical aspects of chemical behaviour can be described through the geometry. This approach is also used in studying chemical reactions in multicomponent systems.
### 1.6. System 1: Single Component Closed System
Consider a *closed* system containing a single component in the absence of an external field (). The energy supplied by the surroundings is derived from the sum of the heat transfer, $`dQ`$, and mechanical work, $`dW_M`$. In this case the entropy production $`d_IS=0`$, and the entropy of the system is given by,
$$dS=d_ES=\frac{dQ}{T}$$
(1.32)
Therefore, in molar form the $`1`$-form of the energy, $`du`$, is given by,
$$du=dQ+dW_M=Tdspdv$$
(1.33)
### 1.7. System 2: Multi-Component Closed System
Consider a multi-component system in which changes in internal energy can occur due to chemical reactions. In this case, the entropy production $`d_IS`$ is given by ,
$$d_IS=\frac{A}{T}d\xi $$
where A is the affinity of the chemical reaction related to the chemical potentials $`\mu _i`$ by $`A=_i\mu _i\nu _i`$. The $`\nu _i`$ are the stoichiometric coefficients and $`\xi `$ is the extent of reaction. Note, for a single chemical reaction the entropy of the system is given by,
$$dS=d_ES+d_IS=\frac{dQ}{T}+\frac{A}{T}d\xi $$
(1.34)
Now introduce another state function, the Gibbs free energy, G, defined by,
$$G=UTS+pV$$
In terms of this function, conservation of energy can be written as ,
$$dG=dQTdSSdT+Vdp=SdT+VdpAd\xi $$
(1.35)
While this expression is central to the present study, useful geometrical tools are also introduced allowing a more general treatment of the case of $`l`$ independent chemical reactions. In this context, $`(1.34)`$ and $`(1.35)`$ can be restated as,
$$dS=\frac{dQ}{T}+\underset{n}{}\frac{A_n}{T}d\xi _n,$$
(1.36)
and
$$dG=SdT+Vdp\underset{n}{}A_nd\xi _n$$
(1.37)
where $`n=1,\mathrm{},l`$ is the number of chemical reactions that occur.
### 1.8. System 3: Open Systems
Consider an open system in the absence of an external field without entropy production. The expression in $`(1.29)`$ can be generalized to consider changes in the number of moles $`N_1`$,…,$`N_r`$ , viz.
$$dS=d_ES=\frac{d\mathrm{\Psi }}{T}\underset{i}{}\frac{\mu _i}{T}dN_i$$
(1.38)
where $`d\mathrm{\Psi }=TdS+_i\mu _idN_i`$ is the energy flow due to heat transfer and exchange of matter. The corresponding Gibbs free energy is given by,
$$dG=SdT+Vdp+\underset{i}{}\mu _idN_i$$
(1.39)
with $`i=1,\mathrm{},r`$, the number of moles of each reaction component.
## 2. System 1: Single component system
Consider a $`7`$-dimensional thermodynamic phase space P of variables
$`(U,(S,T),(V,p),(N,\mu ))`$ with the contact $`1`$-form given by,
$$\theta =dUTdS+pdV\mu dN$$
(2.1)
where $`N`$ is the number of moles of the component and $`\mu `$ its chemical potential. Next, consider a $`3`$-dimensional Legendre submanifold $`𝒮_U(S,V,N)`$ of this system defined by the constitutive relation,
$$U=U(S,V,N)$$
(2.2)
By homogeneity of degree one of the internal energy, consider the molar form of the constitutive relation $`(2.2)`$ and obtain the following constitutive relation for a closed system with a single component ,
$$u=u(s,v)$$
(2.3)
The differential is given by $`(1.33)`$, i.e.
$$du=Tdspdv$$
Introduce the following thermodynamic parameters:
1. $`C_V`$ is the heat capacity at constant volume:
$$C_V=T(\frac{S}{T})_V,$$
(2.4)
2. $`C_p`$ is the heat capacity at constant pressure:
$$C_p=T(\frac{S}{T})_p,$$
(2.5)
3. $`\alpha `$ is the thermal coefficient of expansion:
$$\alpha =\frac{1}{V}(\frac{V}{T})_p$$
(2.6)
4. $`k_T`$ is the isothermal compressibility:
$$k_T=\frac{1}{V}(\frac{V}{p})_T$$
(2.7)
For the molar case the above parameters are represented in lower case type. The following expression relates the above molar parameters:
$$c_pc_v=vT\frac{\alpha ^2}{k_T}$$
(2.8)
Considering the expression in $`(1.14)`$, the Weinhold metric $`\eta _u`$ is defined on the two-dimensional integral surface $`𝒮_u`$ by ,
$$\eta _{ij_u}=\frac{^2u}{x^ix^j}=\frac{1}{c_v}\left(\begin{array}{cc}T& \frac{T\alpha }{k_T}\\ \frac{T\alpha }{k_T}& \frac{c_p}{vk_T}\end{array}\right)$$
(2.9)
It follows, if $`T0`$, the Weinhold metric is degenerate along the curve $`\gamma _\eta `$ given by,
$$(\frac{p}{v})_T=0$$
(2.10)
which is presented in one of two forms: $`p=p(v)`$ or $`T=T(v)`$ .
The main points of focus are as follows :
1) The equilibrium surface is the union of regions where the Weinhold metric has different signature separated by the curve $`\gamma _\eta `$ where the metric is degenerate;
2) The critical point $`(p_c,T_c,v_c)`$ of the system is the extremum of the functions $`p=p(v)`$ and $`T=T(v)`$;
3) Along the curve of degeneracy, $`\gamma _\eta `$, a *first order phase transition* seems to occur;
4) Scalar curvature of the Weinhold metric is *strongly* influenced by parameters related to non-ideal inter-particle interactions within the system.
###### Remark 3.
For a two dimensional state space, the determinant $`det(\eta _\mathrm{\Phi })`$ and the scalar curvature $``$ of a thermodynamic metric $`\eta _\mathrm{\Phi }`$ are inversely related ,
$$=\frac{1}{4det(\eta _\mathrm{\Phi })^2}det\left(\begin{array}{ccc}\eta _{11}& \eta _{11,1}& \eta _{11,2}\\ \eta _{12}& \eta _{12,1}& \eta _{12,2}\\ \eta _{22}& \eta _{22,1}& \eta _{22,2}\end{array}\right).$$
(2.11)
In point 4) above, it was noted that the scalar curvature of the Weinhold metric is strongly related to interparticle interactions in the system. Moreover, as the determinant of the matrix approaches zero, point 3) implies that the system approaches a phase transition at which point the scalar curvature $``$ of the metric goes to infinity. Thus, the connection of points 3) and 4) suggests the following interpretation:
5) If the system approaches a state ”close enough” to the curve of degeneracy, $`\gamma _\eta `$, the scalar curvature of the metric goes to infinity. Physically this is consistent with a relevant increase in inter-particle interactions between the reactant and product species when the system approaches a phase transition.
Such behavior suggests an intriguing relationship between degeneracy, scalar curvature and inter-particle interactions. In particular, this suggests a *geometrical condition* for a phase transition might be degeneracy (or infinite curvature) of the Weinhold metric $`\eta _{ij_u}`$. The following examples support this suggestion.
Example I: Ideal gas.
Given the equation of state for an Ideal gas, $`pv=RT`$, the Weinhold metric, $`\eta _{ij_u}`$, is given by ,
$$\eta _{ij_u}=\frac{1}{c_v}\left(\begin{array}{cc}T& p\\ p& \frac{c_pp}{v}\end{array}\right)$$
(2.12)
In this case $`c_pc_v=R`$, and ,
$$det(\eta _{ij_u})=\frac{p}{vc_v^2}(Tc_ppv)=\frac{pT}{vc_v}=\frac{RT^2}{c_vv^2}>0$$
(2.13)
The above metric is positive definite on the costitutive surface $`𝒮_u`$. Thus, for an Ideal gas, the energy metric is never degenerate except for trivial cases, i.e. $`T=0`$. This lack of degeneracy is consistent with the characteristics of an Ideal Gas (i.e. it does not display a critical point and therefore does not exhibit a phase transition). Moreover scalar curvature of the Weinhold metric is zero, , i.e.
$$_{\eta _{ij_u}}=0$$
(2.14)
Ruppeiner was the first to suggest that zero curvature might be evidence for the absence of inter-particle interactions. Their absence is precisely the case for an ideal gas. Table $`1`$ gives parallels between the geometric and thermodynamic features of an Ideal Gas.
Example II: van der Waals gas
The equation of state for the van der Waals gas is given by,
$$(p+\frac{a}{v^2})(vb)=RT$$
(2.15)
where a and b are positive constants. Expression $`(2.15)`$ provides a more realistic representation of the actual behavior of real (non-ideal) gases by introducing the additional positive constants a and b, characteristic of the particular gas under consideration. The factor $`(vb)`$ indicates the excluded volume of the molecules, while the factor $`\frac{a}{v^2}`$ is the ”interaction” term.
For the van der Waals gas, the Weinhold metric is given by ,
$$\eta _{ij_u}=\left(\begin{array}{cc}\frac{T}{c_v}& \frac{TR}{(vb)c_v}\\ \frac{TR}{(vb)c_v}& (\frac{TR}{(vb)^2}(1+\frac{R}{c_v})\frac{2a}{v^3})\end{array}\right)$$
(2.16)
which is degenerate along the curve $`\gamma _\eta `$ written in the following forms:
$$s=s(v)=c_v[(2+\frac{R}{c_v})\mathrm{ln}(vb)+\mathrm{ln}(\frac{2ac_v}{Rv^3})]$$
$$p(v)=(v2b)\frac{a}{v^3},T(v)=\frac{2a(vb)^2}{Rv^3}$$
Note, in the limit $`a=b=0`$, the ideal case is recovered with the metric in $`(2.12)`$. Taking the derivatives of the last two expressions $`p=p(v)`$ and $`T=T(v)`$ and setting them to zero, the critical point is obtained,
$$(p_c,T_c,v_c)=(\frac{a}{27b^2},\frac{8a}{27bR},3b)$$
(2.17)
Moreover, for the van der Waals gas the scalar curvature of the Weinhold metric $`_{\eta _{ij_u}}`$ is given by ,
$$_{\eta _{ij_u}}=\frac{aRv^3}{c_v(pv^3av+2ab)^2}$$
(2.18)
In general, the scalar curvature $`_{\eta _{ij_u}}0`$ as $`a0`$ and as the system approaches the degeneracy curve, $`_{\eta _{ij_u}}\mathrm{}`$. On the other hand, expression $`(2.18)`$ does not vanish if the parameter $`b0`$. So while the scalar curvature of the Weinhold metric is not related to the excluded volume, it is strongly influenced by the parameter $`a`$, that includes non-ideal interactions in the system. This finding is entirely consistent with the physical behavior of the van der Waals gas which exhibits both a critical point and phase transition. Indeed in the $`(pT)`$ plane, the following solutions are obtained ,
$$p_r^i(T)=\frac{3v_r^i(T)2}{(v_r^i(T))^3}i=1,\mathrm{},3$$
One of the these solutions is the Pressure-Temperature Phase Boundary (Fig.$`1`$).
Table $`2`$ summarizes parallels between geometric and thermodynamic features of the Van der Waals gas.
Example III: Berthelot gas
As a third example consider the Berthelot gas with the equation of state,
$$(p+\frac{a}{Tv^2})(vb)=RT$$
(2.19)
where a and b are positive constants (analogous to the van der Waals gas).
The Weinhold metric for the Berthelot gas is given by ,
$$\eta _{ij_u}=\left(\begin{array}{cc}\frac{T}{c_v}& \frac{1}{c_v}(\frac{RT}{(vb)}+\frac{a}{Tv^2})\\ \frac{1}{c_v}(\frac{RT}{(vb)}+\frac{a}{Tv^2})& \frac{RT^2v^32a(vb)^2}{Tv^3(vb)^2}+\frac{T}{c_v}(\frac{R}{(vb)}+\frac{a}{T^2v^2})^2\end{array}\right)$$
(2.20)
whose degeneracy is ,
$$RT^2v^32a(vb)^2=2p^2v^3(vb)^2aR(v2b)^2=0$$
(2.21)
In analogy to what was done for the van der Waals gas, it follows that
$$(v_c,p_c,T_c)=(3b,\pm (\frac{aR}{216b^3})^{\frac{1}{2}},\pm (\frac{8a}{27Rb})^{\frac{1}{2}})$$
and that the scalar curvature of the Weinhold metric, $`_{\eta _{ij_u}}`$, is given by ,
$$_{\eta _{ij_u}}=2a\frac{\left(T^4v^4Rc_vL(c_v,v)+T^2v^3RaQ(c_v,v)+a^2W(c_v,v)\right)}{c_v^3T^3v(RT^2v^32a(vb)^2)^2},$$
(2.22)
where
$$L(c_v,v)=(2c_vR)v^23c_vbv+c_vb^2$$
$$Q(c_v,v)=Rv^5+3Rbv^43Rb^2v^3+(Rb^3+c_v+R)v^2b(b2v)(R+c_v)$$
and
$$W(c_v,v)=Rv^7+4Rbv^66Rb^2v^5+(2c_v+R+4Rb^3)v^4(8c_v+3R+Rb^3)bv^3+(12c_v+3R)b^2v^2$$
$$(8c_v+R)b^3v+2c_vb^4$$
The scalar curvature, $`_{\eta _{ij_u}}`$, goes to zero as $`a0`$, and is strongly influenced by this parameter which corresponds to non-ideal interactions in the system. Furthermore, $`_{\eta _{ij_u}}\mathrm{}`$ as the system approaches a phase transition. Once again, degeneracy of the Weinhold metric and non-zero scalar curvature are consistent with the characteristic physical behavior of the Berthelot Gas. Parallels displayed in Table 2 for the van der Waals gas are equally applicable to the Berthelot gas.
## 3. System 2: chemical reactions in a closed system
### 3.1. Single chemical reaction
Consider a closed system comprised of r components among which chemical reactions can occur. First, we focus our attention on single chemical reactions and then introduce multi-component reactions. In a closed system, any change in the masses of the components will occur only from a chemical reaction. Thus, denoting the mass of component i by $`m_i`$, with $`i=1,\mathrm{},r`$, the infinitesimal change in mass can be written as ,
$$dm_i=\nu _iM_id\xi $$
(3.1)
where $`M_i`$ is the molar mass of component $`i`$. The principle of conservation of mass for a closed system is expressed as ,
$$dm=\underset{i}{}\nu _iM_id\xi =0$$
(3.2)
with $`m=_im_i`$. The equation $`_i\nu _iM_i=0`$ is referred to as the stoichiometric equation.
Alternatively, rather than the component masses it is more convenient to consider the number of moles $`N_1`$,…,$`N_r`$ involved in the reaction. Since $`\frac{dm_i}{M_i}=dN_i`$, the infinitesimal change in the mole number of the $`i`$ component, can be expressed as
$$dN_i=\nu _id\xi $$
Let $`N_i^0`$ be the number of moles of component $`i`$ in the initial state of the system. When a reaction occurs, as indicated by the stoichiometric coefficients $`\nu _i`$, the variations of the number of moles of each component $`N_i`$ are not independent. This can be expressed as ,
$$\frac{dN_1}{\nu _1}=\mathrm{}=\frac{dN_i}{\nu _i}=\mathrm{}=d\xi $$
(3.3)
where the extent of the reaction, $`\xi `$, is an extensive variable just like the number of moles. Integrating and taking $`\xi =0`$ as the initial state of the system, we obtain ,
$$N_i=N_i^0+\nu _i\xi i=1,\mathrm{},r$$
In this context, the Legendre submanifold $`𝒮_G`$ can be defined by the constitutive relation, $`G=G(T,p,\xi )`$. Restriction of the contact $`1`$-form, $`\theta =dG+SdTVdp+Ad\xi `$ to the submanifold $`𝒮_G`$ provides,
$$dG=SdT+VdpAd\xi $$
(3.4)
Thus, the general metric $`\eta _{ij_G}=\frac{^2G}{x^ix^j}`$, where $`x^i`$ and $`x^j`$ are the extensive variables, is given by,
$$\eta _{ij_G}=\left(\begin{array}{ccc}\frac{C_p}{T}& \alpha V& \mathrm{\Delta }_rS\\ \alpha V& k_TV& \mathrm{\Delta }_rV\\ \mathrm{\Delta }_rS& \mathrm{\Delta }_rV& (\frac{A}{\xi })_{T,p}\end{array}\right)$$
(3.5)
where $`A=\mathrm{\Delta }_rG`$ is the Gibbs free energy of the reaction and $`\mathrm{\Delta }_rS`$ and $`\mathrm{\Delta }_rV`$ are the entropy and volume of reaction, respectively. Here, the affinity and Gibbs free energy of reaction are used interchangebly.
Naturally, as a reaction takes place, the chemical potential of the components varies and so does the affinity of the reaction. At constant temperature and pressure, the system is at equilibrium whenever the affinity $`A=0`$. Since $`(\frac{A}{\xi })_{T,p}0`$, the determinant of the matrix in Eqn.$`(3.5)`$ is given by,
$$det\eta _{ij_G}=\frac{C_vk_TV}{T}(\frac{A}{\xi })_{T,p}+\frac{C_v}{T}(\mathrm{\Delta }_rV)^2+Vk_T[\mathrm{\Delta }_rS\frac{\alpha }{k_T}\mathrm{\Delta }_rV]^20$$
(3.6)
Of primary interest is what type of information is provided by the degeneracy and, in some simple cases, by the scalar curvature of the metric $`\eta _{ij_G}`$. As an example consider the three-dimensional case of the Ideal gas mixture.
Example I: Ideal gas mixture.
Consider a simple reaction in which substance A converts into substance B. Starting with one mole of A, the relative amounts at some later point in the reaction are $`N_A=1\xi `$ and $`N_B=\xi `$ with $`\xi [0,1]`$. Therefore, the Gibbs free energy can be written in terms of the extent of reaction as,
$$G=(1\xi )\mu _A+\xi \mu _B$$
Now, the chemical potential of an ideal gas mixture is given by $`\mu _i=\mu _i^\theta (T)+RT\mathrm{ln}(\frac{p_i}{p^\theta })`$ where $`i=A,B`$ and the superscript $`\theta `$ indicates some standard state at pressure $`p^\theta `$. This reaction resulting in conversion of the ideal component A into the ideal component B is analogous to a homogeneous mixture of two ideal components, in that the two components are mixed but do not interact.
Considering that $`p_A=(1\xi )p`$ and $`p_B=\xi p`$, where $`p`$ is the total pressure, the Gibbs free energy can be written as,
$$G=(1\xi )\mu _A^\theta +\xi \mu _B^\theta +RT\mathrm{ln}(\frac{p}{p^\theta })+RT[(1\xi )\mathrm{ln}(1\xi )+\xi \mathrm{ln}\xi ]$$
where $`RT[(1\xi )\mathrm{ln}(1\xi )+\xi \mathrm{ln}\xi ]=\mathrm{}G_{mix}`$ is the Gibbs free-energy of mixing.
Thus, the metric $`(3.5)`$ becomes,
$$\eta _{ij_G}=\left(\begin{array}{ccc}(1\xi )\frac{d^2\mu _A^\theta }{dT^2}+\xi \frac{d^2\mu _B^\theta }{dT^2}& \frac{R}{p}& \frac{d\mu _B^\theta }{dT}\frac{d\mu _A^\theta }{dT}+R\mathrm{ln}(\frac{\xi }{1\xi })\\ \frac{R}{p}& \frac{RT}{p^2}& 0\\ \frac{d\mu _B^\theta }{dT}\frac{d\mu _A^\theta }{dT}+R\mathrm{ln}(\frac{\xi }{1\xi })& 0& \frac{RT}{\xi (1\xi )}\end{array}\right)$$
(3.7)
where $`\mathrm{\Delta }_rS=\frac{d\mu _B^\theta }{dT}\frac{d\mu _A^\theta }{dT}+R\mathrm{ln}(\frac{\xi }{1\xi })`$, $`\mathrm{\Delta }_rV=0`$ and $`(\frac{A}{\xi })_{T,p}=\frac{RT}{\xi (1\xi )}`$. The determinant in expression $`(3.6)`$ reduces to,
$$det\eta _{ij_G}=\frac{RT}{p^2}([\frac{d}{dT}(\mu _B^\theta \mu _A^\theta )+R\mathrm{ln}(\frac{\xi }{1\xi })]^2\frac{RT}{\xi (1\xi )}[\frac{R}{T}+(1\xi )\frac{d^2\mu _A^\theta }{dT^2}+\xi \frac{d^2\mu _B^\theta }{dT^2}])$$
(3.8)
If the chemical potentials of the two ideal components in the standard state are explicitly known, useful general information could be extrapolated from the Gibbs metric, its degeneracy and its scalar curvature. Since in general this is not the case, our analysis is restricted to the $`2`$-dimensional isothermal case and to the $`1`$-dimensional isothermal-isobaric case.
3.2 Isothermal single chemical reaction.
When the temperature is kept constant during the chemical reaction, $`dT=0`$ and the expression in Eqn. $`(3.4)`$ reduces to,
$$dG=VdpAd\xi $$
(3.9)
Thus, the metric $`(3.5)`$ reduces to,
$$\eta _{ij_G}=\left(\begin{array}{cc}k_TV& \mathrm{\Delta }_rV\\ \mathrm{\Delta }_rV& (\frac{A}{\xi })_p\end{array}\right)$$
(3.10)
with
$$det\eta _{ij_G}=k_TV(\frac{A}{\xi })_p(\mathrm{\Delta }_rV)^2$$
(3.11)
In the case of an ideal gas mixture, the metric of Eqn.$`(3.10)`$ becomes,
$$\eta _{ij_G}=\left(\begin{array}{cc}\frac{RT}{p^2}& 0\\ 0& \frac{RT}{\xi (1\xi )}\end{array}\right)$$
(3.12)
with determinant $`det\eta _{ij_G}=\frac{R^2T^2}{p^2\xi (1\xi )}`$ which is always different than zero (except for trivial values of some thermodynamic parameters). This implies that the Gibbs metric is, in general, never degenerate for an Ideal mixture. Thus, the physical interpretation of this result is that there is no critical behavior displayed by an Ideal Gas Mixture. Moreover, the scalar curvature of the metric $`(3.12)`$ is zero. Indeed, since the inverse of $`\eta _{ij_G}`$ is given by,
$$\eta _G^{ij}=\left(\begin{array}{cc}\frac{p^2}{RT}& 0\\ 0& \frac{\xi (1\xi )}{RT}\end{array}\right)$$
(3.13)
the third derivatives of the Gibbs potential are given by,
$$\eta _{11,_1}=\frac{2RT}{p^3},\eta _{11,_2}=\eta _{12,_1}=\eta _{21,_1}=0$$
(3.14)
$$\eta _{22,_1}=\eta _{12,_2}=\eta _{21,_2}=0,\eta _{22,_2}=\frac{RT(2\xi 1)}{\xi ^2(1\xi )^2}$$
(3.15)
Therefore, the components of the Ricci tensor $`_{ij}`$, $`(1.21)`$ to $`(1.23)`$, for an isothermal ideal mixture are all zero, namely
$$_{ij}=0i,j=1,2.$$
Using the expression in Eqn.$`(1.25)`$,
$$_{\eta _{ij_G}}=0$$
(3.16)
Obviously, this result is consistent with the fact that the two *Ideal* components when mixed do not interact, and it is essentialy consistent with the features of the single-component Ideal case. This strongly suggests that even in the context of chemical reactions in closed systems, non-zero scalar curvature might provide useful information regarding interactions between components. Although beyond the scope of the present study, this is an interesting path to pursue and, due to the complexity of the system, will require the use of numerical mathematics.
3.3 Isothermal-isobaric single chemical reaction
It is interesting to note that in the case of constant temperature and pressure, the change in Gibbs free energy is given by,
$$dG=Ad\xi $$
(3.17)
In this one-dimensional case, important information can be gleaned from examination of the convexity of the Gibbs free energy function. Consider the condition,
$$\frac{d^2G}{d\xi ^2}=\frac{d\mathrm{\Delta }_rG}{d\xi }=0$$
(3.18)
For an ideal gas mixture ,
$$\frac{d^2G}{d\xi ^2}=\frac{RT}{\xi (1\xi )}>0,\xi [0,1],T0$$
(3.19)
which implies that the Gibbs free energy is a convex function of the extent of reaction. An example of such a function is displayed in Fig.$`2`$, (see ), and defines the condition of stability. Initially, the Gibbs free energy decreases. As a reaction proceeds, the Gibbs free energy of the system continues to decrease until it reaches a minimum value. At equilibrium (constant T and p), the Gibbs energy is at the minimum, and $`\mathrm{\Delta }_rG`$ is equal to zero. At greater extents of reaction the Gibbs free energy is greater than zero and increases.
This well-known result suggests for a non-ideal mixture the change in sign and vanishing of the second derivative of the Gibbs function might provide some insight into the stability of the system. In particular, consider the case in which the Gibbs free energy is not a simple convex function of the extent of the reaction everywhere, but rather, first convex, then concave, then convex again, as shown in Fig $`3`$. Points on the curve where changes between convex and concave behavior occur are inflection points which satisfy the expression in Eqn.$`(3.18)`$. Naturally, a shift of the equilibrium state from one local minimum to another constitutes a first order phase transition induced by a change in the extent of reaction . For a chemical reaction, the system tends to approach chemical equilibrium where the forward and reverse reaction rates are the same and concentrations of the reactant and product species do not change with time. When the equilibrium condition is achieved, proportions of the various compounds remain unchanged, and the reaction ceases to progress. Prior to reaching the point of equilibrium, the system fluctuates between different equilibrium microstates. Suppose that the system is confined in a lower (more stable) Gibbs free-energy minimum and, occasionally, a fluctuation may be large enough to push the system over the maximum to the region of higher energy, i.e. a matastable minimum. A small fluctuation can overcome the shallow barrier back to the more stable equilibrium state . Any thermodynamic system, in this case a chemical reaction, tends to eventually reach the lowest minimum in the Gibbs free energy. Naturally, if the ”unstable” barrier is too high or the minima are far apart a shift of the equilibrium from one local minimum state to another is less probable.
Within this picture, the local curvature of the Gibbs free energy is positive for all points except those between the two inflection points. Moreover, the portion of the curve between the minima at the inflection points is said to be locally stable but globally unstable. In this region on the curve metastable states occur which appear to be stable to small perturbations, but mixed configurations at the same extent of reaction represent more stable states with lower free energy. A straight line connecting the two minima corresponds to a phase boundary , i.e. a phase transition from the phase at one minimum to the phase at the other minimum. Positive local curvature fails at the points of inflection. Local stability determines whether, after a small perturbation, a system will return to the original equilibrium state. Here our focus is on examining conditions leading to failure of local stability.
###### Remark 4.
Ideal Mixture. For an isothermal and isobaric ideal mixture,
$$\frac{d\mathrm{\Delta }_rG}{d\xi }=\frac{RT}{\xi (1\xi )}$$
(3.20)
When $`T0`$, consider that
$$\frac{d\xi }{d\mathrm{\Delta }_rG}=\frac{\xi (1\xi )}{RT}=k\xi (1\xi ),k=\frac{1}{RT}>0$$
(3.21)
which is the so-called logistic equation . The process described by this equation has two equilibrium positions, namely $`\xi =0`$ and $`\xi =1`$. Between these two points the field is directed from $`0`$ to $`1`$. As a result the equilibrium position $`\xi =0`$ is unstable (as soon as the reaction proceeds away from $`\xi =0`$ reactants are converted to products). Meanwhile the equilibrium position $`\xi =1`$ is stable. Moreover, integral curves tend asymptotically to the line $`\xi =1`$ as $`\mathrm{\Delta }_rG+\mathrm{}`$ and to the line $`\xi =0`$ as $`\mathrm{\Delta }_rG\mathrm{}`$. Such curves describe the passage from one state (0) to another (1) in an infinite $`\mathrm{\Delta }_rG`$.
### 3.2. Compounds
Consider a generic chemical reaction written as,
$$|\nu _A|A+|\nu _B|B+\mathrm{}|\nu _S|S+|\nu _T|T+\mathrm{}$$
where components on the left side are designated *reactants* while components on the right hand side are *products* . The $`|\nu _i|`$, $`i=A,B,\mathrm{}`$ are stoichiometric coefficients of the reaction. Another formal representation of a chemical reaction which better lends itself to mathematical manipulation is given by ,
$$\nu _SS+\nu _TT+\mathrm{}+\nu _AA+\nu _BB=\underset{i}{}\nu _ii=0$$
(3.22)
The chemical potential of each component is given by ,
$$\mu _i=\mu _i^{}(T,p)+RT\mathrm{ln}\gamma _ix_i=\mu _i^{}(T,p)+RT\mathrm{ln}a_ii=A,B,\mathrm{}$$
(3.23)
where $`a_i`$ is the activity of component $`i`$. Recall that $`x_i=\frac{N_i}{_jN_j}`$, with $`N_i=N_i^0+\nu _i\xi `$ and that $`\gamma _i`$ depends on the extent of reaction. Thus, the Gibbs free-energy of the reaction is given by ,
$$\mathrm{\Delta }_rG=\underset{i}{}\mu _i\nu _i=\underset{i}{}\nu _i\mu _i^{}+RT\mathrm{ln}\underset{i}{}a_i^{\nu _i}=\underset{i}{}\nu _i\mu _i^{}+RT\mathrm{ln}\underset{i}{}\gamma _i^{\nu _i}x_i^{\nu _i}$$
(3.24)
Introducing the quotient of reaction, $`Q_a=_ia_i^{\nu _i}=_i\gamma _i^{\nu _i}_ix_i^{\nu _i}=Q_\gamma Q_c`$, where the subscript $`c`$ denotes concentration, the expression $`(3.24)`$ can be written as ,
$$\mathrm{\Delta }_rG=\mathrm{\Delta }_rG^\theta +RT\mathrm{ln}Q_a$$
(3.25)
where $`\mathrm{\Delta }_rG^\theta =_i\nu _i\mu _i^{}(T,p^\theta )`$ is the standard Gibbs free energy of reaction. It follows that,
$$Q_a(\xi )=e^{\frac{\mathrm{\Delta }_rG(\xi )\mathrm{\Delta }_rG^\theta }{RT}}$$
(3.26)
Naturally, if the system reaches equilibrium, namely $`\mathrm{\Delta }_rG=0`$, the parameter $`Q_a`$ is denoted by $`K_a`$, the equilibrium constant, and the expression in $`(3.26)`$ becomes ,
$$Q_a^{eq}=K_a=e^{\frac{\mathrm{\Delta }_rG^\theta }{RT}}$$
(3.27)
Thus, for an isothermal and isobaric single chemical reaction,
$$\frac{d\mathrm{\Delta }_rG}{d\xi }=RT\frac{d\mathrm{ln}Q_a}{d\xi }=RT[\frac{d\mathrm{ln}Q_c}{d\xi }+\frac{d\mathrm{ln}Q_\gamma }{d\xi }]$$
$$=RT[\underset{i}{}(\frac{\nu _i^2}{N_i})\frac{(_i\nu _i)^2}{_iN_i}+\frac{d}{d\xi }\mathrm{ln}(\underset{i}{}\gamma _i^{\nu _i})]$$
(3.28)
where
$$\frac{d\mathrm{ln}Q_c}{d\xi }=\underset{i}{}(\frac{\nu _i^2}{N_i})\frac{(_i\nu _i)^2}{_iN_i}$$
(3.29)
For simplicity, denote $`\frac{d}{d\xi }\mathrm{ln}(_i\gamma _i^{\nu _i})=W(\xi )`$. Then, the expression in $`(3.28)`$ can be rewritten as,
$$\frac{d^2G}{d\xi ^2}=RT[\underset{i}{}(\frac{\nu _i^2}{N_i})\frac{(_i\nu _i)^2}{_iN_i}+W(\xi )]$$
(3.30)
The expression in $`(3.29)`$ denotes the influence of the relative amounts of reactants and products at each extent of the reaction while $`W(\xi )`$ represents the relative strength of non-ideal (inter-particle) interactions existent between products and reactants. Consequently, at any value of the extent of reaction, there is a ”possible” value of W such that the two mentioned forces exactly balance one another. Thus, at a given certain extent of the reaction determined by the relative amounts of reactants and products, $`W`$ at that point corresponds to the relative strength of non-ideal interactions that must exist between products and reactants for a failure of local stability.
###### Remark 5.
For the ideal gas mixture, $`Q_a=Q_c=\frac{\xi }{1\xi }`$ and therefore, $`\frac{dQ_a}{d\xi }=\frac{1}{(1\xi )^2}>0`$ (see Fig. $`4`$). In this case, $`Q_\gamma =1`$ and the expression in $`(3.30)`$ reduces to,
$$\frac{d^2G}{d\xi ^2}=RT[\underset{i}{}(\frac{\nu _i^2}{N_i})\frac{(_i\nu _i)^2}{_iN_i}]$$
(3.31)
Moreover, when the sum of the stoichiometric coefficients vanishes (i.e the isothermal-isobaric Ideal gas mixture, see $`(3.19)`$), the expression in $`(3.30)`$ reduces further to,
$$\frac{d^2G}{d\xi ^2}=RT\underset{i}{}(\frac{\nu _i^2}{N_i})>0$$
(3.32)
###### Theorem 2.
Let $`T0`$. Then, for an isobaric and isothermal single chemical reaction, $`\frac{d^2G}{d\xi ^2}=0`$ if and only if
$$W(\xi )=\frac{(_i\nu _i)^2}{_iN_i(\xi )}\underset{i}{}\frac{\nu _i^2}{N_i(\xi )}$$
(3.33)
The Gibbs free energy, $`G`$, is a convex function of the extent of reaction whenever,
$$W(\xi )>\frac{(_i\nu _i)^2}{_iN_i(\xi )}\underset{i}{}\frac{\nu _i^2}{N_i(\xi )}$$
(3.34)
and a concave function whenever,
$$W(\xi )<\frac{(_i\nu _i)^2}{_iN_i(\xi )}\underset{i}{}\frac{\nu _i^2}{N_i(\xi )}$$
(3.35)
The curve described by the expression in $`(3.33)`$ is denoted as the curve of phase boundary in the $`W`$-$`\xi `$ plane. Such a curve traces the phase boundary between the convex and concave regions of the Gibbs free energy. In particular, at a fixed value of the extent of reaction, the system is locally stable whenever the value of $`W`$ is less (in absolute value) than the value on the curve of phase boundary. If instead the value of $`W`$ is greater, the system is locally unstable.
In early stages of the reaction, reactant species are far in excess of product species, and $`W`$ is (in absolute value) relatively large. Since a change from reactant phase to product phase is improbable early in the reaction, interactions between products, favoring product formation, must be greater than those between reactants, favoring the reactant phase. $`W`$ indicates the balance of the strengths of the product and reactant interactions required for failure of local stability. As the reaction progresses toward the extremum of the curve of phase boundary, $`\frac{dW}{d\xi }=0`$, the relative difference in strength between the two types of interactions is a minimum. At this critical extent of reaction, a change between reactant and product phases requires the smallest difference between their constituent interactions and is thus most probable. Past this critical point the extent of reaction increases. To achieve local instability the relative strength of the interactions favoring reactants must be increasingly greater than those favoring the products.
Example II: Synthesis Reaction
Consider a simple synthesis reaction in which two or more substances combine to form a more complex substance. For example, $`2`$ moles of di-hydrogen react with $`1`$ mole of oxygen to give $`2`$ moles of water,
$$2H_2+O_22H_2O$$
(3.36)
with
$$N_{H_2}=22\xi N_{O_2}=1\xi andN_{H_2O}=2\xi $$
and the corresponding stoichiometric coefficients given by,
$$\nu _{H_2}=2\nu _{O_2}=1and\nu _{H_2O}=2$$
Then, from $`(3.33)`$,
$$W(\xi )=\frac{6}{\xi (1\xi )(3\xi )}<0$$
(3.37)
A plot of the curve of phase boundary for this reaction is given in Fig.5.
The extremum of this curve defines the ”most probable” transition point between the reactant and product phases. At this point differences in the relative amounts of reactant and product species, and relative differences in the strengths of their non-ideal interactions, are minimal. Taking the derivative with respect to $`\xi `$ (Fig.$`6`$) and setting it to zero yields,
$$\xi =0.4514$$
It follows that $`W=9.5`$. The fact that the critical extent of the reaction is less than 0.5 suggests that the products and associated non-ideal interactions are more strongly favored such that the product phase is preferred even before half the extent of reaction is reached.
Now in analogy, consider the dissociation reaction, i.e. the synthesis reaction in the opposite direction. In this case $`2`$ moles of water split into $`2`$ moles of hydrogen and $`1`$ mole of oxygen. Namely,
$$2H_2O2H_2+O_2$$
(3.38)
Following the same steps as for analysis of the synthesis reaction, the following expression for $`W`$ is obtained,
$$W(\xi )=\frac{6}{\xi (1\xi )(2+\xi )}<0$$
(3.39)
The graph of this curve of phase boundary is shown in Fig.7.
Taking the derivative with respect to $`\xi `$ and setting it to zero (see Fig.$`8`$),
$$\xi =0.5486$$
Note, $`W=9.5`$ is apparently ”invariant” to the direction of the chemical reaction. The critical extent of reaction, $`\xi `$ = 0.5486 for the dissociation reaction, consistent with inferences drawn from analysis of the synthesis reaction. That is, interactions between the synthetic species ($`H_2O`$) are more favorable than those between the individual species ($`H_2`$, $`O_2`$). The sum of the two critical extents of reactions is $`0.4514+0.5486=1`$.
Example III: Single Displacement Reaction
A single displacement reaction is one in which an atom (or ion) of a single compound replaces an atom of another compound. As an example, consider the single displacement in which copper ions in a copper sulfate solution are displaced by zinc, forming zinc sulfate:
$$Zn+CuSO_4Cu+ZnSO_4$$
(3.40)
with,
$$N_{Zn}=1\xi N_{CuSO_4}=1\xi N_{Cu}=\xi andN_{ZnSO_4}=\xi $$
The corresponding stoichiometric coefficients are given by,
$$\nu _{Zn}=1\nu _{CuSO_4}=1\nu _{Cu}=1and\nu _{ZnSO_4}=1$$
Then, the expression in $`(3.33)`$ is given by,
$$W(\xi )=\frac{2}{\xi (1\xi )}<0$$
(3.41)
The graph of this curve of phase boundary for this reaction is shown in Fig.9.
Taking the derivative with respect to $`\xi `$ and setting it to zero (see Fig.$`10`$),
$$\xi =0.5$$
It follows that $`W=8`$. In this case $`\xi =0.5`$ corresponds to the minimum difference between interactions associated with the product and reactant species for a phase transition to occur. Note that, for the displacement reaction, the critical extent of the reaction is $`0.5`$. In this case, the number of products species equals the number of reactants species and the critical extent of the reaction is independent of the direction of the process. This is in contrast to what was found for obtain in the synthesis and dissociation reactions where product and reactant species are not equal and the critical extent of the reaction depends on the reaction direction.
### 3.3. Multi-chemical reaction
Here the Gibbs metric is introduced for a closed system with r chemical species in which $`l`$ independent reactions can occur. The total change of mass, $`dm_i`$, is equal to the sum of the changes resulting from the different reactions, ,
$$dm_i=M_i\underset{n}{}\nu _{in}d\xi _nn=1,..,l$$
(3.42)
The principle of Conservation of Mass can be stated as ,
$$dm=\underset{i}{}\underset{n}{}\nu _{in}M_id\xi _n=0$$
(3.43)
As done previously, consider the number of moles of the components in the system instead of their masses. The change in number of moles of component $`i`$ is given by,
$$dN_i=\underset{n}{}\nu _{in}d\xi _n$$
The Legendre submanifold $`𝒮_G`$ is defined by the constitutive relation $`G=G(T,p,\xi _1,\mathrm{}\xi _l)`$ and the expression in Eqn.$`(3.3)`$ becomes,
$$\frac{dN_{1,n}}{\nu _{1,n}}=\mathrm{}=\frac{dN_{i,n}}{\nu _{i,n}}=\mathrm{}=d\xi _n$$
where $`\xi _n`$ is the extent of the n-th reaction with $`i=1,\mathrm{},r`$ and $`n=1,\mathrm{},l`$, . Now, consider the differential of the Gibbs free-energy,
$$dG=SdT+VdP\underset{n}{}A_nd\xi _n$$
(3.44)
By recalling the expression in $`(1.16)`$, the following metric is obtained,
$$\eta _{ij_G}=\left(\begin{array}{ccccc}\frac{C_p}{T}& \alpha V& \mathrm{\Delta }_rS_1& \mathrm{}& \mathrm{\Delta }_rS_r\\ \alpha V& k_TV& \mathrm{\Delta }_rV_1& \mathrm{}& \mathrm{\Delta }_rV_r\\ \mathrm{\Delta }_rS_1& \mathrm{\Delta }_rV_1& (\frac{A_1}{\xi _1})_{T,p,\xi _2,\mathrm{}\xi _l}& \mathrm{}& (\frac{A_l}{\xi _1})_{T,p,\xi _2,\mathrm{}\xi _l}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \mathrm{\Delta }_rS_r& \mathrm{\Delta }_rV_r& (\frac{A_l}{\xi _1})_{T,p,\xi _2,\mathrm{}\xi _l}& \mathrm{}& (\frac{A_l}{\xi _l})_{T,p,\xi _1,\mathrm{}\xi _{l1}}\end{array}\right)$$
(3.45)
This is the Gibbs metric for a multicomponent thermodynamic system in which $`l`$ independent chemical reactions occur.
## 4. System 3: Open systems
Finally, consider open systems in the absence of external fields. Recall that the contact $`1`$-form $`\theta `$, restricted to the Legendre submanifold $`𝒮_G`$ described by the constitutive relation $`G=G(T,p,N_1,\mathrm{},N_r)`$, gives the following differential, see $`(1.39)`$,
$$dG=SdT+Vdp+\underset{i}{}\mu _idN_i$$
(4.1)
Define a partial molar quantity as ,
$$\overline{X_i}=(\frac{X}{N_i})_{T,p,N_{ji}}$$
(4.2)
Since the Gibbs metric is defined as the Hessian of the Gibbs potential, $`\eta _{ij_G}=\frac{^2G}{x^ix^j}`$, where $`x^i`$ and $`x^j`$ are the extensive variables, the following result is obtained,
$$\eta _{ij_G}=\left(\begin{array}{ccccc}\frac{C_p}{T}& \alpha V& \overline{S_1}& \mathrm{}& \overline{S_r}\\ \alpha V& k_TV& \overline{V_1}& \mathrm{}& \overline{V_r}\\ \overline{S_1}& \overline{V_1}& \overline{\mu _{11}}& \mathrm{}& \overline{\mu _{r1}}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \overline{S_r}& \overline{V_r}& \overline{\mu _{r1}}& \mathrm{}& \overline{\mu _{rr}}\end{array}\right)$$
(4.3)
where $`\overline{\mu _{ik}}=(\frac{\mu _i}{N_k})_{T,p,N_{jk}}`$.
At constant temperature and pressure, conservation of energy $`(4.1)`$ reduces to,
$$dG=\underset{i}{}\mu _idN_i$$
(4.4)
Therefore, the metric in Eqn.$`(4.3)`$ becomes,
$$\eta _{ij_G}=\left(\begin{array}{cccc}\overline{\mu _{11}}& \overline{\mu _{21}}& \mathrm{}& \overline{\mu _{r1}}\\ \overline{\mu _{21}}& \overline{\mu _{22}}& \mathrm{}& \overline{\mu _{r2}}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \overline{\mu _{r1}}& \overline{\mu _{r2}}& \mathrm{}& \overline{\mu _{rr}}\end{array}\right)$$
(4.5)
### 4.1. Ideal solutions
Now, introduce the Gibbs metric for the cases of an ideal and non-ideal solution in an open thermodynamic system without chemical reactions. A component $`i`$ in solution is said to be ideal when its chemical potential is given by, ,
$$\mu _i=\mu _i^{}(T,p)+RT\mathrm{ln}x_i$$
(4.6)
where $`\mu _i^{}(T,p)`$ is the standard state chemical potential which is independent of the composition and $`x_i=\frac{N_i}{_jN_j}=\frac{N_i}{N}`$ is the mole fraction. Note, the total number of moles, N, depends on the number of moles of the r components, namely $`N=N(N_1,\mathrm{},N_r)`$.
Using the expression $`(4.6)`$, the metric for an ideal solution that depends on the temperature and the pressure of the system is obtained,
$$\eta _{ij_G}=\left(\begin{array}{ccccc}\frac{C_p}{T}& \alpha V& (\frac{\mu _1^{}}{T})_{p,N_i}+R\mathrm{ln}x_1& \mathrm{}& (\frac{\mu _r^{}}{T})_{p,N_i}+R\mathrm{ln}x_r\\ \alpha V& k_TV& (\frac{\mu _1^{}}{p})_{T,N_i}& \mathrm{}& (\frac{\mu _r^{}}{p})_{T,N_i}\\ (\frac{\mu _1^{}}{T})_{p,N_i}+R\mathrm{ln}x_1& (\frac{\mu _1^{}}{p})_{T,N_i}& RT(\frac{1}{N_1}\frac{1}{N})& \mathrm{}& \frac{RT}{N}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ (\frac{\mu _r^{}}{T})_{p,N_i}+R\mathrm{ln}x_r& (\frac{\mu _r^{}}{p})_{T,N_i}& \frac{RT}{N}& \mathrm{}& RT(\frac{1}{N_r}\frac{1}{N})\end{array}\right)$$
(4.7)
For an isothermal-isobaric system such a metric becomes,
$$\eta _{ij_G}=\left(\begin{array}{cccc}RT(\frac{1}{N_1}\frac{1}{N})& \frac{RT}{N}& \mathrm{}& \frac{RT}{N}\\ \frac{RT}{N}& RT(\frac{1}{N_2}\frac{1}{N})& \mathrm{}& \frac{RT}{N}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \frac{RT}{N}& \frac{RT}{N}& \mathrm{}& RT(\frac{1}{N_r}\frac{1}{N})\end{array}\right)$$
(4.8)
This metric is degenerate everywhere since $`det(\eta _{ij_G})=0`$ for all values of the number of moles of the r components. This result indicates the Gibbs free energy metric for an open system does not distinguish differences between ideal components. As far as the metric is concerned, due to the explicit lack of inter-component interactions, the system is apparently comprised of a single, indistinguishable, ideal component.
### 4.2. Non-Ideal solutions
For non-ideal systems, the chemical potential of component $`i`$ takes the form, ,
$$\mu _i=\mu _i^{}(T,p)+RT\mathrm{ln}\gamma _ix_i$$
(4.9)
where $`\gamma _i`$ is the activity coefficient of species $`i`$ that depends on temperature, pressure and composition of the solution. Expression $`(4.9)`$ considers deviations from ideality. The partial molar entropy, $`\overline{S_i}`$, and the partial molar volume of component $`i`$, $`\overline{V_i}`$, are given by,
$$\overline{S_i}=(\frac{\mu _i^{}}{T})_{p,N_i}+R[\mathrm{ln}\gamma _ix_i+T(\frac{\mathrm{ln}\gamma _i}{T})_{p,N_i}]$$
(4.10)
and
$$\overline{V_i}=(\frac{\mu _i^{}}{p})_{T,N_i}+RT(\frac{\mathrm{ln}\gamma _i}{p})_{T,N_i}$$
(4.11)
Moreover,
$$\overline{\mu _{ii}}=(\frac{\mu _i}{N_i})_{T,p,N_{ji}}=RT[(\frac{1}{N_i}\frac{1}{N})+(\frac{\mathrm{ln}\gamma _i}{N_i})_{T,p,N_{ji}}]$$
(4.12)
and,
$$\overline{\mu _{ik}}=(\frac{\mu _i}{N_k})_{T,p,N_{jk}}=RT[\frac{1}{N}+(\frac{\mathrm{ln}\gamma _i}{N_k})_{T,p,N_{jk}}]$$
(4.13)
For symplicity, denote $`(\frac{\mathrm{ln}\gamma _i}{N_k})_{T,p,N_{jk}}=\overline{(\mathrm{ln}\gamma _i)}_k`$ and retain the usual notation for partial molar entropy and volume. Substituting the expressions $`(4.10)`$-$`(4.13)`$ into the metric $`(4.3)`$, the general metric of a multicomponent non-ideal solution is obtained,
$$\eta _{ij_G}=\left(\begin{array}{ccccc}\frac{C_p}{T}& \alpha V& \overline{S_1}& \mathrm{}& \overline{S_r}\\ \alpha V& k_TV& \overline{V_1}& \mathrm{}& \overline{V_r}\\ \overline{S_1}& \overline{V_1}& RT[(\frac{1}{N_1}\frac{1}{N})+\overline{(\mathrm{ln}\gamma _1)}_1]& \mathrm{}& RT[\frac{1}{N}+\overline{(\mathrm{ln}\gamma _r)}_1]\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \overline{S_r}& \overline{V_r}& RT[\frac{1}{N}+\overline{(\mathrm{ln}\gamma _r)}_1]& \mathrm{}& RT[(\frac{1}{N_r}\frac{1}{N})+\overline{(\mathrm{ln}\gamma _r)}_r]\end{array}\right)$$
(4.14)
and, in the case of an isobaric-isothermal non-ideal system, the Gibbs metric simplifies to,
$$\eta _{ij_G}=RT\left(\begin{array}{cccc}[(\frac{1}{N_1}\frac{1}{N})+\overline{(\mathrm{ln}\gamma _1)}_1]& [\frac{1}{N}+\overline{(\mathrm{ln}\gamma _2)}_1]& \mathrm{}& [\frac{1}{N}+\overline{(\mathrm{ln}\gamma _r)}_1]\\ [\frac{1}{N}+\overline{(\mathrm{ln}\gamma _2)}_1]& [(\frac{1}{N_2}\frac{1}{N})+\overline{(\mathrm{ln}\gamma _2)}_2]& \mathrm{}& [\frac{1}{N}+\overline{(\mathrm{ln}\gamma _r)}_2]\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ [\frac{1}{N}+\overline{(\mathrm{ln}\gamma _r)}_1]& [\frac{1}{N}+\overline{(\mathrm{ln}\gamma _r)}_2]& \mathrm{}& [(\frac{1}{N_r}\frac{1}{N})+\overline{(\mathrm{ln}\gamma _r)}_r]\end{array}\right)$$
(4.15)
$$=RT[\left(\begin{array}{cccc}(\frac{1}{N_1}\frac{1}{N})& \frac{1}{N}& \mathrm{}& \frac{1}{N}\\ \frac{1}{N}& (\frac{1}{N_2}\frac{1}{N})& \mathrm{}& \frac{1}{N}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \frac{1}{N}& \frac{1}{N}& \mathrm{}& (\frac{1}{N_r}\frac{1}{N})\end{array}\right)+\left(\begin{array}{cccc}\overline{(\mathrm{ln}\gamma _1)}_1& \overline{(\mathrm{ln}\gamma _2)}_1& \mathrm{}& \overline{(\mathrm{ln}\gamma _r)}_1\\ \overline{(\mathrm{ln}\gamma _2)}_1& \overline{(\mathrm{ln}\gamma _2)}_2& \mathrm{}& \overline{(\mathrm{ln}\gamma _r)}_2\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \overline{(\mathrm{ln}\gamma _r)}_1& \overline{(\mathrm{ln}\gamma _r)}_2& \mathrm{}& \overline{(\mathrm{ln}\gamma _r)}_r\end{array}\right)]$$
(4.16)
$$=(\eta _{ij_G})_{ideal}+(\eta _{ij_G})_{deviation}$$
(4.17)
where the subscript ”deviation” indicates the degree of non-ideality.
## 5. Conclusions
Geometrical thermodynamics has been applied to study the behavior of several simple thermodynamic systems. For a single component, closed system intriguing relationships between the geometrical concepts of degeneracy and scalar curvature of the Weinhold metric, and the physical concepts of a phase transition and inter-particle, non-ideal interactions are divulged. For multi-component closed systems the Gibbs metric was presented and the scalar curvature and degeneracy of this metric was determined for a few simple cases and could be indicative of physical behavior. In summary, this study provides convincing examples that this geometrical approach to analysis of thermodynamic systems can be applied to actually divulge important physical and chemical behavior. |
warning/0507/cond-mat0507159.html | ar5iv | text | # Conservation of spin current
## Abstract
The conventional definition of spin-current, namely spin density multiplied by the group velocity, is not a conserved quantity due to possible spin rotations caused by spin-orbit (SO) interaction. However, in a model with spin-spin interactions, rotation of a spin causes a dynamic response of surrounding spins that opposes the rotation. Such a many-body effect restores the spin-current conservation. Here we prove that the non-conservation problem of spin-current can be resolved if a self-consistent spin-spin interaction is included in the analysis. We further derive a spin-conductance formula which partitions spin-current into different leads of a multi-lead conductor.
Recently, considerable interest has been paid to the quantum physics of spin-currentsarma-review . It is believed that a controlled spin-current generation, detection, and usage can provide interesting applications to spintronics. Spin-current generation has been classified as “extrinsic” or “intrinsic”. An extrinsic spin-current is generated by external physical factors and driving forces of the spintronics device, such as optical spin injection achieved experimentallysipe and the various spin pumps studied theoreticallysun-prl . An intrinsic spin-current is generated by physical factors existed inside the spintronic device, notable is that generated by various spin-orbit interactionsdas ; zhang ; sinova ; kato ; shen . In particular, it has been theoretically predicted that non-magnetic systems with spin-orbit interaction and under an external electric field, can generate a spin-current flowing perpendicular to the electric fieldzhang . Such an spin-current is termed “dissipation-less” because the electron motion is perpendicular to the electric field. There are so far extensive theoretical work on spin-current physicshalperin ; ref7 , and some experimental works have also appearedkato ; marcus1 which may provide support to some of the theories.
Despite the increasing literature on spin-current physics, it is recognized that the definition of spin-current itself is still somewhat controversial. If one mimics the definition of charge-current, then a spin-current $`I_s`$ can be defined as the time derivative of spin density $`N_s`$, $`I_s<dN_s/dt>`$, here the bracket is the quantum average. In its simplest form at steady state, such a definition gives a spin-current $`I_s=\frac{\mathrm{}}{2}(I_{}I_{})`$, where $`I_{},I_{}`$ are the charge current with spin-up and spin-down, respectively. Clearly, this is a very intuitive definition of spin-current and is adopted by most of the work in literature. Since both spin and velocity are vectors, the spin-current is a tensor. In systems where there is a spin-orbit interaction, the spin density is not conserved: spins can rotate from their initial orientation due to the interaction. Therefore, spin-current $`I_s`$ becomes a non-conservative quantity. A quantity which is not conserved is difficult to study experimentally and indeed, it is unclear what is even measured if an experimental detection method can be found to measure spin-current. Without knowing what is measured, the definition of spin-current becomes non-unique and there have been several definitions in literature.
Although it is unclear a priori if a measurable spin-current must be conserved, the property of conservation would be nice to have, at least theoretically. Such an issue has been discussed in a recent paper of Sun and Xiesun1 who defined an additional rotational contribution to spin-current on top of the above conventional definition, and the total spin-current is then conserved. However, there is no microscopic theory concerning this problem. In the present work, we further investigate this problem by using the conventional definition of spin-current $`I_s<dN_s/dt>`$, but we include a spin-spin interaction into the Hamiltonian for computing the quantum average $`<\mathrm{}>`$. Our basic idea is the following. The existing problem of spin-current non-conservation was due to SO interaction which rotates the spin away from its original direction. In a spin system with spin-spin interactions, the rotation of one spin causes a dynamic response of surrounding spins through the interaction, and this response opposes the original spin rotation. This many-body effect introduces a new term in the spin-current that balances the spin-current non-conservation. In the end, one obtains a final spin-current in the interacting model (which has a different value than that of non-interacting model) that is conserved. In the following, we demonstrate this idea by proving that $`I_s`$ obtained in a model with dipole-dipole interaction is, indeed, conserved. Importantly, the idea is based on a many-body phenomenon, as long as there is spin-spin interaction, our conclusion should be general. Furthermore, since there are always some dipole interactions between spins, including their contribution is also rather natural. Finally, the conserved spin-current gives a linear spin-conductance, we derive an expression for this spin-conductance for multi-probe device systems.
We consider a coherent multi-probe mesoscopic device which consists a device scattering region that connects to a number of leads extending to far away. In a lead labeled by $`\alpha `$, the spin-current $`I_{s\alpha }<dN_s/dt>`$. Following Ref.jauho , we calculate the time derivative using Heisenberg equation of motion and in steady state, the spin current $`I_s`$ can be written in terms of Green’s functions:
$`I_{s\alpha }`$ $`=`$ $`{\displaystyle }{\displaystyle \frac{dE}{4\pi }}\mathrm{Tr}[\sigma _3(𝐆^r𝚺_\alpha ^<𝚺_\alpha ^<𝐆^a`$ (1)
$`+`$ $`𝐆^<𝚺_\alpha ^a𝚺_\alpha ^r𝐆^<)]`$
where $`𝚺^{r,a}`$ are the retarded and advanced self-energies due to the presence of leads; $`𝐆^{<,r,a}`$ are the lesser, retarded and advanced Green’s functions of the device scattering region, respectively. Note that the trace in the last equation can be reduced as:
$`\mathrm{Tr}\left[\sigma _3(𝐆^r𝚺^<𝚺^<𝐆^a)\right]=\mathrm{Tr}[𝐆^aF𝐆^r𝚺_\alpha ^<]`$
where
$`F[𝐆^a]^1\sigma _3\sigma _3[𝐆^r]^1=[\sigma _3,H_0]𝚺^a\sigma _3\sigma _3𝚺^r.`$
Here $`H_0`$ is the Hamiltonian of the device scattering region without the leads. Hence Eq.(1) can be written as
$`I_{s\alpha }`$ $`=`$ $`{\displaystyle }{\displaystyle \frac{dE}{4\pi }}\mathrm{Tr}[𝐆^<(𝚺_\alpha ^a\sigma _3\sigma _3𝚺_\alpha ^r)`$ (2)
$``$ $`𝐆_\alpha ^<(𝚺^a\sigma _3\sigma _3𝚺^r)+[\sigma _3,H_0]𝐆_\alpha ^<]`$
where $`𝐆_\alpha ^<=i𝐆^r𝚪_\alpha f_\alpha 𝐆^a`$ with $`\mathrm{\Gamma }_\beta `$ the linewidth function.
The total spin-current is
$`I_s={\displaystyle \underset{\alpha }{}}I_{s\alpha }={\displaystyle \frac{S}{t}}`$ (3)
where
$`{\displaystyle \frac{S}{t}}i{\displaystyle \frac{dE}{4\pi }\mathrm{Tr}[\dot{\sigma }_3G^<]}.`$ (4)
with $`\dot{\sigma }_3=i[\sigma _3,H_o]`$. Eq.(3) is just the continuity equation for spin current in the absence of external magnetic field and spin relaxationstiles . The physics of this equation is rather clear: in the presence of any spin-orbit interaction the electron spin precesses due to an internal torque on the spin, hence the total spin current flowing into the scattering region is non-zero and it equals to the rate of spin precession $`S/t`$. This is as if there were some “spin accumulation” in the scattering region. In other words, the total spin-current $`I_s`$ is itself not conserved due to spin precessing.
This situation is reminiscent to the continuity equation of charge current $`I_e`$ in the presence of a time dependent fieldbuttiker , namely
$`I_e{\displaystyle \underset{\alpha }{}}I_{e,\alpha }={\displaystyle \frac{Q}{t}}`$ (5)
Here the total charge current $`I_e`$ is equal to the charge accumulation in the scattering region. The problem of charge current conservation under AC fields was discussed by Büttikerbuttiker who pointed out that the total particle current is not conserved under AC conditions, but the total particle current plus total displacement current is a conserved quantity. The total displacement current is precisely $`Q/t`$ in the above equation. Since displacement current results from induction which is related to electron-electron interaction, Büttiker formulated a current conserving theorybuttiker by including the electron-electron interaction at the mean field level, which naturally deduces the displacement current in each lead. The problem of displacement current partition in multi-probe conductor under nonequilibrium conditions has been reported in Ref.wbg, within the nonequilibrium Green’s function formalism.
The similarity between Eqs.(3) and (5) suggests that the problem of spin-current conservation can also be looked at from an interacting spin point of view. As discussed above, a spin-spin interaction causes a dynamic response inthe system which counter-react on any rotation of a spin. In the following, we demonstrate this many-body physics by investigating the consequence of a self-consistent spin-spin interaction in a device Hamiltonian which contains the Rashba spin-orbit (SO) interaction. Indeed, we prove that the self-consistent spin-spin interaction produces a term that exactly cancels the right hand side of Eq.(3) so that $`_\alpha I_{s\alpha }=0`$.
Our model Hamiltonian is ($`\mathrm{}=1`$):
$`H=H_o+V^{ss}`$
where
$`H_o={\displaystyle \underset{n}{}}(ϵ_n+qU_n)d_n^{}d_n+{\displaystyle \underset{nm}{}}(V_{nm}^{so}d_n^{}d_m+h.c.)`$ (6)
$`H^{ss}={\displaystyle \frac{1}{2}}{\displaystyle \underset{nmij}{}}(V_{nmij}^{ss}d_n^{}d_md_i^{}d_j+h.c.)`$
where $`H^{ss}`$ is the spin-spin interaction and $`n`$,$`m`$,$`i`$, and $`j`$ include spin indices. Here $`V^{so}`$ and $`V^{ss}`$, respectively, are matrix elements of the Rashba SO interaction and the spin-spin interaction. In real space, they are given by
$`V^{so}(x)=\alpha _R\sigma (\widehat{z}\times 𝐤)`$ (7)
andjackson
$`V^{ss}(x,x^{})=(g\mu _B)^2\sigma \sigma ^{}{\displaystyle \frac{1}{|xx^{}|}}.`$
To deal with transport problems, we make a mean field analysis on the spin-spin interaction so that $`H^{ss}`$ becomes
$`H^{ss}={\displaystyle \underset{nm}{}}V_{nm}d_n^{}d_m`$
with
$`V_{nm}={\displaystyle \frac{1}{2}}{\displaystyle \underset{ij}{}}(V_{nmij}^{ss}<d_i^{}d_j>+h.c.)`$
In real space, $`V_{nm}`$ becomes,
$`V(x)=g\mu _B\sigma \varphi _M=g\mu _B\sigma 𝐇`$ (8)
with
$`\varphi _M={\displaystyle _V}{\displaystyle \frac{^{}𝐌(x^{})}{|xx^{}|}}𝑑x^{}`$ (9)
The local magnetic moment $`𝐌`$ is determined by spin density which, in the language of nonequilibrium Green’s functions (NEGF), is given by
$`𝐌(x)=ig\mu _B{\displaystyle \frac{dE}{2\pi }\mathrm{Tr}_s[\sigma 𝐆_{xx}^<]}`$
where $`\mathrm{Tr}_s`$ is the trace over spin space and $`𝐆_{xx}^<`$ is the diagonal element of $`𝐆^<`$. Note that $`𝐇=\varphi _M`$ is the self-consistent effective magnetic field due to the spin-spin interaction. The Green’s function also depends on this effective field through $`V(x)`$,
$`𝐆^r={\displaystyle \frac{1}{EH_0g\mu _B\sigma \varphi _M𝚺^r}}.`$ (10)
From Eq.(9), $`\varphi _M`$ satisfies the following Poisson-like equation
$`^2\varphi _M=4\pi \rho _M`$ (11)
where $`\rho _M=𝐌`$ is the effective density of magnetic chargejackson . These equations form a self-consistent problem.
To solve Eq.(11), we consider a very large volume of space surrounding the device scattering region such that the total magnetic charge inside that volume to be zero. Mathematically, this consideration means:
$`{\displaystyle 𝑑x𝐌(x)}=0\mathrm{or}{\displaystyle \frac{dE}{2\pi }\mathrm{Tr}(\sigma G^<)}=0`$ (12)
where $`\mathrm{Tr}=\mathrm{Tr}_s\mathrm{Tr}_o`$ includes the trace over both spin space and orbit space. Eq.(12) can be achieved if the volume in which we solve Eq.(11) is so large that the effective magnetic field $`𝐇`$ on the surface of that volume is zero.
Using the model Hamiltonian Eq.(6) and applying the Heisenberg equation of motion, we can evaluate
$`\dot{\sigma }_3=i[\sigma _3,H_0]=2i\alpha _R\sigma .`$
This result together with Eq.(12) proves
$$\frac{dE}{2\pi }\mathrm{Tr}(\dot{\sigma }_3G^<)=0.$$
Therefore the “spin accumulation” of Eq.(4) is actually zero. In other words, if the spin-spin interaction is included in the Green’s function, there will be no spin accumulation and the spin current is conserved. We emphasize that in the self-consistent formalism, Eqs.(1,10,11) form the basic set of equations for the spin current conserving theory.
Now we derive a “spin conductance” $`𝒢`$ that corresponds to the spin-current. When the external bias is small, we expand the spin-spin interaction $`V(x)`$ in terms of bias $`v_\alpha `$, $`V(x)=_\alpha u_\alpha v_\alpha `$ where we have introduced the notion of spin dependent characteristic potentialbuttiker ; ma $`u_\alpha `$ which satisfy the gauge invariant condition $`_\alpha u_\alpha =1`$, i.e., the spin current depends only on the difference in external bias. Expanding Eq.(2) in terms of small bias $`v_\beta `$, we find at zero temperature,
$`I_{s\alpha }={\displaystyle \underset{\beta }{}}𝒢_{\alpha \beta }v_\beta `$ (13)
where
$`𝒢_{\alpha \beta }=\mathrm{Tr}[g_{\alpha \beta }({\displaystyle \underset{\gamma }{}}g_{\alpha \gamma })_{xx}u_\beta (x)]`$ (14)
where the matrix $`g_{\alpha \beta }`$ is given by
$`g_{\alpha \beta }`$ $`=`$ $`{\displaystyle \frac{i}{4}}[(𝚺_\alpha ^a\sigma _3\sigma _3𝚺_\alpha ^r)(𝚺^a\sigma _3\sigma _3𝚺^r)\delta _{\alpha \beta }`$ (15)
$`+[\sigma _3,H_0]\delta _{\alpha \beta }]𝐆^r\mathrm{\Gamma }_\beta 𝐆^a+h.c.`$
Note that the first term of Eq.(14) comes from expansion of Fermi distribution function and the second term involving the characteristic potential is due to the spin-spin interaction term in the expansion. Linearizing Eq.(11) we have
$`^2\varphi _\alpha (x)=\kappa \mathrm{Tr}_s[({\displaystyle \underset{\eta \gamma }{}}g_{\eta \gamma })_{xx}u_\alpha (x)({\displaystyle \underset{\gamma }{}}g_{\gamma \alpha })_{xx}]`$ (16)
where $`\kappa =(2\pi g\mu _B)/\alpha _R`$ and $`\varphi _M=_\alpha \varphi _\alpha v_\alpha `$. Note that the expansion over external bias, the spin conductance $`g_{\alpha \beta }`$ in both Eq.(15) and Eq.(16) do not depend on the self-consistent interaction. From Eq.(8) we have $`u_\alpha =g\mu _B\sigma \varphi _\alpha `$. From Poisson like equation Eq.(16) we obtain the spin dependent characteristic potential and the spin conductance can be calculated from Eq.(14).
We emphasize that this conductance guarantees that the linear spin-current of Eq.(13) is conserved. Without the spin-spin interaction, the conductance would be given by only the first term on the right hand side of Eq.(17), and the resulting spin-current would not be conserved. In fact, the sum of the second term of Eq.(17) over space, i.e., the quantity $`\mathrm{Tr}[_\gamma g_{\gamma \alpha }]`$, is exactly equal to the ”spin accumulation” $`S_\alpha /t`$ in the scattering region from lead $`\alpha `$ due to Rashba interaction in the small bias limit. To see this, we find from Eq.(15)
$`{\displaystyle \frac{S_\alpha }{t}}=\mathrm{Tr}[{\displaystyle \underset{\gamma }{}}g_{\gamma \alpha }]={\displaystyle \frac{i}{2}}\mathrm{Tr}\left[[\sigma _3,H_0]𝐆^r\mathrm{\Gamma }_\alpha 𝐆^a\right]`$
so that $`_\alpha S_\alpha /tv_\alpha =S/t`$. This means that the spin-spin interaction puts this contribution into the spin conductance itself automatically, so that the right hand side of Eq.(3) vanishes. Eq.(17) therefore partitions the non-conserving part of the spin-current (i.e. the right hand side of Eq.(3) when there is no spin-spin interaction) into each leads, such that the spin-current becomes conserved.
Since the solution of the Poisson like equation requires numerical calculation, analytically we can avoid this by using a quasi-neutrality approximation, i.e., assuming that the effective density of magnetic charge $`\rho _M(x)=0`$ so that the local magnetic moment $`𝐌`$ is independent of position. Then, we find the spin dependent characteristic potential by setting the right hand side of Eq.(16) to zero and obtain: $`u_\alpha (x)=(_\gamma g_{\gamma \alpha })_{xx}/(_{\eta \gamma }g_{\eta \gamma })_{xx}`$. The conductance is then found to be:
$`𝒢_{\alpha \beta }=\mathrm{Tr}[g_{\alpha \beta }{\displaystyle \frac{(_\gamma g_{\alpha \gamma })(_\gamma g_{\gamma \alpha })}{(_{\eta \gamma }g_{\eta \gamma })}}].`$ (17)
This expression is very similar to that of charge current partition in ac situationsbuttiker ; wbg , and it partitions the total spin-current into each lead $`\alpha `$ so that the total spin-current flowing into the device is conserved.
The above microscopic theory result, Eq.(17), is valid for Rashba SO interaction. For a general SO interaction, a similar expression to Eq.(17) can be derived using a phenomenological argumentbuttiker . To do that, we require two conditions: (i) the total spin current is conserved; (ii) the value of spin current depends only on the difference of external bias. The latter condition means that spin current remains unchanged if external bias at each lead is shifted by the same amount. Now, the unconserved spin current $`I_{s\alpha }^c`$ is given by Eq.(13). The “spin accumulation” $`I_s^dS/t`$ is given by $`_\alpha I_{s\alpha }^c=_\beta (_\alpha 𝒢_{\alpha \beta }^c)v_\beta =I_s^d`$, where we have used $`𝒢_{\alpha \beta }^c\mathrm{Tr}g_{\alpha \beta }`$ for the non-conserved spin conductance. Note that the total “spin accumulation” is due to the contribution $`I_{s\alpha }^d`$ from each lead $`\alpha `$, i.e. $`I_s^d=_\alpha I_{s\alpha }^d`$. Since only the total “spin accumulation” is known, we need to find $`I_{s\alpha }^d`$ by partition the spin current. For this purpose, let $`I_{s\alpha }I_{s\alpha }^c+A_\alpha I_s^d`$, or equivalently
$`𝒢_{\alpha \beta }=𝒢_{\alpha \beta }^cA_\alpha {\displaystyle \underset{\gamma }{}}𝒢_{\gamma \beta }^c`$ (18)
where $`A_\alpha `$ is an unknown to be determined. Condition (i) gives $`_\alpha 𝒢_{\alpha \beta }=0`$, hence we obtain $`_\alpha A_\alpha =1`$. Condition (ii) gives gauge invariance $`_\beta 𝒢_{\alpha \beta }=0`$, hence we obtain $`\mathrm{Tr}[_\beta g_{\alpha \beta }A_\alpha _{\gamma \beta }g_{\gamma \beta }]=0`$ from which we find $`A_\alpha =_\gamma 𝒢_{\alpha \gamma }/_{\gamma \eta }𝒢_{\gamma \eta }`$. Therefore Eq.(18) gives:
$`𝒢_{\alpha \beta }=𝒢_{\alpha \beta }^c{\displaystyle \frac{(_\gamma 𝒢_{\alpha \gamma }^c)(_\gamma 𝒢_{\gamma \alpha }^c)}{(_{\eta \gamma }𝒢_{\eta \gamma }^c)}}.`$ (19)
Eq.(19) has the same form as Eq.(17) which is specific to Rashba SO. We therefore propose that Eq.(17) can serve as a phenomenological theory which conserves spin-current regardless of the detailed spin-orbit interactions. Namely, if we use Eq.(17) to compute spin-conductance, the resulting spin-current from Eq.(13) will always be conserved regardless of which SO interaction is present.
In summary, we have proven that the conventional spin-current $`I_s<n_sv_s>`$ for Rashba interaction becomes a conserved quantity if spin-spin dipole interaction is included. Such a dipole interaction introduces a self-consistent field which correlates spins spatially. For general SO interactions, a phenomenological theory for spin current partition is proposed which conserves spin current, and the resulting spin-conductance has the same form as that derived from the microscopic theory of Rashba interaction.
Acknowledgments. We gratefully acknowledge support by a RGC grant from the SAR Government of Hong Kong under grant number HKU 7044/04P. B.G. W is supported by the grant from NSFC under grant number 90303011 and H.G is supported by NSERC of Canada, FQRNT of Québec and CIAR.
<sup>∗)</sup> Electronic address: jianwang@hkusub.hku.hk |
warning/0507/astro-ph0507248.html | ar5iv | text | # The recurrent ultra-luminous X-ray transient NGC~253~ULX1Based on observations obtained with XMM-Newton, an ESA science mission with instruments and contributions directly funded by ESA Member States and NASA
## 1 Introduction
Ultra-luminous X-ray sources (ULXs) are extra-nuclear compact X-ray sources with luminosities considerably exceeding the Eddington luminosity for stellar mass X-ray binaries of $`2`$$`\times 10^{38}`$ erg s<sup>-1</sup> (Makishima et al., 2000). However their luminosities are still lower than that of active galactic nuclei (AGN).
There are currently four preferred models to explain the luminosities of these objects. The first is that ULXs are intermediate mass black holes (IMBHs: M$`{}_{\mathrm{BH}}{}^{}10^210^5\mathrm{M}_{}`$). However IMBHs are at present not explainable with stellar evolutionary models. The alternatives are stellar-mass black hole X-ray binaries where either photon bubble instabilities allow super-Eddington luminosities (Begelman, 2002), anisotropically emitting X-ray binaries (King et al., 2001), or that ULXs are micro-quasars that are observed down the beam of their relativistic jet (e.g. Reynolds et al., 1997). It is therefore important to increase the sample of ULXs to find arguments that favour or exclude the above models. One attempt was the search for ULXs in 313 nearby galaxies from ROSAT HRI observations by Liu & Bregman (2005, hereafter LB2005). A target of this search was the starburst galaxy NGC~253 where they found 21 X-ray sources but only one of them matched their criteria for an ULX (NGC~253~ULX1). This source is located within, but close to the north-east boundary of the D25 ellipse of NGC~253.
We here report on a more detailed analysis of NGC~253~ULX1 including ROSAT, XMM-Newton and Chandra data, and specifically on the detection of a second outburst in one of the XMM-Newton observations.
## 2 Search for the source in XMM-Newton, Chandra and ROSAT archives
We searched the ROSAT, Chandra and XMM-Newton archive for observations of NGC~253. The results are listed in Table 1. Except for two Chandra observations the position of NGC~253~ULX1 was always in the field of view (FOV). Besides the first detection in ROSAT observation 601111h (LB2005), NGC~253~ULX1 was only visible in XMM-Newton observation 0110900101. These XMM-Newton and ROSAT HRI data are further discussed in Sect. 3 and 4, respectively.
For the remaining observations we determined $`3\sigma `$ upper limits for the count rate. From that we obtained upper limit for fluxes and luminosities (cf Table 1). We used WebPIMMS (v3.6c) with the spectral model we got from the analysis of XMM-Newton observation 0110900101 to determine energy conversion factors. The long term light curve of NGC~253~ULX1 is shown in Fig. 2.
## 3 Detailed analysis of XMM-Newton observation 0110900101
NGC~253~ULX1 was detected for the second time on 2000 December 14 with XMM-Newton. The position of the source was within the FOV of both of the MOS and the PN cameras. We used the latest version of the Science Analysis System (SAS v6.1) to process the obtained data from the PN (thin filter) and the two MOS (medium filter) detectors. Periods of high background were determined and excluded from further analysis. The low background exposure times for PN and for the MOS instruments were 23.0 ks and 24.4 ks, respectively.
We applied the source detection tasks eboxdetect and emldetect only on the data from the PN detector, as the source was positioned far from the optical axis and close to the edge of the FOV on the MOS detectors. The obtained position was then corrected using optical reference coordinates from the USNO B1 catalogue (Monet et al., 2003) of three AGN, identified by Vogler & Pietsch (1999, sources X4, X22, X58). The corrected position in J2000 coordinates is $`\alpha =`$ 00h48m20.11s, $`\delta =25^{}10\mathrm{}10\stackrel{}{.}4`$ with an error in position of 0$`\stackrel{}{.}`$3. The derived position is well within the positional errors given by LB2005 for NGC~253~ULX1.
A foreground star with a B magnitude of $`13`$ (Monet et al., 2003) is located close (15.5″) to the obtained position (Fig. 3). We can rule out that the actual detection of NGC~253~ULX1 in observation 0110900101 was caused by this star, as its proper motion of -9.2 mas/yr in RA$`\mathrm{cos}(`$DEC$`)`$ and -3.6 mas/yr in DEC (Zacharias et al., 2004) is too small to match the detected position of NGC~253~ULX1 with that of the star within the period of observations. Additionally there was no detection of the source in other XMM-Newton observations using the same filter.
We extracted energy spectra for NGC~253~ULX1 for all EPIC detectors. For the PN chip we included source counts from an elliptical region with major and minor axes of 27.6″ and 12.3″ respectively. The background region was a circular source-free region with a radius of 48″ on the same CCD close to the source. For MOS the source extraction region was an ellipse with major (minor) axes of 28.95″ (11.3″) for MOS1 and 31.35″ (15.15″) for MOS2, respectively. The background regions were circles with radii of 68″ and 80″ for MOS1 and MOS2, respectively. After subtracting the background the spectra for each instrument were rebinned to a significance level of $`3\sigma `$.
For the spectral analysis XSPEC 11.3.1 was used. The best-fit parameters from different models provided within XSPEC are listed in Table 2. Using the PN and MOS spectra simultaneously the source spectrum was best fitted with a bremsstrahlung model (Fig. 1). The fit of the multicolour disk blackbody model (diskpn) would also be acceptable. However, we favour the bremsstrahlung model since it is less complex and gives a better $`\chi _{red}^2`$. Except for the bremsstrahlung model the foreground absorption ($`N_H`$) had to be fixed to the Galactic foreground absorption as a lower limit (1.30$`\times 10^{20}`$ cm<sup>-2</sup>, Dickey & Lockman, 1990). If the parameter was free to adjust it converged to unreasonably low values.
From the best fitting spectral model we calculated the source flux and, assuming a distance of 2.58 Mpc (Puche et al., 1991) we derived an unabsorbed luminosity of 5.0$`\times 10^{38}`$ erg s<sup>-1</sup> in the 0.3-10.0 keV band.
In order to study the temporal behaviour of the source a background corrected light curve was created using the tasks evselect and lccorr. The source count rate was constant at about 0.8 ct s<sup>-1</sup> within the errors of about 15% during observation 0110900101.
## 4 Analysis of ROSAT observation 601111h
The first detection of NGC~253~ULX1 was in ROSAT observation 601111h (LB2005). The observation (total exposure time 17.5 ks) is spread over ten observing intervals, with different exposure and waiting time for the individual observations. The source was bright enough to determine luminosities for each of these observation intervals.
We calculated count rates using the EXSAS source detection task detect/sources. To reduce noise we only analysed HRI channel 2-15. We used WebPIMMS (v3.6c) and the spectral model retrieved from the XMM-Newton observation (bremsstrahlung, kT$`=2.24`$keV, $`N_H=1.74`$$`\times 10^{20}`$ cm<sup>-2</sup>) to determine energy conversion factors to obtain the corresponding fluxes and luminosities (see lower panel of Fig. 2). The luminosity averaged over the whole observation (1.43$`\times 10^{39}`$ erg s<sup>-1</sup>) is indicated by the dashed line.
During the observation the source showed significant variability by at least a factor of 2.
## 5 Discussion
We detected the recurrence of NGC~253~ULX1 in the XMM-Newton observation from 2000 December 14. This was the first detection after the outburst in 1997, reported from ROSAT HRI observations by LB2005. In all other observations of NGC~253 the luminosity of the source was below the detection limit. This implies brightness variability by at least a factor of 500. Its fastest change in luminosity $`(\text{L}_{\text{max}}/\text{L}_{\text{min}})`$ exceeds a factor of 71 in 120 days.
The improved position (errors of 0$`\stackrel{}{.}`$24 compared to 4″- 10″) of the source determined in Sect. 3 allowed us to search for optical counterparts. We checked images taken with the Wide Field Imager (WFI) on the MPG-ESO 2.2m Telescope at La Silla in the R- (Fig. 3), I- and B-band (limiting magnitudes 24.2, 22.9 and 24.3, respectively) and images taken with the Galaxy Evolution Explorer (GALEX, a space telescope from NASA observing in the ultraviolet) in the NUV (Fig. 4) and FUV (limiting magnitudes 22 and 23, respectively), but no counterpart could be detected.
With the data discussed in Sect. 3 and 4 we can exclude that NGC~253~ULX1 is either a foreground object or a background AGN based on three arguments: (i) We estimated the $`\mathrm{log}(f_\text{x}/f_{\text{opt}})=\mathrm{log}f_\text{x}+(m_\text{v}/2.5)+5.37>3.2`$ using the flux of the ROSAT detection and a lower limit for $`m_\text{v}`$ of 24.2 (averaging the limiting magnitudes of the R- and the B-WFI images, see above). Following Maccacaro et al. (1988) this value exceeds that expected for galactic sources ($`4.6`$ to $`0.6`$) as well as AGNs ($`1.2`$ to $`+1.2`$). (ii) The variability of NGC~253~ULX1 is by a large factor higher than the typical value observed for AGNs $`(1060)`$. (iii) NGC~253~ULX1 shows a bremsstrahlung spectrum, whereas spectra of AGNs above $`2`$keV are typically fitted by a power law. The recurrent outbursts also exclude that the source is the luminous remnant of a recent supernova, like e.g. SN1993J in M81 (Zimmermann & Aschenbach, 2003).
The X-ray spectrum may indicate that NGC~253~ULX1 is a low mass X-ray binary (LMXB). The X-ray emission in these objects is created in the optically thin boundary layer between the disk and the neutron star and comptonization may dominate the spectral emission (White et al., 1988), leading to a spectrum that can be fitted with a bremsstrahlung model. However the ROSAT HRI peak luminosity of 1.43 $`\times 10^{39}`$ erg s<sup>-1</sup> is very high for typical LMXBs. Other systems that show bremsstrahlung spectra are black hole XRBs, e.g. Cyg X-1 (Sunyaev & Truemper, 1979), LMC X-3 and X1755-33 (White et al., 1988). These systems may contain a high or low mass companion.
An additional argument for a low mass companion comes from the lack of an optical counterpart (see above). High mass X-ray binaries (HMXBs) should be detectable at about 22 to 24 mag, extrapolating V magnitudes from HMXBs in the Magellanic Clouds (Liu et al., 2000). We would have detected an object of this brightness in the WFI data.
The luminosity of a compact object radiating at the Eddington limit is given as $`L_{Edd}=1.5\times 10^{38}(M/M_{})`$ erg s<sup>-1</sup>, when electron scattering dominates the opacity. Luminosities higher than 2$`\times 10^{38}`$ erg s<sup>-1</sup> (corresponding to a $`1.4M_{}`$ object, commonly assumed as the maximum mass of a neutron star) suggests that the compact object is a black hole. According to NGC~253~ULX1’s maximum luminosity of 1.43$`\times 10^{39}`$ erg s<sup>-1</sup> the lower limit for the mass of the black hole is 11 M. Therefore NGC~253~ULX1 is not required to be an IMBH.
Another argument against an IMBH is the temperature of NGC~253~ULX1. Miller et al. (2004a) compared intermediate mass black hole candidate ULXs and stellar mass black holes with respect to luminosity and temperature. If we assume the multicolour disk blackbody model then NGC~253~ULX1’s position in the luminosity-disk temperature diagram (Fig. 2 in Miller et al., 2004a) indicates that NGC~253~ULX1 is not an IMBH, but a stellar mass black hole.
Recently another object was found that, like NGC~253~ULX1, showed also a bremsstrahlung spectrum: X-44 in the Antennae Galaxies (NGC 4038/4039) (Miller et al., 2004b). The temperature of X-44 is $`3.7\pm 0.5`$keV and its luminosity is $`1.0_{0.2}^{+1.3}`$ $`\times 10^{40}`$ erg s<sup>-1</sup>. This temperature is about a factor of 1.5 higher than in NGC~253~ULX1, and the luminosity exceeds the luminosity of NGC~253~ULX1 by a factor of 15 compared to the outburst in 1997.
Another interesting ULX to compare NGC~253~ULX1 with is M101~ULX-1 (Kong et al., 2005). It was the first ULX that like NGC~253~ULX1 has been observed during more than one ultra-luminous outburst. Like many other ULXs the spectrum of M101~ULX-1 is best described with an absorbed blackbody model, but the temperature of $`50160`$eV is rather low. M101~ULX-1 has a peak luminosity of about $`10^{41}`$ erg s<sup>-1</sup> ($`0.37`$keV), and the hardness of its spectrum changed between different observations. We do not know whether the spectrum of NGC~253~ULX1 changed in the two observations, as the XMM-Newton data provided the very first spectrum of the source. During 12 years of observations NGC~253~ULX1 showed two outbursts with an interval of three years. In M101~ULX-1 the two outbursts are only seperated by half a year. On shorter time scales NGC~253~ULX1 showed only one drop in luminosity by a factor of $`2`$ during the ROSAT observation, and in the XMM-Newton observation (exposure time $`8.2`$h) no variability could be detected. M101~ULX-1 on the other hand does show short-time-scale variability. Its luminosity changed by a factor of $`\mathrm{}>10`$ on a time scale of hours. The lack of short time variability of NGC~253~ULX1 argues against the relativistic beaming model, since this would require a very stable jet (Reynolds et al., 1997).
###### Acknowledgements.
We thank G. Szokoly for providing us with the images of NGC~253 from the Wide Field Imager, which are based on observations made with ESO Telescopes at the La Silla and Paranal Observatory. Some of the data presented in this paper were obtained from the Multimission Archive at the Space Telescope Science Institute (MAST). STScI is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS5-26555. Support for MAST for non-HST data is provided by the NASA Office of Space Science via grant NAG5-7584 and by other grants and contracts. The XMM-Newton and the ROSAT project is supported by the Bundesministerium für Bildung und Forschung/Deutsches Zentrum für Luft- und Raumfahrt (BMBF/DLR), the Max-Planck Society and the Heidenhain-Stiftung. M.B. acknowledges support from the International Max Planck Research School on Astrophysics (IMPRS). |
warning/0507/math0507121.html | ar5iv | text | # On the poles of topological zeta functions
*Abstract.— We study the topological zeta function $`Z_{top,f}(s)`$ associated to a polynomial $`f`$ with complex coefficients. This is a rational function in one variable and we want to determine the numbers that can occur as a pole of some topological zeta function; by definition these poles are negative rational numbers. We deal with this question in any dimension. Denote $`𝒫_n:=\{s_0f[x_1,\mathrm{},x_n]:Z_{top,f}(s)\text{has a pole in}s_0\}`$. We show that $`\{(n1)/21/ii_{>1}\}`$ is a subset of $`𝒫_n`$; for $`n=2`$ and $`n=3`$, the last two authors proved before that these are exactly the poles less then $`(n1)/2`$. As main result we prove that* each *rational number in the interval $`[(n1)/2,0)`$ is contained in $`𝒫_n`$.*
1. Introduction
Denef and Loeser created in $`1992`$ a new zeta function, which they called the topological zeta function because of the topological Euler–Poincaré characteristic turning up in it. Roughly said, the topological zeta function $`Z_{top,f}`$ associated to a polynomial $`f`$ is a function containing information we can pick out of each chosen embedded resolution of $`f^1\{0\}𝔸^n`$. They introduced it in \[DL1\] in the following way.
Let $`f`$ be a polynomial in $`n`$ variables over $``$ and let $`h:X𝔸^n`$ be an embedded resolution of $`f^1\{0\}`$. To define $`Z_{top,f}`$ we need some data related to the embedded resolution $`(X,h)`$. Let $`E_i,iS`$, be the irreducible components of $`h^1(f^1\{0\})`$, then denote by $`N_i`$ and $`\nu _i1`$ the multiplicities of $`E_i`$ in the divisor on $`X`$ of $`fh`$ and $`h^{}(dx_1\mathrm{}dx_n)`$, respectively. The couples $`(N_i,\nu _i),iS`$, are called the numerical data of the resolution $`(X,h)`$. For $`IS`$ we denote also $`E_I:=_{iI}E_i`$ and $`E_I^{}:=E_I(_{jI}E_j)`$. Further we write $`\chi ()`$ for the topological Euler–Poincaré characteristic.
Definition.— The local *topological zeta function associated to $`f`$* is the rational function in one complex variable
$$Z_{top,f}(s):=\underset{IS}{}\chi (E_I^{}h^1\{0\})\underset{iI}{}\frac{1}{N_is+\nu _i}.$$
There is a global version replacing $`E_I^{}h^1\{0\}`$ by $`E_I^{}`$. When we do not specify, we mean the local one.
Denef and Loeser proved that every embedded resolution gives rise to the same function, so the topological zeta function is a well-defined singularity invariant (see \[DL1\]). Once the motivic Igusa zeta function was introduced, they proved this result alternatively in \[DL2\] by showing that this more general zeta function specialises to the topological one.
In particular the poles of the topological zeta function of $`f`$ are interesting numerical invariants. Various conjectures relate them to the eigenvalues of the local monodromy of $`f`$, see for example \[DL1\]. The poles are part of the set $`\{\nu _i/N_iiS\}`$; therefore the $`\nu _i/N_i`$ are called the candidate poles. Notice that the poles are negative rational numbers.
A related numerical invariant of $`f`$ at $`0^n`$ is its *log canonical threshold* $`c_0(f)`$ which is by definition
$$sup\{c\text{ the pair }(^n,c\text{div}f)\text{ is log canonical in a neighbourhood of }0\}.$$
It is described in terms of the embedded resolution as $`c_0(f)=\mathrm{min}\{\nu _i/N_i0h(E_i),iS\}`$ (see \[Ko2, Proposition 8.5\]). It was studied in various papers of Alexeev, Cheltsov, Ein, de Fernex, Kollár, Kuwata, M<sup>c</sup>Kernan, Mustaţă, Park, Prokhorov, Reid, Shokurov and others. Especially the sets $`𝒯_n:=\{c_0(f)f[x_1,\mathrm{},x_n]\},`$ with $`n_{>0}`$, show up in interesting conjectures, see \[Al\], \[Ko\], \[Ku\], \[M<sup>c</sup>KP\], \[Pr\] and \[Sh\].
In the context of the topological zeta function, it is natural to study similarly the set
$$𝒫_n:=\{s_0f[x_1,\mathrm{},x_n]:Z_{top,f}(s)\text{ has a pole in }s_0\}.$$
The case $`n=1`$ is trivial: $`𝒫_1=\{1/ii_{>0}\}`$.
From now on we assume that $`n2`$. A more or less obvious lower bound for $`𝒫_n`$ is $`(n1)`$, see \[Se1, Section 2.4\]. In \[SV\], the second and the third author studied the ‘smallest poles’ for $`n=2`$ and $`n=3`$. They showed that $`𝒫_2(\mathrm{},\frac{1}{2})=\{\frac{1}{2}\frac{1}{i}i_{>1}\}`$ and that $`𝒫_3(\mathrm{},1)=\{1\frac{1}{i}i_{>1}\}`$. They expected that this could be generalised to
$`𝒫_n(\mathrm{},{\displaystyle \frac{n1}{2}})=\{{\displaystyle \frac{n1}{2}}{\displaystyle \frac{1}{i}}i_{>1}\},\text{for all }n_{>1}.`$
In particular, they predicted that the lower bound $`(n1)`$ could be sharpened to $`n/2`$. This better bound was recently proven by the second author in \[Se2\]. In this article we verify for all $`n4`$ that $`\{(n1)/21/ii_{>1}\}𝒫_n`$, and as main result we show that *any* rational number in the remaining interval $`[(n1)/2,0)`$ is a pole of some topological zeta function.
Theorem.— *For $`n2`$ we have $`[(n1)/2,0)𝒫_n`$.*
With the Thom-Sebastiani principle \[DL3\], $`x_1^i+x_2^2+\mathrm{}+x_n^2`$ is the obvious candidate to have $`(n1)/21/i`$ as a pole of its associated topological zeta function. It is not clear a priori that this will be true for all $`n`$ and $`i`$. We check this in section $`2`$. For the theorem, however, the key is to find a suitable family of polynomials.
We will put the useful information of the resolution into a diagram, which is called the dual intersection graph. It is obtained as follows. One associates a vertex to each exceptional component in the embedded resolution (represented by a dot) and to each component of the strict transform of $`f^1\{0\}`$ (represented by a circle). One also associates to each intersection an edge, connecting the corresponding vertices. The fact that $`E_i`$ has numerical data $`(N_i,\nu _i)`$ is denoted by $`E_i(N_i,\nu _i)`$.
When the strict transform of $`f^1\{0\}`$ is irreducible, we will denote it by $`E_0`$. Let $`E_i`$ be an exceptional variety and let $`E_j`$, $`jJ`$, be the components that intersect $`E_i`$ in $`X`$. We set $`\alpha _j:=\nu _j(\nu _i/N_i)N_j`$ for $`jJ`$; these numbers appear in the calculation of the residue of $`Z_{top,f}`$ in $`\nu _i/N_i`$.
2. The set $`\{(n1)/21/ii_{>1}\}`$ is a subset of $`𝒫_n`$
*Embedded resolution for $`𝒙_1^𝒊+𝒙_2^2+\mathrm{}+𝒙_𝒏^2=0`$, $`𝒏4`$, with $`𝒊`$ even*
After blowing up $`i/2`$ times in the origin, we get an embedded resolution for $`f`$. We present the dual intersection graph for $`i2`$.
The exceptional variety $`E_{i/2}`$ gives the candidate pole $`(n1)/21/i`$ in which we are interested. If $`i2`$, its residue is
$`{\displaystyle \frac{1}{N_{\frac{i}{2}}}}`$ $`\left(\chi (E_{I_1}^{})+\chi (E_{I_2}^{}){\displaystyle \frac{1}{\alpha _{\frac{i}{2}1}}}+\chi (E_{I_3}^{}){\displaystyle \frac{1}{\alpha _0}}+\chi (E_{I_4}^{}){\displaystyle \frac{1}{\alpha _0\alpha _{\frac{i}{2}1}}}\right),`$
where $`I_1=\{\frac{i}{2}\},I_2=\{\frac{i}{2},\frac{i}{2}1\},I_3=\{\frac{i}{2},0\},I_4=\{\frac{i}{2},\frac{i}{2}1,0\}.`$ The Euler–Poincaré characteristics $`\chi (E_{I_j}^{})`$, $`1j4`$, are put in Table $`1`$. These are easily computed since $`E_{i/2}^{n1}`$, and $`E_{i/21}`$ and $`E_0`$ intersect $`E_{i/2}`$ in a hyperplane and a smooth quadric, respectively.
| $`\chi (E_{I_j}^{})`$ | n odd | n even |
| --- | --- | --- |
| $`j=1`$ | $`1`$ | $`1`$ |
| $`j=2`$ | $`0`$ | $`1`$ |
| $`j=3`$ | $`0`$ | $`2`$ |
| $`j=4`$ | $`n1`$ | $`n2`$ |
*Table 1*
Using that $`\alpha _0=(3n)/21/i`$ and $`\alpha _{i/21}=2/i`$, some easy calculations yield that the residue is non-zero, for all $`n`$, $`n4`$.
When $`i=2`$, we blow up just once in the origin to get an embedded resolution. By using $`\alpha _0=\frac{2n}{2},\chi (E_{I_1}^{})=0(n\text{ even}),\chi (E_{I_1}^{})=1(n\text{ odd}),`$ we conclude that also here the residue is non-zero.
*Embedded resolution for $`𝒙_1^𝒊+𝒙_2^2+\mathrm{}+𝒙_𝒏^2=0`$, $`𝒏4`$, with $`𝒊`$ odd*
After blowing up $`(i+1)/2`$ times in the origin, followed by blowing up once more in $`D:=E_{(i+1)/2}E_{(i1)/2}^{n2}`$, we get an embedded resolution with the following dual intersection graph.
The last exceptional variety has $`(n1)/21/i`$ as candidate pole. The relevant subsets in the computation of the residue are $`I_1=\{\frac{i+3}{2}\},I_2=\{\frac{i+3}{2},0\},I_3=\{\frac{i+3}{2},\frac{i+1}{2}\},I_4=\{\frac{i+3}{2},\frac{i1}{2}\},I_5=\{\frac{i+3}{2},\frac{i1}{2},0\}.`$ Here $`E_{(i+3)/2}`$ is a $`^1`$-bundle over $`D`$. For $`j=2,3,4`$ we have that $`E_{I_j}D`$ and $`E_{I_5}`$ is a smooth quadric. With the Euler–Poincaré characteristics of Table $`2`$ and $`\alpha _0=(3n)/21/i`$, $`\alpha _{(i1)/2}=1/i`$ and $`\alpha _{(i+1)/2}=(n1)/2`$, we find that the residue is non-zero, for all $`n4`$.
| $`\chi (E_{I_j}^{})`$ | n odd | n even |
| --- | --- | --- |
| $`j=1`$ | $`0`$ | $`1`$ |
| $`j=2`$ | $`0`$ | $`1`$ |
| $`j=3`$ | $`n1`$ | $`n1`$ |
| $`j=4`$ | $`0`$ | $`1`$ |
| $`j=5`$ | $`n1`$ | $`n2`$ |
*Table 2*
Throwing together these results we obtain
$$\{\frac{n1}{2}\frac{1}{i}i_{>1}\}𝒫_n.$$
Now that we checked this expectation, we proceed proving the theorem.
*Remark.—* Notice that $`m𝒫_{n1}`$ implies that $`m𝒫_n`$. Indeed, any polynomial $`f`$ in $`n1`$ variables can be considered as a polynomial in $`n`$ variables. An embedded resolution for $`f^1\{0\}^{n1}`$ induces the obvious analogous one for $`f^1\{0\}^n=^{n1}\times `$ and, since $`\chi ()=1`$, the two associated topological zeta functions are equal. From this observation it follows that it is sufficient to prove that $`[(n1)/2,(n2)/2)𝒫_n`$. As we showed in this section that $`(n1)/2`$ is contained in $`𝒫_{n1}`$ and thus in $`𝒫_n`$, we restrict ourselves in the next sections to the subset $`((n1)/2,(n2)/2)`$.
3. The set $`(1/2,0)`$ is a subset of $`𝒫_2`$
Considering how candidate poles look like in the formula of the topological zeta function written in terms of newton polyhedra (see \[DL1\]), the number $`(b+2)/(2a+2b)`$ seems to appear as a candidate pole of the topological zeta function associated to $`f(x,y)=x^a(x^b+y^2)`$, where $`a`$ and $`b`$ are positive integers. An easy computation yields:
Lemma.— *When $`a`$ and $`b`$ run through $`2_{>0}`$, $`a2`$, the quotient $`(b+2)/(2a+2b)`$ takes all rational values in $`(1/2,0)`$.*
Taking the lemma into account, the functions $`f(x,y)=x^a(x^b+y^2)`$, where $`a,b2_{>0}`$ and $`a2`$, could be a pretty nice choice to obtain all desired poles. Easy calculations give the following dual resolution graph for $`f`$.
Because $`E_{b/2}`$ is intersected three times by other components, Theorem $`4.3`$ in \[Ve2\] allows us to conclude that $`(b+2)/(2a+2b)`$ is a pole of $`Z_{top,f}`$.
4. The set $`((n1)/2,(n2)/2)`$ is a subset of $`𝒫_n,n3`$
As this set is a translation by $`1/2`$ of expected poles in dimension $`n1`$, the Thom-Sebastiani principle in \[DL3\] is again the motivation why we consider
$$f(x_1,\mathrm{},x_n)=x_n^2+\mathrm{}+x_3^2+x_1^a(x_1^b+x_2^2),$$
where $`a2_{>0}`$ and $`a2`$, to reach the set $`((n1)/2,(n2)/2)`$.
*Embedded resolution for $`𝒛^2+𝒙^𝒂(𝒙^𝒃+𝒚^2)`$*
Let us first explain in dimension $`3`$ which embedded resolution we choose for $`z^2+x^a(x^b+y^2)`$ ($`a,b2_{>0}`$, $`a2`$). We first blow up in the singular locus $`\{x=z=0\}`$ of $`f`$ and further always in the singular locus of the strict transform; the first $`a/2`$ times this is an affine line and the last $`b/2`$ times it is a point. This is the special case for $`n=3`$ in Table $`3`$.
The dual intersection graph looks as follows.
The candidate pole given by the last exceptional surface, $`E_{(a+b)/2}`$, is equal to
$$\frac{a/2+b+1}{a+b}=\frac{b+2}{2a+2b}\frac{1}{2},$$
and thus covers all rational numbers in $`(1,1/2)`$ if $`a`$ and $`b`$ run over $`2_{>0}`$ and $`a2`$.
*Embedded resolution for $`𝒙_𝒏^2+\mathrm{}+𝒙_3^2+𝒙_1^𝒂(𝒙_1^𝒃+𝒙_2^2)`$, $`𝒏>3`$*
The sequence of blow-ups in Table $`3`$ yields an embedded resolution for
$$f(x_1,\mathrm{},x_n)=x_n^2+\mathrm{}+x_3^2+x_1^a(x_1^b+x_2^2),$$
based on the previous one for $`n=3`$.
number $`i`$ of center blow-up equation strict transform blow-up in relevant chart $`1`$ $`x_1=x_3=x_4=\mathrm{}=x_n=0`$ $`x_n^2+\mathrm{}+x_3^2+x_1^{a2}(x_1^b+x_2^2)`$ $`2`$ $`x_1=x_3=x_4=\mathrm{}=x_n=0`$ $`x_n^2+\mathrm{}+x_3^2+x_1^{a4}(x_1^b+x_2^2)`$ $`\mathrm{}`$ $`\mathrm{}`$ $`\mathrm{}`$ $`a/2`$ $`x_1=x_3=x_4=\mathrm{}=x_n=0`$ $`x_n^2+\mathrm{}+x_3^2+x_1^b+x_2^2`$ $`a/2+1`$ $`(0,0,\mathrm{},0)`$ $`x_n^2+\mathrm{}+x_3^2+x_1^{b2}+x_2^2`$ $`a/2+2`$ $`(0,0,\mathrm{},0)`$ $`x_n^2+\mathrm{}+x_3^2+x_1^{b4}+x_2^2`$ $`\mathrm{}`$ $`\mathrm{}`$ $`\mathrm{}`$ $`(a+b)/2`$ $`(0,0,\mathrm{},0)`$ $`x_n^2+\mathrm{}+x_3^2+1+x_2^2`$
*Table 3*
The dual intersection graph here looks as follows.
Now $`\nu _{(a+b)/2}/N_{(a+b)/2}`$ is equal to
$$\frac{a/2+b+1+((a+b)/2)(n3)}{a+b}=\frac{b+2}{2a+2b}\frac{n2}{2},$$
which covers the interval $`((n1)/2,(n2)/2)`$ when $`a`$ and $`b`$ vary in $`2_{>0}`$ with $`a2`$.
*The rational number $`𝝂_{(𝒂+𝒃)/2}/𝑵_{(𝒂+𝒃)/2}`$ is a pole of $`𝒁_{𝒕𝒐𝒑,𝒇}`$*
For all $`n3`$ and $`f(x_1,\mathrm{},x_n)=x_n^2+\mathrm{}+x_3^2+x_1^a(x_1^b+x_2^2)`$, we calculate the residue of $`Z_{top,f}`$ in $`\nu _{(a+b)/2}/N_{(a+b)/2}`$. Observe that if $`(a+b)/(2+b)`$, the exceptional variety $`E_{(a+b)/(2+b)}`$ induces the same candidate pole as $`E_{(a+b)/2}`$. The other exceptional varieties always give rise to other candidate poles.
The subsets playing a role in the contribution of $`E_{(a+b)/(2+b)}`$ to the residue are $`J_1=\{\frac{a+b}{2+b}\},J_2=\{\frac{a+b}{2+b},\frac{a+b}{2+b}1\},J_3=\{\frac{a+b}{2+b},\frac{a+b}{2+b}+1\},J_4=\{\frac{a+b}{2+b},0\},J_5=\{\frac{a+b}{2+b},\frac{a+b}{2+b}1,0\},J_6=\{\frac{a+b}{2+b},\frac{a+b}{2+b}+1,0\}.`$ Notice that when $`n=3`$, $`E_{(a+b)/(2+b)}`$ does not intersect $`E_0`$.
We have that $`E_{(a+b)/(2+b)}`$ is isomorphic to the cartesian product of $`𝔸^1`$ and the blowing-up of $`^{n2}`$ in a point. It is also easy to describe the whole intersection configuration on $`E_{(a+b)/(2+b)}`$.
| $`\chi (E_{J_k}^{})`$ | n odd | n even |
| --- | --- | --- |
| $`k=1`$ | $`0`$ | $`0`$ |
| $`k=2`$ | $`1`$ | $`0`$ |
| $`k=3`$ | $`1`$ | $`0`$ |
| $`k=4`$ | $`0`$ | $`0`$ |
| $`k=5`$ | $`n3`$ | $`n2`$ |
| $`k=6`$ | $`n3`$ | $`n2`$ |
*Table 4*
With the relevant Euler–Poincaré characteristics of Table $`4`$ and $`\alpha _{(a+b)/(2+b)1}=1/i`$, $`\alpha _{(a+b)/(2+b)+1}=1/i`$, we see that $`E_{(a+b)/(2+b)}`$ does not give any contribution to the residue in $`\nu _{(a+b)/2}/N_{(a+b)/2}`$. Alternatively, this is implied by \[Ve1, Proposition 6.5\]. This means we only have to take the contribution of $`E_{(a+b)/2}`$ into account.
To compute this contribution the relevant subsets for the summation in the formula of the topological zeta function are $`I_1=\{\frac{a+b}{2}\},\text{ }I_2=\{\frac{a+b}{2},\frac{a+b}{2}1\},\text{ }I_3=\{\frac{a+b}{2},0\},\text{ }I_4=\{\frac{a+b}{2},\frac{a+b}{2}1,0\}.`$ The Euler–Poincaré characteristics $`\chi (E_{I_j}^{})`$, $`1j4`$, are the same as those given in Table $`1`$ and we have $`\alpha _0=((n4)a+(n3)b+2)/(2(a+b))`$ and $`\alpha _{(a+b)/21}=(2a)/(a+b)`$.
As the residue then is equal to
$$\frac{(2+3a+2b)(na2ab+nb+2)}{(2+a)(a+b)(na4a+2+nb3b)}\text{ for }n\text{ odd and}$$
$$\frac{(2+b)(na2ab+nb+2)}{(2+a)(a+b)(na4a+2+nb3b)}\text{ for }n\text{ even},$$
we find that $`(\nu _{(a+b)/2})/(N_{(a+b)/2})=(b+2)/(2a+2b)(n2)/2`$ is a pole of $`Z_{top,f}`$.
We conclude that $`((n1)/2,(n2)/2)𝒫_n`$, for all $`n3`$.
5. Some remarks
($`1`$) Instead of achieving this result with the method of resolution of singularities one can find the poles of the topological zeta function of the polynomials
$$x_n^2+\mathrm{}+x_3^2+x_1^a(x_1^b+x_2^2)\text{ and }x_n^2+\mathrm{}+x_2^2+x_1^i$$
with the help of Newton polyhedra. Indeed, we can write down the topological zeta function for these polynomials using the formula of Denef and Loeser in \[DL1\]. For example if $`f(x_1,\mathrm{},x_n)=x_n^2+\mathrm{}+x_3^2+x_1^a(x_1^b+x_2^2)`$, where $`a`$ and $`b`$ are positive even integers and $`a2`$, put $`A:=(a+b)s+1+b/2+(n2)(a+b)/2`$ and $`B:=as+1+(n2)a/2`$. We get
$`Z_{top,f}(s)`$ $`=`$ $`(n1){\displaystyle \frac{b}{2AB}}+{\displaystyle \frac{1}{A}}+(n2){\displaystyle \frac{a}{2B}}`$
$`+{\displaystyle \frac{s}{s+1}}({\displaystyle \underset{d=1}{\overset{n1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n2}{d+1}}\right)({\displaystyle \frac{a}{2B}}+{\displaystyle \frac{b}{2AB}})(2)^d`$
$`+{\displaystyle \underset{d=1}{\overset{n1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n1}{d}}\right){\displaystyle \frac{1}{A}}(2)^d+{\displaystyle \underset{d=1}{\overset{n2}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n2}{d}}\right){\displaystyle \frac{b}{2AB}}(2)^d).`$
Handling the problem in this way leads to the same results. One just has to be careful with the dual cones of some faces, namely those that are not a rational simplicial cone.
($`2`$) With a similar definition of $`𝒫_n`$ in each case, the same results hold for local and global versions of the motivic zeta function, the Hodge zeta function and Igusa’s zeta function. Indeed, the results for the topological zeta function imply the results for those ‘finer’ zeta functions. |
warning/0507/cs0507038.html | ar5iv | text | # Upper Bound on the Number of Vertices of Polyhedra with 0,1-Constraint Matrices.
## 1 Introduction
Let $`A\{0,1\}^{m\times d}`$ be a $`0,1`$-matrix and $`b^m`$ be a given real vector. Consider the polyhedron $`P=P(A,b)=\{x^d:Axb\}`$. Denote by $`EXT(P)`$ the set of vertices of the polyhedron $`P`$. The purpose of this short note is to show the following upper bound on the size of $`EXT(P)`$.
###### Theorem 1
For any $`A\{0,1\}^{m\times d}`$ and any $`b^m`$, the number vertices of the polyhedron $`P=\{x^d:Axb\}`$ is at most $`d!`$.
The bound of Theorem 1 is tight: consider the polyhedron:
$$P=\mathrm{\Pi }_{d1}+_{}=\mathrm{conv}\{\sigma ^d:\sigma \text{ is permutation of }[d]\}+\mathrm{cone}\{𝐞_1,\mathrm{},𝐞_d\},$$
i.e. the polyhedron with all permutation of $`[d]`$ as vertices (known as Permutahedron) and the negative of the unit vectors as extreme directions. Then $`P`$ has $`d!`$ vertices and it is known that $`P`$ can have the following linear description with a $`0,1`$-constraint matrix $`A`$:
$$P=\{x^d:\underset{iS}{}x_i\underset{i=1}{\overset{|S|}{}}(ni+1),\text{for all}S[d]\}.$$
In the next section we give a few preliminaries that will be needed to prove the Theorem 1. In Section 3, we give the proof of the theorem.
## 2 Preliminaries
First, we may assume without loss of generality, by translating $`P=P(A,b)`$ if necessary, that the vector $`b`$ is strictly positive. To deal with the fact that $`P`$ is unbounded, we define $`\widehat{P}=P\{x^d:_{i=1}^dx_iM\}`$, for a sufficiently big $`M`$. Note that $`EXT(\widehat{P})EXT(P)\{x^d:_{i=1}^dx_i=M\}`$, and that $`0\mathrm{int}(P)\mathrm{int}(\widehat{P})`$ since $`b>0`$, where for $`X^d`$, we denote by $`\mathrm{int}(P)`$ the interior of $`P`$. We need the following simple facts, which are well-known and also are not difficult to prove.
###### Observation 2
Let $`Q`$ be a $`d`$-dimensional polytope and $`z\mathrm{int}(Q)`$ be a point in the interior of $`Q`$. Let $`F_1=\mathrm{conv}\{p_1,\mathrm{},p_r\}`$ and $`F_2=\mathrm{conv}\{q_1,\mathrm{},q_s\}`$ be two facets of $`Q`$, where $`p_1,\mathrm{},p_r,q_1,\mathrm{},q_sEXT(Q)`$ are vertices of $`Q`$. Let $`\lambda _1,\mathrm{},\lambda _r`$ and $`\mu _1,\mathrm{},\mu _r`$ be positive real numbers. Then
$$\mathrm{int}(\mathrm{conv}\{z,\lambda _1p_1,\mathrm{},\lambda _rp_r\})\mathrm{int}(\mathrm{conv}\{z,\mu _1q_1,\mathrm{},\mu _rq_s\})=\mathrm{}.$$
###### Observation 3 ()
The volume $`\mathrm{Vol}^d(Q)`$ of any $`d`$-dimensional polytope $`Q`$ with integer vertices is an integer multiple of $`\frac{1}{d!}`$. In particular $`\mathrm{Vol}^d(Q)\frac{1}{d!}`$.
Given a polytope $`Q`$ which contains $`\mathrm{𝟎}`$ as an interior point, the polar polytope $`Q^{}`$ is defined as
$$Q^{}=\{x^d:y^Tx1yQ\},$$
see e.g. . It is known that $`Q^{}`$ also contains $`\mathrm{𝟎}`$ as an interior point, and that the vertices and facets of $`Q`$ are in one-to-one correspondence with the facets and vertices of $`Q^{}`$, respectively.
## 3 Proof of Theorem 1
Consider $`(\widehat{P})^{}`$, the polar of the polytope $`\widehat{P}`$ defined by adding a bounding inequality to the polyhedron $`P=P(A,b)`$ as above. Form polarity, it follows that if $`F=\{x^d:xP,a_i^Tx=b_i\}`$ is a facet of $`P`$, where $`a_i^T`$ is the $`i`$th row of $`A`$ and $`b_i`$ is the $`i`$th component of $`b`$, then the corresponding polar vertex $`v_i`$ in $`EXT((\widehat{P})^{})`$ lies on the ray connecting the origin to $`a_i`$, at a distance $`d(a_i,\mathrm{𝟎})=\frac{1}{b_i}`$ from the origin.
Note that $`(\widehat{P})^{}`$ contains exactly one negative vertex $`u=(\frac{1}{M},\mathrm{},\frac{1}{M})`$ which does not correspond to a facet of $`P`$. It follows then that the vertices of $`P`$ are in one-to-one correspondence with the facets of $`(\widehat{P})^{}`$ that do not contain $`u`$. Consider any two such facets $`F_1`$ and $`F_2`$. Then $`F_1=\mathrm{conv}\{\frac{1}{b_{i_1}}a_{i_1},\mathrm{},\frac{1}{b_{i_r}}a_{i_r}\}`$ and $`F_2=\mathrm{conv}\{\frac{1}{b_{j_1}}a_{j_1},\mathrm{},\frac{1}{b_{j_s}}a_{j_s}\}`$, where $`i_1,\mathrm{},i_r`$ and $`j_1,\mathrm{},j_s`$ are indices from $`[m]`$. Now, since $`\mathrm{𝟎}\mathrm{int}(\widehat{P}^{})`$ and $`F_1`$ and $`F_2`$ are facets of $`\widehat{P}^{}`$, we note by Observation 2 that
$$\mathrm{int}(\mathrm{conv}\{\mathrm{𝟎},a_{i_1},\mathrm{},a_{i_r}\})\mathrm{int}(\mathrm{conv}\{\mathrm{𝟎},a_{j_1},\mathrm{},a_{j_s}\})=\mathrm{}.$$
By Observation 3, we have that the volume $`\mathrm{Vol}^d(S(a_{i_1},\mathrm{},a_{i_r})\})`$ of the polytope $`S(a_{i_1},\mathrm{},a_{i_r})=\mathrm{conv}\{\mathrm{𝟎},a_{i_1},\mathrm{},a_{i_r})`$ arising by ”stretching” the vertices of a face $`F=\mathrm{conv}\{\frac{1}{b_{i_1}}a_{i_1},\mathrm{},\frac{1}{b_{i_r}}a_{i_r}\}`$ of $`\widehat{P}^{}`$ is at least $`1/d!`$. Since all such polytopes $`S(a_{i_1},\mathrm{},a_{i_r})`$ lie completely inside the $`d`$-dimensional hypercube, whose volume is $`1`$, we conclude that there number, which is equal to the number of facets of $`\widehat{P}^{}`$ is at most $`d!`$. This finishes the proof.
It can be easily seen that the above proof can be extended to matrices $`A\{K,\mathrm{},1,`$ $`0,1,\mathrm{},K\}^{m\times d}`$, in which case the corresponding upper-bound on the number of vertices will be $`(2K+1)^dd!`$. In particular, if the number of bits to represent a given integer constraint matrix is $`L`$, then the number of vertices of the polyhedron described by $`A`$ is at most $`(2^L+1)^dd!`$. |
warning/0507/nucl-ex0507024.html | ar5iv | text | # Long range absorption in the scattering of 6He on 208Pb and 197Au at 27 MeV.
## I Introduction
The <sup>6</sup>He nucleus has a weakly bound three-body Borromean n-n-$`\alpha `$ structure and is known to have an extended two neutron distribution Tani96 ; Hansen87 . Scattering of <sup>6</sup>He is therefore of considerable interest in nuclear physics since experiments may demonstrate sensitivity to this underlying exotic structure.
The interest of measuring elastic scattering of <sup>6</sup>He arises from the weakly bound nature of the projectile, which affects the dynamics of the collision. <sup>6</sup>He is bound by less than 1 MeV, so we expect that the coupling to the continuum can significantly modify cross sections in the elastic channel. Korsheninnikov et al. Korsh97 , reported on the measurement of proton elastic scattering by <sup>6</sup>He, <sup>8</sup>He and <sup>11</sup>Li.
One important aspect that we want to investigate is the role of Coulomb dipole polarizability May94 ; May95 , which consists in the effect on the elastic channel of coupling to the break-up channels produced by the Coulomb dipole force. This effect has been investigated for the scattering of other weakly bound nuclei on heavy targets, such as d+ <sup>208</sup>Pb Moro99 , <sup>7</sup>Li + <sup>208</sup>Pb Mar98 , and <sup>11</sup>Li+ <sup>208</sup>Pb May99 . It is found that dipole polarizability gives rise to a significant reduction in the elastic scattering cross sections, which is particularly important for weakly bound nuclei. For the case of <sup>6</sup>He scattering on heavy targets, we expect that the strong Coulomb field generated by the target can distort the <sup>6</sup>He projectile, so that the <sup>4</sup>He core is pushed away from the target nucleus, while the neutrons in the halo remain unaffected.
Another significant aspect is the behaviour of the nuclear optical potential. It is an open question whether the optical model, using reasonable optical potentials, is an adequate approach to describe the elastic scattering of weakly bound nuclei. We are interested in determining the geometry and energy dependence of the optical potential obtained from a fit to the elastic cross sections, including the effect of dipole polarizability, and compare it with the optical potentials obtained for the most similar stable nucleus, which is <sup>6</sup>Li. The elastic scattering of <sup>6</sup>He on <sup>209</sup>Bi has been measured in Notre Dame at energies near Kolataprl and below Kolatarc the Coulomb barrier. The analysis of the elastic scattering and reaction cross sections performed by the authors required the use of large, and energy-dependent, imaginary diffuseness parameters to reproduce the data. This result indicates the presence of a long range absorption mechanism, which could be related to the effect of Coulomb excitation.
We have presented in a previous publication the analysis of elastic scattering of <sup>6</sup>He on <sup>208</sup>Pb at 27 MeV, measured at Louvain la Neuve Kakuee03 . We found evidence that an extremely large diffuseness parameter was required to fit the data. This indicates that long range reaction mechanisms occur when <sup>6</sup>He was scattered on <sup>208</sup>Pb. This is an important result, because it seems to indicate that the elastic scattering induced by exotic nuclei is qualitatively different from the scattering of stable nuclei. Nevertheless, it is important to confirm that this feature (the long range absorption) occurs when <sup>6</sup>He collides with other heavy targets. If this is the case, we could recognise long range absorption as a robust feature of the scattering of <sup>6</sup>He at energies around the barrier, produced by its weakly bound structure, that does not depend strongly on the target properties.
In this work we present new experimental data of the quasi-elastic scattering of <sup>6</sup>He on <sup>197</sup>Au at 27 MeV. The $`1/2^+`$ state of <sup>197</sup>Au at 70 keV was not resolved from the $`3/2^+`$ ground state. We have explored the effect of including explicitly the excitation in a coupled channels calculation, and we find that the quasi-elastic differential cross sections (elastic plus inelastic) in the coupled channels calculation is very similar to the elastic differential cross section in an optical model calculation. The reason for it is that the flux going to the inelastic channels is subtracted from that of the elastic channel, leaving the quasi-elastic differential cross sections unaffected by the coupling.
We perform an analysis of the new set of data in parallel with the analysis of the data of <sup>6</sup>He on <sup>208</sup>Pb at the same energy. Our purpose is to search for evidence of long range absorption in these collisions.
## II Experimental set-up
The experiment was performed using the radioactive beam facility at the Cyclotron Research Centre at Louvain la Neuve in Belgium. The <sup>6</sup>He beam used in this experiment was produced via the <sup>7</sup>Li(p,2p)<sup>6</sup>He reaction using LiF powder target contained in a graphite holder Vervier .
The post-accelerated secondary <sup>6</sup>He beam of 27 MeV energy and intensity of $`3\times `$10<sup>6</sup> ions/s was scattered on a <sup>197</sup>Au target, which was in fact the backing of a <sup>7</sup>LiF target. The thickness of the <sup>197</sup>Au layer was approximately 300 $`\mu g/cm^2`$. The thickness of the <sup>7</sup>LiF layer was 400 $`\mu g/cm^2`$. The reaction products were detected in a detection system consisting of a LEDA and a LAMP type detector described in Tom2000 . LEDA and LAMP silicon strip detector arrays cover two different angular ranges from 6-15 and from 23-72 in the laboratory frame, respectively. The details of the experimental setup has been described elsewhere Kakuee03 .
Both the energy and the time of flight with respect to the cyclotron pulse were recorded for each reaction product. The in-beam energy resolution for silicon strip detectors was around 120-140 keV, depending on the scattering angle, which is mainly due to the beam emittance and beam straggling in the target. The timing information in connection with energy spectra enabled us to unambiguously identify the elastic scattering events.
The elastic scattering of <sup>6</sup>He from <sup>197</sup>Au could be readily separated from <sup>6</sup>He scattered from <sup>7</sup>Li and <sup>19</sup>F at angles greater than 10 degrees. Figure 1 shows the separation of elastic scattering events corresponding to <sup>197</sup>Au from those on <sup>7</sup>Li and <sup>19</sup>F.
There is no evidence of break-up into <sup>4</sup>He at forward angles. At larger angles however there is some evidence for a broad break up component around 2/3 of the elastic peak energy. In the present paper no attempt has been made to extract information on break up events, and only the elastic scattering events are investigated.
Only statistical uncertainties are considered in this analysis. The main source of systematic errors in our setup comes from possible uncertainties in the precise position of the LAMP array with respect to the target. This would lead to uncertainties in the scattering angle, which affect the ratio of the measured cross sections to Rutherford cross sections.
In a previous work Kakuee03 we found that this uncertainty in the position of LAMP affected mainly the relative normalisation of the small angles covered by LEDA and the large angles covered by LAMP. There, we found that minor changes in the positioning of LAMP ($`\pm 3`$ mm), affected the values of the long range absorption. As the small angles covered by LEDA only give information on the global normalisation, and this normalisation, for the intermediate angles (20 to 75 degrees) which are of interest, has systematic uncertainties, we have chosen in this work to neglect the data at small scattering angles, and adjust the normalisation of the LAMP data to the theoretical calculations.
We have also considered the effect of adding the number of counts measured in the different strips of the detector. Due to the geometry of LAMP array, we have six strips which correspond to the same nominal scattering angle. So, one could add the counts of all these strips, to reduce statistical uncertainties, as well as the number of data points. The result of applying this procedure is the “averaged data” set, shown in figure 2. This procedure is adequate to visualise the trend of the data. However, we consider that this procedure could hide the presence of some systematic uncertainties, related to the different solid angles of the sectors of the LAMP array, as well as the different scattering angles corresponding to the strips in different sectors, due to beam misalignment. So, we have also considered the “raw data” set, in which one experimental data point is associated to each strip. In the “raw data” set, the effect of beam misalignment is taken into account Kakuee03 , so that the scattering angles corresponding to strips of different sectors are slightly different from the common nominal scattering angle. The set of data points so obtained, with more data and higher statistical uncertainties, are shown in figure 3. The statistical significance of fits to the “averaged data” set would be the same than that of the “raw data” set, if the difference in counting rates of the different sectors was purely statistical, and the difference between the actual scattering angle of each strip and the common nominal scattering angle was negligible. As this hypothesis is not neccesarily true, we prefer to determine the optical potential parameters from a $`\chi ^2`$ minimisation using the “raw data” set, which is directly related to the experimental measurements. However, to describe the qualitative features of the data is more adequate to use the “averaged data” of figure 2.
## III Theoretical calculations
Dipole Coulomb excitation to break-up states can play an important role in the scattering of weakly bound nuclei. To describe the effect of this reaction mechanism in the elastic scattering, one can make use of a dynamic polarisation potential May94 . The form of the polarisation potential is obtained in a semi-classical framework requiring that the second order amplitude for the dipole excitation-deexcitation process and the first order amplitude associated with the polarisation potential are equal for all classical trajectories corresponding to a given scattering energy. This leads to an analytic formula for the polarisation potential for a single excited state May94 . The expression so obtained can be generalised for the case of excitation to a continuum of break-up states May95 giving rise to the following formula:
$`U_{pol}(r)`$ $`=`$ $`{\displaystyle \frac{4\pi }{9}}{\displaystyle \frac{Z_t^2e^2}{\mathrm{}v}}{\displaystyle \frac{1}{(ra_o)^2r}}`$
$`{\displaystyle _{ϵ_b}^{\mathrm{}}}dϵ{\displaystyle \frac{dB(E1,ϵ)}{dϵ}}(g({\displaystyle \frac{r}{a_o}}1,\xi )+if({\displaystyle \frac{r}{a_o}}1,\xi ),)`$
where *g* and *f* are analytic functions defined as
$`f(z,\xi )`$ $`=`$ $`4\xi ^2z^2\mathrm{exp}(\pi \xi )K_{2i\xi }^{\prime \prime }(2\xi z)`$ (2)
$`g(z,\xi )`$ $`=`$ $`{\displaystyle \frac{P}{\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{f(z,\xi ^{})}{\xi \xi ^{}}}𝑑\xi ^{}`$ (3)
and $`\xi =\frac{ϵa_o}{\mathrm{}v}`$ is the Coulomb adiabaticity parameter corresponding to the excitation energy $`ϵ`$ of the nucleus.
When one compares the Coulomb dipole polarisation potentials for the collisions <sup>6</sup>He + <sup>208</sup>Pb, and <sup>6</sup>He + <sup>197</sup>Au, at the same energy of 27 MeV, one would expect that the larger charge of <sup>208</sup>Pb would produce larger dipole polarisation effects for <sup>6</sup>He + <sup>208</sup>Pb than for <sup>6</sup>He +<sup>197</sup>Au (note the factor $`Z_t^2`$ in eq. III). However, the smaller charge of the <sup>197</sup>Au target also makes the adiabaticity parameter $`\xi `$ smaller. This means that the imaginary part of the polarisation potential at long distances, describing long range absorption due to Coulomb break-up, is actually larger for <sup>197</sup>Au than for <sup>208</sup>Pb.
The elastic differential cross sections are analysed assuming the validity of the optical model. The potential that describes the interaction between <sup>6</sup>He and <sup>208</sup>Pb is the sum of a monopole Coulomb potential, a dipole Coulomb polarisation potential and a phenomenological nuclear potential. The monopole Coulomb potential is determined by the charges of the colliding nuclei. Its only parameter is a Coulomb radius which, when taken in a reasonable range, does not affect significantly the cross sections. The dipole Coulomb polarisation potential describes the effect of coupling the ground state to break-up states in the continuum by the dipole Coulomb force. It is a complex, long range and energy dependent potential, which is completely determined from the values of the B(E1) distribution of <sup>6</sup>He. The phenomenological nuclear potential includes the “direct” term of the nuclear interaction between <sup>6</sup>He and <sup>208</sup>Pb as well as the dynamic effects of nuclear coupling to break-up states, quadrupole Coulomb coupling and Coulomb-nuclear interference terms.
For our analysis we take theoretical values of the B(E1) distribution of <sup>6</sup>He ian ; cobis . This determines completely the dipole Coulomb polarisation potential. The real part of the dipole polarization potential is plotted as the dashed-dotted line in Fig. 6, while the imaginary polarization potential is plotted as the dashed line in Fig 5. For the nuclear potential, as a starting point, we use a Woods-Saxon potential, which was obtained from the optical model analysis of elastic scattering of <sup>6</sup>Li, in mass range of 24-208 and energy range of 13-156 MeV Cook82 . We shall refer to this potential as “Cook” potential. The optical model parameters used are given in table 1.
Our starting consideration is that <sup>6</sup>Li and <sup>6</sup>He have similar structures, and both are weakly bound (by 1.475 MeV and 0.975 MeV, respectively). So, the main qualitative difference between them could be that the dipole Coulomb force can break up <sup>6</sup>He into <sup>4</sup>He + 2n, but it cannot break up <sup>6</sup>Li into <sup>4</sup>He + <sup>2</sup>H. The dipole Coulomb operator, in an $`N=Z`$ nucleus, is an isospin 1 operator. Since <sup>6</sup>Li, <sup>4</sup>He and <sup>2</sup>H, have isospin 0 in their ground states, it is not possible that the dipole Coulomb force breaks up <sup>6</sup>Li into <sup>4</sup>He + <sup>2</sup>H.
This difference between <sup>6</sup>Li and <sup>6</sup>He is explicitly taken into account by means of the dipole Coulomb polarisation potential. In figure 2 we present the “averaged data” measured for the collision of <sup>6</sup>He on <sup>197</sup>Au and <sup>208</sup>Pb at 27 MeV bombarding energy. The solid lines are the global optical model calculations which uses Cook potential Cook82 . These calculations, that do not take into account the effect of dipole polarizability, show a well defined rainbow around 43 degrees, which is clearly absent from the experimental data. The calculations including dipole polarizability (dotted lines) also show a rainbow, but it is less pronounced and it appears at a smaller angle. So, we can conclude that the effect of dipole polarizability is clearly visible in the scattering of <sup>6</sup>He on <sup>197</sup>Au and <sup>208</sup>Pb at 27 MeV, and explains in part the disappearance of the rainbow in the experimental data. However, it is also clear that the optical potential obtained from <sup>6</sup>Li scattering (Cook potential), even supplemented with the Dipole Polarisation Potential, is not adequate to reproduce the scattering of <sup>6</sup>He on <sup>197</sup>Au and <sup>208</sup>Pb at 27 MeV.
We now proceed to modify Cook <sup>6</sup>Li optical potential to fit the <sup>6</sup>He scattering data. The first argument is that the reaction channels produced by <sup>6</sup>He scattering (mostly break-up and neutron transfer), will be different from those of <sup>6</sup>Li. These reaction channels affect mainly the imaginary part of the potential, which describes the loss of flux from the elastic channel. So, we allowed the depth and the diffuseness of the potential to vary. The results of these calculations are shown in tables 1 and 2. The fit of the data is fair, but not perfect, as it can be seen from the values of $`\chi ^2`$. If the value of the depth of the real potential is also fitted, then the fit of the data is very good, as it is shown in tables 1 and 2, as well as in figure 3.
In fitting the data, we have taken into account that the data of <sup>6</sup>He on <sup>208</sup>Pb are much more accurate than those of <sup>6</sup>He on <sup>197</sup>Au. So, whenever possible, we kept the same parameters for the two targets. We only allowed to vary the depth of the imaginary potential, because the targets <sup>208</sup>Pb and <sup>197</sup>Au could lead to different reaction channels. The parameters that fit the experimental data of the reaction <sup>6</sup>He on <sup>197</sup>Au, shown in table 2 are similar to those obtained for the reaction <sup>6</sup>He on <sup>208</sup>Pb, shown in table 1. So are the values of $`\chi ^2`$ obtained in the best fits. We can conclude that both sets of data give a consistent message, indicating the presence of long range reaction mechanisms.
It is interesting to comment on the values of the reaction cross sections obtained in the fits. These are around 1900 mb in all the calculations that fit accurately the elastic data. From the reaction cross sections, we have extracted the contribution due to the imaginary part of the Coulomb polarisation potential. This value is an estimate of the reaction cross section which is due to Coulomb break-up. Note that this value is 318 mb for the <sup>208</sup>Pb target, and 564 mb for the <sup>197</sup>Au target. Naively, one would expect that the target with higher charge would induce more Coulomb break-up. However, the lower charge of <sup>197</sup>Au makes the energy of the collision (27 MeV) higher with respect to the Coulomb barrier, reducing the collision time, and thus producing more Coulomb break-up. Also, it should be mentioned that Coulomb break-up leads indeed to long range absorption. This can be seen from the values of $`L_{dp}`$, which is the average L-value of the reaction cross sections due to coulomb excitation. They are considerably larger than the values of $`L_T`$, which is the average L-value of the total reaction cross sections.
The results of the optical model fits discussed above, which have been performed with the code FRESCO FRESCO , are shown in figure 3, and the sets of optical model parameters are shown in table 1 (<sup>208</sup>Pb) and table 2 (<sup>197</sup>Au). We have also performed calculations varying the value of $`a_i`$, and fitting the value of $`W`$ for each $`a_i`$. The values of $`\chi ^2/n`$ for these fits are plotted in Figure 4. The data on the <sup>208</sup>Pb target indicate clearly the need for large imaginary diffuseness parameters, to obtain good fits $`\chi ^2/n1.5`$. However, the data on the <sup>197</sup>Au target can be fitted with the same accuracy with almost any diffuseness parameter. This is an effect of the larger statistical uncertainties of the <sup>197</sup>Au data. In any case, we can say that the large imaginary diffuseness parameter required to reproduce the data on the <sup>208</sup>Pb target is not inconsistent with the values required to reproduce the data on the <sup>197</sup>Au target. Also, the fits show that the explicit inclusion of the Dipole Polarisation Potential reduce the values of the diffuseness parameters required to fit the data in both cases.
The fits presented in tables 1 and 2, and in figure 4 make use of the “raw data” set, which are shown in figure 3. We have also performed the optical potential fits making use of the “averaged data” set, presented in figure 2. We find that the results for the potentials that produce the best fits are very similar in both cases. This indicates that, in these experiments, there were not systematic differences between the cross sections obtained from the detector strips of the different sectors, and that the difference between actual and nominal scattering angles was not important, for the observable considered.
## IV Discussion
One objective of this analysis is to investigate if, as suggested in our previous paper Kakuee03 , there are evidences of long range mechanisms that lead to the loss of flux in the elastic channel at kinematic conditions that suggest the nuclei are far beyond the strong absorption radius, which, in our case, is approximately $`R_{SA}=11.5fm`$ for both targets.
The other objective is to investigate the role of coulomb dipole polarizability in the scattering of this weakly bound, and hence easily polarizable, nucleus.
### IV.1 Evidence for long range absorption mechanisms
The first evidence for the long range absorption comes from the values of the imaginary diffuseness parameters required to fit the data. The plot of $`\chi ^2`$ (figure 4), for the <sup>6</sup>He + <sup>208</sup>Pb, clearly indicates that the imaginary potential needed to fit the data has a diffuseness considerably larger than that of the real potential. A diffuseness parameter in the range of 1.30 to 1.90 fm is needed to fit the data. This is to be compared to the value of the diffuseness of the real potential, which is $`a_r=0.811`$ fm. We point out that this diffuseness is not as large as the one found in the previous work Kakuee03 , but there the diffuseness of the real potential was allowed to vary, and it also acquired large values. The systematics of our calculations show that the imaginary potential has to be much more diffuse than the real potential to reproduce the data.
The second evidence for the long range absorption comes from the relatively large values of the average reaction angular momenta $`<L>_T`$. In a cutoff model, which is not unreasonable to describe the absorption in heavy ion collisions, the absorption is maximum for $`L+1/2<\lambda `$ and negligible for $`L+1/2>\lambda `$. The relation of the cutoff parameter $`\lambda `$ and the reaction cross section is given by satchler $`\sigma _R=\pi \lambda ^2/k^2`$. Then, in this sharp cutoff model, $`L_T+1/2=2/3\lambda `$. The values of the reaction cross sections in these reactions (see tables 1,2) are about 1900 mb. This leads to values of $`L_T13.5`$, which is considerably smaller than the values shown in the table. This is an indication of the fact that the reaction cross sections extend to up to values of $`L`$ which are well beyond the grazing angular momentum.
The third evidence for the long range absorption comes from the values of the imaginary potential. In figure 5 different families of imaginary potentials which fit reasonably the data are plotted as a function of the distance. Let us focus on the solid lines, which correspond to optical model calculations in which the dipole polarisation potential (DPP) is not explicitly included. We see that the lines cross at 13.25 fm (Au) and 14.0 fm (Pb). This indicates the region of sensitivity to the imaginary potential. Note also that the imaginary potential still has sizeable values at distances as large as 16 fm. In figure 6 we present real potentials that fit the data. These potentials cross between 11 and 12 fm, which is the region of the strong absorption radius (11.5 fm). So, we see that while the real potential is determined around the distance of the strong absorption radius, the imaginary potential is determined at much larger distances. This point should be taken into account when investigating the energy dependence of the real and imaginary optical potentials of exotic nuclei.
It should be noticed that these evidences for long range absorption are deduced from a consistent analysis of the scattering data on <sup>197</sup>Au and <sup>208</sup>Pb targets. The data on the <sup>197</sup>Au target, considered separately, are not sufficiently accurate to determine unambiguously the depth and diffuseness of the imaginary potential. However, the other two evidences for long range absorption (large values in the average reaction angular momenta, and sensitivity of the imaginary potential to large distances) come out clearly from the analysis of the <sup>197</sup>Au data. These features are unaffected by the ambiguities of the imaginary potential parameters, as shown in figure 5.
### IV.2 Role of Coulomb dipole polarizability
Having established the presence of long range reaction mechanisms, we will now discuss the relevance of the Coulomb dipole polarizability in this mechanism. For that purpose, we will consider the calculations in which the DPP is explicitly included. In these calculations, the phenomenological imaginary potential describes the absorption produced by dynamical effects different from pure dipole Coulomb excitation.
First, we investigate the change in the values of the imaginary diffuseness parameters. As shown in table 1, the value of the diffuseness parameter when the DPP is explicitly included is 1.46 fm, to be compared with 1.75 fm when it was omitted. This can also be seen in figure 4, where the values of the diffuseness parameters which reduce the values of $`\chi ^2`$ are definitively smaller when the DPP is included. Note that in the $`\chi ^2`$ plot for <sup>6</sup>He + <sup>197</sup>Au, the relevant magnitude is not the absolute minima. In the case of <sup>6</sup>He + <sup>208</sup>Pb, the best fits obtained had $`\chi ^2/n1.5`$. We can argue that, for the less accurate <sup>6</sup>He + <sup>197</sup>Au data, a fit with $`\chi ^2/n1.5`$ already indicates a good fit of the data. Hence, the plots should be interpreted saying that, when the DPP is not explicitly included, the imaginary diffuseness parameter should be less than 1.90 fm, while if the DPP is explicitly included, it should be less than 1.60 fm. This is compatible with the values obtained with the <sup>6</sup>He + <sup>208</sup>Pb data. The data therefore indicate that when the DPP is included, there is a need of long range imaginary potentials, but the range is not as large as when the DPP is not included.
Second, we investigate the values of the reaction cross sections and average L-values produced by the DPP. As shown in tables 1, 2, a significant fraction of the reaction cross section is due to the coulomb dipole excitation mechanism. In addition, the values of the average angular momenta are very large (24.4 and 26.6 respectively for the <sup>197</sup>Au and <sup>208</sup>Pb targets). These values of the angular momenta correspond to distances of closest approach of 14.6 and 14.8 fm, which are well beyond the strong absorption radius. This indicates that coulomb dipole polarizability is indeed a mechanism that generates long range absorption. This produces a reduction of the elastic cross sections at forward angles, which is associated with the disappearance of the rainbow previously discussed. It should be noticed that the absorption cross section due to the effect of dipole polarizability does not increase with the charge of the target. This is due to the effect of the adiabaticity parameter $`\xi `$, which is larger for the <sup>208</sup>Pb target than for the <sup>197</sup>Au target. This means that, although the coulomb dipole force is weaker in the <sup>197</sup>Au target, it is more effective in producing break-up cross sections, which generate absorption in the elastic channel.
Third, we investigate the values of the potentials as a function of the distance. Consider the dashed lines in figure 5. They represent the imaginary potentials which describe absorption by mechanisms different from the dipole polarizability. They cross at distances of 12.5 (Au) and 13.25 (Pb) fm, which are considerably smaller than the crossings of the full lines, where the DPP is not explicitly considered. So, we see explicitly that a long range absorption mechanism is required in addition to the pure dipole coulomb excitation. However, the range of the imaginary potential associated to this additional mechanism is not as large as the imaginary DPP, which is shown by the dotted line.
The DPP has a real component, attractive, which is shown by the dot-dashed line. This potential has a very long range, but its effect on the scattering is determined mainly by its value at the strong absorption radius. As shown in figure 6, all the calculations that fit the data, either with or without the explicit inclusion of the DPP, have values of the real potential which are about 0.90 MeV at 11.5 fm. However, the long range attraction enhances the absorption by the imaginary potentials. This explains that the sum of the DPP imaginary potential and the phenomenologic potentials given by the dashed lines in figure 5 are smaller than the full lines, which represent the phenomenologic imaginary potentials where DPP is not considered.
We can conclude that the effect Coulomb Dipole Polarizability accounts for an important fraction of the long range absorption needed to describe the elastic scattering of <sup>6</sup>He on <sup>208</sup>Pb and <sup>197</sup>Au at 27 MeV. However, other reaction mechanisms such as nuclear break-up, coulomb-nuclear interference effects or neutron transfer to bound or unbound states can play also an important role.
We consider that a proper understanding of the long range absorption is required. The use of nuclear reactions as an spectroscopic tool to investigate the structure of exotic nuclei requires a deep understanding of the reactions induced by exotic nuclei. This work indicates that simple preconceptions based on the experience of the optical model on stable nuclei, such as the role of the strong absorption radius, are not extrapolatable to the scattering of exotic nuclei.
Experimental measurements of the elastic scattering of <sup>6</sup>He and other exotic nuclei on a variety of targets, along with the measurement of the main reaction channels, would be required. Reaction calculations with a proper treatment of the break-up channels would be needed to understand the role of absorption.
Acknowledgements:
The authors thank the staff at the Cyclotron Research Centre accelerator facility in Louvain-la-Neuve, Belgium, for providing us with an intense and good quality radioactive beam. ORK acknowledges a grant from the Iranian government and support from the University of Sevilla for his stay at this university. PJW, TD, ACS, AL, ADP and WBS would like to acknowledge support from the British EPSRC. AMSB, IMB, MVA, AM, MAGA and JGC would like to acknowledge support from the Spanish MCyT under projects FPA2003-05958 and FPA2005-04460. AMSB acknowledges a research grant from the Spanish MCyT. |
warning/0507/nlin0507026.html | ar5iv | text | # 1 . Introduction.
## 1 . Introduction.
More than one century ago Painlevé and his school discovered six irreducible second-order equations, which determine new transcendental functions. Over a long period of time these functions seemed to have no physical applications, but now these equations are widely used for description of different physical processes. At the present the problem of analysis of the higher analogs to the Painlevé equations has appeared. There is a lot of works, devoted to the solution of this problem . One of the fourth-order analogs of the first Painlevé equation is equation
$$w_{zzzz}+18ww_{zz}+9w_z^2+24w^3=z$$
(1.1)
In papers it was shown that equation (1.1) has properties, that are typical for the Painlevé equations $`P_1÷P_6`$. Equation (1.1) belongs to the class of exactly solvable equations, as it has Lax pair and a lot of other typical properties of the exactly solvable equations. However it does not have the first integrals in the polynomial form, that is one of the features of the Painlevé equations. Equation (1.1) seems to determine new transcendental functions just as equations $`P_1÷P_6`$, although the rigorous proof of the irreducibility of equation (1.1) is now the open problem.
Thereupon the study of all the asymptotic forms and power expansions of equation (1.1) is the important stage of the analysis of this equation, as this fact indirectly confirms the irreducibility of equation (1.1).
This equation does not have the exact solutions, and so it is very important to find the asymptotic forms and the power expansions of the solution of this equation, that is the aim of this work.
Let us find all the power expansions for the solution of equation (1.1) in the form of
$$w(z)=c_rz^r+\underset{s}{}c_sz^s$$
(1.2)
at $`z\mathrm{\hspace{0.17em}0}`$, then $`\omega =1`$, $`s>r`$ and at $`z\mathrm{}`$, then $`\omega =1`$, $`s<r`$.
For that we use the power geometry method .
The outline of this paper is as follows. Section 2 is devoted to the general properties of equation (1.1). In sections 3–5 the expansions near $`z=0`$ are found. In sections 6–9 the power expansion near $`z=\mathrm{}`$ and its exponential additions are obtained.
## 2 . The general properties of equation (1.1).
Let us consider the fourth-order equation (1.1)
$$f(z,w)\stackrel{def}{=}w_{zzzz}+18ww_{zz}+9w_z^2+24w^3z=0$$
(2.1)
For monomials of equation (2.1) we have points $`M_1=(4,1),M_2=(2,2),M_3=(2,2),M_4=(0,3),M_5=(1,0)`$.
The carrier of equation is defined by four points $`Q_1=M_1`$, $`Q_2=M_4`$, $`Q_3=M_5`$ and $`Q_4=M_2=M_3`$. Their convex hull $`\mathrm{\Gamma }`$ is the triangle (fig. 1).
This triangle has apexes $`Q_j(j=1,2,3)`$ and edges $`\mathrm{\Gamma }_1^{(1)}=[Q_3,Q_1],\mathrm{\Gamma }_2^{(1)}=[Q_1,Q_2],\mathrm{\Gamma }_3^{(1)}=[Q_2,Q_3]`$
Outward normal vectors $`N_j(j=1,2,3)`$ of edges $`\mathrm{\Gamma }_j^{(1)}(j=1,2,3)`$ are determined by vectors
$$N_1=(1,5),N_2=(1,2),N_3=(3,1)$$
(2.2)
The normal cones $`U_j^{(1)}`$ to edges $`\mathrm{\Gamma }_j^{(1)}`$ are
$$U_j^{(1)}=\mu N_j,\mu >0,j=1,2,3$$
(2.3)
They and the normal cones $`U_j^{(0)}`$ of apexes $`\mathrm{\Gamma }_j^{(0)}=Q_j(j=1,\mathrm{\hspace{0.17em}2},\mathrm{\hspace{0.17em}3})`$ are represented at fig. 2.
We can choose the basis of the lattice of the carrier of equation (2.1) as
$$B_1=(5,\mathrm{\hspace{0.17em}1}),B_2=(3,\mathrm{\hspace{0.17em}2})$$
(2.4)
Let us study solutions, corresponding to the bounds $`\mathrm{\Gamma }_j^{(d)},d=0,1;j=1,\mathrm{\hspace{0.17em}2},\mathrm{\hspace{0.17em}3}`$ in view of the reduced equations, conforming to apexes $`\mathrm{\Gamma }_j^{(0)}(j=1,\mathrm{\hspace{0.17em}2},\mathrm{\hspace{0.17em}3})`$
$$\widehat{\text{f}}_1^{(0)}\stackrel{def}{=}w_{zzzz}=0$$
(2.5)
$$\widehat{\text{f}}_2^{(0)}\stackrel{def}{=}24w^3=0$$
(2.6)
$$\widehat{\text{f}}_3^{(0)}\stackrel{def}{=}z=0$$
(2.7)
and reduced equations, conforming to edges $`\mathrm{\Gamma }_j^{(1)}(j=1,2,3)`$
$$\widehat{\text{f}}_1^{(1)}\stackrel{def}{=}w_{zzzz}z=0$$
(2.8)
$$\widehat{\text{f}}_2^{(1)}\stackrel{def}{=}w_{zzzz}+18ww_{zz}+9w_z^2+24w^3=0$$
(2.9)
$$\widehat{\text{f}}_3^{(1)}\stackrel{def}{=}24w^3z=0$$
(2.10)
Note, that the reduced equations (2.6) and (2.7) are the algebraic ones. According to they do not have non-trivial power or non-power solutions.
## 3 . Solutions, corresponding to apex $`Q_1`$.
Apex $`Q_1=(4,1)`$ is corresponded to reduced equation (2.5).
Let us find the reduced solutions
$$w=c_rz^r,c_r0$$
(3.1)
for $`\omega (1,r)U_1^{(0)}`$.
Since $`p_1<0`$ in the cone $`U_1^{(0)}`$, then $`\omega =1,z0`$ and the expansions are the ascending power series of $`z`$. The dimension of the bound $`d=0`$, therefor
$$g(z,w)=w^4w^1w_{zzzz}$$
(3.2)
We get the characteristic polynomial
$$\chi (r)\stackrel{def}{=}g(z,z^r)=r(r1)(r2)(r3)$$
(3.3)
Its roots are
$$r_1=0,r_2=1,r_3=2,r_4=3$$
(3.4)
Let us explore all these roots.
The root $`r_1=0`$ is corresponded to vector $`R=(1,0)`$ and vector $`\omega RU_1^{(0)}`$.
We obtain the family $`_1^{(1)}\mathrm{\hspace{0.17em}1}`$ of reduced solutions $`w=c_0`$, where $`c_00`$ is arbitrary constant and $`\omega =1`$. The first variation of equation (2.5)
$$\frac{\delta \widehat{\text{f}}_1^{(0)}}{\delta w}=\frac{d^4}{dz^4}$$
(3.5)
gives operator
$$(z)=\frac{d^4}{dz^4}0$$
(3.6)
Its characteristic polynomial is
$$\nu (k)=z^{4k}(z)z^k=k(k1)(k2)(k3)$$
(3.7)
Equation
$$\nu (k)=0$$
(3.8)
has four roots
$$k_1=0,k_2=1,k_3=2,k_4=3$$
(3.9)
As long as $`\omega =1`$ and $`r=0`$, then the cone of the problem is
$$𝒦=\{k>0\}$$
(3.10)
It contains the critical numbers $`k_2=1,k_3=2`$ and $`k_3=3`$. Expansions for the solutions, corresponding to reduced solution (3.1) can be presented in the form
$$w=c_0+c_1z+c_2z^2+c_3z^3+\underset{k=4}{\overset{\mathrm{}}{}}c_kz^k$$
(3.11)
where all the coefficients are constants, $`c_00,c_1,c_2`$, $`c_3`$ are arbitrary ones and $`c_k(k4)`$ are uniquely defined. Denote this family as $`\text{G}_1^{(0)}1`$. Expansion (3.11) with taking into account eight terms is
$$\begin{array}{c}w\left(z\right)=c_0+c_1z+c_2z^2+c_3z^3\left(\frac{3}{2}c_0c_2+c_{0}^{}{}_{}{}^{3}+\frac{3}{8}c_{1}^{}{}_{}{}^{2}\right)z^4+\\ +\left(\frac{1}{120}\frac{9}{10}c_0c_3\frac{3}{5}c_1c_2\frac{3}{5}c_{0}^{}{}_{}{}^{2}c_1\right)z^5+\\ +\left(\frac{1}{40}c_0c_{1}^{}{}_{}{}^{2}\frac{1}{5}c_{2}^{}{}_{}{}^{2}\frac{9}{20}c_1c_3+\frac{7}{10}c_{0}^{}{}_{}{}^{2}c_2+\frac{3}{5}c_{0}^{}{}_{}{}^{4}\right)z^6+\\ +\left(\frac{3}{10}c_{0}^{}{}_{}{}^{2}c_3\frac{1}{280}c_0+\frac{3}{5}c_0c_1c_2+\frac{3}{5}c_{0}^{}{}_{}{}^{3}c_1\frac{3}{10}c_2c_3+\frac{1}{10}c_{1}^{}{}_{}{}^{3}\right)z^7+\mathrm{}\end{array}$$
Let us explore root $`r_2=1`$. The cone of the problem is $`𝒦=\{k>1\}`$. It contains the critical numbers $`k_2=2,k_3=3`$. The expansion of solution, corresponding to the reduced solution
$$_1^{(1)}2:w=c_1z$$
can be written as
$$w(z)=c_1z+c_2z^2+c_3z^3+\underset{k=4}{\overset{\mathrm{}}{}}c_kz^k$$
(3.12)
where $`c_10,c_2`$ and $`c_3`$ are the arbitrary constants. Denote this family as $`\text{G}_1^{(0)}2`$. The expansion of solutions (3.12) with taking into account seven terms is
$$\begin{array}{c}w(z)=c_1z+c_2z^2+c_3z^3\frac{3}{8}c_{1}^{}{}_{}{}^{2}z^4+\left(\frac{1}{120}\frac{3}{5}c_1c_2\right)z^5\\ \left(\frac{1}{5}c_{2}^{}{}_{}{}^{2}+\frac{9}{20}c_1c_3\right)z^6+\left(\frac{1}{10}c_{1}^{}{}_{}{}^{3}\frac{3}{10}c_2c_3\right)z^7+\mathrm{}\end{array}$$
For root $`r_2=2`$ the cone of the problem is $`𝒦=\{k>2\}`$. The critical number is $`k_3=3`$. The expansion of the solutions, corresponding to the reduced solution
$$_1^{(1)}3:w=c_2z^2$$
takes the form
$$w=c_2z^2+c_3z^3+\underset{k=4}{\overset{\mathrm{}}{}}c_kz^k$$
(3.13)
Denote this family as $`\text{G}_1^{(0)}3`$. Expansion (3.13) with taking into account eight terms is
$$\begin{array}{c}w(z)=c_2z^2+c_3z^3+\frac{1}{120}z^5\frac{1}{5}c_{2}^{}{}_{}{}^{2}z^6\frac{3}{10}c_2c_3z^7\frac{9}{80}c_{3}^{}{}_{}{}^{2}z^8\\ \frac{1}{630}c_2z^9+\left(\frac{2}{75}c_{2}^{}{}_{}{}^{3}\frac{41}{33600}c_3\right)z^{10}+\mathrm{}\end{array}$$
(3.14)
For root $`r_3=3`$ the cone of the problem is $`𝒦=\{k>3\}`$. There is no critical number here. The expansion of solutions, corresponding to the reduced solution, is
$$_1^{(1)}4:w=c_3z^3$$
takes the form
$$w(z)=c_3z^3+\underset{k=4}{\overset{\mathrm{}}{}}c_kz^k$$
(3.15)
Denote this family as $`\text{G}_1^{(0)}4`$. The expansion (3.15) with taking into account four terms is
$$\begin{array}{c}w(z)=c_3z^3+\frac{1}{120}z^5\frac{9}{80}c_{3}^{}{}_{}{}^{2}z^8\frac{41}{33600}c_3z^{10}+\mathrm{}\end{array}$$
(3.16)
The expansions of solutions converge for sufficiently small $`|z|`$. The existence and analyticity of expansions (3.11), (3.12), (3.13) and (3.15) follow from Cauchy theorem.
## 4 . Solutions, corresponding to edge $`\mathrm{\Gamma }_1^{(1)}`$.
Edge $`\mathrm{\Gamma }_1^{(1)}`$ is conformed by the reduced equation
$$\widehat{f}_1^{(1)}(z,y)\stackrel{def}{=}w_{zzzz}z=0$$
(4.1)
Normal cone is
$$U_1^{(1)}=\{\mu (1,5),\mu >0\}$$
(4.2)
Therefor $`\omega =1`$, i.e. $`z0`$ and $`r=5`$. Power solutions are found in the form
$$w=c_5z^5$$
For $`c_5`$ we have
$$c_5=\frac{1}{120}$$
(4.3)
The only power solution is
$$_2^{(1)}1:w=\frac{z^5}{120}$$
(4.4)
Compute the critical numbers. The first variation of (2.8) is
$$\frac{\delta \widehat{f}_1^{(1)}}{\delta w}=\frac{d^4}{dz^4}$$
(4.5)
We get the proper numbers
$$k_1=0,k_2=1,k_3=2,k_4=3$$
(4.6)
The cone of the problem
$$𝒦=\{k>5\}$$
does not consist them.
Solution (4.4) is corresponded to two vector indexes $`\stackrel{~}{Q}_1=(0,1),\stackrel{~}{Q}_2=(5,0)`$. There difference $`B=\stackrel{~}{Q}_1\stackrel{~}{Q}_2=(5,\mathrm{\hspace{0.17em}1})`$ equals to vector $`Q_1Q_2`$. So solution (4.4) is conformed to lattice Z, which consists of points $`Q=(q_1,q_2)=k(3,\mathrm{\hspace{0.17em}2})+m(5,\mathrm{\hspace{0.17em}1})=(3k5l,\mathrm{\hspace{0.17em}2}k+l)`$, where $`k`$ and $`l`$ are whole numbers. Points belong to line $`q_2=1`$, if $`l=12k`$. In this case $`q_1=5+7k`$. As long as the cone of the problem here is $`𝒦=\{k>5\}`$, the set of the carrier of solution expansion K takes the form
$$\text{K}=\{5+7n,n\}$$
(4.7)
Then the expansion of solution can be written as
$$w(z)=z^5\left(\frac{1}{120}+\underset{m=1}{\overset{\mathrm{}}{}}c_{5+7m}z^{7m}\right)$$
(4.8)
Expansion (4.8) with taking into account three terms is
$$w(z)=\frac{z^5}{120}\left(1\frac{13}{31680}z^7+\frac{601}{4911667200}z^{14}+\mathrm{}\right)$$
(4.9)
Equation (4.1) does not have exponential additions and non-power asymptotic forms.
## 5 . Solutions, corresponding to edge $`\mathrm{\Gamma }_2^{(1)}`$.
Edge $`\mathrm{\Gamma }_2^{(1)}`$ is corresponded to the reduced equation
$$\widehat{f}_2^{(1)}(z,w)\stackrel{def}{=}w_{zzzz}10ww_{zz}5w^2+10w^3=0$$
(5.1)
The normal cone is
$$U_2^{(1)}=\{\mu (1,2),\mu >0\}$$
(5.2)
Therefor $`\omega =1`$, i.e. $`z0`$ and $`r=2`$. Hence the solution of equation (5.1) we can find in the form
$$w=c_2z^2$$
(5.3)
For $`c_2`$ we have the determining equation
$$c_2^2+6c_2+5=0$$
(5.4)
Consequently we get
$$c_2^{(1)}=1,c_2^{(2)}=5$$
(5.5)
The reduced solutions are
$$_2^{(1)}1:w=z^2$$
(5.6)
$$_2^{(1)}2:w=5z^2$$
(5.7)
Let us compute the corresponding critical numbers. The first variation is
$$\frac{\delta f_2^{(1)}}{\delta w}=\frac{d^4}{dz^4}+18w_{zz}+18w\frac{d^2}{dz^2}+18w_z\frac{d}{dz}+72w^2$$
(5.8)
Applied to solution (5.6), it produces operator
$$^{(1)}(z)=\frac{d^4}{dz^4}\frac{18}{z^2}\frac{d^2}{dz^2}+\frac{36}{z^3}\frac{d}{dz}\frac{36}{z^4}$$
(5.9)
which is corresponded by the characteristic polynomial
$$\nu (k)=k^46k^37k^2+48k36$$
(5.10)
Equation
$$\nu (k)=0$$
(5.11)
has the roots
$$k_1=3,k_2=1,k_3=2,k_4=6$$
(5.12)
With reference to solution (5.7) variation (5.8) gives operator
$$^{(2)}(z)=\frac{d^4}{dz^4}\frac{90}{z^2}\frac{d^2}{dz^2}+\frac{180}{z^3}\frac{d}{dz}+\frac{1260}{z^4}$$
(5.13)
which is corresponded by the characteristic polynomial
$$\nu (k)=k^46k^379k^2+264k+1260$$
(5.14)
with roots
$$k_1=7,k_2=3,k_3=6,k_4=10$$
(5.15)
The cone of the problem here is
$$𝒦=\{k>2\}$$
(5.16)
Therefor for the reduced solution (5.6) three critical numbers belong to the cone, and there are two critical numbers for the reduced solution (5.7) in the cone of the problem.
The set of the carriers of the solution expansions K can be written as
$$\text{K}=\{2+7n,n\}$$
(5.17)
Sets $`\text{K}(0)`$, $`\text{K}(0,3)`$ and $`\text{K}(0,3,6)`$ are
$$\begin{array}{c}\text{K}(1)=\{2+7n+3m,n,m,n+m\mathrm{\hspace{0.17em}0}\}=\\ =\{2,1,4,5,7,8,10,\mathrm{}\}\end{array}$$
(5.18)
$$\begin{array}{c}\text{K}(1,2)=\{2+7n+3m+4k,n,m,k,m+n+k\mathrm{\hspace{0.17em}0}\}=\\ =\{2,1,2,4,5,6,7,8,\mathrm{}\}\end{array}$$
(5.19)
$$\begin{array}{c}\text{K}(1,2,6)=\{2+7n+3m+4k+8l,n,m,k,l,m+n+k+l\mathrm{\hspace{0.17em}0}\}=\\ =\{2,1,2,4,5,6,7,8,\mathrm{}\}\end{array}$$
(5.20)
In this case the expansion for the solution of equation can be represented as
$$\begin{array}{c}w(z)=\frac{2}{z^2}+\underset{n+m+k+l>0}{}c_{2+7n+3m+4k+8l}z^{2+7n+3m+4k+8l}\end{array}$$
(5.21)
Denote this family as $`G_2^11`$. The critical number $`1`$ does not belong to set K, so the compatibility condition for $`c_1`$ holds automatically and $`c_1`$ is the arbitrary constant. The critical number $`2`$ also does not belong to sets K and $`\text{K}(1)`$, therefor the compatibility condition for $`c_2`$ holds too and $`c_2`$ is the arbitrary constant. But critical number $`6`$ is a member of $`\text{K}(1,2)`$, so it is necessary to verify that the compatibility condition for $`c_6`$ holds and that $`c_6`$ is the arbitrary constant. The calculation shows that in this situation the condition holds and $`c_6`$ is the arbitrary constant too. The three-parameter power expansion of solutions, corresponding to the reduced solution (5.6) takes the form
$$\begin{array}{c}w(z)=\frac{1}{z^2}+c_1z+c_2z^2\frac{3}{4}c_{1}^{}{}_{}{}^{2}z^4\left(\frac{3}{4}c_1c_2+\frac{1}{96}\right)z^5+c_6z^6+\\ +\frac{7}{25}c_{1}^{}{}_{}{}^{3}z^7+\frac{1}{4928}c_1\left(1848c_1c_2+17\right)z^8+\left(\frac{1}{8}c_{2}^{}{}_{}{}^{2}c_1\frac{1}{4}c_1c_6+\frac{1}{448}c_2\right)z^9\\ \left(\frac{1}{156}c_{2}^{}{}_{}{}^{3}+\frac{437}{5200}c_{1}^{}{}_{}{}^{4}+\frac{9}{52}c_2c_6\right)z^{10}+\mathrm{}\end{array}$$
(5.22)
The carrier of power expansion, corresponding to reduced solution (5.7), is formed by the sets
$$\begin{array}{c}\text{K}(6)=\{2+7n+8m,n,m,m+n\mathrm{\hspace{0.17em}0}\}=\\ =\{2,5,6,12,13,14,19,20,21,22,27,28,29,30,33,34,35,36,37,38,40,\mathrm{}\}\end{array}$$
(5.23)
$$\begin{array}{c}\text{K}(6,10)=\{2+7n+8m+12k,n,m,k,m+n+k\mathrm{\hspace{0.17em}0}\}=\\ =\{2,5,6,10,12,13,14,17,18,19,20,21,22,24\mathrm{}\}\end{array}$$
(5.24)
The expansion for solution of equation can be written as
$$\begin{array}{c}w(z)=\frac{5}{z^2}+\underset{n+m+k>0}{}c_{2+7n+8m+12k}z^{2+7n+8m+12k}\end{array}$$
(5.25)
Denote this family as $`G_2^12`$. The critical numbers 6 and 10 do not belong to the set K and the number 10 does not belong to the set $`\text{K}(6)`$. For numbers 6 and 10 the compatibility conditions holds automatically, therefor coefficients $`c_6`$ and $`c_10`$ are the arbitrary constants. The two-parameter expansion of solution, corresponding to the reduced solution (5.7), is
$$\begin{array}{c}w(z)=\frac{5}{z^2}+\frac{1}{480}z^5+c_6z^6+c_{10}z^{10}\frac{1}{3502080}z^{12}\\ \frac{1}{4480}c_6z^{13}\frac{3}{68}c_{6}^{}{}_{}{}^{2}z^{14}\frac{3}{24640}c_{10}z^{17}+\mathrm{}\end{array}$$
(5.26)
According to , the expansions of solutions (5.22) and (5.26) do not have power and exponential additions.
## 6 . Solutions, corresponding to edge $`\mathrm{\Gamma }_3^{(1)}`$.
Edge $`\mathrm{\Gamma }_3^{(1)}`$ is corresponded by the reduced equation
$$\widehat{f}_3^{(1)}(z,w)\stackrel{def}{=}24w^3z=0$$
(6.1)
In this case $`\omega =1`$, i.e. $`z0`$ and $`r=1/3`$. The expansions are the descending power series of $`z`$.
Reduced equation (6.1) has three power solutions
$$_3^{(1)}1:w=\phi ^{(1)}(z)=c_{1/3}^{(1)}z^{1/3},c_{1/3}^{(1)}=\frac{1}{2}\sqrt[3]{\frac{1}{3}}$$
(6.2)
$$_3^{(1)}2:w=\phi ^{(2)}(z)=c_{1/3}^{(2)}z^{1/3},c_{1/3}^{(2)}=\frac{1}{4}(1i\sqrt{3})\sqrt[3]{\frac{1}{3}}$$
(6.3)
$$_3^{(1)}3:w=\phi ^{(3)}(z)=c_{1/3}^{(3)}z^{1/3},c_{1/3}^{(3)}=\frac{1}{4}(1+i\sqrt{3})\sqrt[3]{\frac{1}{3}}$$
(6.4)
The shifted carrier of reduced solutions (6.2) – (6.4) gives a vector
$$B=(\frac{1}{3},1)$$
(6.5)
which equals a third of vector $`Q_2Q_1`$. Therefor we explore the lattice, generated by vectors $`Q_3Q_1`$ and $`B`$. The basis of this lattice is $`(3,2)`$ and $`(1/3,1)`$. We have $`Q=(q_1,q_2)=k(3,\mathrm{\hspace{0.17em}2})+m(\frac{1}{3},1)=(3k+m/3,\mathrm{\hspace{0.17em}\hspace{0.17em}2}km)`$, where $`k`$ and $`m`$ are the whole numbers. At the line $`q_2=1`$ we have $`2km=1`$, wherefrom $`m=2k+1`$ and $`q_1=\frac{(17k)}{3}`$. And so the carrier of solution is
$$𝐊=\left\{k=\frac{17n}{3},n\right\}$$
(6.6)
and the expansions of solutions take the form
$$G_3^{(1)}l:w=\phi ^{(l)}(z)=c_{1/3}^{(l)}z^{1/3}+\underset{n=1}{\overset{\mathrm{}}{}}c_{(17n)/3}^{(l)}z^{(17n)/3}$$
(6.7)
Here $`c_{1/3}^{(l)}`$ can be found from reduced solutions (6.2) – (6.4), coefficients $`c_{(17n)/3}^{(l)}`$ are computed sequentially. The calculating of the coefficient $`c_2`$ gives the result $`c_2=1/24`$. The expansion of solution with taking into account five terms is
$$\begin{array}{c}\phi ^{(l)}(z)=c_{1/3}z^{1/3}+\frac{1}{24}z^2\frac{1925}{46656}\frac{1}{c_{1/3}}z^{13/3}+\\ \\ +\frac{509575}{3359232}\frac{1}{c_{1/3}^{}{}_{}{}^{2}}z^{20/3}\frac{445712575}{362797056}\frac{1}{c_{1/3}^{}{}_{}{}^{3}}z^9+\mathrm{}\end{array}$$
(6.8)
The obtained expansions seem to be divergent ones.
## 7 . Exponential additions of the first level.
Let us find the exponential additions to solutions (6.2)-(6.4). We look for the solutions in the form
$$w=\phi ^{(l)}(z)+u^{(l)},l=1,2,3$$
The reduced equation for the addition $`u^{(l)}`$ is
$$M_l^{(1)}(z)u^{(l)}=0$$
(7.1)
where $`M_l^{(1)}(z)`$ is the first variation at the solution $`w=\phi ^{(l)}(z)`$. As long as
$$\frac{\delta f}{\delta w}=\frac{d^{\mathrm{\hspace{0.17em}4}}}{dz^4}+18w_{zz}+18w\frac{d^2}{dz^2}+18w_z\frac{d}{dz}+72w^2$$
(7.2)
then
$$M_l^{(1)}(z)=\frac{d^{\mathrm{\hspace{0.17em}4}}}{dz^4}+18\phi _{zz}^{(l)}+18\phi ^{(l)}\frac{d^{\mathrm{\hspace{0.17em}2}}}{dz^2}+18\phi _z^{(l)}\frac{d}{dz}+72\phi _{}^{(l)}{}_{}{}^{2}$$
(7.3)
Equation (7.1) takes the form
$$\begin{array}{c}\frac{d^4u^{(l)}}{dz^4}+18\phi _{zz}^{(l)}u^{(l)}+18\phi (l)\frac{d^2u^{(l)}}{dz^2}+18\phi _z^{(l)}\frac{du^{(l)}}{dz}+72\phi _{}^{(l)}{}_{}{}^{2}u^{(l)}=0,\\ l=1,2,3\end{array}$$
(7.4)
$$\zeta ^{(l)}=\frac{d\mathrm{ln}u^{(l)}}{dz}$$
(7.5)
then from (7.5) we have
$$\frac{du^{(l)}}{dz}=\zeta ^{(l)}u^{(l)},\frac{d^2u^{(l)}}{dz^2}=\zeta _z^{(l)}u^{(l)}+\zeta _{}^{(l)}{}_{}{}^{2}u^{(l)}$$
$$\frac{d^3u^{(l)}}{dz^3}=\zeta _{zz}^{(l)}u^{(l)}+3\zeta ^{(l)}\zeta _z^{(l)}u^{(l)}+\zeta _{}^{(l)}{}_{}{}^{3}u^{(l)}$$
$$\frac{d^4u^{(l)}}{dz^4}=\zeta _{zzz}^{(l)}u^{(l)}+4\zeta ^{(l)}\zeta _{zz}^{(l)}u^{(l)}+3\zeta _{z}^{(l)}{}_{}{}^{2}u^{(l)}+6\zeta _{}^{(l)}{}_{}{}^{2}\zeta _z^{(l)}u^{(l)}+\zeta _{}^{(l)}{}_{}{}^{4}u^{(l)}$$
By substituting the derivatives
$$\frac{du^{(l)}}{dz},\frac{d^2u^{(l)}}{dz^2},\frac{d^4u^{(l)}}{dz^4}$$
into the equation (7.4) we get the reduced equation in the form
$$\begin{array}{c}u^{(l)}[\zeta _{zzz}^{(l)}+4\zeta ^{(l)}\zeta _{zz}^{(l)}+3\zeta _{z}^{(l)}{}_{}{}^{2}+6\zeta _{}^{(l)}{}_{}{}^{2}\zeta _z^{(l)}+\\ +\zeta _{}^{(l)}{}_{}{}^{4}+18\phi _{zz}^{(l)}+18\phi ^{(l)}\zeta _z^{(l)}+18\phi ^{(l)}\zeta _{}^{(l)}{}_{}{}^{2}+18\phi _z^{(l)}\zeta ^{(l)}+72\phi _{}^{(l)}{}_{}{}^{2}]=0\end{array}$$
(7.6)
Let us find the power expansions for solutions of equation (7.6). The carrier of equation (7.6) consists of points
$$\begin{array}{c}Q_1=(3,1),Q_2=(2,2),Q_3=(1,3),\\ Q_4=(0,4),Q_5=(\frac{1}{3},2),Q_6=(\frac{2}{3},0),Q_7=(\frac{2}{3},1),\\ Q_8=(\frac{5}{3},0),Q_{5,k}=(\frac{17k}{3},2),Q_{6,k}=(\frac{27k}{3},0),\\ Q_{7,k}=(\frac{2+7k}{3},1),Q_{8,k}=(\frac{5+7k}{3},0),k\end{array}$$
(7.7)
The closing of convex hull of points of the carrier of equation (7.6) is the strip. It is represented at fig. 3.
The periphery of the strip contains edges $`\mathrm{\Gamma }_j^{(1)}(j=1,2,3)`$ with normal vectors $`N_1=(6,1),N_2=(0,1),N_3=(0,1)`$. It should take up edge $`\mathrm{\Gamma }_1^{(1)}`$ only . This edge is corresponded by the reduced equation
$$h_1^{(1)}(z,\zeta )\stackrel{def}{=}\zeta ^4+18\phi ^{(l)}\zeta ^2+72\phi _{}^{(l)}{}_{}{}^{2}=0$$
(7.8)
Wherefrom we have
$$\zeta ^2=3\left(3+(1)^m\right)\phi ^{(l)},m=1,2$$
(7.9)
We obtain twelve solutions of equation (7.8)
$$\begin{array}{c}\zeta ^{(l,m,k)}=g_{1/6}^{(l,m,k)}z^{1/6},l=1,2,3;m,k=1,2\end{array}$$
(7.10)
where
$$\begin{array}{c}g_{1/6}^{(l,m,k)}=(1)^k\sqrt{3\left(3+(1)^m\right)c_{1/3}^{(l)}},\\ l=1,2,3;m,k=1,2\end{array}$$
(7.11)
The reduced equation is algebraic one, so it has no critical numbers. Let us compute the carrier of the expansion for solution of equation (7.6). The shifted carrier of equation (7.6) is contained in a lattice, generated by vectors $`B_1=(\frac{7}{3},0),B_2=(1,1)`$. The shifted carrier of solutions (7.10) gives rise to vector $`B_3=(\frac{1}{6},1)`$. The difference $`B_2B_3=(\frac{7}{6},0)=\frac{1}{2}B_1\stackrel{def}{=}B_4`$. Therefor, vectors $`B_1,B_2`$ and $`B_3`$ generate the same lattice as vectors $`B_2,B_4`$. Points of this lattice can be written as
$$Q=(q_1,q_2)=k(1,1)+m(\frac{7}{6},0)=(k+\frac{7m}{6},k)$$
At the line $`q_2=1`$ we have $`k=1`$, and so $`q_1=1+\frac{7m}{6}`$. As long as the cone of the problem here is $`𝒦=\left\{k<\frac{1}{6}\right\}`$, then the set of the carriers of expansions $`𝐊`$ is
$$𝐊=\left\{\frac{17n}{6},n\right\}$$
(7.12)
The expansion for solution of equation (7.6) takes the form
$$\begin{array}{c}\zeta ^{(l,m,k)}=g_{1/6}^{(l,m,k)}z^{1/6}+\underset{n}{}g_{(17n)/6}^{(l,m,k)}z^{(17n)/6},\\ l=1,2,3;m=1,2;k=1,2\end{array}$$
(7.13)
Coefficients $`g_{1/6}^{(l,m,k)}`$ are determined by expression (7.11). Coefficient $`g_1^{(l,m,k)}`$ takes on a value
$$g_1^{(l,m,k)}=\frac{1}{4}$$
(7.14)
The expansion of solution with taking into account four terms takes the form
$$\begin{array}{c}\zeta ^{(l,m,k)}=g_{1/6}z^{1/6}\frac{1}{4}z^1\frac{7}{288}\frac{\left(17g_{1/6}^{}{}_{}{}^{2}+63c_{1/3}\right)}{g_{1/6}\left(g_{1/6}^{}{}_{}{}^{2}+9c_{1/3}\right)}z^{13/6}\\ \frac{49}{1728}\frac{17g_{1/6}^{}{}_{}{}^{4}+36c_{1/3}g_{1/6}^{}{}_{}{}^{2}+567c_{1/3}^{}{}_{}{}^{2}}{g_{1/6}^{}{}_{}{}^{2}\left(g_{1/6}^{}{}_{}{}^{2}+9c_{1/3}\right)^2}z^{10/3}+\mathrm{}\end{array}$$
(7.15)
In view of (7.5) we can find additions $`u^{(l,m,k)}(z)`$. We have
$$u^{(l,m,k)}(z)=C\mathrm{exp}\zeta ^{(l,m,k)}(z)𝑑z$$
Wherefrom we get
$$\begin{array}{c}u^{(l,m,k)}(z)=C_1z^{1/4}\mathrm{exp}\left[\frac{6}{7}g_{1/6}^{(l,m,k)}z^{7/6}+\underset{n=2}{\overset{\mathrm{}}{}}\frac{6}{7(1n)}g_{(17n)/6}^{(l,m,k)}z^{7(1n)/6}\right]\\ l=1,2,3;m=1,2;k=1,2\end{array}$$
(7.16)
Here $`C_1`$ and farther $`C_2`$ and $`C_3`$ are the arbitrary constants. Addition $`u^{(l,m,k)}(z)`$ near $`z\mathrm{}`$ is the exponentially small one in those sectors of complex plane $`z`$, where
$$Re\left[g_{7/6}^{(l,m,k)}z^{1/6}\right]<0$$
(7.17)
Thus for three expansions $`G_3^{(1)}l`$ we get four one-parameter family of additions $`G_3^{(1)}lG_1^1mk`$, where $`m=1,2`$ and $`k=1,2`$.
## 8 . Exponential additions of the second level.
Let us find exponential additions of the second level $`v^{(p)}`$, i.e. the additions to solutions $`u^{(l,m,k)}(z)`$. The reduced equation for addition $`v^{(p)}`$ is
$$M_p^{(2)}(z)v^{(p)}=0$$
(8.1)
where operator $`M_p^{(2)}`$ is the first variation of (7.6). Equation (8.1) for $`v=v^{(p)}`$ takes the form
$$\begin{array}{c}\frac{d^3v}{dz^3}+4\zeta _{zz}v+4\zeta v_{zz}+6\zeta _zv_z+12\zeta \zeta _zv+\\ +6\zeta ^2v_z+4\zeta ^3v+18\phi ^{(l)}v_z+36\phi ^{(l)}\zeta v+18\phi _z^{(l)}v=0\end{array}$$
(8.2)
Assumed that
$$\frac{d\mathrm{ln}v}{dz}=\xi $$
(8.3)
we have
$$\frac{dv}{dz}=\xi v,\frac{d^2v}{dz^2}=\xi _zv+\xi ^2v,\frac{d^3v}{dz^3}=\xi _{zz}v+3\xi \xi _zv+\zeta ^3v$$
(8.4)
From (8.2) we get equation
$$\begin{array}{c}\xi _{zz}+3\xi \xi _z+\xi ^3+4\zeta _{zz}+4\xi _z\zeta +4\xi ^2\zeta +6\xi \zeta _z+12\zeta \zeta _z+6\xi \zeta ^2+\\ +4\zeta ^3+18\phi ^{(l)}\xi +36\zeta \phi ^{(l)}+18\phi _z^{(l)}=0\end{array}$$
(8.5)
Monomials of equation (8.5) are corresponded by the points
$$\begin{array}{c}M_{0,k}=(\frac{1}{2}\frac{7}{6}k,\mathrm{\hspace{0.17em}0}),M_{1,k}=(\frac{1}{3}\frac{7}{6}k,1),\\ M_{2,k}=(\frac{1}{6}\frac{7}{6}k,\mathrm{\hspace{0.17em}2}),M_3=(0,3)\\ k=0,1,2,\mathrm{}\end{array}$$
(8.6)
The carrier of the equation (8.5) is determined by points of the set (8.6). The convex set forms the strip, which is represented at fig. 4. It should examine edge $`\mathrm{\Gamma }_1^{(1)}`$, which is passing through points
$$\begin{array}{c}Q_0=(\frac{1}{2},0),Q_1=(\frac{1}{3},1),Q_2=(\frac{1}{6},2),Q_3=(0,3)\end{array}$$
(8.7)
The reduced equation, corresponding to this edge, is
$$\begin{array}{c}\xi ^3+4\xi ^2\zeta +6\xi \zeta ^2+4\zeta ^3+18\xi \phi ^{(l)}+36\zeta \phi ^{(l)}=0\end{array}$$
(8.8)
The basis of the lattice, corresponding to the carrier of equation (8.5) is
$$\begin{array}{c}B_1=(1,1),B_2=(\frac{7}{6},0)\end{array}$$
The solution of equation (8.8) takes the form
$$\begin{array}{c}\xi ^{(l,m,k,p)}=r_{1/6}^{(l,m,k,p)}z^{1/6},m,k=1,2;l=1,2,3;p=1,2,3\end{array}$$
(8.9)
where $`r=r_{1/6}^{(l,m,k,p)},p=1,2,3`$ are the roots of the equation
$$\begin{array}{c}r^3+4r^2g_{1/6}^{(l,m,k)}+\left(6g_{1/6}^{(l,m,k)}{}_{}{}^{2}+18c_{1/3}^{(l)}\right)r+4g_{1/6}^{(l,m,k)}{}_{}{}^{3}+\\ +36g_{1/6}^{(l,m,k)}c_{1/3}^{(l)}=0\end{array}$$
(8.10)
Equation (8.10) has the roots
$$\begin{array}{c}r_{1/6}^{(l,m,k,1)}=2g_{1/6}^{(l,m,k)},r_{1/6}^{(l,m,k,2)}=g_{1/6}^{(l,m,k)}+\left(18c_{1/3}^{(l)}g_{1/6}^{(l,m,k)}{}_{}{}^{2}\right)^{1/2}\\ r_{1/6}^{(l,m,k,3)}=g_{1/6}^{(l,m,k)}\left(18c_{1/3}^{(l)}g_{1/6}^{(l,m,k)}{}_{}{}^{2}\right)^{1/2}\end{array}$$
(8.11)
The set of carriers of expansions for solution $`𝐊`$ coincides with (7.12). The expansion of solution for $`\xi ^{(l,m,k,p)}`$ takes the form
$$\begin{array}{c}\xi ^{(l,m,k,p)}=r_{1/6}^{(l,m,k,p)}z^{1/6}+\underset{n=1}{\overset{\mathrm{}}{}}r_{(17n)/6}^{(l,m,k,p)}z^{(17n)/6},\\ l=1,2,3;m=1,2;k=1,2;p=1,2,3\end{array}$$
(8.12)
The computing the coefficient $`r_1^{(l,m,k,p)}`$ gives a result $`r_1^{(l,m,k,p)}=1/6`$. The expansion of solution with taking into account three terms is
$$\begin{array}{c}\xi ^{(l,m,k,p)}=r_{1/6}z^{1/6}+\frac{1}{6}z^1+\frac{7}{72}(34g_{1/6}^{}{}_{}{}^{3}r_{1/6}+\\ +36c_{1/3}r_{1/6}g_{1/6}+36c_{1/3}g_{1/6}^{}{}_{}{}^{2}+63c_{1/3}r_{1/6}^{}{}_{}{}^{2}+17g_{1/6}^{}{}_{}{}^{2}r_{1/6}^{}{}_{}{}^{2}+\\ +567c_{1/3}^{}{}_{}{}^{2}+17g_{1/6}^{}{}_{}{}^{4}\left)\right(g_{1/6}\left)^1\right(g_{1/6}^{}{}_{}{}^{2}+9c_{1/3})^1\\ \left(8g_{1/6}r_{1/6}+3r_{1/6}^{}{}_{}{}^{2}+6g_{1/6}^{}{}_{}{}^{2}+18c_{1/3}\right)^1z^{13/6}+\mathrm{}\end{array}$$
(8.13)
The exponential additions $`v^{(l,m,k,p)}(z)`$ to solutions $`u^{(l,m,k)}(z)`$ are
$$\begin{array}{c}v^{(l,m,k,p)}(z)=C_2z^{1/6}\mathrm{exp}\left[\frac{6}{7}r_{1/6}^{(l,m,k,p)}z^{7/6}+\underset{n=2}{\overset{\mathrm{}}{}}\frac{6}{7(1n)}r_{(17n)/6}^{(l,m,k,p)}z^{7(1n)/6}\right],\\ l=1,2,3;m=1,2;k=1,2;p=1,2,3\end{array}$$
(8.14)
## 9 . Exponential additions of the third level.
Let us compute the exponential additions of the third level $`y^{(s)}`$, i.e. the additions to the solutions $`v^{(l,m,k,p)}(z)`$. The reduced equation for addition $`y^{(s)}`$ is
$$M_s^{(3)}(z)y^{(s)}=0$$
(9.1)
Operator $`M_s^{(3)}`$ is the first variation of (8.5). Equation (9.1) for $`y=y^{(l,m,k,p,s)}`$ takes the form
$$\begin{array}{c}y_{zz}+3\xi _zy+3\xi y_z+3\xi ^2y+4\zeta y_z+8\xi \zeta y+\\ +6\zeta _zy+6\zeta ^2y+18\phi ^{(l)}y=0\end{array}$$
(9.2)
Using the substitute
$$\frac{d\mathrm{ln}y}{dz}=\eta $$
(9.3)
we obtain
$$\frac{dy}{dz}=\eta y,\frac{d^2y}{dz^2}=\eta _zy+\eta ^2y$$
(9.4)
From (9.4) we have equation
$$\begin{array}{c}\eta _z+\eta ^2+3\xi _z+3\xi \eta +3\xi ^2+4\eta \zeta +8\xi \zeta +6\zeta _z+6\zeta ^2+18\phi ^{(l)}=0\end{array}$$
(9.5)
Monomials of equation (9.5) are corresponded by points
$$\begin{array}{c}M_{0,k}=(\frac{1}{3}\frac{7}{6}k,\mathrm{\hspace{0.17em}0}),M_{1,k}=(\frac{1}{6}\frac{7}{6}k,\mathrm{\hspace{0.17em}1}),M_2=(0,\mathrm{\hspace{0.17em}2}),\\ k=0,1,2\mathrm{}\end{array}$$
(9.6)
The carrier of equation (9.5) is formed by points (9.6). The convex set forms the strip, which is represented at fig. 5. It should examine edge $`\mathrm{\Gamma }_1^{(1)}`$, which is passing through points
$$\begin{array}{c}Q_0=(\frac{1}{3},0),Q_1=(\frac{1}{6},1),Q_2=(0,2)\end{array}$$
(9.7)
The reduced equation, corresponding to this edge, is
$$\begin{array}{c}\eta ^2+3\xi \eta +3\xi ^2+4\eta \zeta +8\xi \zeta +6\zeta ^2+18\phi ^{(l)}=0\end{array}$$
(9.8)
The solutions of equation (9.8) takes the form
$$\begin{array}{c}\eta ^{(l,m,k,p,s)}=q^{(l,m,k,p,s)}z^{1/6}\\ l=1,2,3;m,k=1,2;p=1,2,3;s=1,2;\end{array}$$
(9.9)
where $`q^{(l,m,k,p,s)}=q`$ are the roots of equation
$$\begin{array}{c}q^2+\left(3r_{1/6}^{(l,m,k,p)}+4g_{1/6}^{(l,m,k)}\right)q+8r_{1/6}^{(l,m,k,p)}g_{1/6}^{(l,m,k)}+\\ +6g_{1/6}^{(l,m,k)}{}_{}{}^{2}+3r_{1/6}^{(l,m,k,p)}{}_{}{}^{2}+18c_{1/3}^{(l)}=0\end{array}$$
(9.10)
The roots of equation (9.10) are
$$\begin{array}{c}q_{1/6}^{(l,m,k,p,s)}=\frac{3}{2}r_{1/6}^{(l,m,k,p)}2g_{1/6}^{(l,m,k)}+\\ +(1)^{s1}\left(\frac{3}{4}r_{1/6}^{(l,m,k,p)}{}_{}{}^{2}2r_{1/6}^{(l,m,k,p)}g_{1/6}^{(l,m,k)}2g_{1/6}^{(l,m,k)}{}_{}{}^{2}18c_{1/3}^{(l)}\right)^{1/2},\\ l=1,2,3;m,k=1,2;p=1,2,3;s=1,2;\end{array}$$
(9.11)
The basis of the lattice, corresponding to the carrier of equation (9.7), is
$$\begin{array}{c}B_1=(1,1),B_2=(\frac{7}{6},0)\end{array}$$
The set of carriers of expansions for solution $`𝐊`$ coincides with (7.12). The expansion of solution for $`\eta ^{(l,m,k,p,s)}`$ takes the form
$$\begin{array}{c}\eta ^{(l,m,k,p,s)}=q_{1/6}^{(l,m,k,p,s)}z^{1/6}+\underset{n=1}{\overset{\mathrm{}}{}}q_{(17n)/6}^{(l,m,k,p,s)}z^{(17n)/6},\\ l=1,2,3;m=1,2;k=1,2;p=1,2,3;;s=1,2;\end{array}$$
(9.12)
Coefficients $`q_{1/6}^{(l,m,k,p,s)},s=1,2`$ are determined by formulas (9.11). The computing of the coefficient $`q_1^{(l,m,k,p,s)}`$ gives a result $`q_1^{(l,m,k,p,s)}=1/6`$. Exponential additions $`y^{(s,p,l,m,k)}(z)`$ to the solutions $`v^{(l,m,k,p)}(z)`$ are
$$\begin{array}{c}y^{(l,m,k,p,s)}(z)=C_3z^{1/6}\\ \mathrm{exp}\left[\frac{6}{7}q_{1/6}^{(l,m,k,p,s)}z^{7/6}+\underset{n=2}{\overset{\mathrm{}}{}}\frac{6}{7(1n)}q_{(17n)/6}^{(l,m,k,p,s)}z^{7(1n)/6}\right]\\ l=1,2,3;m=1,2;k=1,2;p=1,2,3;s=1,2\end{array}$$
(9.13)
Thus we find three levels of the exponential additions to the expansions for solutions of equation near point $`z=\mathrm{}`$.
Solution $`w(z)`$ at $`z\mathrm{}`$ with taking into account the exponential additions has the expansion
$$\begin{array}{c}w(z)=c_{1/3}^{(l)}z^{1/3}+\frac{1}{24}z^2+\underset{n=2}{\overset{\mathrm{}}{}}c_{(17n)/3}^{(l)}z^{(17n)/3}+\\ +C_1z^{1/4}\mathrm{exp}\{F_1(z)+C_2z^{1/6}\mathrm{exp}\{F_2(z)+C_3z^{1/6}\mathrm{exp}\{F_3(z)\}\}\}\end{array}$$
(9.14)
where $`c_{1/3}^{(l)}`$ can be computed by formulas (6.2), (6.3) and (6.4); $`F_1(z)=F_1^{(l,m,k)}(z)`$, $`F_2(z)=F_2^{(l,m,k,p)}(z)`$ and $`F_3(z)=F_3^{(l,m,k,p,s)}(z)`$, ($`l=1,2,3;m,k=1,2;p=1,2,3;s=1,2`$) are
$$\begin{array}{c}F_1^{(l,m,k)}(z)=\frac{6}{7}g_{1/6}^{(l,m,k)}z^{7/6}+\underset{n=2}{\overset{\mathrm{}}{}}\frac{6}{7(1n)}g_{(17n)/6}^{(l,m,k)}z^{7(1n)/6}\end{array}$$
(9.15)
$$\begin{array}{c}F_2^{(p,l,m,k)}(z)=\frac{6}{7}r_{1/6}^{(l,m,k,p)}z^{7/6}+\underset{n=2}{\overset{\mathrm{}}{}}\frac{6}{7(1n)}r_{(17n)/6}^{(l,m,k,p)}z^{7(1n)/6}\end{array}$$
(9.16)
$$\begin{array}{c}F_3^{(l,m,k,p,s)}(z)=\frac{6}{7}q_{1/6}^{(l,m,k,p,s)}z^{7/6}+\underset{n=2}{\overset{\mathrm{}}{}}\frac{6}{7(1n)}q_{(17n)/6}^{(l,m,k,p,s)}z^{7(1n)/6}\end{array}$$
(9.17)
Coefficients $`g_{1/6}^{(l,m,k)}`$, $`r_{1/6}^{(l,m,k,p)}`$ and $`q_{1/6}^{(l,m,k,p,s)}`$ are defined by formulas (7.11), (8.11) and (9.11). The other coefficients are computed sequentially.
## 10 . Conclusion.
Let us formulate the results of this work.
All the power asymptotic forms for equation (2.1) were found.
We also found all the power expansions, corresponding to these asymptotic forms. We denote the obtained families as $`G_1^{(0)}1`$, $`G_1^{(0)}2`$, $`G_1^{(0)}3`$, $`G_1^{(0)}4`$, $`G_1^{(1)}1`$, $`G_2^{(1)}1`$ and $`G_2^{(1)}2`$ (these expansions converge for sufficiently small $`|z|`$). The existence and analyticity of these expansions follow from Cauchy theorem.
We found three families of expansions near $`z=\mathrm{}`$. That are families $`G_3^{(1)}l(l=1,2,3)`$, described by formulas (6.2), (6.3) and (6.4). For each of these expansions we found four exponential additions $`G_3^{(1)}lG_1^1mk(m,k=1,2)`$ expressed by formula (7.16). For them it was computed exponential additions $`G_3^{(1)}lG_1^1mkG_1^{(1)}p(m,k=1,2;p=1,2,3)`$, and then for them the proper exponential additions $`G_3^{(1)}lG_1^{(1)}mkG_1^{(1)}pG_1^{(1)}s(m,k=1,2;p=1,2,3;s=1,2)`$ were found too.
Families $`G_2^{(1)}l`$ and $`G_2^{(1)}2`$ were first found in the paper . However the structure of expansions $`G_2^{(1)}1`$ and $`G_2^{(2)}2`$ was not discussed earlier. The other families of expansions of solution are found for the first time.
Comparing the power expansions of equation (2.1) with power expansions of Painlevé equations $`P_1÷P_6`$ we note, that they differ. This fact can be interpreted as the additional prove for the hypothesis, that the fourth-order equation (2.1) determines new transcendental functions just as equations $`P_1÷P_6`$. |
warning/0507/cond-mat0507366.html | ar5iv | text | # Non-universal critical behaviour of a mixed-spin Ising model on the extended Kagomé latticeThis work was financially supported under the grants VEGA 1/2009/05 and APVT 20-005204.
## 1 Introduction
Investigation of phase transitions and critical phenomena belongs to the most intensively studied topics in the equilibrium statistical physics. A considerable progress in the understanding of order-disorder phenomena has been achieved by solving planar Ising models that represent valuable exceptions of exactly soluble lattice-statistical models with a non-trivial critical behaviour . Although phase transitions of planar Ising models have already been understood in many respects, there are still a lot of obscurities connected with a criticality of more complicated spin systems exhibiting reentrant transitions, non-universal critical behaviour, tricritical phenomenon, etc. It is worthy to mention, however, that several complicated Ising models can exactly be treated by transforming them to the solvable vertex models. A spin-1/2 Ising model on the union jack (centered square) lattice, which represents a first exactly soluble system exhibiting reentrant transitions , can be for instance reformulated as a free-fermion eight-vertex model . It should be also pointed out that an equivalence with the vertex models have already provided a precise confirmation of the reentrant phenomenon in the anisotropic spin-1/2 Ising models on extended Kagomé lattice and centered honeycomb lattice as well.
Despite the significant amount of effort, there are only few exactly soluble Ising models consisting of mixed spins of different magnitudes, which are usually called also as mixed-spin Ising models. A strong scientific interest focused on the mixed-spin systems arises partly on account of much richer critical behaviour they display compared with their single-spin counterparts and partly due to the fact that they represent the most simple models of ferrimagnets having a wide potential applicability in practice. Using the extended versions of decoration-iteration and star-triangle mapping transformations, Fisher has derived exact solutions of the mixed spin-1/2 and spin-$`S`$ ($`S1`$) Ising models on the honeycomb, diced and decorated honeycomb lattices. Notice that these mapping transformations were later on further generalized in order to account also for the single-ion anisotropy effect. The influence of uniaxial and biaxial single-ion anisotropies have precisely been investigated on the mixed-spin honeycomb lattice as well as on some decorated planar lattices . With exception of several mixed-spin models formulated on the Bethe (Cayley tree) lattices, which can be accurately treated within a discrete non-linear map or an approach based on exact recursion equations , these are the only mixed-spin planar Ising models with generally known exact solutions, yet.
One of the most outstanding findings emerging in the phase transition theory is being a non-universal critical behaviour of some planar Ising models, which is in obvious contradiction with the idea of universality hypothesis . The mixed spin-1/2 and spin-$`S`$ Ising model on the union jack lattice represents very interesting system from this viewpoint as it exhibits a remarkable line of bicritical points that have continuously varying critical indices obeying the weak universality hypothesis . In the present article, we shall investigate a topologically similar mixed spin-1/2 and spin-3/2 Ising model on the extended Kagomé lattice by establishing a mapping correspondence with the staggered and uniform eight-vertex models, respectively. In a certain subspace of interaction parameters, the model under investigation becomes exactly soluble as the staggered eight-vertex model satisfying the free-fermion condition . Even if a non-validity of the free-fermion condition in the rest of parameter space is simply ignored, one still obtains rather reliable estimate of the criticality within free-fermion approximation . Finally, the critical points within another subspace of interaction parameters can be approximated from the relevant solution of the uniform eight-vertex model satisfying the zero-field condition.
The outline of this paper is as follows. In Section 2, a detailed formulation of the model is presented and subsequently, the mapping correspondence that ensures an equivalence with the eight-vertex models will be derived. The most interesting numerical results for a critical behaviour will be presented and particularly discussed in Section 3. Finally, some concluding remarks are drawn in Section 4.
## 2 Formulation
Let us begin by considering the mixed spin-1/2 and spin-3/2 Ising model on the extended Kagomé lattice $``$ schematically illustrated in figure 1. The mixed-spin Kagomé lattice consists of the spin-1/2 (empty) and spin-3/2 (filled circles) atoms placed on the six- and four-coordinated sites, respectively. The total Hamiltonian defined upon the underlying lattice $``$ reads:
$`_{mix}=J{\displaystyle \underset{(i,j)𝒥}{\overset{2N}{}}}S_i\sigma _jJ^{}{\displaystyle \underset{(k,l)𝒦}{\overset{2N}{}}}\sigma _k\sigma _lD{\displaystyle \underset{i=1}{\overset{N/2}{}}}S_i^2,`$ (1)
where $`\sigma _j=\pm 1/2`$ and $`S_i=\pm 1/2,\pm 3/2`$ are Ising spin variables, $`J`$ denotes the exchange interaction between nearest-neighbouring spin-1/2 and spin-3/2 pairs and $`J^{}`$ labels the interaction between the spin-1/2 pairs that are next-nearest-neighbours on the extended Kagomé lattice $``$. Finally, the parameter $`D`$ measures a strength of the uniaxial single-ion anisotropy acting on the spin-3/2 sites and $`N`$ denotes the total number of the spin-1/2 sites.
In order to proceed further with calculation, the central spin-3/2 atoms should be firstly decimated from all faces of extended Kagomé lattice $``$. After the decimation, i.e. after performing a summation over spin degrees of freedom of the spin-3/2 sites (filled circles), the partition function of the mixed-spin system can be rewritten as:
$`𝒵_{mix}={\displaystyle \underset{\{\sigma \}}{}}{\displaystyle \underset{m=1}{\overset{N/2}{}}}\omega _m^𝒜(\sigma _i,\sigma _j,\sigma _k,\sigma _l){\displaystyle \underset{n=1}{\overset{N/2}{}}}\omega _n^{}(\sigma _i,\sigma _j,\sigma _k,\sigma _l).`$ (2)
Above, the summation is performed over all possible spin configurations available at the spin-1/2 sites and the first (second) product is over $`N/2`$ faces having four spin-1/2 sites $`\sigma _i`$, $`\sigma _j`$, $`\sigma _k`$, $`\sigma _l`$ placed in the corners of square plaquettes with (without) a central spin-3/2 site in the middle of these plaquettes (see figure 1). The Boltzmann factors $`\omega ^𝒜(a,b,c,d)`$ and $`\omega ^{}(a,b,c,d)`$ assigned to two different kinds of alternating faces, which constitute the checkerboard lattice, can be defined as:
$`\omega ^𝒜(a,b,c,d)`$ $`=`$ $`2\mathrm{exp}[K^{}(ab+bc+cd+da)/2+\mathrm{\Delta }/4]`$
$`\left\{\mathrm{exp}(2\mathrm{\Delta })\mathrm{cosh}[3K(a+b+c+d)/2]+\mathrm{cosh}[K(a+b+c+d)/2]\right\},`$
$`\omega ^{}(a,b,c,d)`$ $`=`$ $`\mathrm{exp}[K^{}(ab+bc+cd+da)/2],`$ (3)
where $`K=J/(k_\mathrm{B}T)`$, $`K^{}=J^{}/(k_\mathrm{B}T)`$, $`\mathrm{\Delta }=D/(k_\mathrm{B}T)`$, $`k_\mathrm{B}`$ is Boltzmann’s constant, and $`T`$ stands for the absolute temperature.
At this stage, the model under investigation can be rather straightforwardly mapped onto the staggered eight-vertex model defined on a dual checkerboard lattice $`_𝒟`$, since Boltzmann factors $`\omega ^𝒜(a,b,c,d)`$ and $`\omega ^{}(a,b,c,d)`$ are being invariant under the reversal of all four spin variables. Actually, there are maximally eight different spin arrangements giving different Boltzmann weights $`\omega ^𝒜(a,b,c,d)`$ and $`\omega ^{}(a,b,c,d)`$ for each kind of face. Diagrammatic representation of eight possible spin arrangements and their corresponding line coverings of the eight-vertex model is shown in figure 2. If, and only if, the adjacent spins are aligned opposite to each other, then solid lines are drawn on the edges of the dual lattice $`_𝒟`$, otherwise they are drawn as broken lines.
It can be easily understood that eight possible line coverings around each vertex of the dual checkerboard lattice always correspond to two spin configurations, one is being obtained from the other by reversing all spins. Since there is even number of solid (broken) lines incident to each vertex of the dual lattice $`_𝒟`$, the model becomes equivalent to the staggered eight-vertex model.
The Boltzmann weights $`\omega ^𝒜(a,b,c,d)`$ and $`\omega ^{}(a,b,c,d)`$, which correspond to eight possible line coverings emerging at vertices of the dual checkerboard lattice, can directly be calculated from equation (3):
$`\omega _1^𝒜`$ $`=`$ $`2\mathrm{exp}(K^{}/2+\mathrm{\Delta }/4)[\mathrm{exp}(2\mathrm{\Delta })\mathrm{cosh}(3K)+\mathrm{cosh}(K)],`$
$`\omega _2^𝒜`$ $`=`$ $`2\mathrm{exp}(K^{}/2+\mathrm{\Delta }/4)[\mathrm{exp}(2\mathrm{\Delta })+1],`$
$`\omega _3^𝒜`$ $`=`$ $`\omega _4^𝒜=2\mathrm{exp}(\mathrm{\Delta }/4)[\mathrm{exp}(2\mathrm{\Delta })+1],`$
$`\omega _5^𝒜`$ $`=`$ $`\omega _6^𝒜=\omega _7^𝒜=\omega _8^𝒜=2\mathrm{exp}(\mathrm{\Delta }/4)[\mathrm{exp}(2\mathrm{\Delta })\mathrm{cosh}(3K/2)+\mathrm{cosh}(K/2)];`$ (4)
$`\omega _1^{}`$ $`=`$ $`\mathrm{exp}(K^{}/2),\omega _1^{}=\mathrm{exp}(K^{}/2),`$
$`\omega _3^{}`$ $`=`$ $`\omega _4^{}=\omega _5^{}=\omega _6^{}=\omega _7^{}=\omega _8^{}=1.`$ (5)
Unfortunately, there does not exist general exact solution of the staggered eight-vertex model with arbitrary Boltzmann weights $`\omega _i^𝒜`$ and $`\omega _j^{}`$ $`(i,j=18)`$. However, if the weights (4) and (5) satisfy so-called free-fermion condition:
$`\mathrm{\Omega }_1\mathrm{\Omega }_2+\mathrm{\Omega }_3\mathrm{\Omega }_4=\mathrm{\Omega }_5\mathrm{\Omega }_6+\mathrm{\Omega }_7\mathrm{\Omega }_8,`$ (6)
the staggered eight-vertex model then becomes exactly soluble as the free-fermion model solved several years ago by Hsue, Lin and Wu . The expressions which enter into the free-fermion condition (6) can be defined through:
$`\mathrm{\Omega }_1`$ $`=`$ $`\omega _1^𝒜\omega _1^{}+\omega _2^𝒜\omega _2^{},\mathrm{\Omega }_2=\omega _3^𝒜\omega _3^{}+\omega _4^𝒜\omega _4^{},`$
$`\mathrm{\Omega }_3`$ $`=`$ $`\omega _5^𝒜\omega _6^{}+\omega _5^{}\omega _6^𝒜,\mathrm{\Omega }_4=\omega _7^𝒜\omega _8^{}+\omega _7^{}\omega _8^𝒜,`$
$`\mathrm{\Omega }_5\mathrm{\Omega }_6`$ $`=`$ $`\omega _1^𝒜\omega _1^{}\omega _3^𝒜\omega _3^{}+\omega _2^𝒜\omega _2^{}\omega _4^𝒜\omega _4^{}+\omega _5^𝒜\omega _6^{}\omega _7^𝒜\omega _8^{}+\omega _5^{}\omega _6^𝒜\omega _7^{}\omega _8^𝒜,`$
$`\mathrm{\Omega }_7\mathrm{\Omega }_8`$ $`=`$ $`\omega _1^𝒜\omega _1^{}\omega _4^𝒜\omega _4^{}+\omega _2^𝒜\omega _2^{}\omega _3^𝒜\omega _3^{}+\omega _5^𝒜\omega _6^{}\omega _7^{}\omega _8^𝒜+\omega _5^{}\omega _6^𝒜\omega _7^𝒜\omega _8^{}.`$ (7)
It can be readily proved that the free-fermion condition (6) holds in our case just as $`D\pm \mathrm{}`$, or $`T\mathrm{}`$. The restriction to infinitely strong single-ion anisotropy consequently leads to the familiar phase transitions from the standard Ising universality class, because in this case our model effectively reduces to a simple spin-1/2 Ising model on the extended Kagomé lattice. Within the manifold given by the constraint (6), the free-fermion model becomes critical as long as:
$`\mathrm{\Omega }_1+\mathrm{\Omega }_2+\mathrm{\Omega }_3+\mathrm{\Omega }_4=2\text{max}\{\mathrm{\Omega }_1,\mathrm{\Omega }_2,\mathrm{\Omega }_3,\mathrm{\Omega }_4\}.`$ (8)
It is noteworthy, however, that the critical condition (8) yields rather reliable estimate of the criticality within so-called free-fermion approximation even if a non-validity of the free-fermion condition (6) is simply ignored.
Now, we shall establish an approximate mapping between the staggered and uniform eight-vertex models, since the second branch of exact solution is available just for the latter model under the zero-field condition . For this purpose, let us define average Boltzmann weights of the staggered eight-vertex model, which would approximately transform the staggered eight-vertex model into the uniform one:
$`\stackrel{~}{\omega }_i`$ $`=`$ $`\omega _i^𝒜\omega _i^{},(i=18).`$ (9)
Note that the uniform eight-vertex model satisfies the zero-field condition just when its Boltzmann weights are pairwise and symmetrically equal to each other:
$`\stackrel{~}{\omega }_1=\stackrel{~}{\omega }_2,\stackrel{~}{\omega }_3=\stackrel{~}{\omega }_4,\stackrel{~}{\omega }_5=\stackrel{~}{\omega }_6,\stackrel{~}{\omega }_7=\stackrel{~}{\omega }_8.`$ (10)
As we already have $`\stackrel{~}{\omega }_3=\stackrel{~}{\omega }_4`$, $`\stackrel{~}{\omega }_5=\stackrel{~}{\omega }_6`$, and $`\stackrel{~}{\omega }_7=\stackrel{~}{\omega }_8`$, the zero-field case is consequently reached by imposing the condition $`\stackrel{~}{\omega }_1=\stackrel{~}{\omega }_2`$ only, or equivalently:
$`\mathrm{exp}(2\mathrm{\Delta })={\displaystyle \frac{\mathrm{exp}(2K^{})\mathrm{cosh}(K)}{\mathrm{cosh}(3K)\mathrm{exp}(2K^{})}},`$ (11)
According to Baxter’s exact solution , the zero-field eight-vertex model becomes critical on the manifold (10) if:
$`\stackrel{~}{\omega }_1+\stackrel{~}{\omega }_3+\stackrel{~}{\omega }_5+\stackrel{~}{\omega }_7=2\text{max}\{\stackrel{~}{\omega }_1,\stackrel{~}{\omega }_3,\stackrel{~}{\omega }_5,\stackrel{~}{\omega }_7\}.`$ (12)
It is easy to check that $`\stackrel{~}{\omega }_1`$ represents in our case the largest Boltzmann weight, thus, the condition determining the criticality can also be written in this equivalent form:
$`\mathrm{exp}(K_c^{})`$ $`[`$ $`\mathrm{exp}(2\mathrm{\Delta }_c)\mathrm{cosh}(3K_c)+\mathrm{cosh}(K_c)]=`$
$`1`$ $`+`$ $`\mathrm{exp}(2\mathrm{\Delta }_c)+2\mathrm{exp}(2\mathrm{\Delta }_c)\mathrm{cosh}(3K_c/2)+2\mathrm{cosh}(K_c/2),`$ (13)
where $`K_c=J/(k_\mathrm{B}T_c)`$, $`K_c^{}=J^{}/(k_\mathrm{B}T_c)`$, $`\mathrm{\Delta }_c=D/(k_\mathrm{B}T_c)`$, and $`T_c`$ denotes the critical temperature. It should be stressed, nevertheless, that the critical exponents (with exception of $`\delta `$ and $`\eta `$) describing a phase transition of the zero-field eight-vertex model depend on the function $`\mu =2\mathrm{arctan}(\stackrel{~}{\omega }_5\stackrel{~}{\omega }_7/\stackrel{~}{\omega }_1\stackrel{~}{\omega }_3)^{1/2}`$, in fact:
$`\alpha =\alpha ^{}=2{\displaystyle \frac{\pi }{\mu }},\beta ={\displaystyle \frac{\pi }{16\mu }},\nu =\nu ^{}={\displaystyle \frac{\pi }{2\mu }},\gamma ={\displaystyle \frac{7\pi }{8\mu }},\delta =15,\eta ={\displaystyle \frac{1}{4}},`$ (14)
Finally, let us explicitly evaluate the critical exponent $`\beta `$ that determines disappearance of the spontaneous order as the critical temperature is approached from below:
$`\beta ^1={\displaystyle \frac{32}{\pi }}\mathrm{arctan}\left\{{\displaystyle \frac{\mathrm{exp}(2\mathrm{\Delta }_c)\mathrm{cosh}(3K_c/2)+\mathrm{cosh}(K_c/2)}{[\mathrm{exp}(2\mathrm{\Delta }_c)+1]^{3/4}[\mathrm{exp}(2\mathrm{\Delta }_c)\mathrm{cosh}(3K_c)+\mathrm{cosh}(K_c)]^{1/4}}}\right\}.`$ (15)
## 3 Results and discussion
Now, let us turn our attention to a discussion of the most interesting results obtained for the ground-state and finite-temperature phase diagrams. Solid lines displayed in figure 3 represent ground-state phase boundaries separating four distinct long-range ordered phases that emerge in the ground state when $`J>0`$.
Spin order drawn in broken rectangles shows a typical spin configuration within basic unit cell of each phase. As could be expected, a sufficiently strong antiferromagnetic next-nearest-neighbour interaction $`J^{}`$ alters the structure of the ground state owing to a competing effect with the nearest-neighbour interaction $`J`$. Due to a competition between the interactions, the central spins are free to flip within the phases III and IV and thus, these phases exhibit a remarkable coexistence of order and disorder. At last, it is worthwhile to mention that a broken line connecting both triple points depicts a projection of the approximate critical line (13) into the $`J^{}D`$ plane. As this projection crosses zero-temperature plane along the ground-state transition line $`D/J=3/2J^{}/J`$ between the phases I and IV, it is quite reasonable to suspect that this line represent a location of phase transitions between these phases.
Let us investigate more deeply this line of critical points. The critical temperatures calculated from the uniform zero-field eight-vertex model must simultaneously obey both the zero-field condition (11) as well as the critical condition (13). It is easy to check that the former condition necessitates $`1.5<J^{}/J<0.5`$ and $`1.0<D/J<0.0`$. Figure 4(a) displays a projection of this critical line into the $`J^{}T_c`$ plane (the dependence scaled to the left axis) and respectively, a projection into the $`J^{}D`$ plane which is scaled to the right axis. Along this critical line, the critical exponents are expected to vary with interaction parameters as they have to follow the equations (14). For illustration, figure 4(b) shows how the critical index $`\beta `$ changes along the critical line. Apparently, the exponent $`\beta `$ approaches its smallest
possible value $`1/16`$ by reaching both triple points with zero critical temperature, however, it is also quite interesting to ascertain that its greatest value is below the value $`1/8`$ that predicts the universality hypothesis for planar Ising systems .
Before concluding, few remarks should be addressed to a global finite-temperature phase diagram plotted in figure 5, which displays the critical temperature as a function of the ratio $`J^{}/J`$ for several values of the single-ion anisotropy $`D/J`$. Critical boundaries depicted as solid lines represent exact critical points obtained from the free-fermion solution (8) of the staggered eight-vertex model obtained under the constraint (6), which is fulfilled in the limiting cases $`D/J\pm \mathrm{}`$. Dotted critical lines show estimated critical temperatures calculated from the free-fermion approximation simply ignoring a non-validity of the free-fermion condition (6) for any finite value of $`D/J`$. Approximative solution related to the critical points (13) of the uniform zero-field eight-vertex model on the variety (11) is displayed as a rounded broken line.
It is quite obvious from the ground-state phase diagram (figure 3) that a right (left) wing of the displayed critical boundaries corresponds to the phase I (III) if $`D/J>0.0`$, while it corresponds to the phase II (IV) if $`D/J<1.0`$. Actually, the exact as well as approximate critical points resulting from the free-fermion solution correctly reproduce the ground-state boundaries between these phases. When the single-ion anisotropy parameter is selected within the range $`1.0<D/J<0.0`$ (see for instance the curve for $`D/J=0.5`$), however, the critical line obtained from the free-fermion approximation meets at a bicritical (circled) point with the critical line of the equivalent uniform zero-field eight-vertex model as it has been already reasoned by Lipowski and Horiguchi who have solved similar spin system on the union-jack lattice. In such a case, the right and left part (with respect to the bicritical point) of this critical line separate the phases I and IV, respectively, and a line of first-order phase transitions is expected to terminate at this special multicritical point. There are strong indications supporting this concept , actually, the almost straight broken line depicting the zero-field condition (11) should always show a coexistence of these two phases as it starts from a point that determines their coexistence in the ground state. With regard to the aforementioned arguments one may conclude that a coexistence surface between the phases I and IV lies inside the area, which is bounded by the line of bicritical points (rounded broken line) having the non-universal interaction-dependent critical exponents.
Finally, we should remark a feasible appearance of reentrant transitions which can be observed in the critical lines nearby the coexistence points $`D/J=0.0`$ and $`1.0`$. It is quite apparent that the observed reentrance can be explained in terms of the coexistence of a partial order and partial disorder emerging in both the high-temperature reentrant phases III and IV. As a matter of fact, the partial disorder of the spin-3/2 atoms can compensate a loss of entropy that occurs in these phases due to a thermally induced partial ordering of the spin-1/2 atoms what is in a good accordance with a necessary condition conjectured for the appearance of reentrant phase transitions -.
## 4 Concluding Remarks
The work reported in the present article provides a relatively precise information on the critical behaviour of the mixed spin-1/2 and spin-3/2 Ising model on the extended Kagomé lattice by establishing a mapping correspondence with the staggered and uniform eight-vertex models, respectively. The main focus of the present work has been aimed at the examination of the criticality depending on the single-ion anisotropy strength as well as the strength of the competing next-nearest-neighbour interaction. The location of the critical boundaries has accurately been determined from the free-fermion solution of the staggered eight-vertex model and the zero-field solution of the uniform eight-vertex model, respectively, whereas the validity of both mappings is restricted to the certain subspaces of interaction parameters only. In the rest of parameter space, the free-fermion approximation has been used to estimate the critical boundaries as this method should provide meaningful approximation giving rather reliable estimate to the true transition temperatures.
The greatest theoretical interest in this model arises due to the remarkable critical line consisting of bicritical points, which bounds a coexistence surface between two long-range ordered phases. The bicritical points can be characterized by non-universal interaction-dependent critical exponents that satisfy the weak universality hypothesis. Moreover, the same arguments as those suggested by Lipowski and Horiguchi have enabled us to identify the zero-field condition (11) with a location of the first-order transition lines separating these two ordered phases.
It should be remarked that the considered spin system also shows reentrant phase transitions on account of the competition between the nearest- and next-nearest-neighbour interactions. Our results are in agreement with the conjecture stating that the reentrance appears as a consequence of the coexistence of a partial order and disorder, namely, the partial disorder induced among spin-3/2 atoms can compensate the loss of entropy, which occurs on behalf of the partial ordering of the spin-1/2 atoms in both the high-temperature partially ordered phases. |
warning/0507/astro-ph0507596.html | ar5iv | text | # Quantum corrections to the inflaton potential and the power spectra from superhorizon modes and trace anomalies.
## I Introduction
Inflation is a central part of early Universe cosmology passing many observational tests and becoming a predictive scenario scrutinized by current and forthcoming observations. Inflation was introduced to solve several shortcomings of the standard Big Bang cosmologyguth -riottorev . It provides a mechanism for generating scalar (density) and tensor (gravitational wave) perturbationsmukha -mukhanov . A distinct aspect of inflationary perturbations is that these are generated by quantum fluctuations of the scalar field(s) that drive inflation. After their wavelength becomes larger than the Hubble radius, these fluctuations are amplified and grow, becoming classical and decoupling from causal microphysical processes. Upon re-entering the horizon, during the matter era, these classical perturbations seed the inhomogeneities which generate structure upon gravitational collapsemukha -mukhanov . A great diversity of inflationary models predict fairly generic features: a gaussian, nearly scale invariant spectrum of (mostly) adiabatic scalar and tensor primordial fluctuations, making the inflationary paradigm fairly robust. The gaussian, adiabatic and nearly scale invariant spectrum of primordial fluctuations provide an excellent fit to the highly precise wealth of data provided by the Wilkinson Microwave Anisotropy Probe (WMAP)komatsu ; spergel ; kogut ; peiris . Perhaps the most striking validation of inflation as a mechanism for generating *superhorizon* (‘acausal’) fluctuations is the anticorrelation peak in the temperature-polarization (TE) angular power spectrum at $`l150`$ corresponding to superhorizon scaleskogut ; peiris .
The confirmation of many of the robust predictions of inflation by current high precision observations places inflationary cosmology on solid grounds. Forthcoming observations will begin to discriminate among different inflationary models, placing stringent constraints on them. There are small but important telltale discriminants amongst different models: non-gaussianity, a running spectral index for scalar and tensor perturbations, an isocurvature component for scalar perturbations, the ratios for the amplitudes between tensor and scalar modes, etc. Already WMAP reports a hint of deviations from constant scaling exponents (running spectral index) and rules out the purely monomial $`\mathrm{\Phi }^4`$ potentialpeiris .
Amongst the wide variety of inflationary scenarios, single field *slow roll* modelsbarrow ; stewlyth provide an appealing, simple and fairly generic description of inflation. Its simplest implementation is based on a scalar field (the inflaton) whose homogeneous expectation value drives the dynamics of the scale factor, plus small quantum fluctuations. The inflaton potential, is fairly flat during inflation. This flatness not only leads to a slowly varying Hubble parameter, hence ensuring a sufficient number of e-folds, but also provides an explanation for the gaussianity of the fluctuations as well as for the (almost) scale invariance of their power spectrum. A flat potential precludes large non-linearities in the dynamics of the *fluctuations* of the scalar field. The current WMAP data seems to validate the simpler one-field slow roll scenariopeiris . Furthermore, because the potential is flat the scalar field is almost massless, and modes cross the horizon with an amplitude proportional to the Hubble parameter. This fact combined with a slowly varying Hubble parameter yields an almost scale invariant primordial power spectrum. Upon crossing the horizon the phases of the quantum fluctuations freeze out and a growing mode dominates the dynamics, i.e. the quantum fluctuations become classical (see ref.liddle and references therein). Departures from scale invariance and gaussianity are determined by the departures from flatness of the potential, namely by derivatives of the potential with respect to the inflaton. These derivatives can be combined into a hierarchy of dimensionless slow roll parametersbarrow that allow an assessment of the *corrections* to the basic predictions of gaussianity and scale invarianceliddle . The slow-roll approximation has been recently cast as a $`1/N_{efolds}`$ expansionn , where $`N_{efolds}`$ is the number of efolds before the end of inflation when modes of cosmological relevance today first crossed the Hubble radius.
The basic scenario of inflation driven by a scalar field must be interpreted as an *effective* field theoryhector resulting from integrating out heavy degrees of freedom. In particular, in the effective field theory description, the *classical* scalar potential that determines the dynamics of the inflaton, results from integrating out degrees of freedom *much heavier* than the scale of inflation.
Forthcoming observations have the potential of measuring the inflationary potential at least within a span in field amplitude corresponding to the 8-10 e-folds during which wavelengths of cosmological relevance first cross the Hubble radiuslidsey . These observations will measure the full inflaton potential including all possible quantum corrections and not just the classical (tree level) potential. This possibility motivates us to assess the *quantum* corrections to the inflationary potential from fields lighter than the inflaton, since in the effective field theory description, the classical inflaton potential already includes contributions from heavier fields. We focus on light fields since these can exhibit infrared enhanced contributions to the effective potential as discussed in ref.nuestros ; nuestronor .
Our goal is to obtain the *effective* potential that includes the one loop quantum corrections from fields that are *light* during the relevant inflationary stage.
Our program of study focuses on the understanding of quantum aspects of the basic inflationary paradigm. In previous studies we addressed the decay of inflaton fluctuationsnuestros and more recentlynuestronor we focused on the quantum corrections to the equations of motion of the inflaton and the scalar fluctuations during slow roll inflation, from integrating out not only the inflaton fluctuations but also the excitations associated with another scalar field. Since the power spectra of fields with masses $`mH`$ are nearly scale invariant, strong infrared enhancements appear as revealed in these studiesnuestros ; nuestronor . In addition, we find that a particular combination of slow roll parameters which measures the departure from scale invariance of the fluctuations provides a natural infrared regularization.
The small parameter that determines the validity of inflation as an effective *quantum field theory* below the Planck scale is $`H/M_{Pl}`$ where $`H`$ is the Hubble parameter during inflation and therefore the scale at which inflation occurs. The slow roll expansion is in a very well defined sense an *adiabatic* approximation since the time evolution of the inflaton field is slow on the expansion scale. Thus the small dimensionless ratio $`H/M_{Pl}`$, which is required for the validity of an effective field theory (EFT) is logically *independent* from the small dimensionless combinations of derivatives of the potential which ensure the validity of the slow-roll expansion. Present datapeiris indicate a very small amplitude of tensor perturbations which is consistent with $`H/M_{Pl}1`$.
Therefore, in this article we will invoke *two independent* approximations, the effective field theory (EFT) and the slow roll approximation. The former is defined in terms of an expansion in the ratio $`H/M_{Pl}`$, whereas the latter corresponds to an expansion in the (small) slow roll parameters which has recently been identified with an expansion in $`1/N_{efolds}`$n .
It is important to highlight the main differences between slow roll inflation and the post-inflationary stage. During slow roll inflation the dynamics of the scalar field is slow on the time scale of the expansion and consequently the change in the amplitude of the inflaton is small and quantified by the slow roll parameters. The slow roll approximation is indeed an *adiabatic approximation*. In striking contrast to this situation, during the post-inflationary stage of reheating the scalar field undergoes rapid and large amplitude oscillations that cannot be studied in a perturbative expansionreheatnuestro ; ramsey .
Brief summary of results: we obtain the quantum corrections to the inflaton potential up to one loop by including the contributions from scalar and tensor perturbations of the metric as well as one light scalar and one light fermion field coupled generically to the inflaton. Therefore this study provides the most complete assessment of the general backreaction problem up to one loop that includes not only metric perturbations, but also the contributions from fluctuations of other light fields with a generic treatment of both bosonic and fermionic degrees of freedom. Motivated by an assessment of the quantum fluctuations that *could* be of observational interest, we focus on studying the effective inflaton potential during the cosmologically relevant stage of slow roll inflation.
Both light bosonic fields as well as scalar density perturbations feature an infrared enhancement of their quantum corrections which is regularized by slow roll parameters. Fermionic contributions as expected do not feature any infrared enhancement and neither does the graviton contribution to the energy momentum tensor. We find that in slow roll and for light bosonic and fermionic fields there is a clean separation between the super and subhorizon contributions to the quantum corrections from scalar density metric and light bosonic field perturbations. For these fields the superhorizon contribution is of zero order in slow roll as a consequence of the infrared enhancement regularized by slow roll parameters. The subhorizon contribution to the energy momentum tensor from all the fields is completely determined by the trace anomaly of minimally coupled scalars, gravitons and fermionic fields. We find the one loop effective potential to be
$$V_{eff}(\mathrm{\Phi }_0)=V(\mathrm{\Phi }_0)\left[1+\frac{H_0^2}{3(4\pi )^2M_{Pl}^2}\left(\frac{\eta _v4ϵ_v}{\eta _v3ϵ_v}+\frac{3\eta _\sigma }{\eta _\sigma ϵ_v}+𝒯\right)\right]$$
(1)
where $`V(\mathrm{\Phi }_0)`$ is the *classical* inflaton potential, $`\eta _v,ϵ_v,\eta _\sigma `$ slow-roll parameters and $`𝒯=𝒯_\mathrm{\Phi }+𝒯_s+𝒯_t+𝒯_\mathrm{\Psi }=\frac{2903}{20}=145.15`$ is the total trace anomaly from the scalar metric, tensor, light scalar and fermion contributions.
The terms that feature ratios of slow roll parameters arise from superhorizon contributions from curvature and scalar field perturbations. The last term in eq.(1) is independent of slow-roll parameters and is completely determined by the trace anomalies of the different fields. It is the hallmark of the subhorizon contributions.
In the case when the mass of the light bosonic scalar field is much smaller than the mass of the inflaton fluctuations, we find the following result for the scalar curvature and tensor fluctuations including the one-loop quantum corrections,
$`|\mathrm{\Delta }_{k,eff}^{(S)}|^2=|\mathrm{\Delta }_k^{(S)}|^2\left\{1+{\displaystyle \frac{2}{3}}\left({\displaystyle \frac{H_0}{4\pi M_{Pl}}}\right)^2\left[1+{\displaystyle \frac{\frac{3}{8}r(n_s1)+2\frac{dn_s}{d\mathrm{ln}k}}{(n_s1)^2}}+{\displaystyle \frac{2903}{40}}\right]\right\}`$ (2)
(3)
$`|\mathrm{\Delta }_{k,eff}^{(T)}|^2=|\mathrm{\Delta }_k^{(T)}|^2\left\{1{\displaystyle \frac{1}{3}}\left({\displaystyle \frac{H_0}{4\pi M_{Pl}}}\right)^2\left[1+{\displaystyle \frac{1}{8}}{\displaystyle \frac{r}{n_s1}}+{\displaystyle \frac{2903}{20}}\right]\right\},`$ (4)
(5)
$`r_{eff}{\displaystyle \frac{|\mathrm{\Delta }_{k,eff}^{(T)}|^2}{|\mathrm{\Delta }_{k,eff}^{(S)}|^2}}=r\left\{1{\displaystyle \frac{1}{3}}\left({\displaystyle \frac{H_0}{4\pi M_{Pl}}}\right)^2\left[1+{\displaystyle \frac{\frac{3}{8}r(n_s1)+\frac{dn_s}{d\mathrm{ln}k}}{(n_s1)^2}}+{\displaystyle \frac{8709}{20}}\right]\right\}.`$ (6)
The quantum corrections turn out to enhance the scalar curvature fluctuations and to reduce the tensor fluctuations as well as their ratio $`r`$. The quantum corrections are always small, of the order $`\left(\frac{H_0}{M_{Pl}}\right)^2`$, but it is interesting to see that these quantum effects are dominated by the trace anomalies and they correct both scalar and tensor fluctuations in a definite direction. Moreover, it is the tensor part of the trace anomaly which numerically yields the largest contribution.
Quantum trace (conformal) anomalies of the energy momentum tensor in gravitational fields constitute an important aspect of quantum field theory in curved backgrounds, (see for example BD and references therein). In black hole backgrounds they are related to the Hawking radiation. It is interesting to see here that the trace anomalies appear in a relevant cosmological problem and dominate the quantum corrections to the primordial spectrum of curvature and tensor fluctuations.
In section II we compute the effective potential including scalars, gravitons and fermionic fields, in section III we present the quantum corrections to scalar curvature and tensor fluctuations, and in section IV we present our conclusions.
## II The effective potential
In our recent calculation of the quantum corrections to the effective potentialnuestros different expansions appear: the expansion in the effective field theory ratio $`H_0/M_{Pl}`$ where $`H_0`$ is the Hubble parameter during the relevant stage of inflation, and the expansion in slow roll parameters. These expansions are logically different: the slow roll expansion is an *adiabatic* expansion in the sense that the dynamics of the inflaton is slower than the universe expansion, while the (dimensionless) interaction vertices and the loop expansion are determined by the effective field theory parameter $`H_0/M_{Pl}`$nuestros .
During slow roll inflation, the dynamics of the scale factor and the inflaton are determined by the following set of (semi) classical equations of motion
$`H_0^2={\displaystyle \frac{1}{3M_{Pl}^2}}\left[{\displaystyle \frac{1}{2}}(\dot{\mathrm{\Phi }}_0)^2+V(\mathrm{\Phi }_0)\right],`$ (7)
$`\ddot{\mathrm{\Phi }}_0+3H_0\dot{\mathrm{\Phi }}_0+V^{}(\mathrm{\Phi }_0)=0.`$ (8)
where $`M_{Pl}=1/\sqrt{8\pi G}=\mathrm{2.4\hspace{0.33em}10}^{18}`$GeV. Slow roll inflation is tantamount to the statement that the dynamics of the expectation value of the scalar field $`\mathrm{\Phi }_0`$ is slow on the scale of the cosmological expansion. The slow roll approximation is indeed an *adiabatic* approximation in terms of a hierarchy of small dimensionless quantities related to the derivatives of the inflaton potential. Somebarrow ; liddle of these slow roll parameters are given by<sup>1</sup><sup>1</sup>1We follow the definitions of $`\xi _V;\sigma _V`$ in ref.peiris . ($`\xi _V;\sigma _V`$ are called $`\xi _V^2;\sigma _V^3`$, respectively, inbarrow ).
$`ϵ_V={\displaystyle \frac{M_{Pl}^2}{2}}\left[{\displaystyle \frac{V^{^{}}(\mathrm{\Phi }_0)}{V(\mathrm{\Phi }_0)}}\right]^2,\eta _V=M_{Pl}^2{\displaystyle \frac{V^{^{\prime \prime }}(\mathrm{\Phi }_0)}{V(\mathrm{\Phi }_0)}},`$ (9)
$`\xi _V=M_{Pl}^4{\displaystyle \frac{V^{}(\mathrm{\Phi }_0)V^{^{\prime \prime \prime }}(\mathrm{\Phi }_0)}{V^2(\mathrm{\Phi }_0)}},\sigma _V=M_{Pl}^6{\displaystyle \frac{\left[V^{^{}}(\mathrm{\Phi }_0)\right]^2V^{(IV)}(\mathrm{\Phi }_0)}{V^3(\mathrm{\Phi }_0)}}.`$
The slow roll approximationbarrow ; liddle ; lidsey corresponds to $`ϵ_V\eta _V1`$ with the hierarchy $`\xi _V𝒪(ϵ_V^2);\sigma _V𝒪(ϵ_V^3)`$, namely $`ϵ_V`$ and $`\eta _V`$ are first order in slow roll, $`\xi _V`$ second order in slow roll, etc. Recently a correspondence between the slow roll expansion and an expansion in $`1/N_{efolds}`$ has been establishedn with $`ϵ_V,\eta _V1/N_{efolds};\xi _V1/N_{efolds}^2;\sigma _V1/N_{efolds}^3`$, etc.
During slow roll inflation the equation of motion (7)-(8) are approximated by
$`\dot{\mathrm{\Phi }}_0={\displaystyle \frac{V^{}(\mathrm{\Phi }_0)}{3H_0}}+\text{higher order in slow roll},`$ (10)
(11)
$`H_0^2={\displaystyle \frac{V(\mathrm{\Phi }_0)}{3M_{Pl}^2}}\left[1+{\displaystyle \frac{ϵ_V}{3}}+𝒪(ϵ_V^2,ϵ_V\eta _V)\right],`$ (12)
The scale factor is given by
$$C(\eta )=\frac{1}{H\eta (1ϵ_V)}.$$
(13)
In the effective field theory interpretation of inflation, the *classical* inflaton potential $`V(\mathrm{\Phi })`$ should be understood to include the contribution from integrating out fields with masses *much larger* than $`H_0`$. Our goal is to obtain the one loop quantum corrections from *fields that are light during inflation*. Therefore we consider that the inflaton is coupled to a light scalar field $`\sigma `$ and to Fermi fields with a generic Yukawa-type coupling. We take the fermions to be Dirac fields but it is straightforward to generalize to Weyl or Majorana fermions. We also include the contribution to the effective potential from scalar and tensor metric perturbations, thereby considering their *backreaction* up to one loop.
The Lagrangian density is taken to be
$$=\sqrt{g}\left\{\frac{1}{2}\dot{\phi }^2\left(\frac{\stackrel{}{}\phi }{2a}\right)^2V(\phi )+\frac{1}{2}\dot{\sigma }^2\left(\frac{\stackrel{}{}\sigma }{2a}\right)^2\frac{1}{2}m_\sigma ^2\sigma ^2G(\phi )\sigma ^2+\overline{\mathrm{\Psi }}\left[i\gamma ^\mu 𝒟_\mu \mathrm{\Psi }m_fY(\phi )\right]\mathrm{\Psi }\right\}$$
(14)
where $`G(\mathrm{\Phi })`$ and $`Y(\mathrm{\Phi })`$ are generic interaction terms between the inflaton and the scalar and fermionic fields. Obviously this Lagrangian can be further generalized to include a multiplet of scalar and fermionic fields and such case can be analyzed as a straightforward generalization. For simplicity we consider one bosonic and one fermionic Dirac field.
The Dirac $`\gamma ^\mu `$ are the curved space-time $`\gamma `$ matrices and the fermionic covariant derivative is given byweinberg ; BD ; duncan ; casta
$`𝒟_\mu `$ $`=`$ $`_\mu +{\displaystyle \frac{1}{8}}[\gamma ^c,\gamma ^d]V_c^\nu \left(D_\mu V_{d\nu }\right)`$ (15)
$`D_\mu V_{d\nu }`$ $`=`$ $`_\mu V_{d\nu }\mathrm{\Gamma }_{\mu \nu }^\lambda V_{d\lambda }`$
where the vierbein field is defined as
$$g^{\mu \nu }=V_a^\mu V_b^\nu \eta ^{ab},$$
$`\eta _{ab}`$ is the Minkowski space-time metric and the curved space-time matrices $`\gamma ^\mu `$ are given in terms of the Minkowski space-time ones $`\gamma ^a`$ by (greek indices refer to curved space time coordinates and latin indices to the local Minkowski space time coordinates)
$$\gamma ^\mu =\gamma ^aV_a^\mu ,\{\gamma ^\mu ,\gamma ^\nu \}=2g^{\mu \nu }.$$
We will consider that the light scalar field $`\sigma `$ has vanishing expectation value at all times, therefore inflationary dynamics is driven by one single scalar field, the inflaton $`\varphi `$. We now separate the homogeneous expectation value of the inflaton field from its quantum fluctuations as usual by writing
$$\phi (\stackrel{}{x},t)=\mathrm{\Phi }_0(t)+\delta \phi (\stackrel{}{x},t).$$
We will consider the contributions from the quadratic fluctuations to the energy momentum tensor. There are *four* distinct contributions: i) scalar metric (density) perturbations, ii) tensor (gravitational waves) perturbations, iii) fluctuations of the light bosonic scalar field $`\sigma `$, iv) fluctuations of the light fermionic field $`\mathrm{\Psi }`$.
Fluctuations in the metric are studied as usual mukha ; mukhanov ; giova ; hu ; riotto . Writing the metric as
$$g_{\mu \nu }=g_{\mu \nu }^0+\delta ^sg_{\mu \nu }+\delta ^tg_{\mu \nu }$$
where $`g_{\mu \nu }^0`$ is the spatially flat FRW background metric which in conformal time is given by
$$g_{\mu \nu }^0=C^2(\eta )\eta _{\mu \nu },C(\eta )a(t(\eta ))$$
and $`\eta _{\mu \nu }=\text{diag}(1,1,1,1)`$ is the flat Minkowski space-time metric. $`\delta ^{s,t}g_{\mu \nu }`$ correspond to the scalar and tensor perturbations respectively, and we neglect vector perturbations. In longitudinal gauge
$`\delta ^sg_{00}`$ $`=`$ $`C^2(\eta )\mathrm{\hspace{0.33em}2}\varphi `$ (17)
$`\delta ^sg_{ij}`$ $`=`$ $`C^2(\eta )\mathrm{\hspace{0.33em}2}\psi \delta _{ij}`$ (19)
$`\delta ^tg_{ij}`$ $`=`$ $`C^2(\eta )h_{ij}`$
where $`h_{ij}`$ is transverse and traceless and we neglect vector modes since they are not generated in single field inflationmukha ; mukhanov ; giova ; hu ; riotto .
Gauge invariant variables associated with the fluctuations of the scalar field and the potentials $`\varphi ,\psi `$ are constructed explicitly in ref.mukhanov where the reader can find their expressions. Expanding up to quadratic order in the scalar fields, fermionic fields and metric perturbations the part of the Lagrangian density that is quadratic in these fields is given by
$$_Q=_s[\delta \phi ^{gi},\varphi ^{gi},\psi ^{gi}]+_t[h]+_\sigma [\sigma ]+_\mathrm{\Psi }[\overline{\mathrm{\Psi }},\mathrm{\Psi }],$$
where
$`_t[h]={\displaystyle \frac{M_{Pl}^2}{8}}C^2(\eta )_\alpha h_i^j_\beta h_j^i\eta ^{\alpha \beta },`$ (20)
(21)
$`_\sigma [\sigma ]=C^4(\eta )\left\{{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\sigma ^{}}{C}}\right)^2{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\sigma }{C}}\right)^2{\displaystyle \frac{1}{2}}M_\sigma ^2[\mathrm{\Phi }_0]\sigma ^2\right\},`$ (22)
(23)
$`_\mathrm{\Psi }[\overline{\mathrm{\Psi }},\mathrm{\Psi }]=\overline{\mathrm{\Psi }}\left[i\gamma ^\mu 𝒟_\mu \mathrm{\Psi }M_\mathrm{\Psi }[\mathrm{\Phi }_0]\right]\mathrm{\Psi },`$
where the prime stands for derivatives with respect to conformal time and the labels (gi) refer to gauge invariant quantitiesmukhanov . The explicit expression for $`[\delta \phi ^{gi},\varphi ^{gi},\psi ^{gi}]`$ is given in eq. (10.68) in ref.mukhanov . The effective masses for the bosonic and fermionic fields are given by
$`M_\sigma ^2[\mathrm{\Phi }_0]`$ $`=`$ $`m_\sigma ^2+G(\mathrm{\Phi }_0)`$ (24)
$`M_\mathrm{\Psi }[\mathrm{\Phi }_0]`$ $`=`$ $`m_f+Y(\mathrm{\Phi }_0).`$
We will focus on the study of the quantum corrections to the Friedmann equation, for the case in which both the scalar and fermionic fields are light in the sense that during slow roll inflation,
$$M_\sigma [\mathrm{\Phi }_0],M_\mathrm{\Psi }[\mathrm{\Phi }_0]H_0,$$
(26)
at least during the cosmologically relevant stage corresponding to the 50 or so e-folds before the end of inflation.
In conformal time the vierbeins $`V_a^\mu `$ are particularly simple
$$V_a^\mu =C(\eta )\delta _a^\mu $$
(27)
and the Dirac Lagrangian density simplifies to the following expression
$$\sqrt{g}\overline{\mathrm{\Psi }}\left(i\gamma ^\mu 𝒟_\mu \mathrm{\Psi }M_\mathrm{\Psi }[\mathrm{\Phi }_0]\right)\mathrm{\Psi }=C^{\frac{3}{2}}\overline{\mathrm{\Psi }}\left[i\overline{)}M_\mathrm{\Psi }[\mathrm{\Phi }_0]C(\eta )\right]\left(C^{\frac{3}{2}}\mathrm{\Psi }\right)$$
(28)
where $`i\overline{)}`$ is the usual Dirac differential operator in Minkowski space-time in terms of flat space time $`\gamma `$ matrices.
¿From the quadratic Lagrangian given above the quadratic quantum fluctuations to the energy momentum tensor can be extracted.
The effective potential is identified with $`T_0^0`$ in a spatially translational invariant state in which the expectation value of the inflaton field is $`\mathrm{\Phi }_0`$. During slow roll inflation the expectation value $`\mathrm{\Phi }_0`$ evolves very slowly in time, the slow roll approximation is indeed an adiabatic approximation, which justifies treating $`\mathrm{\Phi }_0`$ as a constant in order to obtain the effective potential. The time variation of $`\mathrm{\Phi }_0`$ only contributes to higher order corrections in slow-roll. This is standard in *any* calculation of an effective potential. The energy momentum tensor is computed in the FRW inflationary background determined by the *classical* inflationary potential $`V(\mathrm{\Phi }_0)`$, and the slow roll parameters are also explicit functions of $`\mathrm{\Phi }_0`$. Therefore the energy momentum tensor depends *implicitly* on $`\mathrm{\Phi }_0`$ through the background and *explicitly* on the masses for the light bosonic and fermionic fields given above.
Therefore the effective potential is given by
$$V_{eff}(\mathrm{\Phi }_0)=V(\mathrm{\Phi }_0)+\delta V(\mathrm{\Phi }_0)$$
(29)
where
$$\delta V(\mathrm{\Phi }_0)=T_0^0[\mathrm{\Phi }_0]_s+T_0^0[\mathrm{\Phi }_0]_t+T_0^0[\mathrm{\Phi }_0]_\sigma +T_0^0[\mathrm{\Phi }_0]_\mathrm{\Psi }$$
(30)
$`(s,t,\sigma ,\mathrm{\Psi })`$ correspond to the energy momentum tensors of the quadratic fluctuations of the scalar metric, tensor (gravitational waves), light boson field $`\sigma `$ and light fermion field $`\mathrm{\Psi }`$ fluctuations respectively. Since these are the expectation values of a quadratic energy momentum tensor, $`\delta V(\mathrm{\Phi }_0)`$ corresponds to the *one loop correction* to the effective potential.
### II.1 Light scalar fields
We begin by analyzing the contribution to the effective potential from the light bosonic scalar field $`\sigma `$ because this study highlights the main aspects which are relevant in the case of scalar metric (density) perturbations.
The bosonic Heisenberg field operators are expanded as follows
$$\sigma (\stackrel{}{x},\eta )=\frac{1}{C(\eta )\sqrt{\mathrm{\Omega }}}\underset{\stackrel{}{k}}{}e^{i\stackrel{}{k}\stackrel{}{x}}\left[a_{\sigma ,\stackrel{}{k}}S_\sigma (k,\eta )+a_{\sigma ,\stackrel{}{k}}^{}S_\sigma ^{}(k,\eta )\right]$$
(31)
where $`\mathrm{\Omega }`$ is the spatial volume.
During slow roll inflation the effective mass of the $`\sigma `$ field is given by eq. (24), just as for the inflaton fluctuation. It is convenient to introduce a parameter $`\eta _\sigma `$ defined to be
$$\eta _\sigma =\frac{M_\sigma ^2[\mathrm{\Phi }_0]}{3H_0^2}.$$
(32)
Hence, the statement that the $`\sigma `$ field is light corresponds to the condition
$$\eta _\sigma 1.$$
(33)
This dimensionless parameter plays the same role for the $`\sigma `$ field as the parameter $`\eta _V`$ given by eq. (9) does for the inflaton fluctuation.
The mode functions $`S_\sigma (k,\eta )`$ in eq. (31) obey the following equations up to quadratic ordernuestronor
$$S_\sigma ^{^{\prime \prime }}(k,\eta )+\left[k^2+M_\sigma ^2(\mathrm{\Phi }_0)C^2(\eta )\frac{C^{^{\prime \prime }}(\eta )}{C(\eta )}\right]S_\sigma (k,\eta )=0.$$
Using the slow roll expressions eq.(13) and in terms of $`\eta _\sigma `$, these mode equations become
$$S_\sigma ^{^{\prime \prime }}(k,\eta )+\left[k^2\frac{\nu _\sigma ^2\frac{1}{4}}{\eta ^2}\right]S_\sigma (k,\eta )=0;\nu _\sigma =\frac{3}{2}+ϵ_V\eta _\sigma +𝒪(ϵ_V^2,\eta _\sigma ^2,\eta _V^2,ϵ_V\eta _V).$$
During slow roll inflation $`\mathrm{\Phi }_0`$ is approximately constant, and the slow roll expansion is an *adiabatic* expansion. As usual in the slow roll approximation, the above equation for the mode functions is solved by assuming that $`\mathrm{\Phi }_0`$, hence $`\nu _\sigma `$ are *constant*. This is also the same type of approximation entailed in *every* calculation of the effective potential. Therefore during slow roll, the solution of the mode functions above are
$$S_\sigma (k,\eta )=\frac{1}{2}\sqrt{\pi \eta }e^{i\frac{\pi }{2}(\nu _\sigma +\frac{1}{2})}H_{\nu _\sigma }^{(1)}(k\eta ).$$
This choice of mode functions defines the Bunch-Davis vacuum, which obeys $`a_\stackrel{}{k}|0>_{BD}=0`$. It is important to highlight that there is no unique choice of vacuum or initial state, a recognition that has received considerable attention in the literature, see for exampleinistate ; mottola and references therein. In this study we focus on Bunch-Davis initial conditions since this has been the standard choice to study the power spectra and metric perturbations, hence we can compare our results to the standard ones in the literature, postponing for further study the assessment of different initial states.
The contribution to the effective potential from the light scalar field $`\sigma `$ is given by
$$T_0^0_\sigma =\frac{1}{2}\dot{\sigma }^2+\left(\frac{\sigma }{C(\eta )}\right)^2+M_\sigma ^2[\mathrm{\Phi }_0]\sigma ^2,$$
where the dot stands for derivative with respect to cosmic time. The expectation values are in the Bunch-Davis vacuum state and yield the following contributions
$`{\displaystyle \frac{1}{2}}\left(\dot{\sigma }\right)^2`$ $`=`$ $`{\displaystyle \frac{H_0^4}{16\pi }}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dz}{z}}z^2\left|{\displaystyle \frac{d}{dz}}\left[z^{\frac{3}{2}}H_{\nu _\sigma }^{(1)}(z)\right]\right|^2`$ (34)
$`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\sigma }{C^2(\eta )}}\right)^2`$ $`=`$ $`{\displaystyle \frac{H_0^4}{16\pi }}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dz}{z}}z^5\left|H_{\nu _\sigma }^{(1)}(z)\right|^2`$ (35)
$`{\displaystyle \frac{M_\sigma ^2[\mathrm{\Phi }_0]}{2}}\sigma ^2(\stackrel{}{x},t)`$ $`=`$ $`{\displaystyle \frac{3H_0^2\eta _\sigma }{2}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dk}{k}}𝒫_\sigma (k,t),`$ (36)
where $`𝒫_\sigma (k,t)`$ is the power spectrum of the $`\sigma `$ field, which in terms of the spatial Fourier transform of the field $`\sigma _\stackrel{}{k}(t)`$ is given by
$$𝒫_\sigma (k,t)=\frac{k^3}{2\pi ^2}\left|\sigma _\stackrel{}{k}^2(t)\right|=\frac{H_0^2}{8\pi }(k\eta )^3\left|H_{\nu _\sigma }^{(1)}(k\eta )\right|^2.$$
For a light scalar field during slow roll the power spectrum of the scalar field $`\sigma `$ is nearly scale invariant and the index $`\nu _\sigma 3/2`$. In the exact scale invariant case $`\nu _\sigma =3/2`$,
$$z^3\left|H_{\frac{3}{2}}^{(1)}(z)\right|^2=\frac{2}{\pi }[1+z^2]$$
and the integral of the power spectrum in eq. (36) not only features logarithmic and quadratic *ultraviolet* divergences but also a logarithmic *infrared* divergence. During slow roll and for a light but massive scalar field the quantity
$$\mathrm{\Delta }_\sigma =\frac{3}{2}\nu _\sigma =\eta _\sigma ϵ_V+𝒪(ϵ_V^2,\eta _\sigma ^2,ϵ_V\eta _\sigma ),1$$
is a measure of the departure from scale invariance and provides a natural *infrared regulator*. We note that the contribution from eq. (36) to the effective potential, which can be written as
$$\frac{3H_0^4\eta _\sigma }{16\pi }_0^{\mathrm{}}\frac{dz}{z}z^3\left|H_{\nu _\sigma }^{(1)}(z)\right|^2,$$
is *formally* smaller than the contributions from eqs.(34)-(35) by a factor $`\eta _\sigma 1`$. However, the logarithmic infrared divergence in the exact scale invariant case, leads to a single *pole* in the variable $`\mathrm{\Delta }_\sigma `$ as described in detail in refsnuestros ; nuestronor . To see this feature in detail, it proves convenient to separate the infrared contribution by writing the integral above in the following form
$$_0^{\mathrm{}}\frac{dz}{z}z^3\left|H_{\nu _\sigma }^{(1)}(z)\right|^2=_0^{\mu _p}\frac{dz}{z}z^3\left|H_{\nu _\sigma }^{(1)}(z)\right|^2+_{\mu _p}^{\mathrm{}}\frac{dz}{z}z^3\left|H_{\nu _\sigma }^{(1)}(z)\right|^2.$$
In the first integral we obtain the leading order contribution in the slow roll expansion, namely the pole in $`\mathrm{\Delta }_\sigma `$, by using the small argument limit of the Hankel functions
$$z^3\left|H_{\nu _\sigma }^{(1)}(z)\right|^2\stackrel{z0}{=}\left[\frac{2^{\nu _\sigma }\mathrm{\Gamma }(\nu _\sigma )}{\pi }\right]^2z^{2\mathrm{\Delta }_\sigma }$$
which yields
$$_0^{\mu _p}\frac{dz}{z}z^3\left|H_{\nu _\sigma }^{(1)}(z)\right|^2=\frac{2}{\pi }\left[\frac{1}{2\mathrm{\Delta }_\sigma }+\frac{\mu _p^2}{2}+\gamma 2+\mathrm{ln}(2\mu _p)+𝒪(\mathrm{\Delta }_\sigma )\right],$$
In the second integral for small but fixed $`\mu _p`$, we can safely set $`\mathrm{\Delta }_\sigma =0`$ and by introducing an upper momentum (ultraviolet) cutoff $`\mathrm{\Lambda }_p`$, we finally find
$$_0^{\mathrm{\Lambda }_p}\frac{dz}{z}z^3\left|H_{\nu _\sigma }^{(1)}(z)\right|^2=\frac{1}{\pi }\left[\frac{1}{\mathrm{\Delta }_\sigma }+\mathrm{\Lambda }_{p}^{}{}_{}{}^{2}+\mathrm{ln}\mathrm{\Lambda }_p^2+2\gamma 4+𝒪(\mathrm{\Delta }_\sigma )\right]$$
The simple pole in $`\mathrm{\Delta }_\sigma `$ reflects the infrared enhancement arising from a nearly scale invariant power spectrum. While the terms that depend on $`\mathrm{\Lambda }_p`$ are of purely ultraviolet origin and correspond to the specific regularization scheme, the simple pole in $`\mathrm{\Delta }_\sigma `$ originates in the *infrared* behavior and is therefore independent of the regularization scheme. A covariant regularization of the expectation value $`\sigma ^2(\stackrel{}{x},t)`$ will yield a result which features a simple pole in $`\mathrm{\Delta }_\sigma `$ plus terms which are ultraviolet finite and regular in the limit $`\mathrm{\Delta }_\sigma 0`$. Such regular terms yield a contribution $`𝒪(H^4\eta _\sigma )`$ to eq.(36) and are subleading in the limit of light scalar fields because they do not feature a denominator $`\mathrm{\Delta }_\sigma `$.
Therefore, to leading order in the slow roll expansion and in $`\eta _\sigma 1`$, the contribution from eq.(36) is given by,
$$\frac{M_\sigma ^2[\mathrm{\Phi }_0]}{2}\sigma ^2(\stackrel{}{x},t)=\frac{3H_0^4}{(4\pi )^2}\frac{\eta _\sigma }{\eta _\sigma ϵ_V}+\mathrm{subleading}\mathrm{in}\mathrm{slow}\mathrm{roll}.$$
In the first two contributions given by eqs.(34)-(35) extra powers of momentum arising either from the time or spatial derivatives, prevent the logarithmic infrared enhancements. These terms are infrared finite in the limit $`\mathrm{\Delta }_\sigma 0`$ and their leading contribution during slow roll can be obtained by simply setting $`\nu _\sigma =3/2`$ in these integrals, which feature quartic, quadratic and logarithmic ultraviolet divergences. A covariant renormalization of these two terms will lead to an ultraviolet and an infrared finite contribution to the energy momentum tensor of $`𝒪(H_0^4)`$, respectively. For the term given by eq.(36), the infrared contribution that yields the pole in $`\mathrm{\Delta }_\sigma `$ compensates for the $`\eta _\sigma 1`$ in the numerator, after renormalization of the ultraviolet divergence, the ultraviolet and infrared finite contributions to this term will yield a contribution to the energy momentum tensor of order $`𝒪(H_0^4\eta _\sigma )`$, without the small denominator, and therefore subleading. This analysis indicates that the leading order contributions to the energy momentum tensor for light scalar fields is determined by the infrared pole $`1/\mathrm{\Delta }_\sigma `$ from eq.(36) and the fully renormalized contributions from (34)-(35), namely to leading order in slow roll and $`\eta _\sigma `$
$$T_0^0_\sigma =\frac{3H_0^4}{(4\pi )^2}\frac{\eta _\sigma }{\frac{3}{2}\nu _\sigma }+\frac{1}{2}\dot{\sigma }^2+\left(\frac{\sigma }{C(\eta )}\right)^2_{ren}$$
(37)
In the expression above we have displayed explicitly the pole at $`3/2\nu _\sigma =\eta _\sigma ϵ_V`$.
In calculating the second term (renormalized expectation value) to leading order in eq.(37) we can set to zero the slow roll parameters $`ϵ_V,\eta _V`$ as well as the mass of the light scalar, namely $`\eta _\sigma =0`$. Hence, to leading order, the second term is identified with the $`00`$ component of the renormalized energy momentum tensor for a free massless minimally coupled scalar field in exact de Sitter space time. Therefore we can extract this term from the known result for the renormalized energy momentum tensor for a minimally coupled free scalar boson of mass $`m_\sigma `$ in de Sitter space time with a Hubble constant $`H_0`$ given byBD ; fordbunch ; sanchez
$`T_{\mu \nu }_{ren}`$ $`=`$ $`{\displaystyle \frac{g_{\mu \nu }}{(4\pi )^2}}\left\{m_\sigma ^2H_0^2\left(1{\displaystyle \frac{m_\sigma ^2}{2H_0^2}}\right)\left[\psi \left({\displaystyle \frac{3}{2}}+\nu \right)\psi \left({\displaystyle \frac{3}{2}}\nu \right)+\mathrm{ln}{\displaystyle \frac{m_\sigma ^2}{H_0^2}}\right]+{\displaystyle \frac{2}{3}}m_\sigma ^2H_0^2{\displaystyle \frac{29}{30}}H_0^4\right\},`$ (38)
$`\nu `$ $``$ $`\sqrt{{\displaystyle \frac{9}{4}}{\displaystyle \frac{m_\sigma ^2}{H_0^2}}}.`$ (40)
where $`\psi (z)`$ stands for the digamma function. This expression corrects a factor of two in ref.BD ; dowker . In eq. (6.177) in BD the D’Alambertian acting on $`G^1(x,x^{})`$ was neglected. However, in computing this term, the D’Alambertian must be calculated *before* taking the coincidence limit. Using the equation of motion yields the extra factor 2 and the expression eq.(40). This result eq.(40) for the renormalized energy momentum tensor was obtained by several different methods: covariant point splitting, zeta-function and Schwinger’s proper time regularizationsBD ; dowker .
The simple pole at $`\nu =3/2`$ manifest in eq.(40) coincides precisely with the similar simple pole in eq. (37) as can be gleaned by recognizing that $`m_\sigma ^2=3H^2\eta _\sigma `$ as stated by eq.(32). This pole originates in the term $`m_\sigma ^2<\sigma ^2>`$, which features an infrared divergence in the scaling limit $`\nu _\sigma =3/2`$. All the terms that contribute to the energy momentum tensor with space-time derivatives are infrared finite in this limit. Therefore, from the energy momentum tensor eq.(40) we can extract straightforwardly the leading contribution to the renormalized expectation value in eq.(37) in the limit $`H_0m_\sigma `$, and neglecting the slow roll corrections to the scale factor. It is given by the last term in the bracket in eq. (40). Hence, we find the leading order contribution
$$T_0^0_\sigma =\frac{H_0^4}{(4\pi )^2}\left[\frac{3\eta _\sigma }{\eta _\sigma ϵ_V}\frac{29}{30}+𝒪(ϵ_V,\eta _\sigma ,\eta _V)\right]$$
(41)
The last term is completely determined by the trace anomalyBD ; sanchez ; fordbunch ; duff ; dowker ; hartle ; fujikawa which is in turn determined by the short distance correlation function of the field and the background geometry.
Therefore, we emphasize that in the slow roll approximation there is a clean and unambiguous separation between the contribution from superhorizon modes, which give rise to simple poles in slow roll parameters and that of subhorizon modes whose leading contribution is determined by the trace anomaly and the short distance behavior of the field.
### II.2 Scalar metric perturbations
The gauge invariant energy momentum tensor for quadratic scalar metric fluctuations has been obtained in ref.abramo where the reader is referred to for details. In longitudinal gauge and in cosmic time it is given by
$`T_0^0_s=`$ $`M_{Pl}^2\left[12H_0\varphi \dot{\varphi }3(\dot{\varphi })^2+{\displaystyle \frac{9}{C^2(\eta )}}(\varphi )^2\right]`$ (42)
$`+{\displaystyle \frac{1}{2}}(\dot{\delta \phi })^2+{\displaystyle \frac{(\delta \phi )^2}{2C^2(\eta )}}+{\displaystyle \frac{V^{\prime \prime }(\mathrm{\Phi }_0)}{2}}(\delta \phi )^2+2V^{}(\mathrm{\Phi }_0)\varphi \delta \phi `$
where the condition $`\varphi =\psi `$ valid in scalar field inflation has been used, and the dots stand for derivatives with respect to cosmic time.
In longitudinal gauge, the equations of motion in cosmic time for the Fourier modes aremukhanov ; riotto
$`\ddot{\varphi }_\stackrel{}{k}+\left(H_02{\displaystyle \frac{\ddot{\mathrm{\Phi }}_0}{\dot{\mathrm{\Phi }}_0}}\right)\dot{\varphi }_\stackrel{}{k}+\left[2\left(\dot{H}_02H_0{\displaystyle \frac{\ddot{\mathrm{\Phi }}_0}{\dot{\mathrm{\Phi }}_0}}\right)+{\displaystyle \frac{k^2}{C^2(\eta )}}\right]\varphi _\stackrel{}{k}=0`$ (43)
(44)
$`\ddot{\delta \phi }_\stackrel{}{k}+3H\dot{\delta \phi }_\stackrel{}{k}+\left[V^{\prime \prime }[\mathrm{\Phi }_0]+{\displaystyle \frac{k^2}{C^2(\eta )}}\right]\delta \phi _\stackrel{}{k}+2V^{}[\mathrm{\Phi }_0]\varphi _\stackrel{}{k}4\dot{\mathrm{\Phi }}_0\dot{\varphi }_\stackrel{}{k}=0,`$ (45)
with the constraint equation
$$\dot{\varphi }_\stackrel{}{k}+H_0\varphi _\stackrel{}{k}=\frac{1}{2M_{Pl}}\delta \phi _\stackrel{}{k}\dot{\mathrm{\Phi }}_0.$$
(46)
Just as in the case of the scalar fields, we expect an infrared enhancement arising from superhorizon modes, therefore, following ref.abramo we split the contributions to the energy momentum tensor as those from superhorizon modes, which will yield the infrared enhancement, and the subhorizon modes for which we can set all slow roll parameters to zero. Just as discussed above for the case of the $`\sigma `$ field, since spatio-temporal derivatives bring higher powers of the momenta, we can neglect all derivative terms for the contribution from the superhorizon modes. Therefore, the contribution from superhorizon modes which will reflect the infrared enhancement is extracted fromabramo
$$T_0^0_{IR}\frac{1}{2}V^{\prime \prime }[\mathrm{\Phi }_0]\left(\delta \phi (\stackrel{}{x},t)\right)^2+2V^{}[\mathrm{\Phi }_0]\varphi (\stackrel{}{x},t)\delta \phi (\stackrel{}{x},t).$$
(47)
The analysis of the solution of eq.(43) for superhorizon wavelengths in ref. mukhanov shows that in exact de Sitter space time $`\varphi _\stackrel{}{k}\mathrm{constant}`$, hence it follows that during quasi-de Sitter slow roll inflation for superhorizon modes
$$\dot{\varphi }_\stackrel{}{k}(\mathrm{slow}\mathrm{roll})\times H_0\varphi _\stackrel{}{k}$$
(48)
Therefore, for superhorizon modes, the constraint equation (46) yields
$$\varphi _\stackrel{}{k}=\frac{V^{}(\mathrm{\Phi }_0)}{2V(\mathrm{\Phi }_0)}\delta \phi _\stackrel{}{k}+\mathrm{higher}\mathrm{orders}\mathrm{in}\mathrm{slow}\mathrm{roll}.$$
(49)
Inserting this relation in eq.(45) and consistently neglecting the term $`\dot{\varphi }_\stackrel{}{k}`$ according to eq.(48), we find the following equation of motion for the gauge invariant scalar field fluctuation in longitudinal gauge
$$\ddot{\delta \phi }_\stackrel{}{k}+3H_0\dot{\delta \phi }_\stackrel{}{k}+\left[\frac{k^2}{C^2(\eta )}+3H_0^2\eta _\delta \right]\delta \phi _\stackrel{}{k}=0,$$
(50)
where we have used the definition of the slow roll parameters $`ϵ_V;\eta _V`$ given in eq.(9), and introduced
$$\eta _\delta \eta _V2ϵ_V$$
(51)
This is the equation of motion for a minimally coupled scalar field with mass squared $`3H_0^2\eta _\delta `$ and we can use the results obtained in the case of the scalar field $`\sigma `$ above. The quantum field $`\delta \phi (\stackrel{}{x},t)`$ is expanded as
$$\delta \phi (\stackrel{}{x},\eta )=\frac{1}{C(\eta )\sqrt{\mathrm{\Omega }}}\underset{\stackrel{}{k}}{}e^{i\stackrel{}{k}\stackrel{}{x}}\left[a_{\delta ,\stackrel{}{k}}S_\delta (k,\eta )+a_{\delta ,\stackrel{}{k}}^{}S_\delta ^{}(k,\eta )\right],$$
(52)
where the mode functions are given by
$$S_\delta (k,\eta )=\frac{1}{2}\sqrt{\pi \eta }e^{i\frac{\pi }{2}(\nu _\delta +\frac{1}{2})}H_{\nu _\delta }^{(1)}(k\eta );\nu _\delta =\frac{3}{2}+ϵ_V\eta _\delta =\frac{3}{2}+3ϵ_V\eta _V.$$
(53)
In this case, the slow roll quantity that regulates the infrared behavior is $`\mathrm{\Delta }_\delta \eta _V3ϵ_V`$.
Again we choose the Bunch-Davies vacuum state annihilated by the operators $`a_{\delta ,\stackrel{}{k}}`$. Therefore, the contribution to $`T_0^0`$ from superhorizon modes to lowest order in slow roll is given by
$$T_0^0_{IR}=3H_0^2\left(\frac{\eta _V}{2}2ϵ_V\right)\left[_0^{\mathrm{}}\frac{dk}{k}𝒫_\delta (k,\eta )\right]_{IR}$$
(54)
where the power spectrum of scalar fluctuations is given by
$$𝒫_\delta (k,\eta )=\frac{k^3}{2\pi ^2}\left|\delta \phi _\stackrel{}{k}(t)\right|^2=\frac{H_0^2}{8\pi }(k\eta )^3\left|H_{\nu _\delta }^{(1)}(k\eta )\right|^2$$
(55)
and the subscript $`IR`$ in the integral refers only to the infrared pole contribution to $`\mathrm{\Delta }_\delta `$. Repeating the analysis presented in the case of the scalar field $`\sigma `$ above, we finally find
$$T_0^0_{IR}=\frac{3H_0^4}{(4\pi )^2}\frac{\eta _V4ϵ_V}{\eta _V3ϵ_V}+\mathrm{subleading}\mathrm{in}\mathrm{slow}\mathrm{roll}$$
(56)
For subhorizon modes with wavevectors $`ka(t)H_0`$, the solutions of the equation (43) aremukhanov
$$\varphi _\stackrel{}{k}(t)e^{\pm ik\eta }\dot{\varphi }_\stackrel{}{k}(t)\frac{ik}{a(t)}\varphi _\stackrel{}{k}(t)$$
(57)
For $`ka(t)H_0`$ the constraint equation (46) entails thatabramo
$$\varphi _\stackrel{}{k}(t)\frac{ia(t)}{2M_{Pl}k}\dot{\mathrm{\Phi }}_0\delta \phi _\stackrel{}{k}.$$
(58)
Replacing the expressions eqs.(57)-(58) in eq.(42) yields that all the terms featuring the gravitational potential $`\varphi `$ are suppressed with respect to those featuring the scalar field fluctuation $`\delta \phi `$ by powers of $`H_0a(t)/k1`$ as originally observed in ref.abramo . Therefore the contribution from subhorizon modes to $`T_{0s}^0`$ is given by
$$T_{0s}^0_{UV}\frac{1}{2}(\dot{\delta \phi })^2+\frac{(\delta \phi )^2}{2a^2}$$
(59)
where we have also neglected the term with $`V^{\prime \prime }[\mathrm{\Phi }_0]3H_0^2\eta _V`$ since $`k^2/a^2H_0^2`$ for subhorizon modes. Therefore, to leading order in slow roll we find the renormalized expectation value of $`T_{00s}`$ is given by
$$T_{0s}^0_{ren}\frac{3H_0^4}{(4\pi )^2}\frac{\eta _V4ϵ_V}{\eta _V3ϵ_V}+\frac{1}{2}\dot{\delta \phi }^2+\left(\frac{\delta \phi }{C(\eta )}\right)^2_{ren}$$
(60)
To obtain the renormalized expectation value in eq.(60) one can set all slow roll parameters to zero to leading order and simply consider a massless scalar field minimally coupled in de Sitter space time. This is precisely what we have already calculated in the case of the scalar field $`\sigma `$ above by using the known results in the literature for the covariantly renormalized energy momentum tensor of a massive minimally coupled fieldBD ; fordbunch ; sanchez ; dowker , and we can just borrow the result from eq.(41). We find the following final result to leading order in slow roll
$$T_{0s}^0_{ren}=\frac{H_0^4}{(4\pi )^2}\left[\frac{\eta _V4ϵ_V}{\eta _V3ϵ_V}\frac{29}{30}+𝒪(ϵ_V,\eta _\sigma ,\eta _V)\right]$$
(61)
The last term in eq. (61) is completely determined by the trace anomaly of a minimally coupled scalar field in de Sitter space timeBD ; sanchez ; duff ; dowker .
### II.3 Tensor perturbations
Tensor perturbations correspond to massless fields with two physical polarizations. The quantum fields are written as
$$h_j^i(\stackrel{}{x},\eta )=\frac{1}{C(\eta )M_{Pl}\sqrt{2\mathrm{\Omega }}}\underset{\lambda =\times ,+}{}\underset{\stackrel{}{k}}{}e^{i\stackrel{}{k}\stackrel{}{x}}ϵ_j^i(\lambda ,\stackrel{}{k})\left[a_{\lambda ,\stackrel{}{k}}S_h(k,\eta )+a_{\lambda ,\stackrel{}{k}}^{}S_h^{}(k,\eta )\right],$$
where the operators $`a_{\lambda ,\stackrel{}{k}},a_{\lambda ,\stackrel{}{k}}^{}`$ obey the usual canonical commutation relations, and $`ϵ_j^i(\lambda ,\stackrel{}{k})`$ are the two independent traceless-transverse tensors constructed from the two independent polarization vectors transverse to $`\widehat{𝐤}`$, chosen to be real and normalized such that $`ϵ_j^i(\lambda ,\stackrel{}{k})ϵ_k^j(\lambda ^{},\stackrel{}{k})=\delta _k^i\delta _{\lambda ,\lambda ^{}}`$.
The mode functions $`S_h(k,\eta )`$ obey the differential equation
$$S_h^{^{\prime \prime }}(k,\eta )+\left[k^2\frac{\nu _h^2\frac{1}{4}}{\eta ^2}\right]S_h(k,\eta )=0;\nu _h=\frac{3}{2}+ϵ_V+𝒪(ϵ_V^2,\eta _\sigma ^2,\eta _V^2,ϵ_V\eta _V)$$
(62)
The solutions corresponding to the Bunch-Davies vacuum annihilated by the operators $`a_{\lambda ,\stackrel{}{k}}`$ are
$$S_h(k,\eta )=\frac{1}{2}\sqrt{\pi \eta }e^{i\frac{\pi }{2}(\nu _h+\frac{1}{2})}H_{\nu _h}^{(1)}(k\eta ),$$
(63)
The energy momentum tensor for gravitons only depends on derivatives of the field $`h_j^i`$ therefore its expectation value in the Bunch Davies (BD) vacuum does not feature infrared singularities in the limit $`ϵ_V0`$. The absence of infrared singularities in the limit of exact de Sitter space time, entails that we can extract the leading contribution to the effective potential from tensor perturbations by evaluating the expectation value of $`T_{00}`$ in the BD vacuum in exact de Sitter space time, namely by setting all slow roll parameters to zero. This will yield the leading order in the slow roll expansion.
Because de Sitter space time is maximally symmetric, the expectation value of the energy momentum tensor is given byweinberg ; BD
$$T_{\mu \nu }_{BD}=\frac{g_{\mu \nu }}{4}T_\alpha ^\alpha _{BD}$$
(64)
and $`T_\alpha ^\alpha `$ is a space-time constant, therefore the energy momentum tensor is manifestly covariantly conserved. Of course, in a quantum field theory there emerge ultraviolet divergences and the regularization procedure must be compatible with the maximal symmetry. A large body of work has been devoted to study the trace anomaly in de Sitter space time implementing a variety of powerful covariant regularization methods that preserve the symmetryduff ; BD ; dowker ; hartle ; fujikawa ; sanchez yielding a renormalized value of the expectation value of the $`T_{\mu \nu }_{BD}`$ given by eq. (64). Therefore, the full energy momentum tensor is completely determined by the trace anomaly BD ; sanchez ; duff .
The contribution to the trace anomaly from gravitons has been given in refs.duff ; sanchez ; BD , it is
$$T_\alpha ^\alpha _t=\frac{717}{80\pi ^2}H_0^4$$
(65)
From this result, we conclude that
$$T_0^0_t=\frac{717}{320\pi ^2}H_0^4$$
(66)
This result differs by a numerical factor from that obtained in ref.finelli , presumably the difference is a result of a different regularization scheme.
### II.4 Fermion fields
The Dirac equation in the FRW geometry is given by \[see eq.(28)\],
$$\left[i\overline{)}M_\mathrm{\Psi }[\mathrm{\Phi }_0]C(\eta )\right]\left(C^{\frac{3}{2}}\mathrm{\Psi }(\stackrel{}{x},\eta )\right)=0.$$
(67)
The solution $`\mathrm{\Psi }(\stackrel{}{x},\eta )`$ can be expanded in spinor mode functions as
$$\mathrm{\Psi }(\stackrel{}{x},\eta )=\frac{1}{C^{\frac{3}{2}}(\eta )\sqrt{\mathrm{\Omega }}}\underset{\stackrel{}{k},\lambda }{}e^{i\stackrel{}{k}\stackrel{}{x}}\left[b_{\stackrel{}{k},\lambda }U_\lambda (\stackrel{}{k},\eta )+d_{\stackrel{}{k},\lambda }^{}V_\lambda (\stackrel{}{k},\eta )\right],$$
(68)
where the spinor mode functions $`U,V`$ obey the Dirac equations
$`\left[i\gamma ^0_\eta \stackrel{}{\gamma }\stackrel{}{k}M(\eta )\right]U_\lambda (\stackrel{}{k},\eta )`$ $`=`$ $`0`$ (69)
$`\left[i\gamma ^0_\eta +\stackrel{}{\gamma }\stackrel{}{k}M(\eta )\right]V_\lambda (\stackrel{}{k},\eta )`$ $`=`$ $`0`$ (70)
and
$$M(\eta )M_\mathrm{\Psi }[\mathrm{\Phi }_0]C(\eta )$$
(71)
Following the method of refs.boyarel ; baacke , it proves convenient to write
$`U_\lambda (\stackrel{}{k},\eta )`$ $`=`$ $`\left[i\gamma ^0_\eta \stackrel{}{\gamma }\stackrel{}{k}+M(\eta )\right]f_k(\eta )𝒰_\lambda `$ (72)
$`V_\lambda (\stackrel{}{k},\eta )`$ $`=`$ $`\left[i\gamma ^0_\eta +\stackrel{}{\gamma }\stackrel{}{k}+M(\eta )\right]g_k(\eta )𝒱_\lambda `$ (73)
with $`𝒰_\lambda ;𝒱_\lambda `$ being constant spinorsboyarel ; baacke obeying
$$\gamma ^0𝒰_\lambda =𝒰_\lambda ,\gamma ^0𝒱_\lambda =𝒱_\lambda $$
(74)
The mode functions $`f_k(\eta );g_k(\eta )`$ obey the following equations of motion
$`\left[{\displaystyle \frac{d^2}{d\eta ^2}}+k^2+M^2(\eta )iM^{}(\eta )\right]f_k(\eta )`$ $`=`$ $`0`$ (75)
$`\left[{\displaystyle \frac{d^2}{d\eta ^2}}+k^2+M^2(\eta )+iM^{}(\eta )\right]g_k(\eta )`$ $`=`$ $`0`$ (76)
Neglecting the derivative of $`\mathrm{\Phi }_0`$ with respect to time, namely terms of order $`\sqrt{ϵ_V}`$ and higher, the equations of motion for the mode functions are given by
$`\left[{\displaystyle \frac{d^2}{d\eta ^2}}+k^2{\displaystyle \frac{\nu _+^2\frac{1}{4}}{\eta ^2}}\right]f_k(\eta )`$ $`=`$ $`0`$ (77)
$`\left[{\displaystyle \frac{d^2}{d\eta ^2}}+k^2{\displaystyle \frac{\nu _{}^2\frac{1}{4}}{\eta ^2}}\right]g_k(\eta )`$ $`=`$ $`0`$ (78)
where
$$\nu _\pm =\frac{1}{2}\pm i\frac{M_\mathrm{\Psi }[\mathrm{\Phi }_0]}{H_0}$$
The scalar product of the spinors $`U_\lambda (\stackrel{}{k},\eta ),V_\lambda (\stackrel{}{k},\eta )`$ yields
$`U_\lambda ^{}(\stackrel{}{k},\eta )U_\lambda ^{}(\stackrel{}{k},\eta )`$ $`=`$ $`𝒞^+(k)\delta _{\lambda ,\lambda ^{}}`$ (79)
$`V_\lambda ^{}(\stackrel{}{k},\eta )V_\lambda ^{}(\stackrel{}{k},\eta )`$ $`=`$ $`𝒞^{}(k)\delta _{\lambda ,\lambda ^{}}`$
where
$`𝒞^+(k)`$ $`=`$ $`f_k^{^{}}(\eta )f_k^{}(\eta )+\left(k^2+M^2(\eta )\right)f_k^{}(\eta )f_k(\eta )+iM(\eta )\left(f_k^{}(\eta )f_k^{}(\eta )f_k(\eta )f_k^{^{}}(\eta )\right)`$ (81)
$`𝒞^{}(k)`$ $`=`$ $`g_k^{^{}}(\eta )g_k^{}(\eta )+\left(k^2+M^2(\eta )\right)g_k^{}(\eta )g_k(\eta )iM(\eta )\left(g_k^{}(\eta )g_k^{}(\eta )g_k(\eta )g_k^{^{}}(\eta )\right)`$
are constants of motion by dint of the equations of motion for the mode functions $`f_k(\eta ),g_k(\eta )`$. The normalized spinor solutions of the Dirac equation are therefore given by
$`U_\lambda (\stackrel{}{k},\eta )={\displaystyle \frac{1}{\sqrt{𝒞^+(k)}}}\left[if_k^{}(\eta )\stackrel{}{\gamma }\stackrel{}{k}f_k(\eta )+M(\eta )f_k(\eta )\right]𝒰_\lambda `$ (83)
(84)
$`V_\lambda (\stackrel{}{k},\eta )={\displaystyle \frac{1}{\sqrt{𝒞^{}(k)}}}\left[ig_k^{}(\eta )+\stackrel{}{\gamma }\stackrel{}{k}g_k(\eta )+M(\eta )g_k(\eta )\right]𝒰_\lambda `$
We choose the solutions of the mode equations (77)-(78) to be
$$f_k(\eta )=\sqrt{\frac{\pi k\eta }{2}}e^{i\frac{\pi }{2}(\nu _++\frac{1}{2})}H_{\nu _+}^{(1)}(k\eta ),g_k(\eta )=\sqrt{\frac{\pi k\eta }{2}}e^{i\frac{\pi }{2}(\nu _{}+\frac{1}{2})}H_\nu _{}^{(2)}(k\eta )$$
(85)
We also choose the Bunch-Davies vacuum state such that $`b_{\stackrel{}{k},\lambda }|0>_{BD}=0;d_{\stackrel{}{k},\lambda }|0>_{BD}=0`$. The choice of the mode functions eq.(85) yield the following normalization factors
$$𝒞^+(k)=𝒞^{}(k)=2k^2.$$
The energy momentum tensor for a spin $`1/2`$ field is given byBD
$$T_{\mu \nu }=\frac{i}{2}\left[\overline{\mathrm{\Psi }}\gamma _{(\mu }\underset{\nu )}{\overset{}{𝒟}}\mathrm{\Psi }\right]$$
and its expectation value in the Bunch-Davis vacuum is equal to
$$T_0^0_{BD}=\frac{2}{C^4(\eta )}\frac{d^3k}{(2\pi )^3}\left\{M(\eta )\mathrm{Im}\left[g_k^{}(\eta )g_k^{}(\eta )\right]\right\}$$
where $`M(\eta )`$ and $`g_k(\eta )`$ are given by eqs.(71) and (85), respectively. It is clear that this energy momentum tensor does not feature any infrared sensitivity because the index of the Bessel functions is $`\nu _\pm 1/2`$. Of course this is expected since fermionic fields cannot feature large amplitudes due to the Pauli principle.
A lengthy computation using covariant point splitting regularization yields the following result
$$T_0^0_\mathrm{\Psi }=\frac{11H_0^4}{960\pi ^2}\{1+\frac{120}{11}^2(^2+1)[\mathrm{Re}\psi (2+i)\frac{19}{12}\gamma 2\mathrm{ln}2]\},\frac{M_\mathrm{\Psi }[\mathrm{\Phi }_0]}{H_0}$$
(86)
The first term in the bracket in eq.(86) is recognized as the trace anomaly for fermions and is the only term that survives in the massless limitduff ; hartle ; BD ; fujikawa ; dowker ; sanchez . For light fermion fields, $`1`$, and the leading contribution to the energy momentum tensor is completely determined by the trace anomaly, hence in this limit the contribution to the covariantly regularized effective potential from (Dirac) fermions is given by
$$T_0^0_\mathrm{\Psi }=\frac{11H_0^4}{960\pi ^2}\left[1+𝒪(^2)\right]$$
This result is valid for Dirac fermions and it must be divided by a factor 2 for Weyl or Majorana fermions.
### II.5 Summary
In summary, we find that the effective potential at one-loop is given by,
$$\delta V(\mathrm{\Phi }_0)=\frac{H_0^4}{(4\pi )^2}\left[\frac{\eta _V4ϵ_V}{\eta _V3ϵ_V}+\frac{3\eta _\sigma }{\eta _\sigma ϵ_V}+𝒯_\mathrm{\Phi }+𝒯_s+𝒯_t+𝒯_\mathrm{\Psi }+𝒪(ϵ_V,\eta _V,\eta _\sigma ,^2)\right],$$
where $`(s,t,\sigma ,\mathrm{\Psi })`$ stand for the contributions of the scalar metric, tensor fluctuations, light boson field $`\sigma `$ and light fermion field $`\mathrm{\Psi }`$, respectively. We have
$`𝒯_\mathrm{\Phi }`$ $`=`$ $`𝒯_s={\displaystyle \frac{29}{30}}`$ (87)
$`𝒯_t`$ $`=`$ $`{\displaystyle \frac{717}{5}}`$ (89)
$`𝒯_\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \frac{11}{60}}`$
The terms that feature the *ratios* of combinations of slow roll parameters arise from the infrared or superhorizon contribution from the scalar density perturbations and scalar fields $`\sigma `$ respectively. The terms $`𝒯_{s,t,\mathrm{\Psi }}`$ are completely determined by the trace anomalies of scalar, graviton and fermion fields respectively. Writing $`H_0^4=V(\mathrm{\Phi }_0)H_0^2/[3M_{Pl}^2]`$ we can finally write the effective potential to leading order in slow roll
$$V_{eff}(\mathrm{\Phi }_0)=V(\mathrm{\Phi }_0)\left[1+\frac{H_0^2}{3(4\pi )^2M_{Pl}^2}\left(\frac{\eta _V4ϵ_V}{\eta _V3ϵ_V}+\frac{3\eta _\sigma }{\eta _\sigma ϵ_V}\frac{2903}{20}\right)\right]$$
(90)
There are several remarkable aspects of this result:
i) the infrared enhancement as a result of the near scale invariance of scalar field fluctuations, both from scalar density perturbations as well as from a light scalar field, yield corrections of *zeroth order in slow roll*. This is a consequence of the fact that during slow roll the particular combination $`\mathrm{\Delta }_\sigma =\eta _\sigma ϵ_V`$ of slow roll parameters yield a natural infrared cutoff.
ii) the final one loop contribution to the effective potential displays the effective field theory dimensionless parameter $`H_0^2/M_{Pl}^2`$ confirming our previous studiesnuestros ; nuestronor ,
iii) the last term is completely determined by the trace anomaly, a purely geometric result of the short distance properties of the theory.
## III Quantum Corrections to the Curvature and Tensor Fluctuations
The quantum corrections to the effective potential lead to quantum corrections to the amplitude of scalar and tensor fluctuations.
The scalar curvature and tensor fluctuations in the slow-roll regime are giving by the formulasliddle
$$|\mathrm{\Delta }_k^{(S)}|^2=\frac{1}{8\pi ^2ϵ_V}\left(\frac{H}{M_{Pl}}\right)^2,|\mathrm{\Delta }_k^{(T)}|^2=\frac{1}{2\pi ^2}\left(\frac{H}{M_{Pl}}\right)^2.$$
(91)
where $`H`$ stands for the Hubble parameter and $`ϵ_V`$ is given by eq.(9).
We can include the leading quantum corrections in eq.(91) replacing in it $`H`$ and $`ϵ_V`$ by the corrected parameters $`H_{eff}`$ and $`ϵ_{eff}`$. That is,
$$H_{eff}^2=H_0^2+\delta H^2,ϵ_{eff}=ϵ_V+\delta ϵ_V$$
(92)
with
$$H_{eff}^2=\frac{V_{eff}(\mathrm{\Phi }_0)}{3M_{Pl}^2},ϵ_{eff}=\frac{M_{Pl}^2}{2}\left[\frac{V_{eff}^{^{}}(\mathrm{\Phi }_0)}{V_{eff}(\mathrm{\Phi }_0)}\right]^2,$$
(93)
and where $`V_{eff}(\mathrm{\Phi }_0)`$ is given by eq.(90). We thus obtain,
$`{\displaystyle \frac{\delta H^2}{H_0^2}}={\displaystyle \frac{1}{3}}\left({\displaystyle \frac{H_0}{4\pi M_{Pl}}}\right)^2\left[{\displaystyle \frac{\eta _V4ϵ_V}{\eta _V3ϵ_V}}+{\displaystyle \frac{3\eta _\sigma }{\eta _\sigma ϵ_V}}{\displaystyle \frac{2903}{20}}\right],`$ (94)
(95)
$`{\displaystyle \frac{\delta ϵ_V}{ϵ_V}}={\displaystyle \frac{2}{3}}\left({\displaystyle \frac{H_0}{4\pi M_{Pl}}}\right)^2\left\{{\displaystyle \frac{\xi _V+12ϵ_V\left(2ϵ_V\eta _V\right)}{2\left(\eta _V3ϵ_V\right)^2}}+{\displaystyle \frac{3\eta _\sigma }{\left(\eta _\sigma ϵ_V\right)^2}}\left[\eta _\sigma +\eta _V2ϵ_V\sqrt{2ϵ_V}M_{Pl}{\displaystyle \frac{d\mathrm{log}M_\sigma [\mathrm{\Phi }_0]}{d\mathrm{\Phi }_0}}\right]{\displaystyle \frac{2903}{20}}\right\}`$
Inserting eq.(94) into eqs.(92) and (93) yields after calculation, for the scalar perturbations,
$`|\mathrm{\Delta }_{k,eff}^{(S)}|^2=|\mathrm{\Delta }_k^{(S)}|^2\left[1{\displaystyle \frac{\delta ϵ_V}{ϵ_V}}+{\displaystyle \frac{\delta H^2}{H^2}}\right]=`$ (96)
(97)
$`=|\mathrm{\Delta }_k^{(S)}|^2\{1{\displaystyle \frac{1}{3}}\left({\displaystyle \frac{H_0}{4\pi M_{Pl}}}\right)^2[{\displaystyle \frac{\xi _V+12ϵ_V^2\eta _V^25ϵ_V\eta _V}{\left(\eta _V3ϵ_V\right)^2}}+`$ (98)
(99)
$`+{\displaystyle \frac{3\eta _\sigma }{\left(\eta _\sigma ϵ_V\right)^2}}[\eta _\sigma 3ϵ_V+2\eta _V2\sqrt{2ϵ_V}M_{Pl}{\displaystyle \frac{d\mathrm{log}M_\sigma [\mathrm{\Phi }_0]}{d\mathrm{\Phi }_0}}]{\displaystyle \frac{2903}{20}}]\},`$ (100)
and for the tensor perturbations,
$`|\mathrm{\Delta }_{k,eff}^{(T)}|^2=|\mathrm{\Delta }_k^{(T)}|^2\left[1+{\displaystyle \frac{\delta H^2}{H^2}}\right]=`$ (101)
(102)
$`=|\mathrm{\Delta }_k^{(T)}|^2\left\{1+{\displaystyle \frac{1}{3}}\left({\displaystyle \frac{H_0}{4\pi M_{Pl}}}\right)^2\left[{\displaystyle \frac{\eta _V4ϵ_V}{\eta _V3ϵ_V}}+{\displaystyle \frac{3\eta _\sigma }{\eta _\sigma ϵ_V}}{\displaystyle \frac{2903}{20}}\right]\right\}.`$ (103)
where $`M_\sigma [\mathrm{\Phi }_0]`$ and $`\eta _\sigma `$ are given by eqs.(24) and (32), respectively.
The case when the field $`\sigma `$ is much lighter than the inflaton permit simplifications, since
$$\eta _\sigma \left(\frac{m_\sigma }{m_{inflaton}}\right)^2\eta _V,$$
for $`m_\sigma ^2m_{inflaton}^2`$, we can neglect terms proportional to $`\eta _\sigma `$ in the expressions for $`|\mathrm{\Delta }_k^{(S)}|^2`$ and $`|\mathrm{\Delta }_{k,eff}^{(T)}|^2`$. In this case the quantum corrections to the power spectra obtain a particularly illuminating expression when the slow-roll parameters in eqs.(96)-(101) are written in terms of the CMB observables $`n_s,r`$ and the spectral running of the scalar index using
$`ϵ_V={\displaystyle \frac{r}{16}},\eta _V={\displaystyle \frac{1}{2}}(n_s1+{\displaystyle \frac{3}{8}}r),`$ (104)
(105)
$`\xi _V={\displaystyle \frac{r}{4}}(n_s1+{\displaystyle \frac{3}{16}}r){\displaystyle \frac{1}{2}}{\displaystyle \frac{dn_s}{d\mathrm{ln}k}},\eta _V3ϵ_V={\displaystyle \frac{1}{2}}(n_s1).`$ (106)
We find from eqs.(96)-(101),
$`|\mathrm{\Delta }_{k,eff}^{(S)}|^2=|\mathrm{\Delta }_k^{(S)}|^2\left\{1+{\displaystyle \frac{2}{3}}\left({\displaystyle \frac{H_0}{4\pi M_{Pl}}}\right)^2\left[1+{\displaystyle \frac{\frac{3}{8}r(n_s1)+2\frac{dn_s}{d\mathrm{ln}k}}{(n_s1)^2}}+{\displaystyle \frac{2903}{40}}\right]\right\}`$ (107)
(108)
$`|\mathrm{\Delta }_{k,eff}^{(T)}|^2=|\mathrm{\Delta }_k^{(T)}|^2\left\{1{\displaystyle \frac{1}{3}}\left({\displaystyle \frac{H_0}{4\pi M_{Pl}}}\right)^2\left[1+{\displaystyle \frac{1}{8}}{\displaystyle \frac{r}{n_s1}}+{\displaystyle \frac{2903}{20}}\right]\right\}.`$ (109)
We see that the anomalies contribution $`\frac{2903}{40}=72.575`$ and $`\frac{2903}{20}=145.15`$ presumably dominate both quantum corrections. The other terms are generally expected to be smaller than these large contributions from the anomalies. These anomalous contributions are dominated in turn by the tensor part \[see eq.(89)\]. Only fermions give contributions with the opposite sign. However, one needs at least $`783`$ species of (Dirac) Fermions to compensate for the tensor part.
These quantum corrections also affect the ratio $`r`$ of tensor/scalar fluctuations as follows,
$$r_{eff}\frac{|\mathrm{\Delta }_{k,eff}^{(T)}|^2}{|\mathrm{\Delta }_{k,eff}^{(S)}|^2}=r\left\{1\frac{1}{3}\left(\frac{H_0}{4\pi M_{Pl}}\right)^2\left[1+\frac{\frac{3}{8}r(n_s1)+\frac{dn_s}{d\mathrm{ln}k}}{(n_s1)^2}+\frac{8709}{20}\right]\right\}$$
(110)
We expect this quantum correction to the ratio to be negative as the anomaly contribution dominates: $`\frac{8709}{20}=435.45`$.
Therefore, the quantum corrections enhance the scalar curvature fluctuations while they reduce the tensor fluctuations as well as their ratio $`r`$. The quantum corrections are small, of the order $`\left(\frac{H_0}{M_{Pl}}\right)^2`$, but it is interesting to see that the quantum effects are dominated by the trace anomalies and they correct both fluctuations in a definite direction.
## IV Conclusions
Motivated by the premise that forthcoming CMB observations may probe the inflationary potential, we study its quantum corrections from scalar and tensor metric perturbations as well as those from one light scalar and one light (Dirac) fermion field generically coupled to the inflaton. The reason for this study is that the measurements probe the full effective inflaton potential, namely the classical potential plus its quantum corrections. We have focused on obtaining the quantum corrections to the effective potential during the cosmologically relevant quasi-de Sitter stage of slow roll inflation. Both, scalar metric fluctuations, as well as those from a light scalar field, feature infrared enhancements as a consequence of the nearly scale invariance of their power spectra. A combination of slow roll parameters appropriate for each case provides a natural infrared regularization. We find that to leading order in slow roll, there is a clean and unambiguous separation between the contributions to the effective potential from superhorizon modes of the scalar metric perturbations as well as the scalar field, and those from subhorizon modes. Only the contributions to the total energy momentum tensor from curvature perturbations and the light scalar field feature an infrared enhancement, while those from gravitational waves and fermions do not feature any infrared sensitivity.
In all cases, scalar metric, tensor, light scalar and fermion fields, the contribution from subhorizon modes is determined by the trace anomaly, while the contribution from superhorizon modes, only relevant for curvature and scalar field perturbations, are infrared enhanced as the inverse of a combination of slow roll parameters which measure the departure from scale invariance in each case.
The one loop effective potential to leading order in slow roll is given by eq. (90). The last term, independent of the slow roll parameters, is completely determined by the trace anomalies of scalar, tensor and fermionic fields (therefore solely determined by the space-time geometry), while the first term reveals the hallmark infrared enhancement of superhorizon fluctuations. Using this result we have obtained the one-loop quantum corrections to the amplitude of (scalar) curvature and tensor perturbations in terms of the CMB observables $`n_s,r`$ and $`dn_s/d\mathrm{ln}k`$ and the total trace anomaly $`𝒯`$ of the different fields.
As we anticipated in ref.nuestros ; nuestronor , the strength of the one loop corrections is determined by the effective field theory parameter $`(H_0/M_{Pl})^2`$. While this quantity is observationally of $`𝒪(10^{10})`$, there is an important message in this result: the robustness of slow-roll inflation as well as the reliability of the effective field theory description.
There is a simple interpretation of the above result: in the effective field theory approach, the ‘classical’ inflaton potential $`V(\mathrm{\Phi }_0)`$ includes contributions from integrating out the fields with scales much heavier than the scale of inflation $`H_0`$.
The contribution to the energy momentum from light fields, yield the effective potential, however for fields with mass scales $`H_0`$ the dominant scale in the problem is $`H_0`$ and on dimensional grounds the contribution to the covariantly renormalized energy momentum tensor must be $`H_0^4`$. This argument would fail in the presence of *infrared* divergences, and indeed the mass terms from curvature perturbations and from the scalar field $`\sigma `$ feature an infrared enhancement because of their nearly scale invariant power spectrum. The mass term is of first order in the slow roll, however the infrared enhancement brings about a denominator which is also of first order in slow roll yielding a ratio which is of zeroth order in slow roll. Hence, this remarkable result validates the simple power counting that yields the overall scale $`H_0^4`$ for the one loop correction.
An important bonus of the slow roll approximation is that the contributions from superhorizon and subhorizon modes can be *unambiguously separated* and the latter are completely determined by the trace anomaly, a purely geometrical result which only depends on the short distance (ultraviolet) properties. Quantum trace anomalies of the energy momentun tensor in gravitational fields constitute a nice and important chapter of QFT in curved backgrounds, (see BD and references therein). Our results here show that these trace anomalies dominate the quantum corrections to a relevant cosmological problem: the primordial power spectrum of curvature and tensor fluctuations.
###### Acknowledgements.
D.B. thanks the US NSF for support under grant PHY-0242134, and the Observatoire de Paris and LERMA for hospitality during this work. He also thanks L. R. Abramo for interesting discussions. This work is supported in part by the Conseil Scientifique de l’Observatoire de Paris through an ‘Action Initiative, BQR’. |
warning/0507/hep-th0507244.html | ar5iv | text | # CERN–PH–TH/2005-136IC/2005/042; LMU-ASC 55/05hep-th/0507244 Open string topological amplitudes and gaugino masses
## 1 Introduction
Possible relation between string theory and topological field theories may have profound origin and help in understanding the non-perturbative nature of the underlying fundamental theory. A concrete example is the topological partition function of the twisted Calabi-Yau sigma-model, $`F^{(g,h)}`$, defined on a general Riemann surface of genus $`g`$ having $`h`$ boundaries . In the absence of holes, it is known that $`F^{(g,0)}`$ describes a sequence of higher derivative F-terms, $`𝒲^{2g}`$, in the effective four-dimensional (4d) $`N=2`$ supergravity, obtained upon compactification of Type II superstring on the corresponding Calabi-Yau manifold . Here, $`𝒲`$ is the chiral superfield of the $`N=2`$ gravity multiplet. Although $`F^{(g,0)}`$ is expected to be a holomorphic function of the chiral moduli, there is an anomaly which is captured by a set of recursion relations . This non-holomorphicity is induced by an anomaly of the BRST current in the topological theory, while from the space-time point of view, it is due to the propagation of massless states that lead to non-localities in the effective action .
In this work, we perform a similar study of the genus-0 series $`F^{(0,h)}`$. In the presence of boundaries, $`h0`$, the topological twist involves D-branes, and thus supersymmetry is broken to $`N=1`$. The partition function $`F^{(0,h)}`$ describes a sequence of higher derivative F-terms, $`(\mathrm{Tr}W^2)^{h1}`$, where $`W`$ is the chiral $`N=1`$ gauge superfield . They give rise to amplitudes involving two gauge fields and $`2(h2)`$ gauginos. In particular, we rederive the relation between the scattering amplitudes and the topological partition function by using the Neveu-Schwarz-Ramond formalism, similarly to Ref.. The non-trivial property of this result is that in the physical amplitudes, besides the cancellation of the string oscillator contributions, there is a cancellation of the prefactor corresponding to the non-compact 4d momentum integration. This cancellation may not necessarily extend to higher genus amplitudes (involving additional insertions of closed string fields), although it certainly holds in the low energy effective field theory limit .
The same amplitudes have been computed in the past, in the context of Calabi-Yau compactifications of the heterotic string and were shown to be identical to the partition function of the topological theory obtained by twisting the left-moving supersymmetric sector which has $`N=2`$ superconformal symmetry . However, unlike in the Type II case, the recursion relations describing the holomorphic anomaly of the heterotic topological partition function do not close among themselves. They involve a new class of topological quantities $`F_n^{(0,h)}`$, corresponding to correlation functions of anti-chiral fields, that describe F-terms of the type $`\mathrm{\Pi }^n(\mathrm{Tr}W^2)^{h1}`$, where $`\mathrm{\Pi }`$’s are chiral projections of non-holomorphic functions of $`N=1`$ chiral superfields. Using Type I – heterotic string duality , one expects that the heterotic topological partition function at genus $`h1`$ coincides with $`F^{(0,h)}`$ in the corresponding perturbative limit. On the Type I side, the non closure of the recursion relations among $`F^{(0,h)}`$’s is due to the existence of only one BRST charge on the boundaries, instead of two acting separately on left and right movers.
An important property of $`F^{(0,h)}`$ with $`h3`$ is the breaking of R-symmetry. A particularly interesting term is $`(\mathrm{Tr}W^2)^2`$. Combined with supersymmetry breaking induced by vacuum expectation values (VEV’s) of R-preserving D-term auxiliaries, it can generate Majorana gaugino masses, $`m_{1/2}D^2m_0^4`$ (in string units), where $`m_0`$ is the scalar mass scale. In this work, we present an explicit calculation of $`F^{(0,3)}`$ on a world-sheet with three boundaries, in a simple example of Type IIB compactified on $`T^6`$ with magnetized D9 branes, or equivalently, upon T-duality, of intersecting D6 branes in Type IIA orientifolds . In a vacuum configuration preserving $`N=1`$ supersymmetry, a non trivial $`F^{(0,3)}`$ is generated if one of the three brane stacks associated to the 3 boundaries is displaced away from the intersection of the other two; this implies the breaking of R-symmetry in the above example. On the other hand, if the brane configuration breaks supersymmetry, Majorana gaugino masses are generated. In the limit of small scalar masses, they behave as $`m_{1/2}m_0^4`$ with a coefficient given precisely by the topological partition function $`F^{(0,3)}`$. Similarly, the $`\mathrm{\Pi }\mathrm{Tr}W^2`$ term associated to the topological amplitude $`F_1^{(0,2)}`$ generates, upon supersymmetry breaking, one-loop masses for non-chiral brane fermions, of the same order as the gaugino masses.
Actually, the fact that a non-vanishing gaugino mass is generated at the perturbative order $`(g=0,h=3)`$, which corresponds to Euler characteristic $`1`$, is in agreement with the general arguments of Ref., based on $`N=2`$ superconformal $`U(1)`$ charge conservation. It is the same order as $`(g=1,h=1)`$, which corresponds to the gravitational mediation of supersymmetry breaking from the bulk to the gauge (brane) sector. This contribution depends strongly on the mechanism of supersymmetry breaking in the closed string sector , ignored in the present work. From the effective field theory point of view, the string diagram $`(g=0,h=3)`$ describes a particular gauge mediation by two-loop corrections in the gauge D-brane theory.
The paper is organized as follows.
$``$ In Section 2, we describe the moduli space of the relevant world-sheet of genus zero with $`h`$ boundaries, $`(g=0,h)`$, as obtained by an appropriate involution of the Riemann surface of genus $`g=h1`$ with no boundaries, $`(g=h1,0)`$ . We also discuss the partition function and the so-called correction factors associated to Neumann and Dirichlet boundary conditions .
$``$ In Section 3, we compute a string amplitude involving two gauge fields and $`2(h2)`$ gauginos on a world-sheet of genus 0 with $`h`$ boundaries. We thus extract the coefficient of the F-term $`(\mathrm{Tr}W^2)^{h1}`$ in the effective action, as a function of the closed string moduli parameterizing the Calabi-Yau geometry. For simplicity, we perform the computation for general $`N=1`$ intersecting brane configurations in $`N=2`$ orbifold compactifications of Type II superstrings. The generalization to Calabi-Yau is straightforward using the results of Ref.. We show that this coupling-coefficient is given by the topological partition function $`F^{(0,h)}`$ of the associated twisted Calabi-Yau sigma-model.
$``$ In Section 4, we discuss the holomorphic anomaly and the corresponding recursion relations. We first review the situation in the heterotic theory and then identify the possible contributions to the anomaly coming from various degeneration limits (dividing geodesic and handle) involving massless open and closed string states.
$``$ In Section 5, we consider an explicit example of Type I string compactifications on a factorized $`T^6=(T^2)^3`$ (instead of general orbifolds) with magnetized D9 branes, or equivalently with D6 branes at angles. We then compute $`F^{(0,3)}`$ in the case where the angles are chosen to preserve $`N=1`$ supersymmetry. We show that a non-zero answer is obtained if one of the three brane stacks associated to the three boundaries is moved away from the intersection of the other two, in each of the three tori. For non-parallel branes, this deformation does not change the spectrum, but breaks all R-symmetries. $`F^{(0,3)}`$ is expressed in a compact integral form as a function of $`T^6`$ moduli and Wilson lines.
$``$ In Section 6, we compute the Majorana gaugino masses for arbitrary brane angles that break supersymmetry, originating from the same world-sheet with three boundaries. We show that in the $`N=1`$ supersymmetric limit, they behave as $`m_0^4`$, with $`m_0`$ the scalar mass, with a coefficient given by the topological partition function $`F^{(0,3)}`$, in accordance with the result of Section 5.
$``$ In Section 7, we compute certain matter mass terms that are related to the F-terms appearing in the holomorphic anomaly of the gaugino mass $`F^{(0,3)}`$. These terms involve (twisted) charged fields from brane intersections.
$``$ Section 8 contains some explicit examples. In particular, we consider the case of three stacks of branes forming pairwise three different $`N=2`$ supersymmetric sectors. We can then evaluate the $`F^{(0,3)}`$ integral explicitly, as well as all quantities related to it by the holomorphic anomaly, namely $`F_1^{(0,2)}`$ and $`F_2^{(0,1)}`$.
$``$ Our results are summarized in Section 9 which also contains discussions on possible applications and open problems.
$``$ Finally in the Appendix, we derive the lattice contribution of the twisted bosons to the topological partition functions $`F^{(0,h)}`$ for the case of magnetized D9 branes. This involves periods associated with Prym differentials and to our knowledge a detailed treatment of the zero modes associated with the twisted bosons for the open string case has not been made in the physics literature (for the closed string case see Ref.). The result of this Appendix for the special case $`h=3`$ is used in sections 5 and 6.
## 2 Genus $`g=0`$ World-Sheet with $`h`$ Boundaries From the Involution of $`g=h1`$ Riemann Surface
A planar ($`h1`$)-loop open string world-sheet with $`h`$ boundaries, see Fig. 1,
admits a closed orientable genus $`h1`$ double cover. The embedding is described by an anti-conformal involution, with the boundaries made of its fixed points. In the canonical homology basis, the a-cycles are invariant, $`𝐚_i\overline{𝐚}_i=𝐚_i`$, while the b-cycles have their orientation reverted, $`𝐛_i\overline{𝐛}_i=𝐛_i^1`$, see Fig. 2.
The anti-symplectic involution matrix has the form
$$I=\left(\begin{array}{cc}11& 0\\ 0& 11\end{array}\right),$$
(2.1)
where $`11`$ is the $`(h1)\times (h1)`$ identity matrix. In terms of the a-cycles of the double cover, the boundaries of the $`(g=0,h)`$ world-sheet are $`\alpha _1=𝐚_1,\mathrm{},\alpha _k=𝐚_k𝐚_{k1}^1,\mathrm{},\alpha _h=𝐚_{h1}^1`$. It can be made into a contractible region by cutting along the curves $`\gamma _i`$, $`i=1,\mathrm{},h1`$, as shown in Fig. 1. The corresponding fundamental polygon is shown in Fig. 3.
Under the involution, $`\gamma _i\overline{\gamma }_i`$, such that $`\gamma _i\overline{\gamma }_i^1=_{k=1}^{k=i}𝐛_k`$.
The period matrix $`\tau `$ of the double cover is restricted by the condition of invariance under the involution: $`\overline{\tau }=I(\tau )=\tau `$, thus
$$\tau =it,$$
(2.2)
where $`t`$ is a real, symmetric matrix. Note that the positivity requirement for the period matrix now reads, in particular, $`dett>0`$. The modular transformations that preserve the involution are represented by $`Sp(2g,)`$ matrices of the form
$$M=\left(\begin{array}{cc}(D^1)^T& 0\\ 0& D\end{array}\right).$$
(2.3)
Since $`D`$ and $`D^1`$ must be integer-valued, the “relative modular group” is $`GL(g,)`$. This group essentially interchanges the boundaries, transforming the period matrix as $`\tau D\tau D^T`$. Note that since $`detD=\pm 1`$, the determinant $`dett`$ remains invariant.
The double cover endows the world-sheet with the basis of holomorphic differentials $`\omega _i`$ normalized as
$$_{𝐚_j}\omega _i=\delta _{ij},_{𝐛_j}\omega _i=it_{ij},$$
(2.4)
The involution transforms them into anti-holomorphic differentials
$$\overline{\omega _i(z)}=\overline{\omega }_i(\overline{z});_{\gamma _j}\omega _i\overline{\omega }_i=i\underset{k=1}{\overset{k=j}{}}t_{ik}.$$
(2.5)
Furthermore, the normalization condition (2.4) implies
$$_{\alpha _m}\omega _n=\delta _{mn}\delta _{m(n1)}.$$
(2.6)
For future use, we also need the surface integrals
$$\omega _i\overline{\omega }_j=d\left(^z\omega _i\right)(\overline{\omega }_j\omega _j)$$
(2.7)
which, after using the fundamental polygon with the boundary conditions $`\overline{\omega }_j|_{\alpha _k}=\omega _j|_{\alpha _k}`$, become
$$\omega _i\overline{\omega }_j=\frac{i}{2}\underset{k=1}{\overset{k=h1}{}}_{\alpha _{k+1}}\omega _i_{\gamma _k}(\overline{\omega }_j\omega _j)=\frac{t_{ij}}{2}.$$
(2.8)
In particular, the above result implies $`t_{ii}>0`$ for the diagonal elements of the period matrix.
The coordinates of a string propagating on a circle of radius $`R`$ have the form
$$X(P)=2\pi R\underset{i=1}{\overset{i=g}{}}[L_i_{P_0}^P\omega _i(z)+\overline{L}_i_{P_0}^P\overline{\omega }_i(\overline{z})],$$
(2.9)
where $`P_0`$ is an arbitrary base point. Due to the reality property $`\overline{\omega }_j|_{\alpha _k}=\omega _j|_{\alpha _k}`$, the Neumann (N) and Dirichlet (D) boundary conditions read, respectively,
$$\overline{L}_i=\pm L_i\{\genfrac{}{}{0pt}{}{+\mathrm{for}N}{\mathrm{for}D}$$
(2.10)
A string satisfying Neumann boundary conditions can wind around the boundaries $`\alpha _i`$. Keeping in mind the double cover, it is convenient to parameterize these windings in terms of the winding numbers $`n_i`$ around $`𝐚_i`$:
$$L_i^N=\frac{n_i}{2}.$$
(2.11)
On the other hand, in the Dirichlet directions, the boundary strings cannot wind, but an open string stretching between two different boundaries can wind from one end to the other. If it winds $`m_i`$ times between $`\alpha _i`$ and $`\alpha _{i+1}`$, then
$$L_i^D=i\underset{j}{}(t^1)_{ij}m_j.$$
(2.12)
In terms of the double cover, $`2m_i`$ corresponds to the (even) number of windings around $`𝐛_i`$. The classical action, $`S_{\mathrm{cl}}=\frac{1}{\pi \alpha ^{}}_zX_{\overline{z}}X`$, can be computed by using (2.8):
$$S_{\mathrm{cl}}^N=\frac{\pi R^2}{2}ntn^T,S_{\mathrm{cl}}^D=2\pi R^2mt^1m^T,$$
(2.13)
where $`R`$ is the compactification radius in $`\alpha ^{}`$ units.
The full partition function, in addition to the lattice sum $`Z_{\mathrm{cl}}=e^{S_{\mathrm{cl}}}`$, includes the bosonic determinant factors. These can be obtained by taking the appropriate square root of the corresponding expression for the double cover:
$$(det\mathrm{Im}\tau )^{1/2}|Z_1|^1(R_\mathrm{\Sigma }det\mathrm{Im}\tau )^{1/4}Z_1^{1/2},$$
(2.14)
where $`Z_1`$ is the chiral part of the determinant and $`R_\mathrm{\Sigma }`$ is the “correction factor” depending on the boundary conditions . For a general involution, it has the form
$$I=\left(\begin{array}{cc}\mathrm{\Gamma }& 0\\ \mathrm{\Delta }& \mathrm{\Gamma }\end{array}\right):R_\mathrm{\Sigma }^\mathrm{N}=(R_\mathrm{\Sigma }^\mathrm{D})^1=det(\frac{1\mathrm{\Gamma }}{2}\mathrm{Im}\tau +\frac{1+\mathrm{\Gamma }}{2}(\mathrm{Im}\tau )^1).$$
(2.15)
Thus in our case,
$$R_\mathrm{\Sigma }^\mathrm{N}=(dett)^1\mathrm{and}$$
$$(det\mathrm{Im}\tau )^{1/2}|Z_1|^1\{\genfrac{}{}{0pt}{}{Z_1^{1/2}\mathrm{for}N}{(dett)^{1/2}Z_1^{1/2}\mathrm{for}D}$$
(2.16)
## 3 $`W^{2(h1)}`$ F-terms from Open String Topological Amplitudes
In Type II theory, the F-terms of the form $`𝒲^{2g}`$, where $`𝒲`$ is the chiral $`N=2`$ supergravity multiplet, are determined at genus $`g`$ by the topological partition function $`F_g`$ . In this section, we discuss a similar structure in Type I theory: the open string diagrams with $`h`$ boundaries ($`h1`$ open string loops) generate the effective action terms $`(\mathrm{Tr}ϵ^{\alpha \beta }W_\alpha W_\beta )^{(h1)}`$, where $`W`$ is the familiar chiral (spinorial) gauge field strength superfield and the trace is over gauge indices in the fundamental representation . We will demonstrate this fact by evaluating the amplitudes involving $`2(h2)`$ gauginos and two gauge bosons coupled to $`h`$ boundaries of a genus zero surface, in the Neveu-Schwarz-Ramond formalism, similarly to Ref.. It is very convenient to describe this surface in terms of its $`g=h1`$ double cover introduced in the previous section. Then the computation proceeds by essentially repeating the steps of the original computation in Type II theory.
We are interested in the effective action term $`(\mathrm{Tr}ϵ_{\alpha \beta }W_\alpha W_\beta )^{(h1)}|_\mathrm{F}=(h1)\mathrm{Tr}[(ϵ^{\alpha \beta }\lambda _\alpha \lambda _\beta )^{(h2)}g^{ac}g^{bd}_{ab}_{cd}]+\mathrm{}`$, where $`\lambda `$ are gauginos and $``$ are the self-dual combinations of gauge field strengths. Thus the amplitude under consideration involves $`h2`$ pairs of gauginos at zero momentum, with opposite helicities inside each pair, plus one pair of gauge bosons in the momentum-helicity configuration corresponding to a self-dual gauge field strength. In order to provide the desired gauge charge, helicity and momentum configuration, the gaugino vertex operators are distributed in pairs among $`h2`$ boundaries, the two gauge boson vertices are inserted on a separate boundary while one boundary remains “empty”.
The gaugino vertex operator of definite helicity $`\alpha `$, at zero momentum, in the canonical $`1/2`$ ghost picture, reads:
$$V_\alpha ^{(1/2)}(x)=:e^{\phi /2}S_\alpha S_{int}:,$$
(3.1)
where $`x`$ is a position on the boundary of the world-sheet, $`\phi `$ is the scalar bosonizing the superghost system, and $`S_\alpha `$ ($`S_{int}`$) is the space-time (internal) spin field. Upon complexification of the four space-time fermionic coordinates, $`\psi _I`$ for $`I=1,2`$, and introducing the bosonization scalars $`e^{i\varphi _I}=\psi _I`$, one has
$$S_\alpha =e^{\pm {\displaystyle \frac{i}{2}}\left(\varphi _1+\varphi _2\right)};\alpha =\pm .$$
(3.2)
The gauge boson vertex operator at momentum $`p`$ and polarization $`ϵ`$, in the $`0`$ ghost picture, is:
$$V_{p,ϵ}^{(0)}(x)=:ϵ_\mu (X^\mu +ip\psi \psi ^\mu )e^{ipX}:.$$
(3.3)
The above vertices must be supplemented by the appropriate Chan-Paton factors, which we omit here for simplicity.
In order to balance the ghost charge, we change one half of gaugino vertex operators to $`+1/2`$ ghost picture: this is done by inserting $`h2`$ picture changing operators at the boundaries. Recall that the picture changing operator (PCO) is defined as $`e^\phi T_F`$, where $`T_F`$ is the world-sheet supercurrent. In addition, due to the supermoduli integration, there are the usual $`2g2`$ PCO insertions on the genus $`g`$ double cover that can be realized as $`2(h2)`$ boundary insertions on the open string world-sheet. As a result, the total number of PCO insertions is $`3(h2)`$.
We will see below that the purely bosonic parts of the gauge boson vertex operators do not contribute to the amplitude under consideration. To make the calculation simpler, we choose a kinematical configuration corresponding to $`\psi _1\psi _2`$ from one vertex and to $`\overline{\psi }_1\overline{\psi }_2`$ from the second. The amplitude then becomes
$$A_h=\underset{i=1}{\overset{h2}{}}e^{\phi /2}S_+S_{int}(x_i)\underset{i=1}{\overset{h2}{}}e^{\phi /2}S_{}S_{int}(y_i)\psi _1\psi _2(z)\overline{\psi }_1\overline{\psi }_2(w)\underset{i=1}{\overset{3(h2)}{}}e^\phi T_F(z_i)_h$$
(3.4)
The above expression has exactly the same structure as the left-moving (or right-moving) part of the topological amplitude written in Eq.(3.6) of Ref.. The amplitudes of Ref. describe scattering processes involving graviphotons and gravitons; here, these particles are replaced by gauginos and gauge bosons, respectively. Now the vertex positions $`x_i`$, $`y_i`$, $`z`$ and $`w`$ are integrated over the boundary while the supercurrents are inserted at a priori arbitrary points $`z_i`$ of the boundary. In order to demonstrate a similar, topological nature of $`A_h`$, we can limit our considerations to the case of orbifold compactifications. The twists around b-cycles of the double cover correspond to the brane angles, while the twists around a-cycles belong to the orbifold group of the Type II theory. Thus, the former are fixed by the brane configuration, while the latter are summed over all elements of the orbifold group. Below, we call both types of periodicity conditions as orbifold twists.
In the case of orbifolds, the internal $`N=2`$ SCFT is realized in terms of free bosons and fermions. We consider for simplicity orbifolds realized in terms of $`3`$ complex bosons $`X_I`$ and left- (right-) moving fermions $`\psi _I`$ ($`\stackrel{~}{\psi }_I`$), with $`I=3,4,5`$. Since all vertices involving these fermions are inserted at the boundary, from now on we can identify the left- and right-movers. Let $`h_O`$ be an orbifold twist defined by $`h_O=\{h_I\}`$, and its action on $`X_I`$ is $`X_Ie^{2\pi ih_I}X_I`$ and similarly for $`\psi _I`$. Space-time supersymmetry implies that one can always choose the $`h_I`$’s to satisfy the condition:
$$\underset{I}{}h_I=0.$$
(3.5)
On the genus $`g=h1`$ double cover, we must associate one orbifold twist to each homology cycle $`𝐚_i,𝐛_i`$, for $`i=1,\mathrm{},h1`$. In the following we shall denote by $`\{h_O\}=\{\{h_I\}\}`$ the set of all twists along different cycles. One can bosonize the complex fermions
$$\psi _I=e^{i\varphi _I},\overline{\psi }_I=e^{i\varphi _I}.$$
(3.6)
In terms of these bosons, the internal part of gaugino vertex operators reads
$$S_{int}=e^{{\displaystyle \frac{i}{2}}\left(\varphi _3+\varphi _4+\varphi _5\right)}.$$
(3.7)
Similarly, the internal part of the supercurrent at the boundary becomes
$$T_F^{int}=G^{}+G^+;G^{}=\underset{I=3}{\overset{5}{}}e^{i\varphi _I}X^I.$$
(3.8)
By internal charge conservation, since all $`2(h2)`$ gauginos carry charge $`+1/2`$ for $`\varphi _3`$, $`\varphi _4`$ and $`\varphi _5`$, only the internal parts (3.8) of the supercurrents contribute: $`h2`$ $`T_F`$’s must contribute $`1`$ charge each for $`\varphi _3`$ and similarly $`1`$ for $`\varphi _4`$ and $`\varphi _5`$ each. Thus only $`G^{}`$ contributes in this amplitude.
In order to compute the amplitude 3.4, we repeat the steps leading from Eq.(3.6) to Eq.(3.18) of Ref.. In particular, in order to cast the contribution of world-sheet fermions into a form that makes the summation over their spin structures tractable by using the simplest form of Riemann identity for theta functions, we are led to the following choice of the supercurrent insertion points:
$$\underset{a=1}{\overset{3(h2)}{}}z_a=\underset{i=1}{\overset{h2}{}}y_iz+w+2\mathrm{\Delta },$$
(3.9)
where $`\mathrm{\Delta }`$ is the Riemann $`\theta `$ constant associated to the double cover. Note that this is an allowed gauge choice even if all points are located at the boundary. After summing over all spin structures and using exactly the same bosonization formulae and theta function identities as in section 3 of Ref., we obtain the following expression
$`A_h`$ $`=`$ $`det\omega _i(x_j,z)det\omega _i(y_j,w){\displaystyle \frac{_Idet(X_I\omega _{h_I,i}(u_{j,I}))}{deth_a(z_b)}}`$ (3.10)
$`ϵ_{a_1,\mathrm{},a_{3(h2)}}{\displaystyle \underset{i=1}{\overset{3(h2)}{}}}{\displaystyle (\mu _ih_{a_i})(\text{lattice sum})},`$
where $`X_I`$ are the instanton modes twisted along the homology cycles by a particular set of orbifold twists $`\{h_I\}`$ and $`\omega _{h_I,i}`$, $`i=1,\mathrm{},h2`$, are the differentials associated to the boundary conditions twisted by $`\{e^{2i\pi h_I}\}`$. The above expression is written for one particular partition $`u_{i,I}(i=1,\mathrm{},h2)`$, $`I=3,4,5`$ of the positions of $`T_F`$’s which contribute $`1`$ charge for $`\varphi _3`$, $`\varphi _4`$ and $`\varphi _3`$, respectively. At the end, one must consider all possible partitions $`\{u_{i,I}\}`$ and antisymmetrize; note that as a set $`\{z_a\}=_{i,I}\{u_{i,I}\}`$. Finally, $`h_a`$, $`a=1,\mathrm{},3(h2)`$ are the $`3(h2)`$ quadratic differentials whose determinant appears after a number of technical steps explained in Ref., making use of the bosonization formulae. They also appear as the zero modes of $`b`$-ghosts, contracted with the Beltrami differentials $`\mu _i`$; the corresponding integral gives the measure over moduli space of genus zero surfaces with $`h`$ disconnected boundaries.
There is only one difference between Eq.(3.10) and the analogous Eq. (3.18) of Ref.: the absence of $`(det\mathrm{Im}\tau )^2`$ factor which arises in Type II theory after integrating over the space-time zero modes $`X^\mu `$ (i.e. the four-dimensional momenta). In our case, the correction factor (2.15) for the Neumann boundary condition eliminates this determinant, c.f. Eq.(2.16). Furthermore, for each compact direction, there is a lattice weight $`Z_{cl}=e^{S_{cl}}`$, see Eq.(2.13), which according to Eq.(2.16) should be multiplied by $`(dett)^{1/2}`$ only in the case of Dirichlet boundary conditions.
The result (3.10) is for a fixed partition $`\{u_{i,I}\}`$ of the $`z_a`$’s. As mentioned earlier one must consider all possible partitions and antisymmetrize. Furthermore, $`X_I\omega _{h_I,i}`$ are holomorphic quadratic differentials. Therefore, summing aver all partitions $`\{u_{i,I}\}`$ with the proper antisymmetrization gives:
$$\underset{\{u\}}{}\underset{I}{}det(X_I\omega _{h_I,i}(u_{i,I}))=Bdeth_a(z_b),$$
(3.11)
where $`B`$ is $`z_a`$-independent. In this way, the amplitude (3.10) becomes manifestly independent of the PCO’s insertion points.
Now it remains to integrate the vertex positions $`x_i`$, $`y_i`$, $`z`$ and $`w`$ over all $`h`$ disconnected boundaries $`\alpha `$. For $`z`$ and $`w`$ located at a specific boundary $`\alpha _k`$, each pairing $`(x,y)`$ of opposite helicity gauginos at other boundaries gives rise to a specific Chan-Paton factor and, as a consequence of Eq.(2.6), it picks up only one of the $`(h1)!^2`$ terms from the product $`det\omega _i(x_j,z)det\omega _i(y_j,w)`$. The integral of such a term is simply 1. As a result,
$$A_h=Nh!__hBdet(\mu _ah_b)(\mathrm{lattice}\mathrm{sum}),$$
(3.12)
where $`_h`$ is the moduli space of genus 0 Riemann surfaces with $`h`$ boundaries.<sup>2</sup><sup>2</sup>2Strictly speaking one descends on $`_h`$ after averaging over periodicity conditions of the orbifold group around the a-cycles. The combinatorial factor $`h!`$ arises as follows. First, there are $`h`$ choices of the “empty” boundary, then $`h1`$ choices of the gauge bosons’ boundary and finally, $`(h2)!`$ ways of distributing say positive helicity gauginos – the positions of negative helicity gauginos determines the Chan-Paton factor. The additional factor $`N`$ comes from the “empty” boundary and counts the number of D-branes.
The effective action term that reproduces the amplitude (3.12) is the F-term
$$S_{\mathrm{eff}}=F^{(0,h)}(\mathrm{Tr}ϵ^{\alpha \beta }W_\alpha W_\beta )^{(h1)}|_F,$$
(3.13)
with the coefficient given by the topological partition function
$$F^{(0,h)}=Nh__h\underset{a=1}{\overset{3h6}{}}(\mu _aG^{})(\mathrm{lattice}\mathrm{sum}),$$
(3.14)
where we have used Eqs.(3.11) and (3.12). Here $`G^{}=_IX^I\overline{\psi }^I`$ is the supercurrent in the topological twisted theory. Note that in the latter $`\overline{\psi }^I`$ carry dimension 1 and therefore they have $`h2`$ zero modes $`\omega _{h_I,i}`$ with $`i=1,\mathrm{},h2`$.
## 4 Holomorphic Anomaly and $`\mathrm{\Pi }`$-terms
Type I models considered above are dual to heterotic models. In particular the $`N=1`$ Type I models should be dual to heterotic models based on an $`N=(2,0)`$ world-sheet superconformal field theory. Correspondingly the topological amplitudes giving rise to F-terms $`F^{(0,h)}(\mathrm{Tr}W^2)^{h1}`$ in Type I has a counterpart in heterotic theory. In Ref., such amplitudes were shown to be given by partition functions of topological heterotic theory obtained by twisting the left-moving $`N=2`$ superconformal algebra. Indeed genus $`g`$ partition function of the topological heterotic theory $`F_g^H`$ gives the coupling $`F_g^H(\mathrm{Tr}W^2)^g`$.
The topological partition functions in Type II theories satisfy a holomorphic anomaly equation; their derivatives with respect to anti-holomorphic moduli can be expressed in terms of topological partition functions of lower genera. In the heterotic theory, however, as shown in Ref., anti-holomorphic derivatives of $`F_g^H`$ give rise to a larger class of “topological” quantities $`F_{g,n}^H`$ labeled by genus $`g`$ and insertions of $`n`$ pairs of anti-chiral operators. In the effective field theory, these quantities are associated to F-terms of the form $`(\mathrm{Tr}W^2)^g\mathrm{\Pi }^n`$. Let us briefly recall this difference between the Type II and heterotic topological theories.
The left-moving $`N=2`$ superconformal algebra is generated by the stress tensor $`T`$, a $`U(1)`$ current $`J`$ and the two dimension $`3/2`$ superconformal generators $`G^+`$ and $`G^{}`$ where the superscripts $`\pm `$ refer to their $`U(1)`$ charges. Topological twisting of this algebra amounts to modifying $`TT+\frac{1}{2}J`$. With respect to the new $`T`$ the dimensions of $`G^+`$ and $`G^{}`$ are respectively $`1`$ and $`2`$. One defines the BRST operator for the topological theory $`Q_{BRST}`$ as the contour integral of $`G^+`$ current. $`G^{}`$ having dimension $`2`$ and satisfying the condition $`Q_{BRST}G^{}=T`$ plays the role of the $`b`$ ghost field of the string theory. Thus the measure on the moduli space of a Riemann surface of genus $`g`$ is defined by computing the correlation function of $`3g3`$ $`G^{}`$’s which are each folded with a Beltrami differential. The physical states of the theory are given by the $`Q_{BRST}`$-cohomology; the chiral primaries $`\varphi _i^+`$ having charge $`+1`$ have now zero dimension, they are in $`Q_{BRST}`$-cohomology and can be inserted at punctures on the Riemann surfaces, while the anti-chiral vertex operators (in the 0 picture in string theory) are BRST exact operators $`Q_{BRST}\varphi _{\overline{i}}^{}=G^+\varphi _{\overline{i}}^{}`$. Inserting such an operator in the Riemann surface gives a total holomorphic derivative in the world-sheet moduli $`_t`$, since by deforming the contour $`Q_{BRST}`$ will act on one of the $`G^{}`$ folded with the Beltrami differentials and convert it to $`T`$. Thus the only possible contributions can come from the boundaries of the moduli space of the world-sheet, namely the degenerations of the original Riemann surface along some non-trivial cycle which could be either homologically trivial or non-trivial.
In the homologically trivial degeneration limit, the surface splits into two Riemann surfaces $`\mathrm{\Sigma }_1`$ and $`\mathrm{\Sigma }_2`$ of lower genera, say genus $`g_1`$ and $`g_2`$, with one puncture each at $`P_1`$ and $`P_2`$. The operator $`\varphi _{\overline{i}}^{}`$ will be on one of the surfaces (say $`\mathrm{\Sigma }_1`$). Since in the twisted theory the total charge on the sphere must be +3, it follows that the operators that appear at $`P_1`$ and $`P_2`$ are dual of each other and carry charges +2 and +1, respectively. The Beltrami differentials together with $`G^{}`$ split on the two surfaces according to the $`U(1)`$ charge conditions in the twisted theory. The resulting term is of the form
$$F_{\overline{i},\overline{j}}^{g_1}G^{\overline{j}j}D_jF^{g_2},$$
(4.1)
where $`D`$ is the holomorphic covariant derivative, and $`j`$ denotes the chiral operator carrying charge +1 at $`P_2`$, while the operator at $`P_1`$ carrying charge +2 is related to the anti-chiral operator (with charge $`1`$) labeled by $`\overline{j}`$ by the action of a holomorphic 3-form operator $`\rho (z)`$ carrying charge +3:
$$\varphi _{\overline{j}}^{}\varphi _{\overline{j}}^{++}𝑑z\rho (z)\varphi _{\overline{j}}^{}.$$
(4.2)
Such an operator $`\rho `$ of dimension 0 in the twisted theory exists for all $`N=2`$ superconformal field theories leading to space-time supersymmetry: $`\rho =e^{i\sqrt{3}H}`$ with $`\sqrt{3}H`$ being the $`U(1)`$ current. The metric $`G_{j\overline{j}}`$ appearing in Eq.(4.1) is defined by the inner product $`\varphi _j^+\varphi _{\overline{j}}^{++}`$. In Ref., these new topological objects $`F_{\overline{i},\overline{j}}^g`$ were related to the couplings of the effective action terms $`(\mathrm{Tr}W^2)^g\mathrm{\Pi }`$, denoted symbolically $`F_1^g`$. By further taking anti-holomorphic derivatives of these terms one arrives at generalized topological quantities denoted by $`F_{\overline{i}_1\mathrm{}\overline{i}_n;\overline{j}_1\mathrm{}\overline{j}_n}^g`$ where $`\overline{i}_k`$ denote charge $`1`$ insertions while $`\overline{j}_k`$ denote charge +2 insertions. These quantities were identified with the effective action terms $`(\mathrm{Tr}W^2)^g(\mathrm{\Pi })^n`$.
Similar reasoning for the case of handle degeneration results in a term of the form
$$C^{j\overline{j}}D_jF_{\overline{i},\overline{j}}^{g1}$$
(4.3)
with
$$C^{j\overline{j}}=G^{j\overline{k}}Q_{\overline{k}\mathrm{}}G^\mathrm{}\overline{\mathrm{}}Q_{\overline{\mathrm{}}k}G^{k\overline{j}},$$
(4.4)
where $`Q`$ is the charge operator associated with the gauge field $`W`$ (recall that in the heterotic theory the non-topological right moving part of the gauge vertex contains the charge operator). This is because only $`2g2`$ of the original $`2g`$ gauge superfields sit on the genus $`g1`$ surface obtained from the handle degeneration and the remaining two $`W`$’s sit at the node. As a result, the propagator connecting the two punctures $`P_1`$ and $`P_2`$ comes with the charges of the two $`W`$’s.<sup>3</sup><sup>3</sup>3More generally there can be several different gauge fields and the corresponding topological partition function will carry the labels of the gauge fields. In this case the $`Q`$’s in equation (4.4) will carry the labels of the two gauge fields that sit at the node.
In Type II, we have also the right-moving $`N=2`$ superconformal algebra, which upon twisting gives rise to another BRST operator $`\overline{Q}_{BRST}`$. The anti-chiral vertex operators now (in the (0,0) ghost picture) are of the form $`Q_{BRST}\overline{Q}_{BRST}\varphi ^{}`$ and inserting it gives rise to a double derivative $`_t_{\overline{t}}`$ in the moduli space integrals, where $`t`$ is standard plumbing fixture coordinate describing the degeneration of the surface. To get a non- vanishing contribution at $`t0`$ the integrand must behave as $`\mathrm{ln}(t\overline{t})`$. The latter behavior is obtained only if the operator $`\varphi ^{}`$ approaches the node (the puncture $`P_1`$) and has a non-vanishing structure constant $`C_{\overline{i}\overline{j}\overline{k}}`$ for some $`\overline{k}`$, the integral of its position giving rise to $`\mathrm{ln}(t\overline{t})`$. As a result
$$F_{\overline{i},\overline{j}}^g=e^{2K}C_{\overline{i}\overline{j}\overline{k}}G^{\overline{k}k}D_kF^g,$$
(4.5)
where $`K`$ is the Kähler potential (the appearance of $`e^{2K}`$ can be deduced by matching the Kähler weights). Moreover for the handle degeneration case $`C^{j\overline{j}}`$ in Eq.(4.4) reduces to the inverse metric $`G^{j\overline{j}}`$. This is because the charge operator $`Q`$ is essentially replaced by space-time momenta which are part of the graviphoton field strength in the $`N=2`$ $`(𝒲^2)^g`$ term.
Returning to the case of open strings (Type I) we expect the situation to be similar to the heterotic string since firstly we are dealing with $`N=1`$ theories and secondly Type I and heterotic theories are dual to each other. This can be seen more directly by examining the corresponding topological theories. Indeed on surfaces with boundaries, only the sum of the left and right moving BRST operators is compatible with the boundary conditions. As a result, upon inserting an anti-chiral operator in the topological partition function one can only deform the contour associated with this combination of left and right BRST operators resulting in an integrand which is a total derivative in the world-sheet moduli space. There is no second BRST operator which could give an additional total derivative. To study the possible boundaries of the moduli space (degenerations) let us, for simplicity, restrict ourselves to surfaces of genus 0 with $`h`$ boundaries $`\mathrm{\Sigma }_{(0,h)}`$. In the following we assume $`h3`$. The moduli space of such surface is $`3(h2)`$ real dimensional implying that there are $`3(h2)`$ insertions of the sum of left and right moving $`G^{}`$ that are folded with the Beltrami differentials. This surface therefore carries a (left plus right) $`U(1)`$ charge equal to $`3(h2)`$ which is exactly the number dictated by the $`U(1)`$ anomaly in the twisted theory. Inserting an anti-chiral operator and deforming the contour associated with the corresponding BRST operator converts one of the $`G^{}`$ into the stress-tensor, giving rise to a total derivative in the moduli space. The boundary terms come from the degeneration limits of the surface. There are two degenerations with open string intermediate states analogous to the dividing and handle degeneration cases of heterotic string, and one with intermediate closed string (see Fig. 4).
1. The first degeneration arises when the surface splits into two Riemann surfaces $`\mathrm{\Sigma }_{(0,h_1)}`$ and $`\mathrm{\Sigma }_{(0,h_2)}`$ (where $`h=h_1+h_21`$ and $`h_1,h_22`$) with one puncture each, say $`P_1`$ and $`P_2`$, at the boundaries of the two surfaces. The original anti-chiral insertion $`\mathrm{\Phi }_{\overline{i}}^{}`$ is on one of the components, say $`\mathrm{\Sigma }_{(0,h_1)}`$. Since the total charge (i.e. left plus right $`U(1)`$ charge) on a disk in the twisted theory, as dictated by the $`U(1)`$ anomaly, should be +3, it follows that the open string insertion at $`P_1`$ carries charge +2, while the open string insertion at $`P_2`$ is its metric-dual operator $`\varphi _j^+`$ carrying charge +1. The charge +2 state is obtained now by the action of $`\rho _L+\rho _R`$ on a charge $`1`$ operator. The leftover $`3(h2)1`$ (left plus right) $`G^{}`$’s folded with the Beltrami differentials distribute themselves on $`\mathrm{\Sigma }_{(0,h_1)}`$ and $`\mathrm{\Sigma }_{(0,h_2)}`$. Their numbers on the two surfaces are respectively $`3(h_12)+1`$ and $`3(h_22)+1`$ which is dictated by the $`U(1)`$ charge anomalies on the two surfaces, as well as by the dimension of the moduli space of the corresponding one punctured Riemann surfaces. The resulting term is
$$F_{\overline{i},\overline{j}}^{(0,h_1)}G^{\overline{j}j}D_jF^{(0,h_2)};h_1+h_2=h+1.$$
(4.6)
2. The second open string degeneration is analogous to the handle degeneration of the closed string and results in twice-punctured surface $`\mathrm{\Sigma }_{(0,h1)}`$ with punctures $`P_1`$ and $`P_2`$ on a boundary with the intermediate states at the two punctures carrying charge +2 and +1 respectively. The dimension of this twice-punctured moduli space is $`3(h3)+2`$ which is exactly the number of the leftover $`G^{}`$’s. This is also in agreement with the total $`U(1)`$ charge anomaly. The resulting term is
$$C^{j\overline{j}}D_jF_{\overline{i},\overline{j}}^{(0,h1)}.$$
(4.7)
3. Finally there is also a degeneration with closed string intermediate state when the surface splits into $`\mathrm{\Sigma }_{(0,h_1)}`$ and $`\mathrm{\Sigma }_{(0,h_2)}`$ with one puncture each $`P_1`$ and $`P_2`$ in the interior of the two surfaces. It is easy to see that $`h_1+h_2=h`$ (where now $`h_1,h_21`$). The plumbing fixture coordinate $`\tau `$ is now complex and the boundary corresponds to $`t=|\tau |\mathrm{}`$. Thus the angular part of $`\tau `$ is still to be integrated. The moduli space of the two surfaces with one puncture each in the interior of the surfaces is $`3(h_12)+2`$ and $`3(h_22)+2`$ real dimensional, respectively, implying that as many $`G^{}`$’s are distributed on the two surfaces, respectively. The remaining one $`G^{}`$ is sitting at the node which is folded with the Beltrami differential corresponding to the angular part of $`\tau `$. Now the question is what are the $`U(1)`$ charges carried by the closed string intermediate state at $`P_1`$ and $`P_2`$. Since the $`U(1)`$ anomaly for a sphere (which is the relevant surface for the intermediate closed string propagation) is +6, we conclude that the sum of the $`U(1)`$ charges at $`P_1`$ and $`P_2`$ is +6. If $`\varphi _{\overline{i}}^{}`$ is on the first surface, then there are two possibilities: the charges at $`(P_1,P_2)`$ are $`(+4,+2)`$ or $`(+3,+3)`$. The charges here are the sum of the left and right moving parts. Charge +4 therefore means (left, right) (anti-chiral, anti-chiral) closed string state where the left and right charge +3 operators $`\rho _L`$ and $`\rho _R`$ have been applied on the anti-chiral operators to convert them into charge +2 operators. Charge +2 operator is the (chiral, chiral) state and hence represents the holomorphic derivative with respect to the corresponding target space modulus. Charge +4 and +2 therefore correspond to the antiholomorphic and holomorphic complex structure moduli, respectively, of the target space. Recall that here we work in the Type I description, where the relevant topological twist is the one of B-model. We thus obtain a term similar to (4.6):
$$F_{\overline{i},\overline{j}}^{(0,h_1)}G^{\overline{j}j}D_jF^{(0,h_2)};h_1+h_2=h,$$
(4.8)
where $`j`$ and $`\overline{j}`$ are closed string states.
Finally, charge +3 at a puncture corresponds to the diagonal combination of (chiral, anti-chiral) plus (anti-chiral, chiral) states which are in fact the Kähler moduli of the target space. These moduli are expected to decouple from the topological B-model. In the Type I context, they are complexified with Ramond-Ramond (RR) fields which are associated with continuous shift symmetries in perturbation theory that make the holomorphic dependence on the Kähler moduli trivial. Charge +3 can also include the identity operator corresponding to exchange of dilaton in the parent string theory. Indeed, in the example discussed in the section 8, there will be such a contribution to the holomorphic anomaly from the identity channel in the closed string degeneration.
Note that for $`h=2`$ the first type of degeneration is absent, while the second gives the well known holomorphic anomaly equation for the gauge couplings. The terms arising from the degeneration with closed string intermediate states need to be understood further, however in the following we will focus on the terms coming from the first two types of degenerations that involve open string intermediate states. The new terms that appear in the holomorphic anomaly equations are $`F_n^{(0,h)}`$’s, that is $`F_{\overline{i}_1\mathrm{}\overline{i}_n;\overline{j}_1\mathrm{}\overline{j}_n}^{(0,h)}`$ where $`\overline{i}`$ indices refer to open string insertions of anti-chiral operators carrying charge -1 while $`\overline{j}`$ indices refer to anti-chiral operators carrying charge +2 (by the action of $`\rho `$). In heterotic theory, it was shown that these quantities originate from F-terms of the form $`(\mathrm{Tr}W^2)^h\mathrm{\Pi }^n`$ . In components they include terms like $`^2(\lambda ^2)^{h1}_{k=1}^n(\psi _{\overline{i}_k}.\psi _{\overline{j}_k})`$ among others, where $`\lambda `$ are the gauginos and $`\psi `$ are non-chiral matter fermions. The proof of this statement in the case of open string is identical to the one given for the heterotic case, since in the open string, by going to the double cover, the spin structure sum boils down to just one sector, as in the heterotic string. When supersymmetry is broken by a D-term expectation value, the one-loop term $`\mathrm{\Pi }\mathrm{Tr}W^2`$ gives rise to fermion masses similar to a “Higgs” $`\mu `$-term in the effective field theory. These masses will be studied in section 7.
## 5 $`F^{(0,3)}`$ in Type I with Magnetized D9 Branes
In this section, we consider an explicit example of Type I string compactified on a factorized six-torus $`T^6=(T^2)^3`$, with magnetized D9 branes. We will then compute the topological partition function $`F^{(0,3)}`$ on a world-sheet with three boundaries, in a $`N=1`$ supersymmetric configuration corresponding to an appropriate choice of the magnetic fields. For this purpose, we first recall the main properties of the relevant Riemann surface $`(g=0,h=3)`$.
### 5.1 Properties of the $`(g=0,h=3)`$ Riemann surface
The $`(g=0,h=3)`$ surface $`\mathrm{\Sigma }_{(0,3)}`$ can be obtained from the double torus of genus 2, by applying the world-sheet involution (2.1) exchanging left and right movers, as described in section 2 and shown in Fig. 5.
The period matrix $`\tau `$, invariant under the involution, is purely imaginary:
$$\tau =i\left(\begin{array}{cc}t_{11}& t_{12}\\ t_{12}& t_{22}\end{array}\right),$$
(5.1)
where $`t_{ij}`$ are three real parameters, dual to the sizes of the three holes. They correspond to the proper time variables of the closed string propagation channels.
The “relative modular group”, defined by the modular transformations $`Sp(4,)`$ that preserve the involution (2.1), is the group $`GL(2,)`$ . It transforms the period matrix as:
$$\tau D\tau D^T,$$
(5.2)
where $`D`$ is an arbitrary $`2\times 2`$ matrix with an inverse of integer entries. This group, however, in general transforms the boundaries (fixed under the involution) to linear combinations of the boundaries. The relevant modular transformations are the ones that act at most as permutations of the boundaries. This group is generated by
$$D_1=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)D_2=\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right)$$
(5.3)
corresponding to
$$D_1:t_{11}t_{22}D_2:\left(\begin{array}{c}t_{11}t_{11}+t_{22}+2t_{12}\\ t_{12}t_{12}t_{22}\end{array}\right).$$
(5.4)
It is easy to see that $`D_1^2=D_2^2=(D_1D_2)^3=1`$.
Using the above transformations and the positivity of the period matrix, one can choose as fundamental domain of integration the ordering:
$$\sqrt{t_{11}t_{22}}t_{12}0t_{11}t_{22}<\mathrm{}.$$
(5.5)
The three degeneration limits, described in section 4, correspond to the following boundaries of the fundamental domain:
1. $`t_{12}0`$, corresponding to shrinking the dividing geodesics of genus 2. Then $`\mathrm{\Sigma }_{(0,3)}`$ degenerates to a product of two annuli with one common boundary stretched in the two annuli through a massless open string. Moreover, the period matrix becomes diagonal with $`t_{11}`$ and $`t_{22}`$ the closed string proper times of the two annuli.
2. $`t_{12}\sqrt{t_{11}t_{22}}`$, implying the vanishing of the period matrix determinant. In this limit, $`\mathrm{\Sigma }_{(0,3)}`$ degenerates into an annulus with a massless open string attached at one of its boundaries. It amounts to taking the infrared limit of one of the two gauge loops in the effective field theory. Then $`t_{22}`$ is the proper time of the annulus in the closed string channel while $`t_{11}`$ parameterizes the length of the open string.
3. $`t_{22}\mathrm{}`$, corresponding to pinching an intermediate closed string connecting an annulus with a disk. Then $`t_{11}`$ becomes the proper time of the annulus in the closed string channel and $`t_{12}`$ parameterizes the position of the shrinking hole.
### 5.2 Partition functions
We now discuss the bosonic and fermionic partition functions. Quantum determinants can be obtained by taking the appropriate square root of the corresponding expression on the genus 2 double cover. In the bosonic case, there is a correction factor $`R_\mathrm{\Sigma }`$ that depends on the involution and on the boundary conditions . It is given by Eq.(2.16), as described in section 2. In the case of a compact boson, one has to multiply the quantum determinant with the lattice sum of momenta or winding modes (2.13). The resulting partition function for N boundary conditions reads:
$$Z_B=Z_1^{1/2}Z_{\mathrm{cl}}^{N,D}Z_{\mathrm{cl}}^N=R\underset{\stackrel{}{n}}{}e^{{\displaystyle \frac{\pi R^2}{2}}\stackrel{}{n}t\stackrel{}{n}^T},$$
(5.6)
where $`\stackrel{}{n}=(n_1,n_2)`$ are the winding numbers around the two $`𝐚`$ cycles. In the degeneration limit $`t_{22}\mathrm{}`$, non-vanishing $`n_2`$ windings are exponentially suppressed, and for $`n_2=0`$ one recovers the annulus partition function depending on the closed string proper time $`t_{11}`$. The $`t_{12}`$ dependence, corresponding to the position of the shrinking boundary, drops. Similarly, in the case of D boundary conditions, one has:
$$Z_{\mathrm{cl}}^D=R(dett)^{1/2}\underset{\stackrel{}{m}}{}e^{2\pi R^2\stackrel{}{m}t^1\stackrel{}{m}^T}.$$
(5.7)
The method of section 2 can also be applied to the fermionic determinants giving rise to theta functions. The partition function of a complex fermion depends on 16 spin structures corresponding to the four boundary conditions $`\stackrel{}{a}=(a_1,a_2)`$ and $`\stackrel{}{b}=(b_1,b_2)`$ around the non-trivial cycles $`𝐚_i`$ and $`𝐛_i`$:
$`Z_f\left[{\displaystyle \genfrac{}{}{0pt}{}{\stackrel{}{a}}{\stackrel{}{b}}}\right]`$ $`=`$ $`Z_1^{1/2}\mathrm{\Theta }\left[{\displaystyle \genfrac{}{}{0pt}{}{\stackrel{}{a}}{\stackrel{}{b}}}\right](t)`$
$`=`$ $`Z_1^{1/2}{\displaystyle \underset{\stackrel{}{n}=(n_1,n_2)}{}}e^{\pi \left(\stackrel{}{n}+\stackrel{}{a}\right)t\left(\stackrel{}{n}+\stackrel{}{a}\right)^T+2i\pi \left(\stackrel{}{n}+\stackrel{}{a}\right)\stackrel{}{b}^T}.`$
The partition function involves a sum over spin structures with appropriate coefficients $`c\left[\genfrac{}{}{0pt}{}{\stackrel{}{a}}{\stackrel{}{b}}\right]`$, determined at one loop level. As usually, at higher loops the corresponding coefficients are determined by the factorization properties of the vacuum amplitude. Indeed, by considering the factorization limit $`t_{12}0`$, we find:
$$c\left[\genfrac{}{}{0pt}{}{\stackrel{}{a}}{\stackrel{}{b}}\right]=c_A\left(\genfrac{}{}{0pt}{}{a_1}{b_1}\right)c_A\left(\genfrac{}{}{0pt}{}{a_2}{b_2}\right),$$
(5.9)
where $`c_A`$ denote the corresponding coefficients in the annulus amplitude.
### 5.3 The topological amplitude
We consider now a toroidal compactification of Type I string theory on three factorized tori, $`T^6=_{I=3,4,5}T_I^2`$, and a $`N=1`$ supersymmetric configuration of magnetized D9 branes. Since the string diagram has three boundaries, we consider in general three brane stacks associated to the gauge group $`U(N_a)`$, $`a=1,2,3`$. In each of the three abelian factors, there is an internal magnetic field $`H_a^I`$ with components along the three factorized tori $`T_I^2`$, $`I=3,4,5`$. They are quantized in units of the corresponding areas $`\sqrt{G}_I`$, with $`G_I`$ the determinant of the $`T_I^2`$ metric, according to the Dirac quantization condition $`H_a^I=q_a^I/p_a^I\sqrt{G}_I`$, where $`q_a^I`$ is the respective magnetic flux and $`p_a^I`$ the wrapping number of the $`a`$-th brane around the $`I`$-th 2-torus. Note that for each $`I`$ and $`a`$, the two integers $`p_a^I`$ and $`q_a^I`$ are relatively coprime. By T-dualizing three directions, one from each $`T^2`$, one obtains an equivalent description as a Type IIA orientifold with D6 branes at angles related to the magnetic fields . More precisely the angle of each brane relative, for instance, to the horizontal axis of $`T_I^2`$ is $`\theta _I^a=\mathrm{arctan}H_a^I\alpha ^{}`$. In this representation, the condition for having $`N=1`$ unbroken supersymmetry on the $`a`$-th stack takes the simple form that the sum of the corresponding angles in the three tori should vanish: $`\theta _3^a+\theta _4^a+\theta _5^a=0`$.
Let us compute now the topological partition function $`F^{(0,3)}`$ given by the physical amplitude involving two gauge fields on one of the 3 boundaries and two gauginos on another, associated to the effective F-term interaction $`(\mathrm{Tr}W^2)^2`$. From the above discussion, it is clear that the computation is just a particular case of the general orbifold compactification described in section 3. Indeed, all twists around a cycles are trivial, while the angles $`\theta _I^a`$ related to the magnetic fields play exactly the role of the orbifold twists around the $`𝐛`$ cycles, as mentioned already in section 3. By identifying the first two boundaries $`a=1,2`$ with the two $`𝐚_i`$ cycles, $`i=1,2`$, of the homology basis of the genus 2 Riemann surface, and the third boundary $`a=3`$ with the middle “horizontal” cycle, see Fig. 5, one has the relations:
$$2\pi h_I^1=2(\theta _I^1\theta _I^3);2\pi h_I^2=2(\theta _I^3\theta _I^2),$$
(5.10)
where $`h_I^i`$ are the orbifold twists around the two $`𝐛_i`$ cycles.<sup>4</sup><sup>4</sup>4 Note that in the presence of orientifold planes, there are additional diagrams where one or two boundaries are replaced by crosscaps. Their inclusion is straightforward and do not change the physical implications of our results.
Using the identification (5.10), the calculation goes along the lines of section 3 and the result is given by (3.14). In section 3, we had not specified the lattice sum involved explicitly. In the Appendix, we give the detailed derivation of the lattice sums appearing in $`F^{(0,h)}`$ in the example of magnetized D9 branes. Here we just summarize the result for $`h=3`$. In the T-dual version, the three stacks of D6 branes on the three boundaries, will be parallel to some primitive lattice vectors in each plane, $`\stackrel{}{v}_a^I`$, where $`a`$ labels the three boundaries and $`\stackrel{}{v}^I`$ are the two dimensional lattice vectors in the $`I`$-th plane. Since we are considering factorized torii, it is sufficient to focus on one plane, therefore in the following we will drop the index $`I`$ labeling the different planes. In the final formulae we will reinstate the index $`I`$. In terms of the magnetic fluxes in the original D9 branes, these vectors are $`\stackrel{}{v}_a=(p_aR_1,q_aR_2)`$ where $`R_1`$ and $`R_2`$ are the two radii in that plane in the T-dual theory with intersecting branes. The world-sheet can therefore carry windings $`n_a\stackrel{}{v}_a`$ on the $`a`$-th boundary with integer $`n_a`$. However they satisfy the constraint that the sum of the windings over the three boundaries must vanish since it gives the winding on a homologically trivial cycle:
$$\underset{a=1}{\overset{3}{}}n_a\stackrel{}{v}_a=0.$$
(5.11)
Since we are excluding the case when the three stacks of branes are parallel to each other in one or more planes (otherwise there would be enhanced supersymmetry), only one of the $`n_a`$’s is independent and it spans a sublattice of integers that depends on the magnetic flux data $`p_a`$ and $`q_a`$. Thus we have a one-dimensional sublattice sum labeled by, say, $`n_1`$.
On the world-sheet there is another set of integers which appears because the position of different stacks of branes is defined only modulo transverse lattice vectors. Since not all the branes are parallel, we can choose the intersection of two stacks of branes as the origin of the plane. Then the freedom is only in choosing the transverse position of the third stack of branes. Thus we have again a one dimensional lattice sum. Upon Poisson resummation over this lattice we will get a lattice of momenta along the Dirichlet directions of the 3 stacks of branes. The details are given in the Appendix, but here we just give a simple argument to determine what this lattice sum would be. Let $`\stackrel{}{v}_a`$ be the two dimensional dual of $`\stackrel{}{v}_a`$ and let $`\sqrt{G}`$ be the area of the torus (i.e. $`R_1R_2`$). Then $`\stackrel{}{v}_a/\sqrt{G}`$ is the primitive vector in the intersection of the dual momentum lattice and the Dirichlet direction to the $`a`$-th stack of branes. Thus the boundary state of the $`a`$-th branes would carry momentum vectors $`k_a\stackrel{}{v}_a/\sqrt{G}`$ with $`k_a`$ being integers. Conservation of the total momentum gives the same constraint as (5.11) with $`n_a`$ replaced by $`k_a`$.
$$\underset{a=1}{\overset{3}{}}k_a\stackrel{}{v}_a=0.$$
(5.12)
This results in again a one dimensional sublattice sum labeled by say $`k_1`$.
The modular group $`SL(2,)`$, associated to the Kähler modulus of $`T^2`$, acts on the pair of integers $`(n_1,k_1)`$ in the usual way. Since $`(n_1,k_1)`$ span only a sublattice of integers subject to the constraint (5.11) and (5.12), one might wonder if only a subgroup of the full $`SL(2,)`$ survives. However these two constraints are invariant under the full $`SL(2,)`$ symmetry which implies that the symmetry group is indeed the full $`SL(2,)`$.
To proceed further we need to write a classical solution (before the Poisson resummation) carrying the above winding numbers and the transverse positions. In terms of the complex coordinate of the plane $`Z=X_1+iX_2`$, the boundary conditions imply that $`Z`$ is untwisted along the $`𝐚`$ cycles and twisted by say $`g_1=e^{2i\pi h^1}`$ and $`g_2=e^{2i\pi h^2}`$ along the $`𝐛_1`$ and $`𝐛_2`$ cycles of the genus 2 double cover of the surface $`\mathrm{\Sigma }_{(0,3)}`$. Denoting by $`\{g\}`$ the collection of the twists $`g_1`$ and $`g_2`$, we note that by Riemann-Roch theorem there is only one linearly independent holomorphic twisted differential (Prym Differential) $`\omega _{\{g\}}`$. To write the classical solution we need the twisted holomorphic differentials $`\omega _{\{g\}}`$, $`\omega _{\{g^1\}}`$ and their complex conjugates $`\overline{\omega }_{\{g\}}`$, $`\overline{\omega }_{\{g^1\}}`$. The solution is of the form
$$Z(P)=L_{P_0}^P\omega _{\{g\}}+\stackrel{~}{L}_{P_0}^P\overline{\omega }_{\{g^1\}},$$
(5.13)
where $`L`$ and $`\stackrel{~}{L}`$ are determined in terms of winding number and transverse position data and the normalization condition for the twisted differential. Since $`𝐚`$ cycles are untwisted, we can choose the normalization condition (for convenience)
$$_{𝐚_1}\omega _{\{g\}}=g_1^{1/2}.$$
(5.14)
Note that the integral around $`a_2`$ cycle is not independent since integrating over the trivial cycle $`_{i=1}^2(𝐚_i𝐛_i𝐚_i^1𝐛_i^1)`$ gives the constraint
$$[(1g_1)_{𝐚_1}+(1g_2)_{𝐚_2}]\omega _{\{g\}}=0.$$
(5.15)
By using Fig. 6 (for $`h=3`$), one can evaluate
$$_{\mathrm{\Sigma }_{(0,3)}}\omega _{\{g\}}\overline{\omega }_{\{g\}}=\tau _{\{g\}}\overline{\tau }_{\{g\}};\tau _{\{g\}}=D_{\{g\}}+\frac{1}{2}\frac{(1g_1)(1g_2g_1^1)}{1g_2},$$
(5.16)
where
$$D_{\{g\}}=2i\frac{(g_1g_2)^{\frac{1}{2}}}{1g_2^1}_{𝐛_1𝐛_2𝐛_1^1𝐛_2^1}\omega _{\{g\}}.$$
(5.17)
Note that the individual cycles $`𝐛_i`$ are not closed due to the twists, but the cycle $`𝐛_1𝐛_2𝐛_1^1𝐛_2^1`$ is closed. Due to the symmetry of the double cover under the involution (2.1), $`D`$ turns out to be purely imaginary. By using the identity $`_{\mathrm{\Sigma }_{(0,3)}}\omega _{\{g\}}\omega _{\{g^1\}}=0`$ we have the relation $`\tau _{\{g\}}=\tau _{\{g^1\}}`$.
The final result of the lattice contribution as shown in the Appendix (including all the three planes) is<sup>5</sup><sup>5</sup>5Here and in the Appendix \[from Eq.(A.25) onwards\] we do not keep track of overall factors that are completely moduli- and flux data-independent.
$$Z_{\mathrm{lattice}}=\underset{I}{}p_I\left(\frac{\mathrm{Im}\stackrel{~}{\tau }_{\{g_I\}}}{\mathrm{Im}T}\right)^{1/2}\underset{(n_1^I,k_1^I)}{}e^{2i\pi (n_1^I\alpha _1^I+k_1^I\alpha _2^I)}e^{\frac{i\pi }{\mathrm{Im}T_I}|n_1^IT_I+k_1^I|^2\stackrel{~}{\tau }_{\{g_I\}}},$$
(5.18)
where $`T_I=B_I+i\sqrt{G}_I`$ is the usual Kähler modulus of the torus $`T_I^2`$, with $`B_I`$ the two-index antisymmetric tensor. The parameter $`\stackrel{~}{\tau }_{\{g_I\}}=\frac{|v_1^I|^2}{\mathrm{Im}T_I}\tau _{\{g_I\}}`$ is independent of the modulus $`T_I`$, although it depends on the world-sheet moduli and the flux data. Furthermore, $`\alpha _1`$ is the effective Wilson line along the world-volume of the branes and $`\alpha _2`$ is the effective transverse position of the branes (which is T-dual to the second component of the Wilson line in the D9 brane theory). The integer $`p_I`$ is the smallest positive integer such that
$$p_I\frac{v_3^I.v_1^I}{v_3^I.v_2^I}.$$
(5.19)
Note that the above equation implies
$$p_I\frac{v_2^I.v_1^I}{v_3^I.v_2^I}.$$
(5.20)
From the constraints (5.11) and (5.12) it follows that $`n_1^I`$ and $`k_1^I`$ are arbitrary integer multiples of $`p_I`$. As observed in subsection 5.1 \[preceding Eq.(5.4)\], the relevant modular group is the permutation group $`S_3`$ that permutes the three boundaries. Since we have treated the three boundaries asymmetrically in arriving to Eq.(5.18), it is not manifest that the result is modular invariant. For example, we have normalized the Prym differential along the first boundary and the lattice momenta appearing in the expression refer directly to the boundary state of the brane attached to this boundary. Indeed, under the permutation of the boundaries, $`\tau _{\{g_I\}}`$, $`p_I`$ and $`|v_1^I|`$ transform nontrivially. However, the combination $`p_I|v_1^I|\tau _{\{g_I\}}^{\frac{1}{2}}`$ is invariant. For instance, under the exchange of first and second boundaries,
$$v_1^Iv_2^I,p_Ip_I\left|\frac{v_3^I.v_2^I}{v_3^I.v_1^I}\right|\text{and}\tau _{\{g_I\}}\left|\frac{1g_1^I}{1g_2^I}\right|^2\tau _{\{g_I\}}.$$
(5.21)
The latter can be seen by the fact that under the exchange of $`𝐚_1`$ and $`𝐚_2`$ cycles, $`\omega _{\{g_I\}}`$ \[normalized by Eq.(5.14) around $`𝐚_1`$ cycle\], is transformed to $`\frac{1g_1^I}{1g_2^I}\omega _{\{g_I\}}`$, where we used the constraint (5.15). The statement that $`p_I|v_1^I|\tau _{\{g_I\}}^{\frac{1}{2}}`$ is invariant now follows from the relation
$$\left|\frac{1g_1^I}{1g_2^I}\right|=\left|\frac{v_3^I.v_1^I}{v_3^I.v_2^I}\right|\frac{|v_2^I|}{|v_1^I|}.$$
(5.22)
As a result, $`Z_{\mathrm{lattice}}`$ as given in (5.18) is invariant under the permutation group.
In the topological partition function, we have also the insertion of $`G^{}`$’s which are folded with the Beltrami differentials. As mentioned earlier, each $`\overline{\psi }^I`$ field, being of dimension 1 in the topological theory, has now only one zero mode (for each plane $`I`$). It will therefore be replaced by $`\omega _{\{g_I^1\}}`$. Similarly $`Z_I`$ will be replaced by the zero mode
$$Z_I=\frac{v_I}{\mathrm{Im}T_I}(n_1^IT_I+k_1^I)\omega _{\{g_I\}}.$$
(5.23)
Integrating the fermion zero modes we get
$$\underset{a=1}{\overset{3}{}}\mu _aG^{}=\underset{I=1}{\overset{3}{}}\frac{v_I}{\mathrm{Im}T_I}(n_1^IT_I+k_1^I)det\mu _ah_J;h_J=\omega _{\{g_J\}}\omega _{\{g_J^1\}}.$$
(5.24)
We can further evaluate the determinant above by noting that under the deformation of the complex structure represented by the Beltrami differential $`\mu _a`$, the twisted $`(1,0)`$ form $`\omega _{\{g_I\}}`$ picks up a $`(0,1)`$ form given by
$$\omega _{\{g_I\}}\omega _{\{g_I\}}+C\overline{\omega }_{\{g_I^1\}};C\overline{\omega }_{\{g_I^1\}}=\mu _a\omega _{\{g_I\}}\mathrm{modulo}\mathrm{an}\mathrm{exact}\mathrm{form}.$$
(5.25)
As a result
$`{\displaystyle \mu _a\omega _{\{g_I\}}\omega _{\{g_I^1\}}}`$ $`=`$ $`C{\displaystyle \overline{\omega }_{\{g_I^1\}}\omega _{\{g_I^1\}}}`$ (5.26)
$`=`$ $`C(\tau _{\{g_I\}}\overline{\tau }_{\{g_I\}}).`$
Now we will relate this to the variation of the twisted $`\tau _{\{g_I\}}`$. The moduli dependent part of the latter, as seen from Eq.(5.16) and (5.17), is proportional to the ratio of the periods $`_{𝐛_1𝐛_2𝐛_1^1𝐛_2^1}\omega _{\{g_I\}}/_{𝐚_1}\omega _{\{g_I\}}`$. From the variation of the twisted differential given in (5.25) we find
$`\delta _a\tau _{\{g_I\}}`$ $`=`$ $`C(\overline{D}_{\{g_I^1\}}D_{\{g_I\}})`$ (5.27)
$`=`$ $`C(\tau _{\{g_I\}}\overline{\tau }_{\{g_I\}})`$
$`=`$ $`{\displaystyle \mu _a\omega _{\{g_I\}}\omega _{\{g_I^1\}}},`$
where in the second equality we have used Eqs.(5.16) and (5.17) and the fact that $`\tau _{\{g_I\}}=\tau _{\{g_I^1\}}`$. The determinant appearing in (5.24) therefore just gives the Jacobian of the transformation from the moduli to $`\tau _{\{g_I\}}`$. If the latter, for $`I=3,4,5`$ are independent functions on the moduli space, then the Jacobian is non-vanishing and the resulting measure of integration on the moduli space becomes $`_Id\tau _{\{g_I\}}`$.
In the topological twisted theory the non-zero modes of the bosons and fermions cancel leaving only the correction factor coming from the boundary conditions on the bosons discussed in section 2 and summarized in Eq.(2.16). For the twisted bosons, $`\tau `$ in (2.16) is replaced by the corresponding twisted $`\tau `$. The resulting $`\mathrm{det}\tau `$ factor cancels the one appearing in lattice partition function (5.18).
The above calculations were done in the D6 brane theory. By T-dualizing we can go to the magnetized D9 brane theory. This amounts to replacing $`T_IU_I`$ and $`\mathrm{Im}U_I\mathrm{Im}T_I=\sqrt{G_I}`$. Here, $`U`$ is the usual complex structure modulus of the torus $`T^2`$, given in terms of its metric $`G_{ij}`$, $`i,j=1,2`$: $`U=(G_{12}+i\sqrt{G})/G_{11}`$. The resulting topological partition function becomes:
$$F^{(0,3)}=3N\underset{I=3}{\overset{5}{}}d\stackrel{~}{\tau }_{\{g_I\}}p^If[U_I,\stackrel{}{\alpha }_I;i\stackrel{~}{\tau }_{\{g_I\}}],$$
(5.28)
where the integration variables
$$\stackrel{~}{\tau }_{\{g_I\}}=\sqrt{G_I}(p_1^I)^2\left[1+(H_1^I)^2\right]\tau _{\{g_I\}}.$$
(5.29)
Here $`H_a^I=\frac{q_a^I}{p_a^I\sqrt{G_I}}`$ is the magnetic field on the $`a`$-th brane stack in the $`I`$-th plane. The function $`f`$ is given by:
$$f(U,\stackrel{}{\alpha },l)=\frac{1}{\mathrm{Im}U}\underset{n_1,n_2}{}(n_1+n_2\overline{U})e^{2i\pi \stackrel{}{n}\stackrel{}{\alpha }}e^{\frac{\pi }{\mathrm{Im}U}|n_1+n_2U|^2l},$$
(5.30)
where the lattice sum of integers $`\stackrel{}{n}^I`$ for each $`I`$ satisfy the conditions
$$p_1^I\stackrel{}{n}_I+p_2^I\stackrel{}{n}_I^{}+p_3^I\stackrel{}{n}_I^{\prime \prime }=q_1^I\stackrel{}{n}_I+q_2^I\stackrel{}{n}_I^{}+q_3^I\stackrel{}{n}_I^{\prime \prime }=0$$
(5.31)
for some integer vectors $`\stackrel{}{n}_I^{}`$ and $`\stackrel{}{n}_I^{\prime \prime }`$. $`\stackrel{}{\alpha }_I`$ is the Wilson line in the $`I`$-th plane and is given in terms of a certain linear combination of the Wilson lines on the three brane stacks:
$$\stackrel{}{\alpha }^I=\stackrel{}{\alpha }_1^I\frac{p_3^Iq_1^Iq_3^Ip_1^I}{p_3^Iq_2^Iq_3^Ip_2^I}\stackrel{}{\alpha }_2^I+\frac{p_2^Iq_1^Iq_2^Ip_1^I}{p_3^Iq_2^Iq_3^Ip_2^I}\stackrel{}{\alpha }_3^I,$$
(5.32)
where the subscripts $`1,2,3`$ refer to the three different stacks of branes on the three boundaries of the world-sheet. Obviously, the result vanishes for $`\stackrel{}{\alpha }_I=0`$ for any $`I`$, due to the $`\stackrel{}{n}_I\stackrel{}{n}_I`$ symmetry of the lattice sum. Note that non-trivial $`\stackrel{}{\alpha }_I`$, although breaks the corresponding R-symmetry, does not change the spectrum. The positive integers $`p_I`$ appearing in (5.28) satisfy (5.19) which can be reexpressed in terms of the flux data as $`p_I`$ being the smallest positive integer such that
$$p_I\frac{p_3^Iq_1^Iq_3^Ip_1^I}{p_3^Iq_2^Iq_3^Ip_2^I}p_I\frac{p_2^Iq_1^Iq_2^Ip_1^I}{p_3^Iq_2^Iq_3^Ip_2^I}.$$
(5.33)
It is worth making some comments on the topological partition function (5.28). The first comment is regarding modular invariance. From the constraint (5.31) it follows that $`n_1^I`$ and $`n_2^I`$ are arbitrary integer multiples of $`p_I`$. The discussion following Eq.(5.19) then implies that (5.28) is invariant under the permutation of the boundaries.
The second comment is regarding the target space duality properties of (5.28). We have already mentioned that it has full $`SL(2,)_U`$ symmetry despite the fact that the sum is over a sublattice defined by the constraints (5.31). This is because the constraints are $`SL(2,)_U`$ covariant. As for the monodromy under shifts of the Wilson lines we note that the constraints imply that $`\stackrel{}{n}_Ip^I\stackrel{}{Z}_I`$ where $`\stackrel{}{Z}`$ is a two-dimensional vector with arbitrary integer components. It follows then that $`F^{(0,3)}`$ is invariant under shifts $`\stackrel{}{\alpha }_I\stackrel{}{\alpha }_I+\stackrel{}{Z}_I/p_I`$.
The topological partition function (5.28) appears deceptively simple as products of three separate integrals, however, the domain of integration over $`\stackrel{~}{\tau }_{\{g_I\}}`$ is very complicated and in general mixes the three integrals. There are some comments worth making about $`F^{(0,3)}`$:
* Taking derivatives with respect to anti-holomorphic moduli should result in a total derivative in the world-sheet moduli space, as follows from the general discussion of section 4. The anti-holomorphic moduli in the present case are the complex structure closed string moduli $`\overline{U}_I`$ and the complexified Wilson lines
$$A_I=\alpha _2^IU_I\alpha _1^I.$$
(5.34)
The physical quantities are invariant under $`A_IA_I+1/p_I`$ and $`A_IA_I+U_I/p_I`$. The derivatives with respect to the anti-holomorphic moduli can be easily evaluated using the identities:
$`{\displaystyle \frac{f}{\overline{A}}}`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{Im}U}}{\displaystyle \frac{Z}{l}};Z={\displaystyle \underset{n_1,n_2}{}}e^{2i\pi \stackrel{}{a}\stackrel{}{n}}e^{\frac{\pi }{\mathrm{Im}U}|n_1+n_2U|^2l}`$
$`{\displaystyle \frac{f}{\overline{U}}}`$ $`=`$ $`{\displaystyle \frac{1}{4\mathrm{I}\mathrm{m}U}}{\displaystyle \frac{(lf)}{l}}{\displaystyle \frac{\mathrm{Im}A}{(\mathrm{Im}U)^2}}{\displaystyle \frac{Z}{l}},`$ (5.35)
which follow from the form of the function $`f`$ in Eq.(5.30).
* The dependence on the Kähler moduli of the tori, namely $`\sqrt{G_I}`$, should drop out because the complexification of these moduli involve the Ramond-Ramond $`B`$ fields. Since the latter are associated with continuous shift symmetries in perturbation theory, the dependence on them should be trivial.
## 6 Gaugino Masses
In this section, we compute the Majorana gaugino masses in the general non-supersymmetric case with the sum of the brane angles different from zero. Using their relation to the orbifold twists (5.10), we parameterize the supersymmetry breaking as
$$\underset{I=3}{\overset{5}{}}h_I=2ϵ,$$
(6.1)
where for simplicity we dropped the cycle indices. The parameter $`ϵ`$ fixes the scalar masses and in the weak field limit $`ϵ0`$, it is proportional to the expectation value of the corresponding abelian D-term. In this limit, the gaugino masses are determined from the effective F-term $`(\mathrm{Tr}W^2)^2`$ and should therefore be identified with the topological partition function $`F^{(0,3)}`$.
The gaugino vertices in the $`1/2`$ ghost picture are given in Eq.(3.1) and are inserted on one of the three boundaries of the $`(g=0,h=3)`$ Riemann surface. Moreover one has to insert three picture changing operators $`e^\phi T_F`$, which we also choose to be on the boundaries. One of them is needed to change the ghost picture of one gaugino to $`+1/2`$, while the other two arise from the integration over the supermoduli:
$$m_{1/2}=g_s^2[dt]dxdy𝒜;𝒜=V_+^{(1/2)}(x)V_{}^{(1/2)}(y)\underset{i=1}{\overset{3}{}}e^\phi T_F(z_i),$$
(6.2)
where $`g_s`$ is the string coupling and its power takes into account the normalization of the gaugino kinetic terms on the disk. The moduli integration is over the fundamental domain (5.5), while $`x`$ and $`y`$ are integrated over the gaugino boundary. The dependence on the positions $`z_i`$ of the picture changing operators is a gauge artifact and should disappear from the physical amplitude.
By internal charge conservation, since both gauginos carry charge $`+1/2`$ for $`\varphi _3`$, $`\varphi _4`$ and $`\varphi _5`$, only the internal part of the world-sheet supercurrents (3.8) contributes; each $`T_F^{int}`$ should provide $`1`$ charge for $`\varphi _3`$, $`\varphi _4`$ and $`\varphi _5`$, respectively. The amplitude (6.2) then becomes:
$`𝒜`$ $`=`$ $`e^{\phi /2}(x)e^{\phi /2}(y){\displaystyle \underset{I=3}{\overset{5}{}}}e^\phi (u_I){\displaystyle \underset{I=1}{\overset{2}{}}}e^{i\varphi _I/2}(x)e^{i\varphi _I/2}(y)`$ (6.3)
$`\times {\displaystyle \underset{I=3}{\overset{5}{}}}e^{i\varphi _I/2}(x)e^{i\varphi _I/2}(y)e^{i\varphi _I}(u_I),`$
where $`\{u_I\}`$, $`I=3,4,5`$, is a permutation of $`\{z_i\}`$, $`i=1,2,3`$, and an implicit summation over all permutations is understood. Performing the contractions for a given spin structure $`s`$, one finds:
$`𝒜_s`$ $`=`$ $`{\displaystyle \frac{\theta _s^2(\frac{1}{2}(xy))\underset{I=3}{\overset{5}{}}\theta _{s,h_I}(\frac{1}{2}(x+y)u_I)X_{h_I}(u_I)}{\theta _s(\frac{1}{2}(x+y)_{I=3}^5u_I+2\mathrm{\Delta })}}`$ (6.4)
$`\times {\displaystyle \frac{\sigma (x)\sigma (y)}{_{I<J}^{3,4,5}E(u_I,u_J)_{I=3}^5\sigma ^2(u_I)}}\times {\displaystyle \frac{Z_2}{Z_1^4_{I=3}^5Z_{1,h_I}}}Z_{lat},`$
where $`\theta _s`$ is the genus-two theta-function of spin structure $`s`$, $`E`$ is the prime form, $`\sigma `$ is a one-differential with no zeros or poles and $`\mathrm{\Delta }`$ is the Riemann $`\theta `$-constant. $`Z_{1,h_I}`$ is the (chiral) determinant of the $`h_I`$-twisted $`(1,0)`$ system, and $`Z_2`$ is the chiral non-zero mode determinant of the $`(2,1)`$ b-c ghost system. Finally, $`Z_{lat}`$ stands for all zero-mode parts of space-time and internal coordinates, while an implicit summation over lattice momenta should be performed, taking into account also the $`X_I`$ factors in (6.4).
In order to perform an explicit sum over spin structures, one should choose the positions of the picture changing operators to satisfy a condition that makes the argument of the $`\theta `$-function in the denominator equal to $`(xy)/2`$, so that one $`\theta _s`$ factor simplifies. The resulting relation however,
$$\underset{i=1}{\overset{i=3}{}}z_i=y+2\mathrm{\Delta },$$
(6.5)
is not allowed. To bypass this difficulty, we insert in the amplitude (6.2) another vertex of an open string Wilson line associated to the $`I=3`$ internal plane, in the $`1`$ ghost picture:
$$V^{(1)}=:e^\phi \psi _3(w):,$$
(6.6)
accompanied by a forth picture changing operator at the boundary position $`z_4`$. Obviously, we should also perform a third integration over its position $`w`$. The new vertex brings another unit of charge for $`\varphi _3`$, and thus, two of the four supercurrents should provide $`1`$ charge for $`\varphi _3`$. As we will demonstrate below, one can now choose an appropriate gauge condition which allows to perform the spin structure sum and show that the amplitude can be written as the variation with respect to the Wilson line of the original one (6.4), evaluated formally using the “forbidden” gauge choice (6.5).
Indeed, the correlators in (6.3) now become:
$`𝒜`$ $`=`$ $`e^{\phi /2}(x)e^{\phi /2}(y)e^\phi (w){\displaystyle \underset{I=3}{\overset{6}{}}}e^\phi (u_I){\displaystyle \underset{I=1}{\overset{2}{}}}e^{i\varphi _I/2}(x)e^{i\varphi _I/2}(y)`$
$`\times `$ $`e^{i\varphi _3/2}(x)e^{i\varphi _3/2}(y)e^{i\varphi _3}(w){\displaystyle \underset{I=3,6}{}}e^{i\varphi _3}(u_I){\displaystyle \underset{I=4,5}{}}e^{i\varphi _I/2}(x)e^{i\varphi _I/2}(y)e^{i\varphi _I}(u_I),`$
where $`\{u_I\}`$, $`I=3,\mathrm{},6`$, is a permutation of $`\{z_i\}`$, $`i=1,2,3,4`$, and an implicit summation over all permutations is understood. Performing the contractions, one finds that the expression (6.4) is modified as:
$`𝒜_s`$ $`=`$ $`{\displaystyle \frac{\theta _s^2(\frac{1}{2}(xy))\theta _{s,h_3}(\frac{1}{2}(x+y)u_3+wu_6)\underset{I=4,5}{}\theta _{s,h_I}(\frac{1}{2}(x+y)u_I)}{\theta _s(\frac{1}{2}(x+y)+w_{I=3}^6u_I+2\mathrm{\Delta })}}`$
$`\times {\displaystyle \frac{\sigma (x)\sigma (y)\sigma ^2(w)}{_{I<J}^{3,4,5}E(u_I,u_J)_{I=3}^6\sigma ^2(u_I)}}{\displaystyle \frac{E(w,u_4)E(w,u_5)}{E(u_6,u_4)E(wu_6,u_5)}}`$
$`\times {\displaystyle \frac{Z_2Z_{lat}}{Z_1^4_{I=3}^5Z_{1,h_I}}}{\displaystyle \underset{I=3,6}{}}X_{3,h_3}(u_I){\displaystyle \underset{I=4,5}{}}X_{I,h_I}(u_I).`$
To perform the spin structure sum, we use the allowed gauge condition:
$$\underset{i=1}{\overset{i=4}{}}z_i=y+w+2\mathrm{\Delta }.$$
(6.9)
The sum over spin structures converts the first factor in (6) to:
$$\theta _ϵ(x\mathrm{\Delta })\theta _{h_3+ϵ}(u_3+u_6w\mathrm{\Delta })\theta _{h_4+ϵ}(u_4\mathrm{\Delta })\theta _{h_5+ϵ}(u_5\mathrm{\Delta }),$$
(6.10)
where we used Eq.(6.1).
We now use three bosonization formulae. First,
$`\theta _{h_3+ϵ}(u_3+u_6w\mathrm{\Delta }){\displaystyle \frac{E(u_3,u_6)Z_1^{1/2}}{E(u_3,w)E(u_6,w)}}{\displaystyle \frac{\sigma (u_3)\sigma (u_6)}{\sigma (w)}}`$ (6.11)
$`=\overline{\psi }_{h_3+ϵ}(u_3)\overline{\psi }_{h_3+ϵ}(u_6)\psi _{h_3ϵ}(w)Z_{1,h_3+ϵ},`$
where $`\psi `$ and $`\overline{\psi }`$ are conformal fields of dimension zero and one, respectively, twisted as indicated by their subscripts. Then
$$\theta _{h_I+ϵ}(u_I\mathrm{\Delta })\sigma (u_I)Z_1^{1/2}=\omega _{h_I+ϵ}(u_I)Z_{1,h_I+ϵ},$$
(6.12)
where $`\omega _h`$ is an abelian differential twisted by $`h`$. Finally,
$$\theta _ϵ(\underset{I=3}{\overset{6}{}}u_Iw3\mathrm{\Delta })\frac{\underset{I<J}{}E(u_I,u_J)}{_{I=3}^6E(u_I,w)}\frac{\underset{I=3}{\overset{6}{}}\sigma ^3(u_I)}{\sigma ^3(w)}Z_1^{1/2}=\underset{I=3}{\overset{6}{}}b_ϵ(u_I)c_ϵ(w)Z_{2,ϵ},$$
(6.13)
where the correlator and the non-zero mode determinant of the b-c ghost system in the r.h.s. are twisted according to the subscripts.
Multiplying the amplitude (6) by $`\theta _ϵ(y\mathrm{\Delta })/\theta _ϵ(_{I=3}^6u_Iw3\mathrm{\Delta })`$ which is equal to the identity by our gauge condition (6.9), and using the spin structure sum (6.10) and the bosonization formulae (6.11) - (6.13), we obtain:
$`𝒜`$ $`=`$ $`\theta _ϵ(x\mathrm{\Delta })\theta _ϵ(y\mathrm{\Delta }){\displaystyle \underset{I=3}{\overset{5}{}}}\left({\displaystyle \frac{Z_{1,h_I+ϵ}}{Z_{1,h_I}}}\right){\displaystyle \frac{\overline{\psi }_{h_3ϵ}(u_3)\overline{\psi }_{h_3+ϵ}(u_6)\psi _{h_3ϵ}(w)}{_{I=3}^6b_ϵ(u_I)c_ϵ(w)}}`$ (6.14)
$`\times {\displaystyle \underset{I=4,5}{}}\omega _{h_I+ϵ}(u_I)\sigma (x)\sigma (y){\displaystyle \frac{Z_2}{Z_{2,ϵ}}}{\displaystyle \frac{1}{Z_1^3}}{\displaystyle \underset{I=3,6}{}}X_{3,h_3}(u_I){\displaystyle \underset{I=4,5}{}}X_I(u_I).`$
Replacing $`\omega _{h_I+ϵ}(u_I)`$ by $`\overline{\psi }_{h_I+ϵ}(u_I)`$ and using the bosonization formula (6.12) twice, one finds:
$`𝒜={\displaystyle \frac{\underset{i=1}{\overset{4}{}}G_ϵ^{}(z_i)\psi _{h_3ϵ}(w)}{_{i=1}^4b_ϵ(z_i)c_ϵ(w)}}{\displaystyle \underset{I=3}{\overset{5}{}}}\left({\displaystyle \frac{Z_{1,h_I+ϵ}}{Z_{1,h_I}}}\right){\displaystyle \frac{Z_2}{Z_{2,ϵ}}}\left({\displaystyle \frac{Z_{1,ϵ}}{Z_1}}\right)^2\omega _ϵ(x)\omega _ϵ(y),`$ (6.15)
where $`G_ϵ^{}`$ is the internal $`N=2`$ supercurrent twisted by $`ϵ`$.
One can now show that the ratio of correlation functions in (6.15) is independent of the positions $`z_i`$:
$$\frac{𝒩}{𝒟}\frac{\underset{i=1}{\overset{4}{}}G_ϵ^{}(z_i)\psi _{h_3ϵ}(w)}{_{i=1}^4b_ϵ(z_i)c_ϵ(w)}=B_ϵX_{3,h_3}\omega _{h_3}(w),$$
(6.16)
with $`B_ϵ`$ a constant to be determined. Indeed, as a function of $`w`$, both numerator $`𝒩`$ and denominator $`𝒟`$ have first order poles when $`wz_i`$ with residues $`_{i=1}^3G_ϵ^{}(z_i)X_{3,h_3}`$ and $`_{i=1}^3b_ϵ(z_i)`$, when for instance $`wz_4`$. These are equal, up to a $`z_i`$-independent multiplicative factor $`B_ϵ`$. Consider now the combination $`𝒩B_ϵX_{3,h_3}\omega _{h_3}(w)𝒟`$. This is holomorphic in $`w`$ twisted 0-form and therefore vanishes identically. Thus, $`B_ϵ`$ is given by:
$$B_ϵ\omega _{h_3}(w)=\frac{\underset{i=1}{\overset{3}{}}G_ϵ^{}(z_i)}{_{i=1}^3b_ϵ(z_i)}.$$
(6.17)
Note that this is reduced to the same result as the one that would be obtained using the “forbidden” gauge condition (6.5), with an additional differentiation with respect to the Wilson line (6.6) associated to the presence of $`X_{3,h_3}\omega _{h_3}(w)`$ in (6.16).
It follows that the gaugino mass (6.2) is given by:
$$m_{1/2}=g_s^2[dt]B_ϵ\underset{I=3}{\overset{5}{}}\left(\frac{Z_{1,h_I+ϵ}}{Z_{1,h_I}}\right)\frac{Z_2}{Z_{2,ϵ}}\left(\frac{Z_{1,ϵ}}{Z_1}\right)^2,$$
(6.18)
where we performed the boundary integrals over $`x`$ and $`y`$, using the canonical normalization over the $`𝐚`$-cycles $`_𝐚\omega _ϵ=1`$, since $`ϵ`$ is a twist around the $`𝐛`$-cycles. In the limit of small supersymmetry breaking scale $`ϵ0`$, one has $`Z_{1,ϵ}ϵZ_1`$, and thus
$$m_{1/2}g_s^2ϵ^2[dt]B_0=g_s^2ϵ^2F^{(0,3)}.$$
(6.19)
## 7 $`\mathrm{\Pi }`$-terms and Matter Fermion Masses
In the previous section, we studied the (Majorana) gaugino mass terms generated via D-term supersymmetry breaking. They originate from the supersymmetric F-term $`(\mathrm{Tr}W^2)^2`$ once the auxiliary D-component of the vector superfield acquires a non-zero VEV along a magnetized $`U(1)`$ group factor. The corresponding topological coupling is given by $`F^{(0,3)}`$, with the three boundaries of the world-sheet attached (in a T-dual picture) to three different stacks of D6 branes intersecting at angles in the internal compactification space. In order to get a non-zero answer, we had to further break the discrete R-symmetry by turning on suitable Wilson lines. In the following we analyze the structure of the terms appearing in the holomorphic anomaly equation for $`F^{(0,3)}`$. The are of the form $`\mathrm{\Pi }\mathrm{Tr}W^2`$, where $`\mathrm{\Pi }`$ is a chiral projection of a non-holomorphic function of chiral superfields. Generically, its lower component includes a fermion bilinear of the form $`\overline{\mathrm{\Psi }}_{\overline{i}}\overline{\mathrm{\Psi }}_{\overline{j}}`$. Hence, upon D-term supersymmetry breaking, they will also induce some fermion masses, in this case of Dirac type.
In order to examine the holomorphic anomaly of $`F^{(0,3)}`$, we take an anti-holomorphic derivative $`_{\overline{i}}`$, with respect to an open string Wilson line $`\overline{i}`$. From the analysis of section 4, one finds contributions from three possible degeneration limits. The two open string degenerations, (1) and (2), involve $`F_{\overline{i};\overline{j}}^{(0,2)}`$, where $`\overline{j}`$ denotes an intermediate anti-chiral open string state. In this section, we restrict our discussion to brane configurations that do not involve parallel stacks in any of the three internal planes. This means that in the supersymmetric limit there are no $`N=2`$ supersymmetric sectors in the one-loop partition function.<sup>6</sup><sup>6</sup>6In the next section however, we will consider examples with such $`N=2`$ sectors. In the absence of such sectors, the dividing degeneration (1) does not contribute because $`\overline{j}`$ is an untwisted operator and the corresponding annulus diagrams vanish. In the handle degeneration limit (2) however, since the open string propagating through the handle stretches between two non-parallel brane stacks, it is necessarily twisted. The third degeneration limit (3) involves intermediate closed strings which, in $`T^6`$ compactifications, are always untwisted. Hence, the corresponding annulus diagrams also vanish.
Thus, the only quantity that appears in the holomorphic anomaly of $`F^{(0,3)}`$ arises from Eq.(4.7), and involves $`D_jF_{\overline{i};\overline{j}}^{0,2}`$, where $`j`$ and $`\overline{j}`$ label chiral and anti-chiral twisted open string states with $`U(1)`$ charges +1 and +2 respectively. These are therefore bi-fundamental states and represent open strings stretched between two different stacks of D6 branes (say $`a=1`$ and $`a=2`$) that are not parallel to each other. If the intersection of these two stacks of D6 branes preserves supersymmetry then the twist angles $`h_I`$ in the three internal planes (tori), $`I=3,4,5`$, satisfy $`_Ih_I=1`$ and the corresponding twisted states are massless. In the topological theory
$$D_jF_{\overline{i};\overline{j}}^{(0,2)}=𝑑t^1𝑑t^2𝑑x\left\{\underset{m=1}{\overset{2}{}}\mu _m(G_L^{}+G_R^{})\right\}V_j(u)V_{\overline{i}}(x)V_{\overline{j}}(y)^𝒯$$
(7.1)
with
$`V_j^𝒯`$ $`=`$ $`:\sigma e^{i_Ih_I\varphi _I}:`$
$`V_{\overline{j}}^𝒯`$ $`=`$ $`:\overline{\sigma }e^{i_I(1h_I)\varphi _I}:`$ (7.2)
$`V_{\overline{i}}^𝒯`$ $`=`$ $`:e^{i\varphi _3}:`$
where, for concreteness, we chose $`\overline{i}`$ to indicate the anti-holomorphic derivative with respect to the Wilson line in the $`I=3`$ plane. $`\sigma `$ and $`\overline{\sigma }`$ denote the bosonic twisted fields. We have included the superscript $`𝒯`$ to indicate that these are the vertices in the topological theory. Note that while $`V_{\overline{i}}^𝒯`$ has twisted $`U(1)`$ charge $`1`$ and dimension 1 (and hence its position $`x`$ is integrated on the world-sheet boundary), the operators $`V_j^𝒯`$ and $`V_{\overline{j}}^𝒯`$ have dimension 0 and charges +1 and +2 respectively. The world-sheet moduli space, labeled by $`t^m`$, with the corresponding Beltrami differentials $`\mu _m`$, is two dimensional, one being the usual modulus associated with the annulus and the second being the relative distance between $`u`$ and $`y`$ on one of the boundaries. The open string state going through the loop (i.e. annulus) is itself twisted by $`\stackrel{~}{h}_I`$ in the three planes since the two boundaries of the annulus sit on two different stacks of D6 branes (say $`a=1`$ and $`a=3`$). We assume here that the combined system preserves supersymmetry so that $`_I\stackrel{~}{h}_I=0`$ mod integers. Note that as the open string propagating in the annulus crosses one of the twist operator insertions it becomes an open string stretched between the stacks 2 and 3 and when it crosses the second twisted operator it becomes again the string stretched between 1 and 3.
To compute this correlation function we can go to the torus double cover of the annulus and take the appropriate square root. The result is
$`D_jF_{\overline{i};\overline{j}}^{(0,2)}`$ $`=`$ $`{\displaystyle 𝑑t^1𝑑t^2𝑑x\underset{m=1}{\overset{2}{}}d^2z_m\mu _m(z_m)\frac{_I\theta _{\stackrel{~}{h}_I}[yY_I+h_I(uy)]}{\eta ^3_JE(Y_J,u)^{h_J}E(Y_J,y)^{1h_J}}}`$ (7.3)
$`\times E(x,y)^{1_Kh_K^2}ϵ^{bc}X_4(z_b)X_5(z_c)\sigma (u)\overline{\sigma }(y),`$
where $`Y_3=x`$, $`Y_4=z_b`$ and $`Y_5=z_c`$. Here, $`E(x,y)`$ denotes the genus one prime form
$$E(x,y)=\frac{\theta _1(xy)}{\theta _1^{}(0)}=\frac{\theta _1(xy)}{2\pi \eta ^3}.$$
(7.4)
Note that in the above equation, $`\mathrm{}`$ denotes the unnormalized correlator including the partition functions of the internal bosons. The physical string amplitude computed by the above quantity is
$$(^+)^2\overline{\mathrm{\Psi }}_{\overline{i}}\overline{\mathrm{\Psi }}_{\overline{j}}\mathrm{\Phi }_j,$$
(7.5)
where $`^+`$ is the self-dual field strength and $`\mathrm{\Phi }`$ is the twisted scalar. In Ref. this computation has been done for the heterotic string and shown to give rise to the topological amplitude above. The methods of Ref. extend trivially to the open string case, as it can be seen by going to the double cover of the world-sheet, where the spin structure sum involves only one sector (say left-moving sector) exactly as in the heterotic theory. In the following, we go directly to the broken supersymmetry case and derive the above topological term in the limit of supersymmetry restoration, in analogy with the gaugino masses.
We now compute directly the mass term for the fermions $`\mathrm{\Psi }`$ when supersymmetry is broken via a VEV of the auxiliary D component of the gauge vector superfield. Specifically, we take $`_I\stackrel{~}{h}_I=ϵ0`$ mod integer, while keeping the supersymmetry condition on $`h_I`$, namely $`_Ih_I=1`$. In other words the $`h`$-twisted sector representing strings stretched between the D6 brane stacks 1 and 2 does not break supersymmetry but the presence of stack 3 breaks it. This corresponds to the situation when $`\mathrm{\Psi }`$ and its superpartners are massless, but supersymmetry is broken by the presence of the other boundary of the annulus associated to the stack 3. Note that the tree-level mass matrix does not mix $`h`$-twisted fermions with the Wilson line fermions. However we will show below that at one loop the closed string exchange between the stack 3 and the intersection of 1 and 2 gives rise to such a mass term.
The amplitude in question is the annulus three point function of open string states:
$$M_{\overline{i}\overline{j}j}=\overline{\mathrm{\Psi }}_{\overline{i}}\overline{\mathrm{\Psi }}_{\overline{j}}\mathrm{\Phi }_j.$$
(7.6)
The vertex operators for the fermions in the $`(1/2)`$ picture and the scalar in the $`(1)`$ picture are
$`V_j(u)`$ $`=`$ $`:ce^\phi \sigma e^{i_Ih_I\varphi _I}:`$
$`V_{\overline{j}}(y)`$ $`=`$ $`:ce^{\frac{\phi }{2}}e^{\frac{\varphi _1\varphi _2}{2}}\overline{\sigma }e^{i_I(\frac{1}{2}h_I)\varphi _I}:`$ (7.7)
$`V_{\overline{i}}(x)`$ $`=`$ $`:e^{\frac{\phi }{2}}e^{\frac{\varphi _1+\varphi _2}{2}}e^{i\frac{\varphi _3+\varphi _4+\varphi _5}{2}}:.`$
We have inserted the bosonic ghosts $`c`$ at the vertex $`V_j`$ and $`V_{\overline{j}}`$ so we are treating the surface as twice-punctured annulus with the associated two moduli (the modulus of the annulus and the relative position between these two vertices). The vertex $`V_{\overline{i}}`$ is dimension 1 and has to be integrated. The total superghost charge of the three vertices is $`2`$ and therefore we need to insert two picture changing operators $`e^\phi T_F`$. Thus the amplitude (7.6) becomes
$$M_{\overline{i}\overline{j}j}=dt^1dt^2dx\underset{m=1}{\overset{2}{}}[\mu _m(b_L+b_R)]e^\phi T_F(z_1)e^\phi T_F(z_2)V_j(u)V_{\overline{i}}(x)V_{\overline{j}}(y.$$
(7.8)
The above amplitude should be independent of the positions $`z_1`$ and $`z_2`$ of the picture changing operators. Since the total $`\varphi _4`$ and $`\varphi _5`$ charges of the three vertices is +1 each, it follows that the only relevant terms in the picture changing operators are
$$e^\phi T_Fe^\phi \underset{I=4,5}{}e^{i\varphi _I}X_I.$$
(7.9)
The mass matrix becomes
$`M_{\overline{i}\overline{j}j}`$ $`=`$ $`{\displaystyle 𝑑t^1𝑑t^2𝑑x\underset{m=1}{\overset{2}{}}[\mu _m(b_L+b_R)]c(u)c(y)}`$ (7.10)
$`\times {\displaystyle \underset{s}{}}{\displaystyle \frac{\theta _s(xy)^2}{\theta _s(uz_a+\frac{x+y}{2})}}[{\displaystyle \underset{I}{}}\theta _{s+\stackrel{~}{h}_I}({\displaystyle \frac{x+y}{2}}+h_I(uy)Y_I)]`$
$`\times \eta ^8(R_𝒜\text{Im}\tau )^1{\displaystyle \frac{_JE(Y_J,y)^{h_J}E(Y_J,u)^{1h_J}}{E(y,u)^{_Kh_K^2}E(x,y)E(x,u)E(z_b,z_c)}}`$
$`\times ϵ^{bc}X_4(z_b)X_5(z_c)\sigma (u)\overline{\sigma }(y),`$
where $`Y_I`$ are as in Eq.(7.3). The above formula takes into account the annulus correction factor $`R_𝒜`$ for the four spacetime bosons, c.f. Eq.(2.15), and the $`(\text{Im}\tau )^1`$ factor due to their zero modes; here, $`\tau `$ denotes the usual untwisted modulus of the torus double cover. The power of the $`\eta `$ function is determined as follows: $`\eta ^4`$ comes from spacetime bosons, $`\eta ^5`$ come from the 5 real scalars that are the bosonization of 10 fermions (spacetime as well as internal), and finally $`\eta `$ comes from the bosonization of superghost. In order to perform the spin structure sum over $`s`$, we choose the following gauge condition for the positions of the picture changing operators:
$$z_1+z_2uy=0.$$
(7.11)
With this gauge choice the theta function in the denominator coming from the superghosts cancels with one of the theta function coming from the space-time fermions. After summing over spin structures we obtain:
$`M_{\overline{i}\overline{j}j}`$ $`=`$ $`{\displaystyle 𝑑t^1𝑑t^2𝑑x\underset{m=1}{\overset{2}{}}(\mu _m(b_L+b_R))c(u)c(y)}`$ (7.12)
$`\theta _ϵ(0)[{\displaystyle \underset{I}{}}\theta _{\stackrel{~}{h}_Iϵ}(yY_I+h_I(uy))]`$
$`\times \eta ^8(R_𝒜\text{Im}\tau )^1{\displaystyle \frac{_JE(Y_J,y)^{h_J}E(Y_J,u)^{1h_J}}{E(y,u)^{_Kh_K^2}E(x,y)E(x,u)E(z_b,z_c)}}`$
$`\times ϵ^{bc}X_4(z_b)X_5(z_c)\sigma (u)\overline{\sigma }(y).`$
In this equation and in the following $`\theta _ϵ`$ or $`\theta _{\stackrel{~}{h}_Iϵ}`$ denote the odd theta function twisted by $`ϵ`$ or $`\stackrel{~}{h}_Iϵ`$. Note that these twists are purely along the a-cycle (b-cycle) in the open (closed) string channel.
In the next step, we use the bosonization formula for $`b,c`$ system twisted by $`ϵ`$:
$$b(z_1)b(z_2)c(u)c(y)_ϵ=\frac{\theta _ϵ(z_1+z_2uy)E(z_1,z_2)E(u,y)}{\eta _{m=1}^2E(z_m,u)E(z_m,y)}.$$
(7.13)
Here again, $`\mathrm{}`$ is the unnormalized correlator, i.e. the complete result of the four-point function in the $`b,c`$ CFT that also includes its non-zero mode determinant. On the r.h.s., there is $`1/\eta `$ factor because the bosonization of $`b,c`$ is one real scalar. Then we can rewrite the mass term in the form
$`M_{\overline{i}\overline{j}j}`$ $`=`$ $`{\displaystyle 𝑑t^1𝑑t^2𝑑x\underset{m=1}{\overset{2}{}}[\mu _m(b_L+b_R)]c(u)c(y)H_ϵ(z_1,z_2,u,x,y)}`$ (7.14)
$`\times {\displaystyle \frac{\theta _ϵ(0)^2}{\eta ^6}}(R_𝒜\text{Im}\tau )^1,`$
where
$$H_ϵ(z_1,z_2,u,x,y)=\frac{_{a=1}^2G^{}(z_a)V_{\overline{i}}^𝒯(x)V_{\overline{j}}^𝒯(y)V_j^𝒯(u)_ϵ^𝒯}{b(z_1)b(z_2)c(u)c(y)_ϵ}.$$
(7.15)
Here we have used the gauge condition (7.11). Furthermore in Eq.(7.15), the numerator is the correlation function in the internal topological theory twisted by $`ϵ`$; in particular, this means that in Eq.(7.3), the functions $`\theta _{\stackrel{~}{h}_I}`$ are replaced by $`\theta _{\stackrel{~}{h}_Iϵ}`$. One can now argue that $`H`$ does not depend on $`z_m`$. As a function of $`z_1`$, both the numerator and denominator have a first order pole each at $`u`$ and $`y`$ and a first order zero at $`z_2`$. Each of them must have one more zero (since the corresponding line bundle has zero Chern class) but it must be in the same position as both the line bundles are characterized by the same twist $`ϵ`$ (i.e. the same point in the Jacobian torus). Since $`H`$ as a function of $`z_1`$ is a section of the trivial line bundle and has no zero or pole, it must be constant. A similar argument applies for $`z_2`$. Therefore
$$H_ϵ(z_1,z_2,u,x,y)=B_ϵ(u,x,y).$$
(7.16)
Now let us take the $`ϵ0`$ limit and compute the leading term in $`ϵ`$. Since in the closed string channel the twist $`ϵ`$ is purely along the b-cycle, in this limit $`\theta _ϵ(0)ϵ\eta ^3`$. The correction factor for spacetime bosons with Neumann boundary conditions yields $`(R_𝒜\text{Im}\tau )^1=1`$, see Eq.(2.16). The leading contribution is therefore of order $`ϵ^2`$, and at this order we can set $`ϵ=0`$ in $`B_ϵ(u,x,y)`$. Finally $`B_{ϵ=0}`$ multiplied by the ghost correlators involving the Beltrami differentials effectively replaces $`\mu b`$ by $`\mu G^{}`$ which gives the topological quantity $`D_jF_{\overline{i};\overline{j}}^{(0,2)}`$. The final result, to order $`ϵ^2`$, therefore is
$$M_{\overline{i}\overline{j}j}=ϵ^2D_jF_{\overline{i};\overline{j}}^{(0,2)}.$$
(7.17)
Eq.(7.17) determines the Yukawa type coupling (7.6) involving two anti-chiral fermions and one chiral boson. Note that this coupling cannot be derived from a superpotential and is not allowed by supersymmetry. Here, it was induced by a supersymmetry breaking D-term. If the twisted scalar $`\mathrm{\Phi }_j`$ acquires a VEV, it generates a Dirac mass term mixing the bi-fundamental fermions with a Wilson line fermion. In fact, such a VEV breaks also the gauge group, generating further mass mixing of the bi-fundamental fermions with gauginos through gauge Yukawa couplings. Since the right hand side term of Eq.(7.17) appears in the holomorphic anomaly of $`F^{(0,3)}`$, which in turn gives the gaugino mass at this order, we conclude that the corresponding fermion mass matrix elements are given by the holomorphic anomaly of the gaugino mass term.
One can further ask what happens when one takes anti-holomorphic derivative with respect to some moduli say $`_{\overline{i_2}}`$ of $`F_{\overline{i_1};\overline{j_1}}^{(0,2)}`$. From the analysis of section 4, after anti-symmetrizing in $`\overline{i}_1`$ and $`\overline{i}_2`$, we find that the result again comes from various degeneration limits. In particular, the degeneration corresponding to open-string intermediate states gives rise to $`F_{\overline{i_1}\overline{i_2};\overline{j_1}\overline{j_2}}^{(0,1)}`$ associated to a $`\mathrm{\Pi }^2`$ term in the effective action at the disk level. For instance the indices $`\overline{i}_1`$ and $`\overline{i}_2`$ can refer to some open string moduli fields, while $`\overline{j}_1`$ and $`\overline{j}_2`$ can refer to bi-fundamental open string states. This coupling can be evaluated by a 6-point function on a disk, involving two pairs of twist-antitwist fields and the two moduli fields, corresponding to $`D_{j_1}D_{j_2}F_{\overline{i_1}\overline{i_2};\overline{j_1}\overline{j_2}}^{(0,1)}`$.
## 8 A Simple Example
In this section, we present a simple toroidal example in which the topological partition function $`F^{(0,3)}`$, as well as all lower order quantities appearing in the holomorphic anomaly equations, can be computed either explicitly or by using some symmetry arguments.
Our starting point is a configuration in which every two brane stacks of the three boundaries intersect nontrivially only in two out of the three internal planes and are parallel in the remaining one, a configuration different from the one considered in the previous section. Furthermore, to avoid the enhancement of supersymmetry to $`N=2`$, and thus the vanishing of $`F^{(0,3)}`$, the plane in which the branes are parallel must be different in every of the three possible pairs. First we choose the horizontal axis in each plane along the stack $`a=3`$. Then, we pick stacks 2 and 3 to be parallel in the plane $`I=3`$, stacks 1 and 3 to be parallel in the plane $`I=4`$ and stacks 1 and 2 to be parallel in the plane $`I=5`$. This is described by the brane angles:
$$\theta _I^3=0\theta _4^1=\theta _3^2=0\theta _5^1=\theta _5^2\theta _3^1+\theta _5^1=0\theta _4^2+\theta _5^2=0,$$
(8.1)
where the last two relations follow from space-time supersymmetry. According to Eq.(5.10), the corresponding orbifold twists along the two b-cycles are:
$$(g_3^1,g_3^2)=(e^{4\pi i\theta },1)(g_4^1,g_4^2)=(1,e^{4\pi i\theta })(g_5^1,g_5^2)=(e^{4\pi i\theta },e^{4\pi i\theta }),$$
(8.2)
where the angle $`\theta =\theta _3^1`$ is an arbitrary parameter.
Now the constraint (5.31) on the lattice sum is solved trivially, leading for each plane to a summation over two unrestricted integers of the corresponding two-dimensional momentum lattice depending on the complex structure modulus $`U_I`$ and Wilson line $`A_I`$. Indeed, when two branes are parallel within a plane $`I`$, say the stacks $`a=1`$ and $`a=2`$, the corresponding magnetic fluxes $`p_a^I/q_a^I`$ are equal and since $`(p_a^I,q_a^I)`$ are relatively prime, one has $`p_1^I=p_2^I`$ and $`q_1^I=q_2^I`$. Eq.(5.31) then requires $`p_1^I(\stackrel{}{n}_I+\stackrel{}{n}_I^{})+p_3^I\stackrel{}{n}_I^{\prime \prime }=0`$ and $`q_1^I(\stackrel{}{n}_I+\stackrel{}{n}_I^{})+q_3^I\stackrel{}{n}_I^{\prime \prime }=0`$, implying $`\stackrel{}{n}_I+\stackrel{}{n}_I^{}=\stackrel{}{n}_I^{\prime \prime }=0`$ since the third stack must have non-trivial intersection with the other two. One is then left with a summation over two unrestricted integers, defined for instance by the vector $`\stackrel{}{n}_I`$. Note that the the physical Wilson line $`\stackrel{}{\alpha }`$ is given by the difference $`\stackrel{}{\alpha }^I=\stackrel{}{\alpha }_1^I\stackrel{}{\alpha }_2^I`$ and corresponds in the T-dual picture to the relative distance between the two parallel brane stacks.
Finally, there is an $`SL(2,)_I`$ action on each $`U_I`$ and $`A_I`$: $`U_I(aU_I+b)/(cU_I+d)`$ and $`A_IA_I/(cU_I+d)`$. The modular weights of $`F_n^{(0,h)}`$ are determined by their Kähler weights $`n+h2`$ . Thus $`F^{(0,3)}`$ transforms with weight 1 under each $`SL(2,)_I`$ and is also monodromy invariant under $`A_IA_I+1`$ and $`A_IA_I+U_I`$. Furthermore, it vanishes when $`A_I=1/2,U_I/2,(U_I+1)/2`$. By taking derivatives with respect to $`A_I`$ or $`\overline{A}_I`$ and setting $`A_I=0`$ one finds that the zeroes are of first order. On the other hand, for $`A_I=0`$, the two stacks become coincident and there are additional massless states. As a result, $`F^{(0,3)}`$ is singular at the origin and acquires a first order pole instead of a zero as in the other points. A simple ansatz for $`F^{(0,3)}`$, satisfying all properties above, is:
$$F^{(0,3)}=f_3\underset{I}{}H(U_I,A_I);H(U,A)=\frac{\theta ^{}(A)}{\theta (A)}+2i\pi \frac{\mathrm{Im}A}{\mathrm{Im}U},$$
(8.3)
where $`f_3`$ is a numerical constant and the prime denotes differentiation with respect to $`A`$. Actually, this expression, which will be verified subsequently by studying the holomorphic anomalies, also suggests that for this special brane configuration described by the orbifold twists (8.2), the integration domain of the three twisted world-sheet moduli $`\tau _{\{g_I\}}`$ is factorized into three independent integrals over the positive real line, so that each one can be performed explicitly yielding the result (8.3):
$$_0^{\mathrm{}}𝑑lf(U,\stackrel{}{a},l)=\frac{1}{\pi }\underset{n_1,n_2}{}^{}\frac{e^{2i\pi \stackrel{}{n}\stackrel{}{a}}}{n_1+n_2U}=\frac{1}{\pi }H(U,A),$$
(8.4)
where the function $`f`$ is given in Eq.(5.30). An appropriate $`SL(2,)`$ invariant regularization of the above sum leads to the r.h.s. part of the equation, with $`H`$ acquiring a non-holomorphic dependence as in (8.3).
Our strategy to prove Eq.(8.3) will be to compute the holomorphic anomaly as discussed in section 4 and show that the latter is reproduced by (8.3). Holomorphic ambiguity is then fixed by requirement of target space $`SL(2,)_I`$ and $`A_I`$ monodromy properties. Taking a derivative of (8.3) with respect to an anti-holomorphic open string Wilson line, for instance $`\overline{A}_3`$, one finds
$$_{\overline{A}_3}F^{(0,3)}=\frac{\pi f_3}{\mathrm{Im}U_3}\underset{I=4,5}{}H(U_I,A_I).$$
(8.5)
On the other hand, from the general discussion of section 4 on holomorphic anomaly, the contribution comes from various degeneration limits of the surface. To analyze this we need to study the behavior of the twisted $`\tau _{\{g_I\}}`$ in the three degeneration limits as shown in Fig. 4. The case that we are considering is characterized by the twists (8.2). It is more convenient to normalize the three twisted differentials $`\omega _{\{g_I\}}`$ along the periods shown in Fig. 6 of the Appendix, so that $`_{𝐚_2}\omega _{\{g_3\}}=_{𝐚_1}\omega _{\{g_4\}}=_{𝐚_1𝐚_2}\omega _{\{g_5\}}=1`$. The twisted $`\tau _{\{g_I\}}`$ are then defined as $`\tau _{\{g_3\}}=_{𝐛_2}\omega _{\{g_3\}}`$, $`\tau _{\{g_4\}}=_{𝐛_1}\omega _{\{g_4\}}`$ and $`\tau _{\{g_5\}}=_{𝐛_1𝐛_2}\omega _{\{g_5\}}`$. Note that for the twists (8.2), $`𝐛_2`$, $`𝐛_1`$ and $`𝐛_1𝐛_2`$ are closed cycles for $`\omega _{\{g_I\}}`$ for $`I=3,4,5`$ respectively.
Now let us consider the three degeneration limits shown in Fig. 4.
1. In the first one where the genus 2 double cover degenerates along the dividing geodesic corresponding to an open string intermediate state between two annuli whose double covers have cycles $`(𝐚_1,𝐛_1)`$ and $`(𝐚_2,𝐛_2)`$, $`\omega _{\{g_4\}}`$ and $`\omega _{\{g_3\}}`$ degenerate to the untwisted differentials of the two torii respectively while $`\omega _{\{g_5\}}`$ degenerates to twisted differentials on the two torii having first order pole at the node. Thus while $`\tau _{\{g_4\}}`$ and $`\tau _{\{g_3\}}`$ are finite (they are just the untwisted moduli associated with the two torii), $`\tau _{\{g_5\}}`$ becomes infinite. In fact in terms of the plumbing fixture coordinate, say $`t`$, $`\tau _{\{g_5\}}`$ goes as $`i\mathrm{ln}|t|`$. The two other such degenerations are obtained by permutations where either $`\tau _{\{g_3\}}`$ or $`\tau _{\{g_4\}}`$ goes to infinity keeping the remaining two finite. Taking derivative with respect to $`\overline{A}_3`$ and using Eq.(5.35), we note that the relevant such degeneration limit comes from $`\tau _{\{g_3\}}`$ going to infinity. The resulting contribution to the holomorphic anomaly is
$$_{\overline{A}_3}F^{(0,3)}=\underset{I=4,5}{}F_{\overline{A}_3;\overline{A}_I}^{(0,2)}G^{\overline{A}_IA_I}_{A_I}F^{(0,2)}.$$
(8.6)
2. In the second degeneration limit which results in a twice punctured annulus with an open string intermediate state, one of the $`\tau _{\{g_I\}}`$ goes to zero keeping the remaining two finite. For instance if in Figure 6, we move the $`𝐚_2`$ cycle near the real axis (i.e. the third boundary), then $`𝐛_2`$ cycle shrinks to zero so that $`\tau _{\{g_3\}}`$ goes to zero, but $`\tau _{\{g_4\}}`$ and $`\tau _{\{g_5\}}`$ remain finite. $`\tau _{\{g_4\}}`$ becomes the usual untwisted modulus $`\tau `$ associated with the resulting annulus while $`\tau _{\{g_5\}}`$ is a function of $`\tau `$ and the separation between the two punctures. In our case, however, this degeneration limit does not contribute so long as every pair of brane stacks is separated (by suitable Wilson line) in the plane in which they are parallel to each other. This gives masses to intermediate open strings that are stretched between different pairs of stacks and hence this degeneration limit is exponentially suppressed.
3. In the third degeneration limit with intermediate closed string states is obtained by shrinking one of the $`𝐚`$ cycles in Fig. 6. For instance if we shrink $`𝐚_2`$ cycle to a point $`P`$ $`\tau _{\{g_3\}}`$ clearly goes to infinity (going as $`i\mathrm{ln}|t|`$ with $`t`$ being the corresponding plumbing fixture coordinate). In this limit $`\omega _{\{g_4\}}`$ becomes untwisted differential on the remaining torus (double cover of the annulus with boundaries $`𝐚_1`$ and $`𝐚_3`$) and hence becomes the usual untwisted modulus of the annulus in the closed string channel. $`\omega _{\{g_5\}}`$ on the other hand remains a twisted differential on the torus with single poles at the points $`P`$ and its image $`P^{}`$. As a result $`\tau _{\{g_5\}}`$ goes to infinity as $`i\mathrm{ln}|t|`$. This degeneration limit does not contribute to the holomorphic anomaly due to the fact that two of the $`\tau _{\{g_I\}}`$ go to infinity simultaneously in this limit. This is because the integrand in (5.28) appears with one momentum each from each plane through $`f`$ defined in (5.30). Explicitely this factor is $`_I(n_1^I+n_2^I\overline{U}_I)`$. The momentum in one of the planes ($`I=3`$) disappears when one takes the derivative with respect to $`\overline{A}_3`$ due to the identity (5.35), however the other two momentum factors for $`I=4,5`$ still exist in the integrand. Since in this closed string degeneration limit $`\tau _{\{g_3\}}`$ as well as another one of the $`\tau _{\{g_I\}}`$ for $`I=4`$ or $`I=5`$ go to infinity, we conclude that this limit is exponentially suppressed.
Thus the holomorphic anomaly is entirely given by the first degeneration limit (8.6). $`F_{\overline{A}_3;\overline{A}_I}^{(0,2)}`$ appearing in this equation comes from the effective action term $`\mathrm{\Pi }\mathrm{Tr}W^2`$ and we will show in the following that it is given by:
$$F_{\overline{A}_I;\overline{A}_J}^{(0,2)}=\frac{f_2^{IJ}}{\mathrm{Im}U_I\mathrm{Im}U_J}H(U_K,A_K);IJKI,$$
(8.7)
where $`f_2^{IJ}`$ are numerical constants.
The function $`F^{(0,2)}`$ appearing in Eq.(8.6) is the one loop gauge kinetic function providing the threshold corrections to gauge couplings $`\mathrm{Tr}W^2`$ . Its Wilson line dependence receives contributions only from $`N=2`$ supersymmetric sectors and reads:
$$_{A_I}F^{(0,2)}=b_I_{A_I}\mathrm{\Delta }(U_I,A_I);\mathrm{\Delta }(U,A)=\mathrm{ln}|\theta _1(A)|^22\pi \frac{(\mathrm{Im}A)^2}{\mathrm{Im}U},$$
(8.8)
where $`b_I`$ are numerical coefficients related to the $`N=2`$ beta-functions and $`_A\mathrm{\Delta }=H`$. Using (8.7) and the Wilson line metric $`G_{A_I\overline{A}_I}=1/\mathrm{Im}U_I`$, one can then identify (8.6) with (8.5), implying $`\pi f_3=_{I=4,5}f_2^{3I}b_I`$. Note that for our choice of angles (8.1)-(8.2), placing at the two boundaries of the annulus two different brane stacks one obtains an $`N=2`$ sector associated to the plane where the two stacks are parallel.
We now compute $`F_{\overline{A}_I;\overline{A}_J}^{(0,2)}`$ from the four-point annulus amplitude involving two gauge fields and two anti-chiral fermions $`\overline{\chi }_I`$, $`\overline{\chi }_J`$, that belong to the corresponding Wilson line supermultiplets. The gauge boson vertices are given in Eq.(3.3), while the $`\overline{\chi }_I`$ fermion vertex at zero momentum, in the $`1/2`$ ghost picture, is given by:
$$V_{\overline{\chi }_I^\alpha }^{(1/2)}=:e^{\phi /2}S_\alpha S_I:;S_I=e^{\frac{i}{2}(\varphi _I\varphi _J\varphi _K)}IJKI,$$
(8.9)
where we use the notation of section 3. Choosing as fermion vertices, say $`I=3`$ and $`J=4`$, one should also insert a picture changing operator $`e^\phi T_F`$, from which only the term $`e^\phi e^{i\varphi _5}\overline{X}^5`$, from the supercurrent (3.8), contributes. Evaluating the correlator at the quadratic order in external momenta and performing the spin-structure sum is straightforward. All non-zero mode determinants cancel and one is left over with a summation over the lattice momenta of the torus $`I=5`$, associated to the $`N=2`$ sector defined by the first two brane stacks. The result is:
$$F_{\overline{A}_I;\overline{A}_J}^{(0,2)}=\frac{f_2^{IJ}}{\mathrm{Im}U_I\mathrm{Im}U_J}_0^{\mathrm{}}𝑑lf(U_K,\stackrel{}{a}_K;l),$$
(8.10)
where the function $`f`$ is defined in Eq.(5.30). The integral can be performed as in (8.4), leading to (8.7). These results prove our initial ansatz (8.3). In fact, the same expression can be obtained by integrating the holomorphic anomaly equation (8.6) by using the explicit forms for $`F_{\overline{A}_I;\overline{A}_J}^{(0,2)}`$ and $`F^{(0,2)}`$ and $`SL(2,)`$ symmetry.
Now we turn to the holomorphic anomaly of $`F_{\overline{A}_I;\overline{A}_J}^{(0,2)}`$. Taking an anti-holomorphic derivative of (8.7) with respect to $`\overline{A}_K`$, one finds:
$$_{\overline{A}_K}F_{\overline{A}_I;\overline{A}_J}^{(0,2)}=\frac{\pi f_2^{IJ}}{\mathrm{Im}U_I\mathrm{Im}U_J\mathrm{Im}U_K}.$$
(8.11)
On the other hand, using the integral form (8.10) and the identities (5.35), one obtains a contribution only from the origin of the lattice at the boundary $`l=\mathrm{}`$. Since this corresponds to the infrared limit in the closed string channel, one concludes that in this case it is the closed string degeneration limit that contributes and the holomorphic anomaly equation becomes:
$$_{[\overline{A}_K}F_{\overline{A}_I;\overline{A}_J]}^{(0,2)}=F_{\overline{A}_K\overline{A}_I;\overline{A}_J\overline{\mathrm{\Phi }}}^{(0,1)}D_\mathrm{\Phi }F^{(0,1)},$$
(8.12)
where $`\mathrm{\Phi }`$ is a closed string modulus. Actually, $`F^{(0,1)}`$ must be the (tree) gauge kinetic function on the disk, implying that $`\mathrm{\Phi }`$ is the string dilaton which couples on the disk but not on the torus. The first factor in the r.h.s. of (8.12) must then correspond to a $`\mathrm{\Pi }^2`$ term on the disk. It can be computed independently by a four-point amplitude involving two anti-chiral fermions, say $`I`$ and $`J`$, and two scalars, the Wilson line $`\overline{A}_K`$ and the dilaton, at a level quadratic in the external momenta.
Let us close the section by giving a possible form of $`F^{(0,3)}`$ in the case of non-parallel brane stacks. In this case, $`F^{(0,3)}`$ is monodromy invariant under $`A_IA_I+1/p_I`$ and $`A_IA_I+U_I/p_I`$, where $`p_I`$ is defined in Eq.(5.33). Furthermore it vanishes when any of the $`p_IA_I=0,1/2,U_I/2,(U_I+1)/2`$. Unlike the previous case, now there is no pole at the origin because at this point no new massless states emerge. As before, by taking derivatives with respect to $`A_I`$ or $`\overline{A}_I`$ and setting for instance $`A_I=0`$, one finds that the zeroes are of first order. We assume the factorized form:
$$F^{(0,3)}=\underset{I}{}_{A_I}G,$$
(8.13)
where $`G`$ is real modular invariant function. A possible ansatz is
$$G=\underset{I}{}G_I;G_I=(\mathrm{Im}U_I)e^{2\pi \frac{\mathrm{Im}p_I^2A_I^2}{\mathrm{Im}U_I}}[c_o|\theta _1(p_IA_I;U_I)|^2+c_e\underset{e}{}|\theta _e(p_IA_I;U_I)|^2]$$
(8.14)
where subscripts $`e`$ and $`o`$ refer to even and odd spin structures respectively and $`c_o`$ and $`c_e`$ are real $`SL(2,)_I`$ invariant functions of $`U_I`$ which may also depend on the brane angles. We could try to check this ansatz by computing the holomorphic anomaly of $`F^{(0,3)}`$ as before. Taking derivative with respect to $`\overline{A}_3`$ one again gets contribution only from the degeneration limits. This time however, in the first as well as the third degeneration limits all the three twisted $`\tau _{\{g_I\}}`$ go to infinity due to the fact that the resulting torii are twisted in all planes. These limits therefore are exponentially suppressed. The only contribution comes from the second degeneration limit where the resulting twice punctured annulus comes with the insertion of twisted states at the punctures. This is the quantity we considered in section 7. However an explicit computation of this term, although in principle possible, is considerably more difficult and we shall not pursue it here.
## 9 Concluding Remarks
To summarize, in this work we discussed the identification of the topological partition function $`F^{(0,h)}`$ of the two-dimensional $`N=2`$ twisted Calabi-Yau $`\sigma `$-model (in the orbifold limit), on bordered genus-zero world-sheets with $`h`$ boundaries, with the moduli-dependent couplings associated to the F-terms $`(\mathrm{Tr}W^2)^{h1}`$, where $`W`$ is the familiar gauge superfield appearing in the effective four-dimensional $`N=1`$ supersymmetric effective action of the Type I string theory compactified on the same six-dimensional orbifold.
We then studied the holomorphic anomaly equation for the violation of the expected holomorphicity of $`F^{(0,h)}`$. The anomaly equation involves also the correlation functions $`F_n^{(0,h1)}`$ associated to F-terms of the form $`\mathrm{\Pi }^n(\mathrm{Tr}W^2)^{h2}`$, where $`\mathrm{\Pi }`$ denotes a generic chiral projection of a non-holomorphic functions of chiral superfields. This system is similar to the heterotic string case studied in the past and related to it by S-duality. A remaining open problem is the integrability of this equation, which can be shown, as in the heterotic string, only in the absence of the handle degeneration limit. Its non-integrability in the general case suggests that there may be still more “topological” quantities missing, that are required for the closure of the full topological theory on world-sheets with boundaries.
Another interesting open problem is the generalization of the possible physical interpretation in terms of actual string amplitudes of the topological partition function $`F^{(g,h)}`$ for higher genus $`g>0`$. A naive extension of our results, following the methods presented in this work, to world-sheets with handles and boundaries simultaneously, fails because the zero-mode contributions of space-time coordinates do not cancel, and thus, the measure does not become topological.
An important phenomenological application of our results is in theories with supersymmetry broken by the VEVs of auxiliary D-components of gauge vector supermultiplets. Such a breaking appears, in particular, in the presence of non-trivial internal magnetic fields, or equivalently in the T-dual description, of branes intersecting at angles, and can be studied directly at the string level. Then the effective operators $`(\mathrm{Tr}W^2)^2`$ and $`\mathrm{\Pi }\mathrm{Tr}W^2`$ generate fermion masses and the associated topological functions become physically observable in the mass spectrum. The resulting Majorana gaugino masses are given by $`F^{(0,3)}`$ at two loops, while the matter fermion masses are of Dirac type and are given by the topological quantity $`F_{\overline{i};\overline{j}}^{(0,2)}`$, appearing at one loop in the holomorphic anomaly equation of $`F^{(0,3)}`$. They are both of order $`m_0^4/M_S^3`$ for $`m_0<M_S`$, where $`m_0`$ is the supersymmetry breaking scale and $`M_S`$ is the string mass. We presented simple examples of toroidal string compactifications with intersecting branes, where both quantities can be computed explicitly as functions of the closed string geometric moduli and of open string Wilson lines.
The precise implementation of these results in a complete string framework, including the Standard Model and a consistent mechanism of supersymmetry breaking and moduli stabilization, remains of course an open issue. For instance, an important question is the effect of closed string back-reaction on the supersymmetry breaking induced in the open string sector by turning on internal magnetic fields, or equivalently by the presence of D-branes at angles. A possible application is in the recently proposed scenario of split supersymmetry . Our results then provide a natural mechanism to break R-symmetry by string effects and generate gaugino and higgsino masses at the TeV scale, when squark and slepton masses are at high energies, of the order of $`10^{13}`$ GeV, and $`M_S`$ is near the unification scale of $`10^{16}`$ GeV. Alternatively, our results can be used to generate gaugino and non-chiral fermion masses in the context of low string mass scale and large extra dimensions , when supersymmetry is broken at the string scale by appropriate configurations of branes and orientifolds .
## Acknowledgments
We would like to thank M.S.Narasimhan for valuable discussions on Prym differentials. This work was supported in part by the European Commission under the RTN contracts MRTN-CT-2004-503369 and MEXT-CT-2003-509661, in part by the INTAS contract 03-51-6346, and in part by the CNRS PICS # 2530. The research of T.R.T. is supported in part by the U.S. National Science Foundation Grant PHY-0242834. K.S.N. and T.R.T. thank the CERN theory division and I.A. thanks ICTP for hospitality, during multiple visits. T.R.T. is grateful to Dieter Lüst, Stephan Stieberger and to Arnold Sommerfeld Center for Theoretical Physics at Ludwig Maximilians University in Munich for their kind hospitality. Any opinions, findings, and conclusions or recommendations expressed in this material are those of the authors and do not necessarily reflect the views of the National Science Foundation.
## Appendix
## Appendix A Lattice Contribution to the Amplitudes
In this Appendix, we derive the lattice sums involved in the topological partition function $`F^{(0,h)}`$. Since we are considering product of three planes (tori) with fluxes, it is sufficient to focus on one plane. In order to compute the lattice contribution, it is more convenient to go to the T-dual theory where the D9 branes become D6 branes whose world-volumes span lines in each of the 3 internal planes. These lines are parallel to some lattice vectors. They are at angles given by the fluxes in the corresponding plane. Let $`\stackrel{}{v}_i`$ be a primitive lattice vector parallel to the $`i`$-th brane and let $`\stackrel{}{w}_i`$ be such that ($`\stackrel{}{v}_i,\stackrel{}{w}_i)`$, for each $`i`$, spans the two-dimensional lattice describing the torus. This in particular means that $`v_i.w_iv_i^\mu w_i^\nu ϵ_{\mu \nu }=\pm T_2`$ where $`T_2=\sqrt{G}`$ is the area of the basic cell of the lattice (in the following we choose the orientation of $`w_i`$ so that the sign in this equation is plus). We can use a complex coordinate $`Z`$ to describe the plane. Let $`v_i`$ and $`w_i`$ be the complex numbers representing the lattice vectors $`\stackrel{}{v}_i`$ and $`\stackrel{}{w}_i`$, respectively. We are considering a world-sheet of genus zero and $`h`$ boundaries. For this bordered surface, we define the $`𝐚_i`$ cycles, $`i=1,2,\mathrm{},h`$, as the boundaries $`\alpha _i`$ shown in Fig. 1 (of course, these cycles are not independent: they satisfy $`_{i=1}^h𝐚_i=\mathrm{𝟏}`$).<sup>7</sup><sup>7</sup>7These cycles should not be confused with the basis of $`𝐚`$-cycles used in Section 2 and depicted in Fig. 2. Let $`P_i`$ be a point on the $`𝐚_i`$ cycle. Since $`𝐚_i`$ lies on the $`i`$-th brane, $`Z(P_i)`$ takes values
$$Z(P_i)=x_i(P_i)v_i+(m_i+y_i)w_i$$
(A.1)
where $`x_i`$ and $`m_i`$ are arbitrary real and integer numbers, respectively, and $`y_i`$ is a fixed real number such that $`|y_i|1/2`$ which describes the relative transverse position of the $`i`$-th brane.<sup>8</sup><sup>8</sup>8In the D9 brane theory $`y_i`$ is one of the components of the Wilson line which after T-duality represents the transverse position of the brane. The second component of the Wilson line remains a Wilson line on the D6 brane. We are considering the case where not all the branes are parallel. Let us assume that the $`(h1)`$-th and the $`h`$-th branes are not parallel to each other. We can choose a complex coordinate $`Z`$ for the plane such that the intersection of these two branes is at $`Z=0`$ and furthermore we can rotate the coordinate so that the $`h`$-th brane lies on the real axis. With this choice of coordinate, we have set $`v_h`$ a real number and $`m_{h1}=m_h=y_{h1}=y_h=0`$. The remaining integers $`m_i`$ for $`i=1,\mathrm{},h2`$ are arbitrary and result in a ($`h2`$)-dimensional lattice sum. However, one can easily see using the fact that $`(v_i,w_i)`$ span the lattice for each $`i`$, that this change of coordinates results in the shift
$$Z(P_i)=x_i(P_i)v_i+(m_i+\stackrel{~}{y}_i\frac{v_h.v_i}{v_h.v_{h1}}m_{h1})w_i;i=1,\mathrm{},h2$$
(A.2)
where $`\stackrel{~}{y}_i`$ are the effective Wilson lines in the Dirichlet direction (or transverse positions in the T dual picture) and are given by
$$\stackrel{~}{y}_i=y_iy_{h1}\frac{v_h.v_i}{v_h.v_{h1}}y_h\frac{v_{h1}.v_i}{v_{h1}.v_h}.$$
(A.3)
Note that the ratio $`\frac{v_h.v_i}{v_h.v_{h1}}`$ is a rational number. Let $`p`$ be the smallest integer such that $`p`$ times this rational number, for each $`i`$, is integer. Then $`m_{h1}=0,1,\mathrm{},p1`$ give the different conjugacy classes of the lattice sum over $`m_i`$.
As one goes with points $`P_i`$ around the cycles $`𝐚_i`$, the functions $`x_i(P_i)`$ can shift by integers: $`x_ix_i+n_i`$. The integers $`n_i`$ denote the winding of the string on the boundaries. Not all the integers $`n_i`$ are independent, however. This is because $`_{i=1}^h𝐚_i`$ is homotopically trivial, giving rise to the condition
$$\underset{i=1}{\overset{h}{}}n_iv_i=0.$$
(A.4)
Since this is a complex equation with not all the $`v_i`$ parallel, there are only $`h2`$ independent integers, say $`n_i`$ for $`i=1,\mathrm{},h2`$. It is important to note that $`n_i`$ actually span only a sublattice of integers such that there exist $`n_{h1}`$ and $`n_h`$ satisfying Eq.(A.4). Since $`(v_h,w_h)`$ generate the lattice, $`v_i=p_iv_h+q_iw_h`$, $`i=1,\mathrm{},h1`$, for some coprime integers $`(p_i,q_i)`$. Therefore choosing $`n_h=_{i=1}^{h1}n_ip_i`$ the above constraint on $`n_i`$ becomes equivalent to the constraint
$$\underset{i=1}{\overset{h1}{}}n_i(v_h.v_i)\underset{i=1}{\overset{h1}{}}(1g_i^1)n_iv_i=0,$$
(A.5)
where $`g_i=e^{2i\theta _i}`$ with $`\theta _i=arg(v_i)`$ and we have used the fact that, for our choice of coordinate $`Z`$, $`v_h`$ is real.
The boundary conditions on the branes can be conveniently imposed by going to the double cover of $`\mathrm{\Sigma }_{(0,h)}`$ which is a genus $`g=h1`$ Riemann surface, with an anti-analytic $`_2`$ involution that keeps the $`𝐚_i`$ cycles fixed and takes $`𝐛_i`$ cycles to $`𝐛_i^1`$, see Section 2. Then going around $`𝐛_i`$ cycle $`dZg_idZ`$, while around $`𝐚_i`$ cycles $`dZ`$ is invariant. $`_2`$ involution which takes a point $`P`$ to $`P^{}`$ acts on $`Z`$ as follows. Let $`P_0`$ be a base point on the $`𝐚_h`$ cycle lying on the $`h`$-th brane. By our choice of coordinate, $`Z(P_0)`$ is real. We can define $`Z(P)=Z(P_0)+_{P_0}^P𝑑Z`$. $`Z(P)`$ is not single-valued – it depends on the homotopy class $`H(P_0,P)`$ of the path chosen. $`_2`$ involution acts as $`Z(P^{})=\overline{Z}(P)`$ where $`Z(P^{})`$ is defined through a path that is in the homotopy class $`GH(P_0,P^{})`$, with $`G`$ being the $`_2`$ action on the homotopy class. Clearly, $`Z`$ on the $`𝐚_h`$ cycle is real. Now consider a point $`P_i`$ on the $`𝐚_i`$ cycle. Then $`Z(P_i)`$ is given by
$$Z(P_i)Z(P_0)=_{P_0}^{P_i}𝑑Z=_{b_i}𝑑Z+g_i_{P_0}^{P_i}𝑑\overline{Z}=_{b_i}𝑑Z+g_i(\overline{Z}(P_i)Z(P_0)).$$
(A.6)
The general solution of this equation is precisely of the form (A.1).
Now we would like to write the classical solutions for $`Z`$ for a given set of integers $`m_i`$, $`n_i`$ and the transverse positions $`y`$. On genus $`h1`$ surface there are $`h2`$ linearly independent holomorphic twisted (Prym) differentials $`\omega _{j,\{g\}}`$ for $`j=1,\mathrm{},h2`$, that are twisted by $`g_i`$ around $`𝐛_i`$ cycles and untwisted around $`𝐚_i`$; here, the subscript $`\{g\}`$ denotes the collection of twists $`\{g_i\}`$. Let us choose a marking for the surface as shown in Fig. 6.<sup>9</sup><sup>9</sup>9Fig. 6 is just the double cover of Fig. 1 and the corresponding contractible surface obtained by cutting along the lines shown in Fig. 6 is just the double cover of Fig. 3. In the following, for notational simplicity, we will choose the $`𝐚`$-cycles to be the $`h1`$ of the boundaries. This is different from the convention followed in Section 2 but it has the advantage of directly giving the lattice momenta of the boundary states.
One gets the genus $`h1`$ surface by gluing $`𝐚_i`$ and $`𝐚_i^1`$ cycles together. The real axis is the boundary $`𝐚_h`$ which sits on the $`h`$-th brane while the remaining $`𝐚_i`$ cycles sit on the $`i`$-th brane. The $`_2`$ involution takes the upper half-plane to the lower half-plane. Let
$$A_{jk}=_{𝐚_k}\omega _{j,\{g\}},B_{jk}=_{𝐛_k}\omega _{j,\{g\}}.$$
(A.7)
Here the $`𝐛_k`$ cycle is defined as follows. If $`P_k`$ is a point on the $`𝐚_k`$ cycle and $`P_k^{}`$ is its $`_2`$ image on the $`𝐚_k^1`$ cycle, then the $`𝐛_k`$ cycle is the line from $`P_0`$ to $`P_k`$ which is identified with $`P_k^{}`$ (on a different sheet due to $`g_k`$ twist) via gluing and then from $`P_k^{}`$ to $`P_0`$ in this different sheet. Note that $`𝐛_k`$ is not a closed contour when $`g_k`$ is non-trivial. As a result, the integrals of closed twisted differentials will depend on the choice of the base point $`P_0`$. Let $`b_{jk}=_{P_0}^{P_k}\omega _{j,\{g\}}`$ on the path indicated in the Fig. 6, then
$$B_{jk}=b_{jk}g_k\overline{b}_{jk}.$$
(A.8)
In particular, this implies that $`g_k^{\frac{1}{2}}B_{jk}`$ is purely imaginary. As noted above, when $`g_k`$ is non-trivial, $`B_{jk}`$ depends on the choice of the base point $`P_0`$.
Integrating around trivial cycle $`_{i=1}^{h1}(𝐚_i𝐛_i𝐚_i^1𝐛_i^1)`$ gives the condition
$$\underset{i=k}{\overset{h1}{}}(1g_k^1)A_{jk}=0.$$
(A.9)
Since at least one of the $`g_i`$ is not identity (say $`g_{h1}`$), we can always eliminate one of the $`A_{jk}`$’s (say $`A_{j(h1)}`$). We can furthermore normalize $`\omega _{j,\{g\}}`$ so that
$$A_{jk}=g_j^{\frac{1}{2}}\delta _{jk},\mathrm{for}k=1,\mathrm{},h2.$$
(A.10)
Then the above equation implies
$$A_{j(h1)}=g_{h1}^{\frac{1}{2}}\frac{\mathrm{sin}(\theta _j)}{\mathrm{sin}(\theta _{h1})}.$$
(A.11)
We will also need the holomorphic differentials $`\omega _{j\{g^1\}}`$ that are twisted oppositely to $`\omega _{j\{g\}}`$. We will denote the corresponding periods by $`A_{jk}^{}`$ and $`B_{jk}^{}`$ and $`b_{jk}^{}`$. They satisfy all the above equations with $`g_k`$ replaced by $`g_k^1`$.
Given any two closed twisted differentials $`\rho `$ and $`\rho ^{}`$ that are twisted by $`\{g\}`$ and $`\{g^1\}`$ respectively, one can evaluate, by using the standard methods for the marking shown in Fig. 6, the following integral:
$$_\mathrm{\Sigma }\rho \rho ^{}=\underset{k=1}{\overset{h1}{}}[B_kA_k^{}A_kB_k^{}+\underset{\mathrm{}=1}{\overset{k}{}}(1g_kg_{\mathrm{}}^1)A_{\mathrm{}}A_k^{}],$$
(A.12)
where $`A_k=_{𝐚_k}\rho `$, $`A_k^{}=_{𝐚_k}\rho ^{}`$, $`B_k=_{𝐛_k}\rho `$ and $`B_k^{}=_{𝐛_k}\rho ^{}`$.
The $`(1,1)`$ forms $`\omega _{j,\{g\}}\overline{\omega }_{k,\{g\}}`$ are untwisted and therefore can be integrated on the Riemann surface. Using (A.12) and the periods (A.8), (A.10) and (A.11), we find
$$_\mathrm{\Sigma }\omega _{j,\{g\}}\overline{\omega }_{k,\{g\}}=\tau _{jk}\overline{\tau }_{kj},$$
(A.13)
where
$$\tau _{jk}=D_{jk}C_{jk}$$
(A.14)
and $`C`$ is a purely imaginary symmetric $`(h2)\times (h2)`$ matrix which does not depend on the world-sheet moduli. It is given by
$$C_{jk}=i\frac{\mathrm{sin}(\theta _k)\mathrm{sin}(\theta _j\theta _{h1})}{\mathrm{sin}(\theta _{h1})},\mathrm{for}jk.$$
(A.15)
On the other hand, the matrix $`D`$ depends on the world-sheet moduli and is given in terms of the periods $`B_{jk}`$ as follows:
$`D_{jk}`$ $`=`$ $`{\displaystyle \frac{g_k^{\frac{1}{2}}}{2i\mathrm{sin}(\theta _{h1})}}{\displaystyle _{b_kb_{h1}b_k^1b_{h1}^1}}\omega _{j,\{g\}}`$ (A.16)
$`=`$ $`g_k^{\frac{1}{2}}B_{jk}{\displaystyle \frac{\mathrm{sin}(\theta _k)}{\mathrm{sin}(\theta _{h1})}}g_{h1}^{\frac{1}{2}}B_{j(h1)}.`$
It is important to note that while $`B_{jk}`$ depends on the base point $`P_0`$, $`D_{jk}`$ does not. This is due to the fact that the cycle $`𝐛_k𝐛_{h1}𝐛_k^1𝐛_{h1}^1`$ is a closed contour even when $`g_k`$ is non-trivial. Due to Eq.(A.8), $`D`$ is also purely imaginary, however it is in general not symmetric. As a consequence the right hand side of equation (A.13) is purely imaginary and
$$\tau _{jk}\overline{\tau }_{kj}=2(\tau ^S)_{jk},$$
(A.17)
where $`\tau ^S`$ is the symmetric part of $`\tau `$.
Similarly, we can define the corresponding quantities for oppositely twisted differentials and the corresponding periods $`\tau _{jk}^{}`$ satisfy the same equations with $`g_k`$ replaced by $`g_k^1`$. Finally, we have the relation
$$0=_\mathrm{\Sigma }\omega _{j,\{g\}}\omega _{k,\{g^1\}}=\tau _{jk}\tau _{kj}^{}.$$
(A.18)
In the untwisted case, a similar equation implies that the period matrix is symmetric, however in the twisted case it only says that the transposed $`\tau `$ is equal to $`\tau ^{}`$.
Using the twisted differentials, one can write the classical solution for $`Z`$ as
$$dZ=\underset{j=1}{\overset{h2}{}}[L_j\omega _{j,\{g\}}+\stackrel{~}{L}_j\overline{\omega }_{j,\{g^1\}}].$$
(A.19)
Integrals around $`𝐚_k`$ cycles should give the windings $`n_kv_k`$, which implies that $`_{a_k}𝑑Z=n_kv_k`$. Furthermore, $`_{P_0}^{P_k}𝑑Z=(m_k+y_k)w_k`$, due to the boundary conditions (A.1). This implies
$$_{𝐛_k}𝑑Z=(m_k+y_k)(w_kg_k\overline{w}_k)=2i(m_k+y_k)g_k^{\frac{1}{2}}\frac{T_2}{|v_k|},$$
(A.20)
where in the second equality we used the fact that $`w_k.v_k=T_2`$. Here, $`y_k`$ are shifted by $`y_h\frac{w_h.v_i}{T_2}`$, due to the fact that we are choosing our coordinates so that the $`h`$-th brane lies on the real axis. We can then solve for $`L`$ and $`\stackrel{~}{L}`$ as
$`L`$ $`=`$ $`{\displaystyle \frac{N}{2}}+{\displaystyle \frac{1}{2\tau ^S}}[(\tau ^A+C)N+2M]`$ (A.21)
$`\stackrel{~}{L}`$ $`=`$ $`{\displaystyle \frac{N}{2}}{\displaystyle \frac{1}{2\tau ^S}}[(\tau ^A+C)N+2M],`$ (A.22)
where $`\tau ^S`$ and $`\tau ^A`$ are the symmetric and anti-symmetric parts of $`\tau `$, respectively, and $`L`$, $`\stackrel{~}{L}`$, $`N`$ and $`M`$ are $`(h2)`$-dimensional column vectors whose components are respectively $`L_j`$, $`\stackrel{~}{L}_j`$, $`N_j=n_j|v_j|`$ while $`M_j`$ is given by
$$M_j=i\frac{T_2}{|v_j|}[(m_j+y_j)\frac{v_h.v_j}{v_h.v_{h1}}(m_{h1}+y_{h1})].$$
(A.23)
The classical action is
$`S`$ $`=`$ $`i\pi (L^T\tau ^SL+\stackrel{~}{L}^T\tau ^S\stackrel{~}{L})+2\pi i{\displaystyle \underset{j=1}{\overset{h2}{}}}n_jW_j`$ (A.24)
$`=`$ $`i\pi [N^T\tau ^SN+\{2M+(\tau ^A+C)N\}^T(\tau ^S)^1\{2M+(\tau ^A+C)N\}]`$
$`+2\pi i{\displaystyle \underset{j=1}{\overset{h2}{}}}n_jW_j,`$
where $`W_j`$ is the Wilson line on the world-volume of the $`j`$-th brane.
Poisson resummation over $`m_j`$ for $`j=1,\mathrm{},h2`$ gives rise to the following instanton contribution
$`Z_{\mathrm{inst}}`$ $`=`$ $`[{\displaystyle \underset{j=1}{\overset{h2}{}}}{\displaystyle \frac{|v_j|}{T_2}}](\mathrm{det}\tau ^S)^{\frac{1}{2}}{\displaystyle \underset{n_j,k_j,m_{h1}}{}}e^{i\pi (N^T\tau ^SN+\frac{1}{4}K^T\tau ^SK)\pi K^T\tau ^AN}`$ (A.25)
$`\times e^{2\pi i_{j=1}^{h2}(k_j\stackrel{~}{y}_j+n_jW_j)}e^{2\pi i_{j=1}^{h2}k_j\frac{v_h.v_j}{v_h.v_{h1}}m_{h1}}e^{\pi K^TCN},`$
where $`K`$ is a column vector with components $`K_j=\frac{|v_j|}{T_2}k_j`$ and $`\stackrel{~}{y}_j`$ is the effective Wilson line (A.3) in the Dirichlet direction. The sum over $`n_j`$ ($`j=1,\mathrm{},h2`$) is over integers that satisfy the constraint (A.4) while the sum over $`k_j`$ ($`j=1,\mathrm{},h2`$) is over all integers. Finally, the sum over $`m_{h1}`$ ranges from $`0`$ to $`p1`$, where $`p`$ is the smallest integer such that $`p\frac{v_h.v_j}{v_h.v_{h1}}`$ is integer for all $`j`$. The latter sum reduces the $`k_j`$ sums to a sublattice satisfying the constraint
$$\underset{j=1}{\overset{h2}{}}k_j\frac{v_h.v_j}{v_h.v_{h1}}=\mathrm{integer}.$$
(A.26)
This constraint is exactly as the one for the $`n_j`$ and from the above discussion it follows that this implies that there exist integers $`k_a`$, $`a=1,\mathrm{},h`$, such that $`_ak_av_a=0`$. In fact, $`K_a=k_av_a/T_2`$ is a vector in the dual lattice and describes the momentum in the Dirichlet direction for the $`a`$-th brane (i.e. in the direction perpendicular to $`v_a`$).
The phase $`e^{\pi K^TCN}`$ is independent of the world-sheet moduli and of the target space moduli. One can show using the constraints (A.4) and the similar one for $`k_j`$, that $`iK^TCN`$ is an integer implying that this phase is $`\pm 1`$.
We can also include a $`B`$ field on the plane. The contribution of the $`B`$ field to the instanton action is
$$S_B=\frac{2\pi B}{T_2}N^T(M+\frac{1}{2}CN).$$
(A.27)
Including this contribution (after the Poisson resummation) in Eq.(A.25), one finally obtains<sup>10</sup><sup>10</sup>10Here and in Section 5 \[from Eq.(5.18) onwards\] we do not keep track of overall factors that are completely moduli- and flux data-independent.
$`Z_{\mathrm{inst}}`$ $`=`$ $`p[{\displaystyle \underset{j=1}{\overset{h2}{}}}{\displaystyle \frac{|v_j|}{T_2}}](\mathrm{det}\stackrel{~}{\tau }^S)^{\frac{1}{2}}{\displaystyle \underset{n_j,k_j}{}}e^{\frac{i\pi }{T_2}(n^T\overline{T}+\frac{1}{2}k^T)\stackrel{~}{\tau }(nT+\frac{1}{2}k)}`$ (A.28)
$`\times e^{2\pi i(k_j\stackrel{~}{y}_j+n_jW_j)}e^{\pi K^TCN},`$
where $`n`$ and $`k`$ are $`(h2)`$-dimensional column vectors with integer entries $`n_j`$ and $`k_j`$ satisfying the constraints (A.4) and (A.26), respectively, and
$$\stackrel{~}{\tau }_{jk}\frac{|v_j||v_k|}{T_2}\tau _{jk}.$$
(A.29)
It is worth noting that $`\tau _{jk}`$ depend only on the world-sheet moduli and the twist angles. In turn, these angles depend on the target space modulus $`U`$, but not on $`T`$. Since $`\frac{|v_j||v_k|}{T_2}`$ does not depend on $`T`$ as well, it follows that $`\stackrel{~}{\tau }`$ also only depends on $`U`$, but not on $`T`$. Thus the $`T`$-dependence of the partition function is explicit in (A.28). The $`U`$ dependence, however, is implicit and appears through $`\stackrel{~}{\tau }`$.
Let us now consider the degeneration limit of $`\tau `$ when $`𝐚_1`$ cycle vanishes. In Fig. 6, this amounts to shrinking $`𝐚_1`$ and $`𝐚_1^1`$ cycles to points $`z_1`$ and its image, respectively. In this limit, the usual period matrix $`t`$ given by the untwisted differentials degenerates as $`t_{11}i\mathrm{}`$. Since the twisted $`\omega _1`$ will degenerate to a twisted differential with two simple poles at $`z_1`$ and its image, with residues $`g_1^{\frac{1}{2}}`$ and $`g_1^{\frac{1}{2}}`$, respectively. Then Eqs. (A.14), (A.15), (A.16) and (A.8) imply that $`\tau _{11}t_{11}`$ modulo finite terms. Thus the $`t_{11}`$ dependence of the partition function is
$$q_1^{\frac{1}{2}[n_1\stackrel{}{v}_1+(k_1+Bn_1)\frac{\stackrel{}{v}_1}{T_2}]^2},q_1=e^{2i\pi t_{11}}.$$
(A.30)
This agrees with the fact that the boundary state describing the brane parallel to the vector $`\stackrel{}{v}_1`$ involves precisely the lattice states $`[n_1\stackrel{}{v}_1+(k_1+Bn_1)\frac{\stackrel{}{v}_1}{T_2}]`$. As a further check, one might ask what happens when we take the limit of shrinking the $`𝐚_{h1}`$ cycle to a point $`z_{h1}`$. Since in our treatment the $`(h1)`$-th boundary played a special role, it is not immediately obvious that the corresponding degeneration would have a correct interpretation in terms of the boundary state. In this limit all the basis elements of twisted differentials $`\omega _{j,\{g\}}`$ develop simple poles at $`z_{h1}`$ and its image, with residues $`\frac{\mathrm{sin}(\theta _j)}{\mathrm{sin}(\theta _{h1})}g_{h1}^{\pm \frac{1}{2}}`$, respectively. Then Eqs. (A.14), (A.15), (A.16) and (A.8) imply that the leading behavior of the twisted period matrix is given by
$$\tau _{jk}\frac{\mathrm{sin}(\theta _j)\mathrm{sin}(\theta _k)}{\mathrm{sin}(\theta _{h1})^2}t_{(h1)(h1)}$$
(A.31)
and the $`q_{h1}t_{(h1)(h1)}`$ dependence of the partition function is
$`q_{h1}^{\frac{1}{2}_{j,k=1}^{h2}\frac{\mathrm{sin}(\theta _j)\mathrm{sin}(\theta _k)}{\mathrm{sin}(\theta _{h1})^2}\frac{|v_j||v_k|}{T_2^2}(n_jT+k_j)(n_k\overline{T}+k_k)}`$ (A.32)
$`=q_{h1}^{\frac{1}{2}[n_{h1}\stackrel{}{v}_{h1}+(k_{h1}+Bn_{h1})\frac{\stackrel{}{v}_{h1}}{T_2}]^2},`$ (A.33)
where we used the constraints (A.4) and (A.26) for $`n_j`$ and $`k_j`$. This again agrees with the boundary state of the $`(h1)`$-th brane.
Finally, let us consider the case when all the $`h1`$ $`𝐚_j`$-cycles are shrunk to points $`z_j`$ for $`j=1,\mathrm{},h1`$. Then the partition function must describe a $`(h1)`$-point function on the disk whose boundary sits on the $`h`$-th brane. The vertices at the $`h1`$ points are boundary states on the $`j`$-th brane, namely
$$V(z_j,n_j,k_j)=e^{in_j\stackrel{}{v}_j.(\stackrel{}{X}_L\stackrel{}{X}_R)+\frac{i}{T_2}(k_j+Bn_j)\stackrel{}{v}_j.(\stackrel{}{X}_L+\stackrel{}{X}_R)}(z_j)$$
(A.34)
for integers $`n_j`$ and $`k_j`$ satisfying the constraints (A.4), (A.26). To see this, we note that, in this limit, $`\omega _i`$ become
$$\omega _i(z)=\frac{1}{2\pi i}[\frac{e^{i\theta _i}}{zz_i}\frac{e^{i\theta _i}}{z\overline{z}_i}\frac{\mathrm{sin}(\theta _i)}{\mathrm{sin}(\theta _{h1})}(\frac{e^{i\theta _{h1}}}{zz_{h1}}\frac{e^{i\theta _{h1}}}{z\overline{z}_{h1}})].$$
(A.35)
Using this we can compute $`\tau `$ and, after some algebra and taking care of the logarithmic branch cuts, we indeed find that the result is
$$\underset{n_j,k_j}{}\underset{j=1}{\overset{h1}{}}q_j^{\mathrm{\Delta }(n_j,k_j)}V(z_j,n_j,k_j).$$
(A.36)
So far, we have discussed the contribution of the classical solutions to the partition function. In the topological amplitude discussed in this paper we actually need the correlation function
$$\underset{j=1}{\overset{h2}{}}[\overline{\psi }Z(w_j)+\overline{\psi }Z(w_j^{})]_h,$$
(A.37)
where $`w_j^{}`$ is the image of $`w_j`$. In the topological theory $`\overline{\psi }`$ are dimension one fields and have $`h2`$ zero modes and therefore are replaced by the $`h2`$ twisted differentials $`\omega _i`$. Similarly plugging in the classical solution for $`Z`$ (A.19), (A.22) and doing the Poisson resummation over the integers $`m_j`$, we find that the world-sheet supercurrents in the correlation function are replaced by $`det(\mathrm{\Omega }_i(w_j))`$, where
$$\mathrm{\Omega }_i=\underset{j=1}{\overset{h2}{}}\frac{1}{T_2}(n_jT+k_j)v_j(\omega _{j,g}\omega _{i,g^1}+\mathrm{c}.\mathrm{c}.).$$
(A.38)
The above discussion has been limited to the case of magnetized D9 branes on $`T^6`$, and as a result the boundaries ($`𝐚`$-cycles) are untwisted. This allowed us to normalize the Prym differentials along the $`𝐚_i`$-cycles for $`i=1,\mathrm{},h2`$. One can generalize the above to the case of orbifolds or fractional branes where some (or all) of the a-cycles are also twisted. The twisted $`𝐚_i`$ cycles are now not closed, however, we can choose instead the closed cycles $`𝐚_i𝐛_i𝐚_i^1𝐛_i^1`$ to normalize the Prym differentials. This analysis has been carried out in Ref. in the context of closed strings. It is straightforward, to generalize the method of this reference to the case of open strings. |
warning/0507/cond-mat0507419.html | ar5iv | text | # Patterns and Collective Behavior in Granular Media: Theoretical Concepts
## I Introduction
### I.1 Preliminary remarks
Granular materials are ubiquitous in our daily lives and basic to many industries. Yet understanding their dynamic behavior remains a major challenge in physics, see for review Duran (1999); Jaeger *et al.* (1996); Kadanoff (1999); de Gennes (1999); Gollub and Langer (1999); Nedderman (1992); Ristow (1999); Rajchenbach (2000); Ottino and Khakhar (2000). Granular materials are collections of discrete macroscopic solid grains with sizes large enough that Brownian motion is irrelevant (energy of 1 mm grain moving with typical velocity of 1 cm/sec exceeds the thermal energy at least by 10 orders of magnitude). Since thermodynamic fluctuations do not play a role, for granular systems to remain active they have to gain energy either from shear or vibration and are thus far from equilibrium. External volume forces (gravity, electric and magnetic fields) and flows of interstitial fluids such as water or air may also be used to activate the grains. When subjected to a large enough driving force, a granular system may exhibit a transition from a granular solid to a liquid and various ordered patterns of grains may develop. Understanding fundamentals of granular materials draws upon and gives insights into many fields at the frontier of modern physics: plasticity of solids, fracture and friction; complex systems from equilibrium such as colloids, foams, suspensions, and biological self-assembled systems. Moreover, particulate flows are central to a large number of industries including the chemical, pharmaceutical, food, metallurgical, agricultural and construction industries. Beyond these industrial applications, particle laden-flows are widespread in nature, for example dune migration, erosion/deposition processes, landslides, underwater gravity currents and coastal geomorphology, etc.
From a theoretical point of view, it is sometimes useful to employ an analogy between granular matter and ordinary condensed matter and to regard the grains as the equivalent of (classical) atoms. However, this analogy is far from complete, since the dissipative nature of grain interactions is the source of many differences between the two “kinds” of matter. In particular, dissipation is responsible for the fact that most states of granular matter are metastable. The typical macroscopic size of the grains renders thermal fluctuations negligible and most standard thermodynamic concepts inapplicable. Whereas the behavior of dilute granular systems (rapid granular gases) can often be explained using the framework of kinetic theory (see e.g. Brilliantov and Pöschel (2004)), the quantitative theory of dense granular assemblies is far less developed.
In recent years several comprehensive reviews and monographs have appeared on the subject of granular physics, see Jaeger *et al.* (1996); Duran (1999); Ristow (2001); Aradian *et al.* (2002); Rajchenbach (2000); Brilliantov and Pöschel (2004); Kudrolli (2004); Ottino and Khakhar (2000). Yet in most of them the focus has been on actual phenomena and experiments rather than on theoretical concepts and approaches to the problems of granular physics. Furthermore, the scope of granular physics has become so broad that we chose to limit ourselves with reviewing the recent progress in a subfield of granular pattern formation leaving out many interesting and actively developing subjects. We loosely define pattern formation as a dynamical process leading to the spontaneous emergence of nontrivial spatially non-uniform structure which is weakly dependent on initial and boundary conditions. According to our working definition, we include in the scope of the review the patterns in thin layers of vibrated grains (Sec. IV,V), patterns in gravity-driven flows (Sec. VI), granular stratification and banding (Sec. VII), as well as a multitude of patterns found in granular assemblies with complex interactions (Sec. VIII). Before delving into details of theoretical modelling of these pattern-forming systems, we present a brief overview of the relevant experimental findings and main theoretical concepts (Sec. II and III).
### I.2 Fundamental microscopic interactions
Probably the most fundamental microscopic property of granular materials is irreversible energy dissipation in the course of interaction (collision) between the particles. For the case of so-called dry granular materials, i.e. when the interaction with interstitial fluid such as air or water is negligible, the encounter between grains results in dissipation of energy while total mechanical momentum is conserved. In contrast to the interaction of particles in molecular gases, the collisions of macroscopic grains is generally inelastic. There are several well-accepted models addressing the specifics of energy dissipation in the course of collision, see for details e.g. Brilliantov and Pöschel (2004). The simplest case corresponds to in-deformable (hard) frictionless particles with fixed restitution coefficient $`0<e<1`$ characterizing the fraction of energy lost in the course of collision. The relation between the velocities after the collisions ($`𝐯_{1,2}^{}`$) and before the collision ($`𝐯_{1,2}`$) for two identical spherical particles is given by
$$𝐯_{1,2}^{}=𝐯_{1,2}\frac{1+e}{2}[𝐧_{12}(𝐯_1𝐯_2)]𝐧_{12}.$$
(1)
Here $`𝐧_{12}`$ is the unit vector pointed from the center of particle 1 to the center of particle 2 at the moment of collision. The case of $`e=1`$ corresponds to the elastic collisions (particles exchange their velocities) and $`e=0`$ characterizes fully inelastic collisions. For $`0<e<1`$ the total energy loss is of the form
$$\mathrm{\Delta }E=\frac{1e^2}{4}|𝐧_{12}(𝐯_1𝐯_2)|^2.$$
Modelling collisions between particles by a fixed restitution coefficient is very simple and intuitive, however this approximation can be questionable in certain cases. For example, approximation of granular media by a gas of hard particles with fixed $`e`$ often yields non-physical behavior such as inelastic collapse McNamara and Young (1996): divergence of the number of collisions in a finite time, see Subsec. IV.1. In fact, the restitution coefficient is known to depend on the relative velocities of colliding particles and approaches unity as $`|𝐯_1𝐯_2|0`$. This dependence is captured by the visco-elastic modelling of particle collision (see e.g. Ramirez *et al.* (1999)). For non-spherical grains the restitution coefficient may also depend on the point of contact Goldsmith (1964).
Tangential friction forces play an important role in the dynamics of granular matter, especially in dense systems. Friction forces are hysteretic and history dependent (the contact between two grains can be either stuck due to dry friction or sliding depending on the history of interaction). This strongly nonlinear behavior makes the analysis of frictional granular materials extremely difficult. In the majority of theoretical studies, the simplest Coulomb law is adopted: friction is independent on sliding velocity as long as tangential force exceeds the certain threshold Walton (1993). However, the main problem is represented by calculation of the static friction forces. It is well known that frictional contact forces among solid particles exhibit indeterminacy in case of multiple contacts per particle because there are less force balance constraints than stress components (see, e.g. McNamara *et al.* (2004); Unger *et al.* (2005)). To resolve this indeterminacy in simulations, various approximate algorithms have been proposed. In soft particle molecular dynamics simulations the most widely used approach to calculating friction forces is the spring-dashpot model Cundall and Strack (1979); Schäffer *et al.* (1996). Another approach is taken in the contact dynamics method. By assuming that all particles are rigid and treating all contacting particles as performing instantaneous collisions (even those which are in fact in persistent contact), one can compute the contact forces generated during these collisions based on local force balance and impenetrability of the particles constraints (see Moreau (1994); Brendel *et al.* (2004)).
Viscous drag forces due to interaction with interstitial fluid often affect the dynamics of granular materials. Gas-driven particulate flows is an active research area in the engineering community, see e.g. Jackson (2000). Fluid-particle interactions are also involved in many geophysical processes, e.g. dune formation Bagnold (1954). Whereas interaction of small individual particles with the fluid is well-understood in terms of Stokes law, collective interaction and mechanical momentum transfer from particles to fluid remains an open problem. Various phenomenological constitutive equations are used in the engineering community to model fluid-particulate flows, see e.g. Duru *et al.* (2002).
Finally, small particles can acquire electric charge of magnetic moment. In this situation fascinating collective behavior emerge due to competition between short-range collisions and long-range electromagnetic forces, see e.g. Aranson *et al.* (2000); Blair *et al.* (2003a); Sapozhnikov *et al.* (2003). Effects of complex inter-particle interactions on pattern formation in granular systems will be discussed in Sec. VIII.
## II Overview of dynamic behavior in granular matter
In this Section we give a brief overview of the main experiments illustrating the dynamical behavior of granular media and the phenomena to be discussed in greater depth in the following Sections. We classify the experiments according to the way energy is injected into the system: vibration, gravity, or shear.
### II.1 Pattern formation in vibrated layers
Quasi-two-dimensional sub-monolayers of grains subjected to vertical vibration exhibit a surprizing bimodal regime characterized by a dense cluster of closely packed almost immobile grains surrounded by gas of agitated particles, Olafsen and Urbach (1998), Fig. 1. This clustering transition occurs when the magnitude of vibration is reduced (the system is “cooled down”) which is reminiscent to the clustering instability observed in non-driven (freely cooling) gas of inelastic particles discovered by Goldhirsch and Zanetti (1993), Fig. 2. Detailed consideration of clustering phenomena in sub-monolayer systems is given in Sec. IV.
Multilayers of granular materials subject to vertical vibration exhibit spectacular pattern formation. In a typical experimental realization a layer of granular material about 10-30 particle diameters thick is energized by precise vertical vibration produced by an electromagnetic shaker. Depending on experimental conditions, plethora of patterns can be observed, from stripes and squares to hexagons and interfaces, see Fig. 3. While the first observations of patterns in vibrated layers were made more than two centuries ago by Chladni (1787) and Faraday (1831), the current interest in these problems was initiated by Douady *et al.* (1989); Fauve *et al.* (1989) and culminated in the discovery by Umbanhowar *et al.* (1996) of a remarkable localized object, oscillon, Fig. 4. Detailed consideration of these observations and their modelling efforts is given in Sec. V.
In another set of experiments pattern formation was studied in a horizontally vibrated system, see e.g. Ristow (1997); Tennakoon *et al.* (1998); Liffman *et al.* (1997). While there are certain common features, such as sub-harmonic regimes and instabilities, horizontally vibrated systems do not show richness of behavior typical for the vertically vibrated systems, and nontrivial flow regimes are typically localized near the walls. When the granular matter is polydisperse, vertical or horizontal shaking often leads to segregation. The most well-known manifestation of this segregation is the so-called “Brazil nut” effect when large particles float to the surface of a granular layer under vertical shaking Rosato *et al.* (1987). Horizontal shaking is also known to produce interesting segregation band patterns oriented orthogonally to the direction of shaking Mullin (2000, 2002) (see Fig. 5).
### II.2 Gravity-driven granular flows
Gravity-driven systems such as chute flows and sandpiles often exhibit nontrivial patterns and spatio-temporal structures. Possibly the most spectacular are avalanches observed in the layers of granular matter if the inclination exceeds the critical angle (static angle of repose). Avalanches were a subject of continued research for many decades, however only recently it was established that the avalanche shape depends sensitively on the thickness of the layer and the inclination angle: triangular downhill avalanches in thin layers and balloon-shaped avalanches in thicker layers which expand both uphill and downhill, see Fig. 6 and Daerr and Douady (1999); Daerr (2001). Gravity-driven granular flows are prone to a variety of non-trivial secondary instabilities in granular chute flow: fingering Pouliquen *et al.* (1997), see Fig. 7, longitudinal vortices in rapid chute flows Forterre and Pouliquen (2001); Börzsönyi and Ecke (2005), see Fig. 8, long modulation waves Forterre and Pouliquen (2003), and others.
Rich variety of patterns and instabilities has also been found in underwater flows of granular matter: transverse instability of an avalanche fronts, fingering, pattern formation in the sediment behind the avalanche, etc. (see Daerr *et al.* (2003); Malloggi *et al.* (2005a, b)). Whereas certain pattern forming mechanisms are specific to the water-granulate interaction, one also finds striking similarities with the behavior of “dry” granular matter.
### II.3 Flows in rotating cylinders
Energy is often supplied into a granular system through the shear which is driven by the moving walls of the container. One of the most commonly used geometries for this class of systems is a horizontal cylinder rotated around its axis, or rotating drum. Rotating drums partly filled with granular matter are often used in chemical engineering for mixing and separation of particles. Flows in rotating drums recently became a subject of active research in the physics community. For not too high rotating rates the flow regime in the drum is separated into an almost solid-body rotation in the bulk of the drum and a localized fluidized layer near the free surface (Fig. 9). Slowly rotating drums exhibit oscillations related to the gradual increase of free surface angle to the static angle of repose and subsequent fast relaxation to a lower dynamic repose angle via an avalanche. Transition to steady flow is observed for the higher rotation rate Rajchenbach (1990). Scaling of various flow parameters with the rotation speed (e.g. the width of the fluidized layer etc) and development of correlations in “dry” and “wet” granular matter was recently studied by Tegzes *et al.* (2002, 2003).
Rotating drums are typically used to study size segregation in binary mixtures of granular materials. Two types of size segregation can be distinguished: radial and axial. Radial segregation is a relatively fast process and occurs after a few revolutions of the drum. As a result of radial segregation larger particles are expelled to the periphery and a core of smaller particles is formed in the bulk Metcalfe *et al.* (1995); Khakhar *et al.* (1997); Metcalfe and Shattuck (1998); Ottino and Khakhar (2000), see Fig. 10.
Axial segregation, occurring in the long drums, happens on a much longer time scale (hundreds of revolutions). As a result of axial segregation, bands of segregated materials are formed along the drum axis Zik *et al.* (1994); Hill and Kakalios (1994, 1995), see for illustration Fig. 11. The segregated bands exhibit slow coarsening behavior. Even more surprisingly, under certain conditions axial segregation patterns show oscillatory behavior and travelling waves Choo *et al.* (1997); Fiodor and Ottino (2003). Possible mechanisms leading to axial segregation are discussed in Sec. VII.
### II.4 Grains with complex interactions
Novel collective behaviors emerge when the interactions between the grains have additional features caused by shape anisotropy, interstitial fluid, magnetization or electrical charge, etc. In this situation short-range collisions, the hallmark of “traditional” granular systems, can be augmented by long-range forces.
Remarkable patterns including multiple rotating vortices of nearly vertical rods are observed in the system of vibrated rods by Blair *et al.* (2003a), see Fig. 12. The rods jump on their ends slightly tilted and drift in the direction of the tilt.
Mechanically Blair and Kudrolli (2003b) or electrostatically Snezhko *et al.* (2005) driven magnetic grains exhibit formation of long chains, isolated rings or interconnecting networks, see Fig. 13. In this situation magnetic dipole-dipole interaction augments hard-core collisions.
Ordered clusters and nontrivial dynamic states were observed by Voth *et al.* (2002); Thomas and Gollub (2004) in a small system of particles vibrated in liquid (Fig. 14). It was shown that fluid-mediated interaction between particles in a vibrating cavity leads to both long-range attraction and short-range repulsion. A plethora of nontrivial patterns including rotating vortices, pulsating rings, chains, hexagons etc was observed by Sapozhnikov *et al.* (2003a) in the system of conducting particles in dc electric field immersed in poor electrolyte (Fig. 15). The nontrivial competition between electrostatic forces and self-induced electro-hydrodynamic flows determines the structure of emerging pattern.
Granular systems with complex interactions serve as a natural bridge to seemingly different systems such as foams, dense colloids, dusty plasmas, ferrofluids and many others.
## III Main theoretical concepts
Physics of granular media is a diverse and eclectic field incorporating many different concepts and ideas, from hydrodynamics to the theory of glasses. Consequently, many different theoretical approaches have been proposed to address observed phenomena.
### III.1 Kinetic theory and hydrodynamics
Kinetic theory deals with the equations for the probability distributions functions describing the state of granular gas. The corresponding equations, similar to Boltzmann equations for rarefied gases, can be rigorously derived for the dilute gas of inelastically colliding particles with fixed restitution coefficient, although certain generalizations are known, Goldstein and Shapiro (1995); Jenkins and Zhang (2002). Kinetic theory is formulated in terms of the Boltzmann-Enskog equation for the probability distribution function $`f(𝐯,𝐫,t)`$ to find the particles with the velocity $`𝐯`$ at point $`𝐫`$ at time $`t`$. In the simplest case of identical frictionless spherical particles of radius $`d`$ with fixed restitution coefficient $`e`$ it assumes the following form
$$(_t+(𝐯_1))f((𝐯_1,𝐫_1,t)=I[f]$$
(2)
with the binary collision integral $`I[f]`$ in the form
$`I`$ $`=`$ $`d^2{\displaystyle }d𝐯_2{\displaystyle }d𝐧_{12}\mathrm{\Theta }(𝐯_{12}𝐧_{12})|𝐯_{12}𝐧_{12}|\times `$ (3)
$`[\chi f(𝐯_{\mathrm{𝟏}}^{}{}_{}{}^{\prime \prime },𝐫_1,t)f(𝐯_{\mathrm{𝟐}}^{}{}_{}{}^{\prime \prime },𝐫_1d𝐧{}_{1}{}^{}2,t)`$
$`f(𝐯_1,𝐫_1,t)f(𝐯_2,𝐫_1+d𝐧_{12},t)]`$
where $`\chi =1/e^2`$, $`\mathrm{\Theta }`$ is theta-function, and pre-collision velocities $`v_{1,2}`$ and “inverse collision” velocities $`v_{1,2}^{\prime \prime }`$ are related as follows
$$𝐯_{1,2}^{\prime \prime }=𝐯_{1,2}\frac{1+e}{2e}[𝐧_{12}(𝐯_1𝐯_2)]𝐧_{12}$$
(4)
(cf. Eq. (1)). This equation is derived with the usual “molecular chaos” approximation which implies that all correlations between colliding particles are neglected. One should keep in mind, however, that in dense granular systems this approximation can be rather poor due to excluded volume effects and inelasticity of collisions introducing velocity correlations among particles (see, for example, Brilliantov and Pöschel (2004)).
Hydrodynamic equations are obtained by truncating the hierarchy of moment equations obtained from the Boltzmann equation (2) via an appropriately modified Chapman-Enskog procedure (see, e.g., Jenkins and Richman (1985); Brey *et al.* (1998); Garzó and Dufty (1999)). As a result, a set of continuity equations for mass, momentum and fluctuation kinetic energy (or “granular temperature”) is obtained. However, in contrast to conventional hydrodynamics, the applicability of granular hydrodynamics is often questionable because typically there is no separation of scale between microscopic and macroscopic motions <sup>1</sup><sup>1</sup>1Except the case of almost elastic particles with the restitution coefficient $`e1`$, see e.g. Tan and Goldhirsch (1998).
The mass, momentum and energy conservation equations in granular hydrodynamics have the form
$`{\displaystyle \frac{D\nu }{Dt}}`$ $`=`$ $`\nu 𝐮,`$ (5)
$`\nu {\displaystyle \frac{D𝐮}{Dt}}`$ $`=`$ $`𝝈+\nu 𝐠,`$ (6)
$`\nu {\displaystyle \frac{DT}{Dt}}`$ $`=`$ $`𝝈:\dot{𝜸}𝐪\epsilon ,`$ (7)
where $`\nu `$ is the filling fraction (the density of granular material normalized by the density of grains), $`𝐮`$ is the velocity field, $`T=(\mathrm{𝐮𝐮}𝐮^2)/2`$ is the granular temperature, $`D/Dt=_t+(𝐮)`$ is the material derivative, $`𝐠`$ is the gravity acceleration, $`\sigma _{\alpha \beta }`$ is the stress tensor, $`𝐪`$ is the energy flux vector, $`\dot{\gamma }_{\alpha \beta }=_\alpha u_\beta +_\beta u_\alpha `$ is the strain rate tensor, and $`\epsilon `$ is the energy dissipation rate. Eqs. (5)-(7) are structurally similar to the Navier-Stokes equations for conventional fluids except for the last term in the equation for granular temperature $`\epsilon `$ which accounts for the energy loss due to inelastic collisions.
These three equations have to be supplemented by the constitutive relations for the stress tensor $`𝝈`$, energy flux $`𝐪`$, and the energy dissipation rate $`\epsilon `$. For dilute systems, a linear relations between stress $`𝝈`$ and strain rate $`\dot{𝜸}`$ is obtained,
$`\sigma _{\alpha \beta }`$ $`=`$ $`[p+(\mu \lambda )\text{Tr}\dot{𝜸}]\delta _{\alpha \beta }\mu \dot{𝜸}_{\alpha \beta },`$ (8)
$`𝐪`$ $`=`$ $`\kappa T.`$ (9)
In the kinetic theory of two-dimensional gas of slightly inelastic hard disks by Jenkins and Richman (1985), these equations are closed with the following equation of state
$$p=\frac{4\nu T}{\pi d^2}[1+(1+e)G(\nu )],$$
(10)
and the expressions for the shear and bulk viscosities
$`\mu `$ $`=`$ $`{\displaystyle \frac{\nu T^{1/2}}{2\pi ^{1/2}dG(\nu )}}\left[1+2G(\nu )+\left(1+{\displaystyle \frac{8}{\pi }}\right)G(\nu )^2\right],`$ (11)
$`\lambda `$ $`=`$ $`{\displaystyle \frac{8\nu G(\nu )T^{1/2}}{\pi ^{3/2}d}},`$ (12)
the thermal conductivity
$$\kappa =\frac{2\nu T^{1/2}}{\pi ^{1/2}dG(\nu )}\left[1+3G(\nu )+\left(\frac{9}{4}+\frac{4}{\pi }\right)G(\nu )^2\right],$$
(13)
and the energy dissipation rate
$$\epsilon =\frac{16\nu G(\nu )T^{3/2}}{\pi ^{3/2}d^3}(1e^2).$$
(14)
The radial pair distribution function $`G(\nu )`$ for a dilute 2D gas of elastic hard disks can be approximated by the formula Song *et al.* (1989)
$$G_{CS}(\nu )=\frac{\nu (17\nu /16)}{(1\nu )^2}$$
(15)
(this is a two-dimensional analog of the famous Carnahan-Starling formula Carnahan and Starling (1969) for elastic spheres). This formula is expected to work for densities roughly below 0.7. For high density granular gases, this function has been calculated using free volume theory by Buehler *et al.* (1951),
$$G_{FV}=\frac{1}{(1+e)\left[(\nu _c/\nu )^{1/2}1\right]}$$
(16)
where $`\nu _c0.82`$ is the density of the random close packing limit. Luding (2001) proposed a global fit
$$G_L=G_{CS}+(1+\mathrm{exp}((\nu \nu _0)/m_0))^1)(G_{FV}G_{CS})$$
with empirically fitted parameters $`\nu _00.7`$ and $`m_010^2`$. However, even with this extension, the continuum theory comprised of Eqs.(5)-(14) cannot describe the force chains which transmit stress via persistent contacts remaining in the dense granular flows, as well as the hysteretic transition from solid to static regimes and coexisting solid and fluid phases.
The granular hydrodynamics is probably the most universal (however not always the most appropriate) tool for modelling large-scale collective behavior in driven granular matter. Granular hydrodynamics equations in the form (5),(6),(7) and their modifications are widely used in the engineering community to describe a variety of large-scale granular flows, especially for design of gas-fluidized bed reactors Gidaspow (1994). In the physics community granular hydrodynamics is used to understand various instabilities in relatively small-scale flows, such as flow past obstacle Rericha *et al.* (2002), convection Livne *et al.* (2002a, b), floating clusters Meerson *et al.* (2003), longitudinal rolls Forterre and Pouliquen (2002, 2003), patterns in vibrated layers Bougie *et al.* (2005) and others. However, Eqs. (5)-(7) are often used far beyond their applicability limits, viz. dilute flows. Consequently, certain parameters and constitutive relations need to be adjusted heuristically in order to accommodate observed behavior. For example, Bougie *et al.* (2002) had to introduce artificial non-zero viscosity in Eq. (6) for $`\nu 0`$ in order to avoid artificial blowup of the solution. Similarly, Losert *et al.* (2000) introduced the viscosity diverging as density approaches the close packed limit as $`(\nu \nu _c)^\beta `$ with $`\beta 1.75`$ being the fitting parameter in order to describe the structure of dense shear granular flows.
### III.2 Phenomenological models
A generic approach to the description of dense granular flows was suggested by Aranson and Tsimring (2001, 2002) who proposed to treat the shear stress mediated fluidization of granular matter as a phase transition. For this purpose an order parameter characterizing the local state of granular matter and the corresponding phase field model were introduced. According to the model, the order parameter has its own relaxation dynamics and defines the static and dynamic contributions to the shear stress tensor. This approach is discussed in more details in Sec. VI.1.1.
Another popular approach is based on the two-phase description of granular flow, one phase corresponding to rolling grains and the other phase to static ones. This approach, so-called the BCRE model, was suggested by Bouchaud *et al.* (1994, 1995) for description of surface gravity driven flows. The BCRE model has direct relation to depth-averaged hydrodynamic equations (so-called Saint-Venant model) popular in the engineering community. Note that BCRE and Saint-Venant models can be derived in a certain limit from the more general order parameter model mentioned above, for detail see Sec. VI.1.2.
Many pattern-forming systems are often described by generic amplitude equations such as Ginzburg-Landau or Swift-Hohenberg equations Cross and Hohenberg (1993); Aranson and Kramer (2002). This approach allows to explain many generic features of patterns, however in any particular system there are peculiarities which need to be taken into account. This often requires modifications to be introduced into the generic models. This approach was taken by Tsimring and Aranson (1997); Aranson and Tsimring (1998); Aranson *et al.* (1999a); Venkataramani and Ott (1998); Crawford and Riecke (1999) in order to describe patterns in a vibrated granular layer. Details of these approaches can be found in Sec. V.4.
In addition, a variety of tools of statistical physics are applied to diverse phenomena occurring in granular systems. For example, celebrated theory of Lifshitz and Slyozov (1958) developed for coarsening phenomena in equilibrium systems was successfully applied to coarsening of clusters in granular systems Aranson *et al.* (2002), see Section VIII.3.
### III.3 Molecular dynamics simulations
Realistic simulation of granular matter consisting of thousands of particles remains a challenge for physics and computer science. Due to simplicity of microscopic interaction laws (at least for “dry” and non-cohesive granular matter) and relatively small number of particles in granular flows as compared to atomic and molecular systems, the molecular dynamics simulations or discrete element models have a potential to address adequately many phenomena occurring in the granular systems.
There exist three fundamentally different approaches, so-called soft particles simulation method; event driven algorithm and the contact dynamics method for rigid particles. For the review on various molecular dynamics simulation methods we recommend Rapaport (1995); Luding (2004); Pöschel and Schwager (2005).
In the soft particle algorithm, all forces acting on a particle either from walls or other particles or external forces are calculated based on the positions of the particles. Once the forces are found, the time is advanced by the explicit integration of the corresponding Newton equations of motion. Various models are used for calculating normal and tangential contact forces. In majority of implementations, the normal contact forces are determined from the particle overlap $`\mathrm{\Delta }_n`$ which is defined as the difference of the distance between the centers of mass of two particles and the sum of their radii. The normal force $`𝐅_n`$ is either proportional to $`\mathrm{\Delta }_n`$ (linear Hookian contact) or proportional to $`\mathrm{\Delta }^{3/2}`$ (Hertzian contact). In the spring-dashpot model, additional dissipative force proportional to the normal component of the relative velocity is added to model inelasticity of grains. A variety of approaches are used to model tangential forces, the most widely accepted of them being Cundall-Strack algorithm Cundall and Strack (1979), in which the tangential contact is modelled by a dissipative linear spring whose force $`𝐅_t=k_t𝚫_tm/2\gamma _t𝐯_t`$ (here $`𝚫_t`$ is the relative tangential displacement and $`𝐯_t`$ is the relative tangential velocity, $`k_t,\gamma _t`$ are model constants). It is truncated when its ratio to the normal force $`|𝐅_t|/|𝐅_n|`$ reaches the friction coefficient $`\mu `$ according to the Coloumb law. Soft-particles methods are relatively slow and used mostly for the analysis of dense flows when generally faster event-driven algorithms are not applicable, see e.g. Silbert *et al.* (2002a, b, c); Landry *et al.* (2003); Volfson *et al.* (2003a, b).
In the event-driven algorithm, the particles are considered infinitely rigid and move freely (or driven by macroscopic external fields) in the intervals between (instantaneous) collisions. The algorithm updates velocities and positions of the two particles involved in a binary collision (in the simplest frictionless case, according to Eq. (1)), and then finds the time of the next collision and velocities and positions of all particles at that time according to Newton’s law. Thus, the time is advanced directly from one collision to the next, and so variable time step is dictated by the interval between the collisions. While event-driven methods are typically faster for dilute rapid granular flows, they become impractical for dense flows where collisions are very frequent and furthermore particles develop persistent contacts. As a related numerical problem, event-driven methods are known to suffer from so-called “inelastic collapse” when the number of collisions between particles diverges in finite time McNamara and Young (1996). There are certain modifications to this method which allow to circumvent this problem by introducing velocity dependent restitution coefficient (see, e.g. Bizon *et al.* (1998a)), but still event-driven methods are mostly applied to rapid granular flows, see e.g. Ferguson *et al.* (2004); McNamara and Young (1996); Khain and Meerson (2004); Nie *et al.* (2002).
Contact dynamics is a discrete element method like soft-particles and event-driven ones, with the equations of motion integrated for each particle. Similarly to event-driven algorithm and unlike soft-particles method, particle deformations are suppressed by considering particles infinitely rigid. The contact dynamics method considers all contacts occurring within a certain short time interval as simultaneous, and computes all contact forces by satisfying simultaneously all kinematic constraints imposed by impenetrability of the particles and the Coulomb friction law. Imposing kinematic constraints requires contact forces (constraint forces) which cannot be calculated from the positions and velocities of particles alone. The constraint forces are determined in such a way that constraint-violating accelerations are compensated. For comprehensive review on the contact dynamics see Brendel *et al.* (2004).
Sometimes different molecular dynamics methods are often applied to the same problem. Lois *et al.* (2005); Staron *et al.* (2002); Radjai *et al.* (1998) applied contact dynamics methods and Silbert *et al.* (2002a, b); Volfson *et al.* (2003a, b) used soft-particles technique for the analysis of instabilities and constitutive relations in dense granular systems. Patterns in vibrated layers were studied by event-driven simulations by Bizon *et al.* (1998a); Moon *et al.* (2003) and by soft particles molecular dynamics simulations by Prevost *et al.* (2004); Nie *et al.* (2000).
## IV Patterns in sub-monolayers. Clustering, Coarsening and Phase Transitions
### IV.1 Clustering in Freely Cooling Gases
Properties of granular gases are dramatically different from the properties of molecular gases due to inelasticity of collisions between the grains. This leads to the emergence of correlation between colliding particles and violation of the molecular chaos approximation. This in turn gives rise to various pattern-forming instabilities. Perhaps the simplest system exhibiting nontrivial pattern formation in the context of granular matter is freely cooling granular gas: isolated system of inelastically colliding particles. The interest to freely cooling granular gases was triggered by the discovery of clustering by Goldhirsch and Zanetti (1993): spontaneously forming dense clusters emerge as a result of instability of initially homogeneous cooling state, see Fig. 2. This instability, which can be traced in many other granular systems, has a very simple physical interpretation: local increase of the density of granular gas results in the increase in the number of collisions, and, therefore, further dissipation of energy and decrease in the granular temperature. Due to proportionality of pressure to the temperature, the decrease of temperature will consequently decrease local pressure, which, in turn, will create a flux of particles towards this pressure depression, and further increase of the density. This clustering instability has interesting counterparts in astrophysics: clustering of self-gravitating gas Shandarin and Zeldovich (1989) and “radiative instability” in optically thin plasmas Meerson (1996) resulting in interstellar dust condensation.
According to Goldhirsch and Zanetti (1993), the initial stage of clustering can be understood in terms of the instability of a homogeneously cooled state described by the density $`\nu `$ and granular temperature $`T`$. This state is characterized by zero hydrodynamic velocity $`v`$, and the temperature evolution follows from the energy balance equation
$$_tTT^{3/2}$$
(17)
which results in the Haff’s cooling law $`Tt^2`$ Haff (1983). However, the uniform cooling state becomes unstable in large enough systems masking Haff’s law. The discussion of the linear instability conditions can be found e.g. in Babic (1993); Brilliantov and Pöschel (2004). For the case of particles with fixed restitution coefficient $`e`$, the analysis in the framework of linearized hydrodynamics equations (5)-(7) yields the critical wavenumber $`k^{}`$ for the clustering instability
$$k^{}\sqrt{1e^2}$$
(18)
As one sees, the length scale of the clustering instability diverges in the limit of elastic particles $`e1`$.
The clustering instability in a system of grains with constant restitution coefficient results in the inelastic collapse discussed in the previous Section. Whereas the onset of clustering can be well-understood in the framework of granular hydrodynamics (see, e.g. Babic (1993); Hill and Mazenko (2003); Goldhirsch (2003); Brilliantov and Pöschel (2004)), certain subtle features (e.g. scaling exponents for temperature) are only assessed within molecular dynamics simulation because the hydrodynamic description often breaks in dense cold clusters. One recent theoretical approach to the description of the late stages of clustering instability consists in introducing additional regularization into the hydrodynamic description due to the finite size of particles Nie *et al.* (2002); Efrati *et al.* (2005).
Nie *et al.* (2002) argued that cluster formation and coalescence in freely cooling granular gases can be heuristically described by the Burgers equation for hydrodynamic velocity $`v`$ with random initial conditions:
$$_tv+vv=\mu _0^2v$$
(19)
where $`\mu _0`$ is effective viscosity (which is different from the shear viscosity in hydrodynamic equations). In this context clustering is associated with the formation of shocks in the Burgers equation. Perhaps not surprising, a very similar approach was applied for description of the gas of “sticky” particles for the description of the large-scale matter formation in the Universe Gurbatov *et al.* (1985); Shandarin and Zeldovich (1989).
Meerson and Puglisi (2005) conducted molecular dynamics simulations of the clustering instability of a freely cooling dilute inelastic gas in a quasi-one-dimensional setting. This problem was also examined in the framework of granular hydrodynamics by Efrati *et al.* (2005). It was observed that, as the gas cools, stresses become negligibly small, and the gas flows only by inertia. Hydrodynamic description reveals a finite-time singularity, as the velocity gradient and the gas density diverge at some location. The molecular dynamics studies show that finite-time singularities, intrinsic in such flows, are arrested only when close-packed clusters are formed. It was confirmed that the late-time dynamics and coarsening behavior are describable by the Burgers equation (19) with vanishing viscosity $`\mu _0`$. Correspondingly, the average cluster mass grows as $`t^{2/3}`$ and the average velocity decreases as $`t^{1/3}`$. Due to the clustering long-term temperature evolution is $`Tt^{2/3}`$ which is different from Haff’s law $`Tt^2`$ derived for the spatially-homogeneous cooling. Efrati *et al.* (2005) argue that flow by inertia represents a generic intermediate asymptotic of unstable free cooling of dilute granular gases consistent with the Burgers equation (19) description of one-dimensional gas of “sticky particles” suggested by Nie *et al.* (2002).
While there is a qualitative similarity between Burgers shocks and clusters in granular materials at least in one dimension, the applicability of the Burgers equation for the description of granular media is still an open question, especially in two and three dimensions. The main problem is that the Burgers equation can be derived from the hydrodynamic equations only in one dimensional situation, in two and three dimensions the Burgers equation assumes zero vorticity, which possibly oversimplifies the problem and may miss important physics. In fact, molecular dynamics simulations illustrate the development of large-scale vortex flows in the course of clustering instability Catuto and Marconi (2004); van Noije and Ernst (2000).
Finally, Das and Puri (2003) proposed a phenomenological description of the long-term clusters evolution in granular gases. Using the analogy between clustering in granular gases and phase-ordering dynamics in two-component mixtures, Das and Puri (2003) postulated generalized Cahn-Hilliard equations for the evolution of density $`\nu `$ and complex velocity $`\psi =v_x+iv_y`$ (see e.g. Bray (1994))
$`_t\nu `$ $`=`$ $`(^2)^m\left[\nu \nu ^3+^2\nu \right]`$ (20)
$`_t\psi `$ $`=`$ $`(^2)^m\left[\psi |\psi |^2\psi +^2\psi \right]`$ (21)
with $`m0+`$ which characterize globally-conserved dynamics of $`\nu `$ and $`\psi `$ similar to that considered in Sec. VIII.3.1. Das and Puri (2003) argue that this choice is most appropriate due to the non-diffusive character of particles motion and is consistent with the observed morphology of clusters. While it might be very challenging to derive Eqs. (20),(21) from the first principles or to deduce them from hydrodynamic equations, the connection to phase-ordering dynamics is certainly deserves further investigation.
### IV.2 Patterns in Driven Granular Gases
Discovery of the clustering instability stimulated a large number of experimental and theoretical studies, even experiments in low gravity conditions Falcon *et al.* (1999). Since “freely cooling granular gas” is difficult to implement in the laboratory, most experiments were performed in the situation when the energy is injected in the granular system in one or another way. Kudrolli *et al.* (1997) studied two-dimensional granular assemblies interacting with a horizontally vibrating (or “hot”) wall. In agreement with granular hydrodynamics, maximum gas density occurs opposite to the vibrating wall, see Fig. 16. The experimental density distributions are consistent with the modified hydrodynamic approach proposed by Grossman *et al.* (1997). Livne *et al.* (2002a, b); Khain and Meerson (2002); Khain *et al.* (2004a), studied the dynamics of granular gases interacting with a hot wall analytically using granular hydrodynamic theory for rigid disks in the formulation of Jenkins and Richman (1985) and predicted a novel phase-separation or van der Waals-type instability of the one-dimensional density distribution. This instability, reproduced later by molecular dynamics simulations Argentina *et al.* (2002) is different from the usual convection instability as it occurs without gravity and is driven by the coarsening mechanism. Simulations indicated a profound role of fluctuations. One may expect that noise amplification near the instability thresholds in granular systems will be very important due to non-macroscopic number of grains. In the context of phase-separation instability Meerson *et al.* (2004) raised the non-trivial question of the origin of giant fluctuations and break-down of hydrodynamic description in granular systems near the threshold of instability (see also Goldman *et al.* (2004); Bougie *et al.* (2005) on the effect of fluctuations in multilayers). Remarkably, for the granular gas confined between two oscillating walls Khain and Meerson (2004) predicted on the basis of event-driven simulations a novel oscillatory instability for the position of the dense cluster. These predictions, however, have not yet been confirmed experimentally, most likely due to the relatively small aspect ratio of available experimental cells.
Olafsen and Urbach (1998) pioneered experiments with sub-monolayers of particles subject to vertical vibration<sup>2</sup><sup>2</sup>2Sub-monolayer implies less than 100% percent coverage by particles of the bottom plate.. Their studies revealed a surprising phenomenon: formation of a dense closely-packed cluster co-existing with dilute granular gas, see Fig. 1. The phenomenon bears a strong resemblance to the first-order solid/liquid phase transition in equilibrium systems. Similar experiments by Losert *et al.* (1999) discovered propagating fronts between gas-like and solid-like phases in vertically vibrated sub-monolayers. Such fronts are expected in extended systems in the vicinity of the first order phase transition, e.g. solidification fronts in supercooled liquids. Prevost *et al.* (2004) performed experiments with vibrated granular gas confined between two plates. Qualitatively similar phase coexistence was found. The cluster formation in vibro-fluidized sub-monolayers shares many common features with processes in freely cooling granular gases because it is also caused by the energy dissipation due to inelasticity of collisions. However, there is a significant difference: the instability described in Subsec. IV.1 is insufficient to explain the phase separation. A very important additional factor is bistability and co-existence of states due to the nontrivial density dependence of the transfer rate of particle’s vertical to horizontal momentum. Particles in a dense closed-packed cluster likely obtain less horizontal momentum than in a moderately dilute gas because in the former particle vibrations are constrained to the vertical plane by interaction with neighbors. In turn, in a very dilute gas the vertical to horizontal momentum transfer is also inhibited due to lack of particle collisions. Another factor here is that vibration is not fully equivalent to the interaction with a heat bath. It is well known that even a single particle interacting with a periodically vibrating plate exhibits coexistence of dynamic and static states Losert *et al.* (1999).
There were several simulation studies of clustering and phase coexistence in vibrated granular submonolayers. Nie *et al.* (2000); Prevost *et al.* (2004) reproduced certain features of cluster formation and two-phase co-existence by means of large-scale three-dimensional molecular dynamics simulations. Since realistic three-dimensional simulations are still expensive and extremely time-consuming, simplified modelling of the effect of a vibrating wall by a certain multiplicative random forcing on individual particles was employed by Cafiero *et al.* (1999). While the multiplicative random forcing is an interesting theoretical idea, it has to be used with caution as it is not guaranteed to reproduce subtle details of particle dynamics, especially the sensitive dependence of the vertical to horizontal momentum transfer as the function of the density.
### IV.3 Coarsening of clusters
One of the most intriguing questions in the context of phase coexistence in vibrofluidized granular sub-monolayers is a possibility of Ostwald-type ripening and coarsening of clusters similar to that observed in equilibrium systems Lifshitz and Slyozov (1958, 1961). In particular, the scaling law for the number of macroscopic clusters is of special interest because it gives a deep insight into the similarity between equilibrium thermodynamic systems and non-equilibrium granular systems. The experiments Olafsen and Urbach (1998); Losert *et al.* (1999); Prevost *et al.* (2004); Sapozhnikov *et al.* (2003) demonstrated emergence and growth of multiple clusters but did not have sufficient aspect ratio to address the problem of coarsening in a quantitative way.
Nevertheless, as it was suggested by Aranson *et al.* (2000), statistical information on out-of-equilibrium Ostwald ripening can be obtained in a different granular system: electrostatically driven granular media. This system permits one to operate with extremely small particles and obtain a very large number of macroscopic clusters. In this system the number of clusters $`N`$ decays with time as $`N1/t`$. This law is consistent with interface-controlled Ostwald ripening in two dimensions, see Wagner (1961). Whereas mechanisms of energy injections are different, both vibrofluidized and electrostatically-driven systems show similar behavior: macroscopic phase separation, coarsening, transition from two- to three-dimensional cluster growth, etc Sapozhnikov *et al.* (2003). In Aranson *et al.* (2002) the theoretical description of granular coarsening was developed in application to the electrostatically driven grains, however we postpone the description of this theory to Sec. VIII.3. We anticipate that a theory similar to that formulated in Aranson *et al.* (2002) can be applicable to mechanically fluidized granular materials as well. The main difference there is the physical mechanism of energy injection which will possibly affect the specific form of the conversion rate function $`\varphi `$ in Eq. (69) in Sec. VIII.3.
## V Surface waves and patterns in vibrated multilayers of granular materials
### V.1 Chladni patterns and heaping
Driven granular systems often manifest collective fluid-like behavior: shear flows, convection, surface waves, and pattern formation (see e.g. Jaeger *et al.* (1996)). Surprisingly, even very thin (less than ten) layers of sand under excitation exhibit pattern formation which is quite similar (however with some important differences) to the corresponding patterns in fluids. One of the most fascinating examples of these collective dynamics is the appearance of long-range coherent patterns and localized excitations in vertically-vibrated thin granular layers.
Experimental studies of vibrated layers of sand have a long and illustrious history, beginning from the seminal works by Chladni (1787) and Faraday (1831) in which they used a violin bow and a membrane to excite vertical vibrations in a thin layer of grains. The main effect observed in those early papers, was “heaping” of granular matter in mounds near the nodal lines of the membrane oscillations. This behavior was immediately (and correctly) attributed to the “acoustic streaming”, or nonlinear detection of the nonuniform excitation of grains by membrane modes. One puzzling result by Chladni was that a very thin powder would collect at the anti-nodal regions where the amplitude of vibrations is maximal. As Faraday demonstrated by evacuating the container, this phenomenology is caused by the role of air permeating the grains in motion. Evidently, the interstitial gas becomes important as the terminal velocity of a free fall $`v_t=\nu gd^2/18\mu `$ becomes of the order of the plate velocity, and this condition is fulfilled for $`1020`$ $`\mu m`$ particles on a plate vibrating with frequency 50 Hz and acceleration amplitude $`g`$.
In subsequent years the focus of attention was diverted from dynamical properties of thin layers of vibrated sand, and only in the last third of the 20th century physicists returned to this old problem equipped with new experimental capabilities. The dawn of the new era was marked by the studies of heaping by Jenny (1964). In subsequent papers Walker (1982); Dinkelacker *et al.* (1987); Evesque *et al.* (1989); Douady *et al.* (1989); Laroche *et al.* (1989), more research has been performed of heaping with and without interstitial gas, with somewhat controversial conclusions on the necessity of ambient gas for heaping (see, e.g. Evesque (1990)). Eventually, after more careful analysis Pak *et al.* (1995) concluded that heaping indeed disappears as the pressure of the ambient gas tends to zero or the particle size increases. This agreed with numerical molecular dynamics simulations Taguchi (1992); Gallas *et al.* (1992a, b, c); Gallas and Sokołowski (1993); Luding *et al.* (1994) which showed no heaping without interstitial gas effects. Recent studies of deep layers ($`50<N<200`$) of small particles ($`10<d<200\mu m`$) by Falcon *et al.* (1999a); Duran (2000, 2001) showed a number of interesting patterns and novel instabilities caused by interstitial air. In particular, Duran (2001) observed formation of isolated droplets of grains after periodic taping similar to the Rayleigh-Taylor instability in ordinary fluids.
Jia *et al.* (1999) proposed a simple model for heap formation which is motivated by these experiments. In a discrete lattice version of the model, the decrease in local density due to vibrations is modelled by the random creation of empty sites in the bulk. The bulk flow is simulated by the dynamics of empty sites, while the surface flow is modelled by rules similar to the sandpile model (see Sec. VI.2). This model reproduced both convection inside the powder and the heap formation for sufficiently large probability of empty site formation (which mimics the magnitude of vibration). Jia *et al.* (1999) also proposed the continuum model which has a simple form of a nonlinear reaction-diffusion equation, for the local height of the sandpile
$$_th=D^2h+\mathrm{\Omega }h\beta h^2.$$
(22)
However this model is perhaps too generic and lacks the specific physics of the heaping process.
### V.2 Standing wave patterns
While heaping may or may not appear depending on the gas pressure and the particle properties at small vertical acceleration, at higher vertical acceleration patterns of standing waves emerge in thin layers. They were first reported by Fauve *et al.* (1989); Douady *et al.* (1989) in a quasi two-dimensional geometry. These waves oscillated at the half of the driving frequency, which indicates the sub-harmonic resonance characteristic for parametric instability. This first observation spurred a number of experimental studies of standing waves in thin granular layers in two and three dimensional geometries Melo *et al.* (1994, 1995); Umbanhowar *et al.* (1996); Clément *et al.* (1996); Aranson *et al.* (1999b); Mujica and Melo (1998). Importantly, these studies were performed in evacuated containers, which allowed to obtain reproducible results not contaminated by heaping. Fig. 3 shows a variety of regular patterns observed in vibrated granular layers under vibration Melo *et al.* (1994). As a result of these studies, the emerging picture of pattern formation appears as follows.
The particular pattern is determined by the interplay between driving frequency $`f`$ and acceleration of the container $`\mathrm{\Gamma }=4\pi ^2𝒜f^2/g`$ ($`𝒜`$ is the amplitude of oscillations, $`g`$ is the gravity acceleration) Melo *et al.* (1994, 1995). The layer of grains remains flat for $`\mathrm{\Gamma }<2.4`$ more-less independent of driving frequency. At higher $`\mathrm{\Gamma }`$ patterns of standing waves emerge. At small frequencies $`f<f^{}`$ (for experimental conditions of Melo *et al.* (1995) , $`f^{}`$ 45 Hz) the transition is subcritical, leading to the formation of square wave patterns, see Fig.3b. For higher frequencies $`f>f^{}`$ the selected pattern is quasi-one-dimensional stripes (Fig. 3a), and the transition becomes supercritical. In the intermediate region $`ff^{}`$, localized excitations (oscillons, Fig.4) and various bound states of oscillons (Fig.3f) were observed within the hysteretic region of the parameter plane. Both squares and stripes, as well as oscillons, oscillate at the half of the driving frequency, which indicates the parametric mechanism of their excitation. The wavelength of the cellular patterns near the onset scales linearly with the depth of the layer and diminishes with the frequency of vibration Umbanhowar and Swinney (2000). The frequency corresponding to the strip-square transition was shown to depend on the particle diameter $`d`$ as $`d^{1/2}`$. This scaling suggests that the transition is controlled by the relative magnitude of the energy influx from the vibrating plate $`f^2`$ and the gravitational dilation energy $`gd`$. At higher acceleration ($`\mathrm{\Gamma }>4`$), stripes and squares become unstable, and hexagons appear instead (Fig. 3c). Further increase of acceleration at $`\mathrm{\Gamma }4.5`$ converts hexagons into a domain-like structure of flat layers oscillating with frequency $`f/2`$ with opposite phases. Depending on parameters, interfaces separating flat domains, are either smooth or “decorated” by periodic undulations (Fig. 3e). For $`\mathrm{\Gamma }>5.7`$ various quarter-harmonic patterns emerge. The complete phase diagram of different regimes observed in a three-dimensional container is shown in Fig. 17. For even higher acceleration ($`\mathrm{\Gamma }>7`$) the experiments reveal surprising phase bubbles and spatio-temporal chaos oscillating approximately at one fourth the driving frequency Moon *et al.* (2002).
Subsequent investigations revealed that periodic patterns share many features with convective rolls in Rayleigh-Bénard convection, for example skew-varicose and cross-roll instabilities de Bruyn *et al.* (1998).
### V.3 Simulations of vibrated granular layers
The general understanding of the standing wave patterns in thin granular layers can be gained by the analogy with ordinary fluids. The Faraday instability in fluids and corresponding pattern selection problems have been studied theoretically and numerically in great detail (see e.g. Zhang and Viñals (1997)). The primary mechanism of instability is the parametric resonance between the spatially uniform periodic driving at frequency $`f`$ and two counter-propagating gravity waves at frequency $`f/2`$. However, this instability in ordinary fluids leads to a supercritical bifurcation and square wave patterns near offset, and as a whole the corresponding phase diagram lacks the richness of the granular system. Of course this can be explained by the fact that there are many qualitative differences between granular matter and fluids, such as presence of strong dissipation, friction and the absence of surface tension in the former. Interestingly, localized oscillon-type objects were subsequently observed in vertically vibrated layers of non-Newtonian fluid Lioubashevski *et al.* (1999), and stipe patterns were observed in highly viscous fluid Kiyashko *et al.* (1996). The theoretical understanding of the pattern formation in a vibrated granular system presents a challenge, since unlike fluid dynamics there is no universal theoretical description of dense granular flows analogous to the Navier-Stokes equations. In the absence of this common base, theoretical and computational efforts in describing these patterns followed several different directions. Aoki *et al.* (1996) were first to perform molecular dynamics simulations of patterns in the vibrated granular layer. They concluded that grain-grain friction is necessary for pattern formation in this system. However, as noted by Bizon *et al.* (1997), this conclusion is a direct consequence of the fact that the algorithm of Aoki *et al.* (1996), which is based on the Lennard-Jones interaction potential and velocity dependent dissipation, leads to the restitution coefficient of particles approaching unity for large collision speeds rather than decreasing according to experiments.
Bizon *et al.* (1998a, b) performed event-driven simulations of colliding grains on a vibrated plate assuming constant restitution (see also Luding *et al.* (1996) for earlier two-dimensional event-driven simulations). It was demonstrated that even without friction, patterns do form in the system, however only supercritical bifurcation to stripes is observed. It turned out that friction is necessary to produce other patterns observed in experiments, such as squares and $`f/4`$ hexagons. Simulations with frictional particles reproduced the majority of patterns observed in experiments and many features of the bifurcation diagram (with the important exception of the oscillons). Bizon *et al.* (1998a) set out to match an experimental cell and a numerical system, maintaining exactly the same size container and sizes and the number of particles. After fitting only two parameters of the numerical model, Bizon *et al.* (1998a) were able to find a very close quantitative agreement between various patterns in the experimental cell and patterns in simulations throughout the parameter space of the experiment (frequency of driving, amplitude of acceleration, thickness of the layer), see Fig. 18.
Shinbrot (1997) proposed a model which combined ideas from molecular dynamics and continuum modelling. Specifically, the model ignored vertical component of particle motion and assumed that impact with the plate adds certain randomizing horizontal velocity to the individual particles. The magnitude of the random component being added at each impact served as a measure of impact strength. After the impact particles were allowed to travel freely in the horizontal plane for a certain fraction of a period after which they inelastically collide with each other (a particle acquires momentum averaged over all particles in its neighborhood). This model did reproduce a variety of patterns seen in experiments (stripes, squares, and hexagons) for various values of control parameters (frequency of driving and impact strength), however it did not describe some of the experimental phenomenology (localized objects as well as interfaces), besides it also produced a number of intricate patterns not seen in experiments.
### V.4 Continuum theories
The first continuum models of pattern formation in vibrating sand were purely phenomenological. In the spirit of weakly-nonlinear perturbation theories Tsimring and Aranson (1997) introduced the complex amplitude $`\psi (x,y,t)`$ of sub-harmonic oscillations of the layer surface, $`h=\psi \mathrm{exp}(i\pi ft)+c.c.`$. The equation for this function on the symmetry grounds in the lowest order was written as
$$_t\psi =\gamma \psi ^{}(1i\omega )\psi +(1+ib)^2\psi |\psi |^2\psi \nu \psi .$$
(23)
Here $`\gamma `$ is the normalized amplitude of forcing at the driving frequency $`f`$. The linear terms in Eq. (23) can be obtained from the complex growth rate for infinitesimal periodic layer perturbations $`h\mathrm{exp}[\mathrm{\Lambda }(k)t+ikx]`$. Expanding $`\mathrm{\Lambda }(k)`$ for small $`k`$, and keeping only two leading terms in the expansion $`\mathrm{\Lambda }(k)=\mathrm{\Lambda }_0\mathrm{\Lambda }_1k^2`$ gives rise to the linear terms in Eq. (23), where $`b=Im\mathrm{\Lambda }_1/Re\mathrm{\Lambda }_1`$ characterizes ratio of dispersion to diffusion and parameter $`\omega =(Im\mathrm{\Lambda }_0+\pi f)/Re\mathrm{\Lambda }_0`$, characterizes the frequency of the driving.
The only difference between this equation and the Ginzburg-Landau equation for the parametric instability Coullet *et al.* (1990) is the coupling of the complex amplitude $`\psi `$ to the “slow mode” $`\nu `$ which characterizes local dissipation in the granular layer ($`\nu `$ can be interpreted as coarse-grained layer’s number density). This slow mode obeys its own dynamical equation
$`_t\nu `$ $`=`$ $`\alpha (\nu |\psi |^2)+\beta ^2\nu .`$ (24)
This equation describes re-distribution of the averaged thickness due to the diffusive flux $`\nu `$, and an additional flux $`\nu |\psi |^2`$ is caused by the spatially nonuniform vibrations of the granular material. This coupled model was used by Tsimring and Aranson (1997); Aranson and Tsimring (1998) to describe the pattern selection near the threshold of the primary bifurcation. The phase diagram of various patterns found in this model is shown in Fig. 19. At small $`\alpha \nu \beta ^1`$ (which corresponds to low frequencies and thick layers), the primary bifurcation is subcritical and leads to the emergence of square patterns. For higher frequencies and/or thinner layers, transition is supercritical and leads to roll patterns. At intermediate frequencies stable localized solutions of Eqs.(23),(24) corresponding to isolated oscillons and a variety of bound states were found in agreement with experiment. The mechanism of oscillon stabilization is related to the oscillatory asymptotic behavior of the tails of the oscillon (see Fig. 20), since this underlying periodic structure provides pinning for the circular front forming the oscillon. Without such pinning, the oscillon solution could only exist at a certain unique value of a control parameter (e.g. $`\gamma `$), and would either collapse or expand otherwise.
Let us note that stable localized solutions somewhat resembling oscillons have recently been found in the nonlinear Schrödinger equation with additional linear dissipation and parametric driving Barashenkov *et al.* (2002).
Phenomenological model (23),(24) also provides a good description of patterns away from the primary bifurcation - hexagons and interfaces Aranson *et al.* (1999a). In high-frequency limit the slow mode dynamics can be neglected ($`\nu `$ becomes enslaved by $`\psi `$), and the dynamics can me described by a single parametric Ginzburg-Landau equation (23).
It is convenient to shift the phase of the complex order parameter via $`\stackrel{~}{\psi }=\psi \mathrm{exp}(i\varphi )`$ with $`\mathrm{sin}2\varphi =\omega /\gamma `$. The equations for real and imaginary part $`\stackrel{~}{\psi }=A+iB`$ are:
$`_tA`$ $`=`$ $`(s1)A2\omega B(A^2+B^2)A+^2(AbB),`$
$`_tB`$ $`=`$ $`(s+1)B(A^2+B^2)B+^2(B+bA),`$ (25)
where $`s^2=\gamma ^2\omega ^2`$. At $`s<1`$, Eqs. (25) has only one trivial uniform state $`A=0,B=0`$, At $`s>1`$, two new uniform states appear, $`A=\pm A_0,B=0,A_0=\sqrt{s1}`$. The onset of these states corresponds to the period doubling of the layer flights sequence, observed in experiments Melo *et al.* (1994, 1995) and predicted by the simple inelastic ball model Melo *et al.* (1994, 1995); Mehta and Luck (1990). Signs $`\pm `$ reflect two relative phases of layer flights with respect to container vibrations.
Weakly-nonlinear analysis reveals that the uniform states $`\pm A_0`$ lose their stability with respect to finite-wavenumber perturbations at $`s<s_c`$, and the nonlinear interaction of growing modes leads to hexagonal patterns. The reason for this is that the non-zero base state $`A=\pm A_0`$ lacks the up-down symmetry $`\psi \psi `$ and the corresponding amplitude equations contains quadratic terms which are known to favor hexagons close to onset (see, e.g. Cross and Hohenberg (1993)). In the regime when the uniform states $`A=\pm A_0,B=0`$ are stable, there is an interface solution connecting these two asymptotic states. This interface may exhibit transversal instability which leads to decorated interfaces (see experimental Fig. 3e). Due to symmetry, the interfaces are immobile, however breaking the symmetry of driving can lead to interface motion. This symmetry breaking can be achieved by additional subharmonic driving at frequency $`f/2`$. The interface will move depending on the relative phases of $`f`$ and $`f/2`$ harmonics of driving. This interface drift was predicted in Aranson *et al.* (1999a) and observed in the subsequent work Aranson *et al.* (1999b). As it was noted by Aranson *et al.* (1999b) (see also later work by Moon *et al.* (2003)), moving interfaces can be used to separate granular material of different sizes. The stability and transition between flat and decorated interfaces was studied theoretically and experimentally by Blair *et al.* (2000). It was shown that non-local effects are responsible for the saturation of transverse instability of interfaces. Moreover, new localized solutions (“superoscillons”) were found for large accelerations. In contrast with conventional oscillons existing on the flat background oscillating with driving frequency $`f`$, i.e. in our notation $`\psi =0`$, the superoscillons exist on the background of the flat period-doubled solution $`\psi 0`$.
Another description of the primary pattern-forming bifurcation was done by Crawford and Riecke (1999) in the framework of the generalized Swift-Hohenberg equation
$`_t\psi `$ $`=`$ $`R\psi (_x^2+1)^2\psi +b\psi ^3c\psi ^5+\epsilon [(\psi )^3]`$ (26)
$``$ $`\beta _1\psi (\psi )^2\beta _2\psi ^2^2\psi .`$
Here the (real) function $`\psi `$ characterizes the amplitude of the oscillating solution, so implicitly it is assumed that the whole pattern always oscillates in phase. Terms proportional to $`\epsilon `$ have been added to the standard Swift-Hohenberg equation first introduced for description of convective rolls (see, e.g. Cross and Hohenberg (1993)) since they are known to favor square patterns, and extended fifth-order local nonlinearity allowed to simulate subcritical bifurcation for $`R<0`$. This equation also describes both square and stripe patterns depending on the magnitude of $`\epsilon `$ and for negative $`R`$ has a stable oscillon-type solution.
Even more generic approach was taken by Venkataramani and Ott (1998, 2001) who argued that the spatio-temporal dynamics of patterns generated by parametric forcing can be understood in the framework of a discrete-time, continuous space system which locally exhibits a sequence of period-doubling bifurcations and whose spatial coupling operator selects a certain spatial scale. In particular they studied the discrete-time system
$$\xi _{n+1}(𝐱)=[M(\xi _n(𝐱))]$$
(27)
where local mapping $`M(\xi )`$ is described by a Gaussian map
$$M(\xi )=\stackrel{~}{r}\mathrm{exp}[(\xi 1)^2/2]$$
and the linear spatial operator $``$ has an azimuthally symmetric Fourier transform
$$f(k)=\text{sign}[k_c^2k^2]\mathrm{exp}[k^2(1k^2/2k_0^2))/2].$$
Here $`k`$ is the wavenumber, $`k_c,k_0`$ are two inverse length scales characterizing the spatial coupling, and $`\stackrel{~}{r}`$ describes the amplitude of forcing. While this choice of the spatial operator appears rather arbitrary, it leads to a phase diagram on the plane $`(k_c/k_0,r)`$ which is similar to the experimental one.
Several authors Cerda *et al.* (1997); Eggers and Riecke (1999); Park and Moon (2002) attempted to develop a quasi two-dimensional fluid-dynamics-like continuum description of the vibrated sand patterns. These models deal with mass and momentum conservation equations which are augmented by specific constitutive relations for the mass flux and pressure. Cerda *et al.* (1997) assumed that during impact particles acquire horizontal velocities proportional to the gradient of local thickness, then during the flight that move freely with these velocities and redistribute mass, and during the remainder of the cycle the layer diffusively relaxes on the plate. The authors found that a flat layer is unstable with respect to square pattern formation, however the transition is supercritical. In order to account for the subcritical character of the primary bifurcation to square patterns, the authors postulated the existence of a certain critical slope (related to the repose angle) below which the free flight initiated by the impact does not occur. They also observed the existence of localized excitations (oscillons and bound states), however they appeared only as transients in the model. Park and Moon (2002) generalized this model by explicitly writing the momentum conservation equation and introducing the equation of state for the hydrodynamic pressure which is proportional to the square of the velocity divergence. This effect provides saturation of the free-flight focusing instability and leads to a squares-to-stripes transition at higher frequencies which was missing in the original model Cerda *et al.* (1997). By introducing multiple free-flight times and contact times Park and Moon (2002) were also able to reproduce hexagonal patterns and superlattices.
Full three-dimensional continuum simulations based on the granular hydrodynamics equations (5),(6),(7) were performed by Bougie *et al.* (2005). Quantitative agreement was found between this description and event-driven molecular dynamics simulations and experiments in terms of the wavelength dependence on the vibration frequency (Fig.21) although the authors had to introduce a certain regularization procedure in the hydrodynamic equations in order to avoid artificial numerical instabilities for $`\nu 0`$. Since standard granular hydrodynamics does not take into account friction among particles, the simulations only yielded stripe pattern, in agreement with earlier molecular dynamics simulations. Furthermore, the authors found a small but systematic difference ($``$10%) between the critical value of plate acceleration in fluid-dynamical and molecular dynamics simulations which could be attributed to the role of fluctuations near the onset. Proper account of inter-particle friction and fluctuations within the full hydrodynamics description still remains an open problem (see more on that in Section VI).
Fluctuations are expected to play a significantly greater role in granular hydrodynamics than in usual fluids, because the total number of particles involved in the dynamics per characteristic spatial scale of the problem is many orders of magnitude smaller than the Avogadro number. The apparatus of fluctuating hydrodynamics which was developed in particular for description of transition to rolls in Rayleigh-Bénard convection Swift and Hohenberg (1977), has been recently applied to the granular patterns Goldman *et al.* (2004); Bougie *et al.* (2005). The Swift-Hohenberg theory is based on the equation for the order parameter $`\psi `$,
$$_t\psi =[ϵ(^2+k_0^2)^2]\psi \psi ^3+\eta (𝐱,t),$$
(28)
where $`ϵ`$ is the bifurcation parameter, $`k_0`$ is the wavenumber corresponding to the most unstable perturbations, and $`\eta `$ is the Gaussian $`\delta `$ correlated noise term with intensity $`F`$. The Swift-Hohenberg theory predicts that noise offsets the bifurcation value of the control parameter from the mean-field value $`ϵ_{MF}=0`$ to the critical value $`ϵ_cF^{2/3}`$. Furthermore, the Swift-Hohenberg theory describes the transition to the linear regime which is expected to work far away from the bifurcation point for small noise intensity when the magnitude of noise-excited modes scales as $`|ϵϵ_c|^{1/2}`$, while the time coherence of fluctuations and the amplitude of spectral peaks decays as $`|ϵϵ_c|^1`$. Fitting the Swift-Hohenberg equation (28) to match the transition in vibrated granular layer, Goldman *et al.* (2004); Bougie *et al.* (2005) found a good agreement with molecular dynamics simulations and experiments (see, e.g., Fig.22). Interestingly, the magnitude of the fitted noise term in Eq.(28) $`F3.510^3`$ turned out to be an order of magnitude greater than for convective instability in a fluid near a critical point Oh and Ahlers (2003). This discrepancy could stem from the fact that the Swift-Hohenberg theory, developed for ordinary fluids, is formally valid for the second-order phase transition, whereas in granular system the transition to square patterns is of the first order type. Consequently, the nonlinear terms can be important near the transition point and may distort the scaling for the noise amplitude.
There have been attempts to connect patterns in vibrated layers with the phenomenon of granular “thermoconvection”. Since high-frequency vibration in many aspects is similar to “hot” wall, it was argued that one should expect granular temperature gradients, density inversion, and, consequently convection instability similar to that observed in heated from below liquid layers. The theoretical analysis based on granular hydrodynamic equations (5),(6),(7) supports the existence of a convective instability in a certain range of parameters He *et al.* (2002); Khain and Meerson (2003). Multiple convection roles were observed in molecular dynamics simulations Sunthar and Kumaran (2001); Paolotti *et al.* (2004). However, the experiments are not conclusive enough Wildman *et al.* (2001). In particular it appears very hard to discriminate between convection induced by vibration and convective flows induced by walls, see e.g. Pak and Behringer (1993); Garcimarin *et al.* (2002).
Vibrated bottom plate is not the only way to induce parametric patterns in thin granular layers. Li *et al.* (2003) demonstrated that periodically modulated airflow through a shallow fluidized bed also produces interesting patterns in the granular layer which oscillate at half the driving frequency (Fig. 23). While the physical mechanism of interaction between the airflow and grains is quite different from the collisional energy transfer in vibrated containers, phenomenological models based on the principal symmetry of the problem should be able to describe the gas-driven granular layer as well. In case of the parametric Ginzburg-Landau model Eq. (23), the order parameter would correspond to the amplitude of the subharmonic component of the surface deformation, and the driving term would be related to the amplitude of the flow modulation. Moreover, variations of the mean flow rate act similar to the variations of the gravitational acceleration in the mechanical system, which may give an additional means to control the state of the system.
## VI Patterns in gravity-driven dense granular flows
In this Section we overview theoretical models for various pattern-forming instabilities in dense gravity-driven granular flows.
### VI.1 Avalanches in thin granular layers
Gravity-driven particulate flows are a common occurrence in nature (dune migration, erosion/deposition processes, land slides, underwater gravity currents and coastal geomorphology) and in various industrial applications having to do with handling granular materials, including their storage, transport, and processing. One of the most spectacular (and often very dangerous) forms of gravity-driven granular flows is the avalanche. Avalanches occur spontaneously when the slope of the granular material exceeds a certain angle (static angle of repose) or they can be initiated at somewhat smaller angle by applying a finite perturbation. Laboratory studies of avalanches are often carried out in rotating drums (see below) or in a chute geometry when a layer of sand is titled at a certain fixed angle. Daerr and Douady (1999) conducted experiments with a thin layer of granular matter on sticky (velvet) inclined plane, see Fig. 24. Surprising diversity of avalanche behavior was observed in this seemingly simple system: triangular avalanches developed in thin layers ($`h`$ is the layer thickness) and for small inclination angles $`\varphi `$, whereas in thicker layers or steeper angles $`\varphi `$ the avalanches assumed balloon shaped with upper edge of the avalanche propagating up-hill, see Fig. 6. According to Rajchenbach (2002a, 2003) the rear front of the balloon-like avalanche propagates uphill with the velocity roughly one half of the downhill velocity of the head front, and the velocity of the head is also two times larger than the depth-averaged flow velocity. The stability diagram is outlined in Fig. 24: a granular layer is stable below solid line (so-called $`h_{stop}`$ limit according to Pouliquen (1999)), spontaneous avalanching was observed above the dashed line. Between dashed and solid lines the layer exhibits bistable behavior: finite perturbation can trigger an avalanche, otherwise the layer remains stable. The dotted line with $`\times `$-symbols indicates the transition between triangular and balloon avalanches.
#### VI.1.1 Partially fluidized flows
The avalanche dynamics described above is an example of a wide class of partially fluidized granular flows. In such flows part of grains flows past each other while other grains maintain static contacts with their neighbors. The description of such flows still represents a major challenge for the theory. In particular, one is faced with the problem of constructing the constitutive relation for the stress tensor $`𝝈`$. In dense quasi-static flows a significant part of the stresses is transmitted through quasi-static contacts between particles as compared with short collisions in dilute flows.
Stimulated by the non-trivial avalanche dynamics in experiments by Daerr and Douady (1999); Daerr (2001), Aranson and Tsimring (2001, 2002) suggested a generic continuum description of partially fluidized granular flows. According to this theory, the ratio of the static part $`𝝈^s`$ to the fluid part $`𝝈^f`$ of the full stress tensor is controlled by the order parameter $`\rho `$. The order parameter is scaled in such a way that in granular solid $`\rho =1`$ and in well developed flow (granular liquid) $`\rho 0`$. On the “microscopic level” the order parameter is defined as a fraction of the number of static (or persistent) contacts of the particles $`Z_s`$ to total number of the contacts $`Z`$, $`\rho =Z_s/Z`$ within a mesoscopic volume which is large with respect to the particle size but small compared with characteristic size of the flow.
Due to a strong dissipation in dense granular flows the order parameter $`\rho `$ is assumed to obey purely relaxational dynamics controlled by the Ginzburg-Landau-type equation for the generic first order phase transition,
$$\frac{D\rho }{Dt}=D^2\rho \frac{F(\rho ,\delta )}{\rho }.$$
(29)
Here $`D`$ is the diffusion coefficient. $`F(\rho ,\delta )`$ is a free energy density which is postulated to have two local minima at $`\rho =1`$ (solid phase) and $`\rho =0`$ (fluid phase) to account for the bistability near the solid-fluid transition.
The relative stability of the two phases is controlled by the parameter $`\delta `$ which in turn is determined by the stress tensor. The simplest assumption consistent with the Mohr-Coloumb yield criterion is to take it as a function of $`\varphi =\mathrm{max}|\sigma _{mn}/\sigma _{nn}|`$, where the maximum is sought over all possible orthogonal directions $`m`$ and $`n`$ (we consider here only two-dimensional formulation of the model, an objective three dimensional generalization was recently proposed by Gao *et al.* (2005)). Furthermore, there are two angles which characterize the fluidization transition in the bulk of granular material, an internal friction angle $`\mathrm{tan}^1\varphi _1`$ such that if $`\varphi >\varphi _1`$ the static equilibrium is unstable, and the “dynamic repose angle” $`\mathrm{tan}^1\varphi _0`$ such that at $`\varphi <\varphi _0`$, the “dynamic” phase $`\rho =0`$, is unstable. Values of $`\varphi _0`$ and $`\varphi _1`$ depend on microscopic properties of the granular material, and in general they do not coincide. Aranson and Tsimring (2001, 2002) adopted the simple algebraic form of the control parameter $`\delta `$,
$$\delta =(\varphi ^2\varphi _0^2)/(\varphi _1^2\varphi _0^2).$$
(30)
Order parameter equation (29) has to be augmented by boundary conditions. While this is a complicated issue in general, a simple but meaningful choice is to take no-flux boundary conditions at free surfaces and smooth walls, and solid phase condition $`\rho =1`$ near sticky or rough walls.
For the flow of thin granular layers on inclined planes Eqs. (6), (29) can be simplified. Using the no-slip boundary condition at the bottom and no-flux condition at the top of the layer and fixing the lowest-mode structure of the order parameter in the direction perpendicular to the bottom of the chute ($`z=0`$, see Inset to Fig. 25), $`\rho =1A(x,y)\mathrm{sin}(\pi z/h)`$, $`h`$ is the local layer thickness, $`A(x,y)`$ is slowly-varying function, one obtains equations governing the evolution of thin layer Aranson and Tsimring (2001, 2002):
$`_th`$ $`=`$ $`\alpha _x(h^3A)+{\displaystyle \frac{\alpha }{\varphi }}\left(h^3Ah\right),`$ (31)
$`_tA`$ $`=`$ $`\lambda A+_{}^2A+{\displaystyle \frac{8(2\delta )}{3\pi }}A^2{\displaystyle \frac{3}{4}}A^3`$ (32)
where $`_{}^2=_x^2+_y^2`$, $`\lambda =\delta 1\pi ^2/4h^2`$, $`\alpha 0.12\mu ^1g\mathrm{sin}\overline{\phi }`$, $`\mu `$ is the shear viscosity, $`\overline{\phi }`$ is the chute inclination, $`\varphi =\mathrm{tan}\overline{\phi }`$. Control parameter $`\delta `$ includes correction due to change in the local slope $`\delta =\delta _0+\beta h_x`$, $`\beta 1/(\varphi _1\varphi _0)1.53`$ depending on the value of $`\varphi `$. The last term in Eq. (31) is also due to change of local slope angle $`\phi `$ and is obtained from the expansion $`\phi \overline{\phi }+h_x`$. This term is responsible for the saturation of the avalanche front slope (without it the front would be arbitrarily steep). While it was not included in original publications Aranson and Tsimring (2001, 2002), this term is important for large wavenumber cut-off of long-wave instability observed by Forterre and Pouliquen (2003), see Sec. VI.3. Numerical and analytic solutions of Eqs. (31),(32) exhibit strong resemblance with experiment: triangular avalanches in thin layers and balloon-like avalanches in thicker layers, see Fig. 26. The corresponding phase diagram agrees quantitatively with an experimental one having only one fitting parameter (viscosity $`\mu `$), Fig. 25.
In subsequent work Aranson and Tsimring (2002), this theory was generalized to other dense shear granular flows including flows in rotating drums, two- and three-dimensional shear cells, etc. The model also was tested in soft-particle molecular dynamics simulations Volfson *et al.* (2003a, b, 2004). Orpe and Khakhar (2005) used the partial fluidization model of Aranson and Tsimring (2001, 2002); Volfson *et al.* (2003a) for the description of velocity profiles three-dimensional shear flows in a rotating drum. The comparison between experimental data and theory shows that the partial fluidization model describes reasonably well entire velocity profile and the flow rheology, however experimental methods for independently estimating the order parameter model are needed. Gao *et al.* (2005) recently developed an objective (coordinate system independent) formulation of the partial fluidization theory which allows for the straightforward generalization to three-dimensional systems.
#### VI.1.2 Two-phase flow approach of granular avalanches
Another approach treating near-surface granular flows as two-phase systems was developed by a number of authors, see e.g. Bouchaud *et al.* (1994, 1995); Boutreux *et al.* (1998); Douady *et al.* (1999); Mehta (1994); Khakhar *et al.* (2001) and many others. For review on recent models of surface flows see Aradian *et al.* (2002). All these models distinguish rolling and static phases of granular flow described by the set of coupled equations for the evolution of thicknesses of both phases, $`R`$ and $`h`$, respectively. The phenomenological theory by Mehta (1994); Bouchaud *et al.* (1994, 1995) (often called BCRE theory) provides an intuitive description of the flow. In shallow granular layers, even simpler depth-averaged granular hydrodynamic equations (so-called Saint-Venant models) often provides quite accurate description, see Savage and Hutter (1989); Douady *et al.* (1999); Khakhar *et al.* (2001); Lajeunesse *et al.* (2004).
In the most general and compact form the BCRE theory can be represented by a pair of equations for evolution of $`R`$ and $`h`$,
$`_th`$ $`=`$ $`\mathrm{\Gamma }(h,R)`$ (33)
$`_tR`$ $`=`$ $`v_d_xR\mathrm{\Gamma }(h,R)`$ (34)
where $`\mathrm{\Gamma }(h,R)`$ is the exchange term, or a conversion rate between rolling and static grains, and $`v_d`$ is the downhill grain velocity. Physical meaning of the BCRE model is very simple: Eq. (33) expresses the increase in the height due to deposition of rolling grains, and Eq. (34) describes advection of rolling fraction by the flow with velocity $`v_d`$ and depletion due to conversion to static fraction. The limitations and generalizations of the BCRE model are discussed by Boutreux *et al.* (1998); Aradian *et al.* (2002).
Douady *et al.* (2002) applied the following two-phase model to describe avalanches in thin granular layers:
$`_th+2U_xh`$ $`=`$ $`{\displaystyle \frac{g}{\overline{\mathrm{\Gamma }}}}\left(\mathrm{tan}\varphi \mu (h)\right)`$ (35)
$`_t\zeta +2U_xh`$ $`=`$ $`0`$ (36)
where $`\zeta =R+h`$ is the position of free surface, $`U`$ is depth-averaged velocity of the flow. In addition to BCRE model Eqs. (35),(36) include two phenomenological functions: $`\overline{\mathrm{\Gamma }}`$ characterizes the mean velocity gradient of a single bead on incline, and $`\mu (h)`$ describes depth-dependent friction with the bottom. According to Douady *et al.* (2002) a three-dimensional version of Eqs. (35),(36) describes transition from triangular to uphill avalanches, however details of the transition depend sensitively on the choice of functions $`\overline{\mathrm{\Gamma }}`$ and $`\mu (h)`$.
Depth-averaged description in the form of Eqs. (35),(36) was used by Börzsönyi *et al.* (2005) to address the difference between shapes of avalanches for sand and glass particles in a chute flow. The authors reduced Eqs. (35),(36) to the modified Burgers equation
$$_th+a(h)_xh=\mu (h)_x^2h$$
(37)
where function $`a(h)h^{3/2}`$ and effective viscosity $`\mu \sqrt{h}`$. This description connects avalanches with the “Burgers” shocks. Eq. (37) implies that all avalanches will eventually decay, in contrast to experiments indicating that only small avalanches decay whereas large avalanches grow and/or form stationary waves Daerr and Douady (1999); Daerr (2001a). This discrepancy is likely due to the fact that reduction of the full model (35),(36) to the single equation (37) does not take into account the bistable nature of granular flows.
While two-phase description of granular flow is simple and rather intuitive, it can be problematic when a clear-cut separation between rolling and static phases is absent, especially near the onset of motion. The order parameter approach can be more appropriate in this situation. Furthermore, the two-phase equations can be derived from the partial fluidization model described in the Sec. VI.1.1 as a sharp-interface limit of the continuum order parameter model Aranson and Tsimring (2002).
#### VI.1.3 Avalanche shape
On the basis of simple kinematic considerations Rajchenbach (2002b) suggested an analytic expression for the shape of triangular and balloon-like avalanches. For the balloon-like avalanches the shape is given by the envelope of the expanding circles with the center drifting downhill:
$$x^2+\left(y2\overline{v}t+\frac{5}{2}\overline{v}\tau \right)^2=\left(\frac{1}{2}\overline{v}\tau \right)^2,\mathrm{\hspace{0.33em}\hspace{0.33em}0}<\tau <t$$
(38)
where $`\overline{v}`$ is the velocity of the rear front. For the triangular avalanches the shape is given by the envelope of dilating ellipses
$$\left(\frac{\overline{v}x}{2v_{}}\right)^2+\left(y\frac{3}{2}\overline{v}t\right)^2=\left(\frac{1}{2}\overline{v}t\right)^2.$$
(39)
Here $`v_{}`$ is perpendicular velocity. While these heuristic relations are consistent with experimental observation, see Fig. 27, their connection to continuum dynamical models of granular flows remains to be understood.
### VI.2 Statistics of avalanches and sandpile model
It is well known that in real sandpiles avalanches can vary widely in size. The wide distribution of scales in real avalanches stimulated Bak *et al.* (1987) to introduce a “sandpile cellular automaton” as a paradigm model for the self-organized criticality, the phenomenon which occurs in slowly driven non-equilibrium spatially extended systems when they asymptotically reach a critical state characterized by a power-law distribution of event sizes. The set of rules which constitute the sandpile model is very simple. Unit size “grains” are dropped one by one on a one-dimensional lattice in random places and form vertical stacks. If a local slope (the difference between heights of two neighboring stacks) exceeds a certain threshold value, a grain hops from the higher to the lower stack. This may trigger an “avalanche” of subsequent hops until the sandpile returns to the stable state. After that another grain is dropped and the relaxation process repeats. The size of an avalanche is determined by the number of grains set into motion by adding a single grain to a sandpile. This model asymptotically reaches a critical state in which the mean angle is equal to the critical slope, and avalanches have a universal power-law distribution of sizes, $`P(s)s^\alpha `$ with $`\alpha 1.5`$.
The relevance of this model and its generalizations to the real avalanches is still the matter of debate. The sandpile model by Bak *et al.* (1987) is defined via a single repose angle and so its asymptotic behavior has the properties of the critical state for a second-order phase transition. Real sandpiles are characterized by two angles of repose and thus exhibit features of the first-order phase transition. Moreover, concept of self-organized criticality is related to a power-law distribution of avalanche sizes, thus reliable experimental verification of self-organized criticality requires accumulation of very large statistics of avalanche events and a large-scale experimental setup. Finite size effects should strongly affect the power-law behavior.
Experiments with avalanches in slowly rotating drums Jaeger *et al.* (1989); Rajchenbach (2000) and chute flows Lemieux and Durian (2000) do not confirm the scale-invariant distribution of avalanches. In other experiments with large mono-disperse glass beads dropped on a conical sandpile Costello *et al.* (2003) claimed existence of the self-organized criticality with $`\alpha 1.5`$. Characteristics of the size distribution depended on the geometry of the sandpile, physical and geometrical properties of grains, and the way the grains are dropped on the pile, contrary to the universal concept of self-organized critical behavior. Self-organized criticality was also claimed in the avalanche statistics in three-dimensional pile of anisotropic grains (long rice), however a smaller scaling exponent $`\alpha 1.2`$ was measured for the avalanche size distribution Aegerter *et al.* (2004). Interestingly, rice piles were observed to demonstrate roughening dynamics of their surface as the distribution of active sites in the self-organized critical state shows a self-affine structure with the fractal exponent $`d_B=1.85`$ Aegerter *et al.* (2004). This is consistent with the theoretically predicted mapping between self-organized criticality and roughening observed for example in Kardar-Parizi-Zhang model Paczuski and Boettcher (1996).
One can argue that real sandpiles should not exhibit self-organized criticality in a strict sense due to hysteresis and the existence of two critical repose angles. However, since the difference between the angles is relatively small, one cannot exclude power-law type behavior in the finite range of avalanche sizes. This circumstance possibly explains significant scatter in experimental results and scaling exponents for avalanche size distribution and the dependence on grain shape and material properties.
### VI.3 Instabilities in granular chute flows
Granular chute flows exhibit a variety of pattern-forming instabilities, including fingering Pouliquen *et al.* (1997); Malloggi *et al.* (2005a), longitudinal vortices Forterre and Pouliquen (2001, 2002); Börzsönyi and Ecke (2005), long surface waves Forterre and Pouliquen (2003), segregation and stratification Gray and Hutter (1997a); Makse *et al.* (1997b), etc.
Pouliquen *et al.* (1997) studied experimentally a granular chute flow on a rough inclined plane. Experiments performed with polydisperse sand particles demonstrated fingering instability of the front propagating down the slope, similar to that observed in fluid films flowing down inclined plane Zhou *et al.* (2005); Troian *et al.* (1989). However, similar experiments with smooth monodisperse glass beads exhibited no instability. The authors argued that the instability was due to a flow-induced size segregation in a polydisperse granular matter. The segregation indeed was found near the avalanche front. However, similar experiments Malloggi *et al.* (2005a, b) showed a fingering front instability without a significant size segregation. Thus, the question of the mechanism of fingering instability is still open.
Experiments by Forterre and Pouliquen (2001, 2002); Börzsönyi and Ecke (2005) show the development of longitudinal vortices in rapid chute flows, see Fig. 8. The vortices develop for large inclination angles and large flow rates in the regime of accelerating flow when the flow thickness decreases and the mean flow velocity increases along the chute. Forterre and Pouliquen (2001) proposed an explanation of this phenomenon in terms of granular “thermoconvection”. Namely, rapid granular flow has a high shear near the rough bottom which leads to the local increase of granular temperature and consequently creates a density inversion. In turn, the density inversion trigger a convection instability similar to that in ordinary fluids. The critical instability wavelength $`\lambda _C`$ is determined by the depth of the layer $`h`$ (in experiment $`\lambda _c3h`$).
In a subsequent work Forterre and Pouliquen (2002) studied the formation of longitudinal vortices and the stability of granular chute flows in the framework of granular hydrodynamics Eqs.(5)-(7). The inverse density profile appears when a heuristic boundary condition at the bottom relating slip velocity and heat flux is introduced. Steady-state solution of Eqs.(5)-(7) indeed yields an inverse density profile (Fig. 28) which turns out to be unstable with respect to short-wavelength perturbations for large flow velocities, see Fig. 29. While the linear stability analysis captured many important features of the phenomenon, there are still open questions. The stability analysis was performed for the steady flow whereas the instability occurs in the regime of accelerating flow. Possibly due to this assumption the linear stability analysis yielded oscillatory instability near the onset of vortices, whereas for the most part, the vortices appear to be steady. Another factor which is ignored in the theory is the air drag. The high flow velocity in the experiment (about 1-2 m/sec) is of the order of the terminal velocity of an individual grain in air, and therefore air drag may affect the granular flow.
Forterre and Pouliquen (2003) presented an experimental study of the long-surface-wave instability developing in granular flows on a rough inclined plane, Fig. 30. This instability was known from previous studies Savage (1979); Davies (1990), however no precise characterization of the instability had been performed. Forterre and Pouliquen (2003) measured the threshold and the dispersion relation of the instability by imposing a controlled perturbation at the entrance of the flow and measuring its evolution down the slope, see Fig. 31. The results are compared with the prediction of a linear stability analysis conducted in the framework of depth-averaged Saint-Venant-type equations similar to those described in Sec. VI.1.2:
$`_th`$ $`+`$ $`_x(uh)=0`$ (40)
$`_t(uh)`$ $`+`$ $`\alpha _x(u^2h)=\left(\mathrm{tan}\theta \mu (u,h)_xh\right)gh\mathrm{cos}(\theta )`$
where $`h`$ is local thickness, $`\theta `$ is inclination angle, $`u`$ is depth-averaged flow velocity, $`\mu (h,u)`$ is a function describing effective depth and velocity dependent bottom friction, $`\alpha O(1)`$ is a constant determined by the velocity profile within the layer. According to Forterre and Pouliquen (2003), the instability is similar to the long-wave instability observed in classical fluids but with characteristics that can dramatically differ due to the specificity of the granular rheology. The theory is able to predict quantitatively the stability threshold and the phase velocity of the waves but fails to describe the observed cutoff of the instability at high wavenumbers. Most likely, one needs to include higher order terms, such as $`_x^2h`$ in the first Eq. (40) in order to account for the cutoff.
The order parameter theory based on Eqs. (31),(32) also reproduces the long-surface wave instability. Furthermore, linearizing Eqs. (31),(32) near the steady flowing solution $`A=A_0+\stackrel{~}{a}\mathrm{exp}[\sigma t+ikx],h=h_0+\stackrel{~}{h}\mathrm{exp}[\sigma t+ikx]`$, after simple algebra one obtains that the growth rate of linear perturbations $`\sigma `$ is positive only in a band restricted by some critical wavenumber and only in the vicinity of $`h_{stop}`$, see Fig. 32. With the increase of $`h`$, i.e. the granular flux, the instability disappears, in agreement with experiments. The nonlinear saturation of the instability results in the development of a sequence of avalanches, which is generally non-periodic, see Fig. 33. The structure shows slow coarsening due to merging of the avalanches. This instability is a possible candidate mechanism of the formation of inhomogeneous deposit structure behind the front of an avalanche.
Conway *et al.* (2003) studied free-surface waves in granular chute flows near a frictional boundary. The experiments showed that the sub-boundary circulation driven by the velocity gradient plays an important role in the pattern formation, suggesting a similarity between wave generation in granular and fluid flows.
A Kelvin-Helmholtz-like shear instability in chute flows was observed by Goldfarbs *et al.* (2002), when two streams of sand flowing on an inclined plane with different velocities were in side-by-side contact with each other. For sufficiently high chute angles and shear rates the interface remains flat. The instability of the interface develops when the chute angle and/or the shear rate is reduced. This instability has been reproduced in soft-particle molecular dynamics simulations by Ciamarra *et al.* (2005) who also observed that in a polydisperse medium this instability leads to grain segregation (See below Sect. VII).
### VI.4 Pattern-forming instabilities in rotating cylinders
Granular media in rotating horizontal cylinders (drums) often show behavior similar to chute flows. For very small rotation rates (as defined by small Froude number $`Fr=\omega ^2R/g`$, where $`\omega `$ is the angular velocity of drum rotation and $`R`$ its radius), well separated in time avalanches occur when the slope of the free surface exceeds a certain critical angle $`\theta _c`$ whereby diminishing this angle to a smaller static repose angle $`\theta _s`$ Rajchenbach (1990); Jaeger *et al.* (1989); Tegzes *et al.* (2002, 2003). The difference between $`\theta _c`$ and $`\theta _s`$ is usually a few degrees. At an intermediate rotation speed, a continuous flow of sand emerges instead of discrete avalanches through a hysteretic transition, similar to the transition in chute flows at large rates of grain deposition Lemieux and Durian (2000). In the bulk, the granular material rotates almost as a solid body with some internal slipping. As moving grains reach the free surface they slide down within a thin near-surface layer Zik *et al.* (1994) (see sketch in Fig. 9). The surface has a nearly flat shape; the arctangent of its average slope defines the so-called dynamic angle of repose $`\theta _d`$.
There are various models addressing the nature of the transition from discrete avalanches to the continuum flow. Linz and Hänggi (1995) proposed a phenomenological model based on a system of equation for the angle of repose $`\varphi `$ and mean flow velocity $`v`$
$`\dot{v}`$ $`=`$ $`g(\mathrm{sin}\varphi k(v)\mathrm{cos}\varphi ))\chi (\varphi ,v)`$
$`\dot{\varphi }`$ $`=`$ $`\overline{\omega }av`$ (41)
where $`\overline{\omega }`$ is the rotation frequency of the drum, $`k(v)=b_0+b_2v^2`$ is the velocity dependent friction coefficient, $`\chi (\varphi ,v)`$ is some cut-off function, and $`a,b_0,b_2`$ are parameters of the model. Despite the simplicity, the model yields qualitatively correct transition from discrete avalanches to continuous flow with the increase of rotation rate $`\overline{\omega }`$, and also predicts logarithmic relaxation of the free surface angle in the presence of vibration.
The transition from avalanches to flow naturally arises in the framework of the partial fluidization theory, Aranson and Tsimring (2002). In this case one can derive a system of coupled equations for the parameter $`\delta `$ (which is related to the surface local angle $`\varphi `$, see Eq. (30)) and the width of fluidized layer $`z_0`$,
$`_tz_0`$ $`=`$ $`_s^2z_0+F(z_0,\delta )\overline{v}_sz_0`$
$`_t\delta `$ $`=`$ $`\overline{\omega }+_s^2J`$ (42)
where $`s`$ is the coordinate along the slope of the granular surface inside the drum, $`J=f(z_0)`$ is the downhill flux of grains, $`\overline{v}`$ is averaged velocity in flowing layer, and functions $`F,f`$ and $`v_0`$ are derived from Eq. (29). This model bears resemblance to the BCRE-type models of surface granular flows which were applied to rotating drums by Khakhar *et al.* (1997); Makse (1999).
Eqs. (42) exhibit stick-slip type oscillations of the surface angle for slow rotation rates and a hysteretic transition to a steady flow for larger rates. Eqs. (42) yields the following scaling for the width of the flowing layer $`z_0`$ in the middle of the drum vs rotation frequency: $`z_0\overline{\omega }^{2/3}`$, which is consistent with experiment Tegzes *et al.* (2002, 2003). After integration over $`s`$ Eqs. (42) can be reduced to a system of two coupled equations for averaged drum angle $`\delta `$ and averaged flow thickness $`z_0`$ somewhat similar to the model of Linz and Hänggi (1995).
Granular flows in long rotating drums under certain conditions also exhibit fingering instability Shen (2002); Fried *et al.* (1998). Similarity between fingering in rotating drums and chute flows Forterre and Pouliquen (2002) suggests that mechanisms described in the Section VI.3 can be responsible for this effect, see also Section VII.
## VII Models of granular segregation
One of the most fascinating features of heterogeneous (i.e., consisting of different distinct components) granular materials is their tendency to segregate under external agitation rather than to mix, as one would expect from the naive entropy consideration. This property is ubiquitous in Nature (see, e.g. Iverson (1997)) and has important technological implications Cooke *et al.* (1976). In fact, some aspects of segregation of small and large particles can be understood on equilibrium thermodynamics grounds Asakura and Oosawa (1958). Since the excluded volume for small particles around large ones becomes smaller when large grains clump together, separated state possesses lower entropy. However, granular systems are driven and strongly dissipative, and this simple equilibrium argument can only be applied qualitatively. The granular segregation is more widespread than it would be dictated by thermodynamics. In fact, any variation in mechanical properties of particles (size, shape, density, surface roughness, etc.) may lead to their segregation. At least for bi-disperse rapid dilute flows the granular segregation can be rigorously treated in the framework of kinetic theory of dissipative gases, see Subsec. III.1. Jenkins and Yoon (2002) employed kinetic theory for a binary mixture for spheres or disks in gravity and derived a simple segregation criterion based on the difference of partial pressures for each type of particles due to the difference in size and/or mass.
Segregation has been observed in most flows of granular mixtures, including granular convection Knight *et al.* (1993), hopper flows Makse *et al.* (1997b); Gray and Hutter (1997a); Samadani *et al.* (1999); Samadani and Kudrolli (2001), flows in rotating drums Zik *et al.* (1994); Hill (1997); Choo *et al.* (1997), and even in freely cooling binary granular gases Catuto and Marconi (2004). Segregation among large and small particles due to shaking has been termed “Brazil nut effect” Rosato *et al.* (1987). The phenomenon of granular segregation was discovered long time ago, and several “microscopic” mechanisms have been proposed to explain its nature, including inter-particle collisions Brown (1939), percolation Williams (1976), and others. In certain cases, separation of grains produces interesting patterns. For example, if a binary mixture of particles which differ both in size and in shape is poured down on a plane, a heap which consists of thin alternating layers of separated particles is formed Gray and Hutter (1997a); Makse *et al.* (1997b), see Fig. 34. Rotating of mixtures of grains with different sizes in long drums produces well separated bands of pure mono-disperse particles Zik *et al.* (1994); Hill (1997); Choo *et al.* (1997); Chicarro *et al.* (1997), Fig. 11. In this Section we only address models of pattern formation due to segregation (stratification and banding), without discussing other manifestations of granular segregation.
### VII.1 Granular stratification
Granular stratification occurs when a binary mixture of particles with different physical properties is slowly poured on a plate Gray and Hutter (1997a); Makse *et al.* (1997b); Koeppe *et al.* (1998). More specifically, it occurs when larger grains have additionally larger roughness resulting in a larger repose angle and the flux of falling particles is small enough to cause intermittent avalanches down the slopes of the heap. The basic mechanism of stratification is related to the avalanches acting as kinetic sieves Savage (1988, 1993). During an avalanche, voids are continuously being created within flowing near-surface layer, and small particles are more likely to fall into them. This creates a downward flux of smaller particles which is compensated by the upward flux of larger particles in order to maintain a zero total particle flux across the flowing layer. Other models of granular segregation in a thin flowing layer Khakhar *et al.* (1997); Dolgunin *et al.* (1998); Khakhar *et al.* (1999) lead to a similar result. Each avalanche leads to the formation of a new pair of layers in which the grains of different sorts are separated (see Fig. 34). This pair of layers grows from the bottom of the pile by upward propagation of a kink at which small particles are stopped underneath large ones. However, when the larger particles were smooth and small particles were rough, instead of stratification only large scale segregation with small particles near the top and large particles near the bottom was observed.
Makse *et al.* (1997b, a) proposed a cellular automata model which generalized the classical sandpile model Bak *et al.* (1987) (see Section VI.2). In this model, a sandpile is built on a lattice, and rectangular grain have identical horizontal size but different heights (see Fig. 35a). Grains are released at the top of the heap sequentially, and they are allowed to roll down the slope. A particle would become rolling if the local slope (defined as the height difference between neighboring columns) exceeds the repose angle. To account for difference in grain properties, four different repose angles $`\theta _{\alpha \beta }`$ were introduced for grains of type $`\alpha `$ rolling on a substrate of type $`\beta `$ ($`\alpha ,\beta \{1,2\}`$ where 1 and 2 stand for small and large grains, respectively). Normally, $`\theta _{21}<\theta _{12}`$ because of the geometry (small grains tend to get trapped by large grains), and one-component repose angles usually lie within this range, $`\theta _{21}<\theta _{11},\theta _{22}<\theta _{12}`$. However the ratio of $`\theta _{11},\theta _{22}`$ depends on the relative roughness of the grains. For $`\theta _{21}<\theta _{11}<\theta _{22}<\theta _{12}`$ (large grains are more rough), the model yields stratification in agreement with experiment (Fig. 35b). If, on the other hand, $`\theta _{22}<\theta _{11}`$ (which corresponds to smaller grains being more rough), the model yields only large-scale segregation: large particles collect at the bottom of the sandpile.
This physical model can also be recast in the form of continuum equations Boutreux and de Gennes (1996); Makse *et al.* (1997a) which generalize the single-species BCRE model of surface granular flows Bouchaud *et al.* (1994) (see Section VI):
$`_tR_\alpha `$ $`=`$ $`v_\alpha _xR_\alpha +\mathrm{\Gamma }_\alpha ,`$ (43)
$`_th`$ $`=`$ $`{\displaystyle \underset{\alpha }{}}\mathrm{\Gamma }_\alpha ,`$ (44)
where $`R_\alpha (x,t),v_\alpha `$ are the thickness and velocity of rolling grains of type $`\alpha `$, $`h(x,t)`$ is the instantaneous profile of the sandpile, and $`\mathrm{\Gamma }_\alpha `$ characterizes interaction between the rolling grains and the substrate of static grains. In the same spirit as in the discrete model, the interaction function $`\mathrm{\Gamma }_\alpha `$ is chosen in the form
$$\mathrm{\Gamma }_\alpha =\{\begin{array}{c}\gamma _\alpha [\theta _l\theta _\alpha (\varphi _\beta )]R_\alpha \hfill \\ \gamma _\alpha \varphi _\alpha [\theta _l\theta _\alpha (\varphi _\beta )]R_\alpha \hfill \end{array}.$$
(45)
Here $`\varphi _\alpha (x,t)`$ is the volume fraction of grains of type $`\alpha `$, and $`\theta _l=_xh`$ is the local slope of the sandpile. This form of the interaction terms implies that the grains of type $`\alpha `$ become rolling if the local slope exceeds the repose angle $`\theta _\alpha (\varphi _\beta )`$ for this type on a surface with composition $`\varphi _\beta (x,t)`$. Assuming that the generalized repose angles $`\theta _\alpha (\varphi _\beta )`$ are linear functions of the concentration
$`\theta _1(\varphi _2)=(\theta _{12}\theta _{11})\varphi _2+\theta _{11},`$ (46)
$`\theta _2(\varphi _2)=(\theta _{12}\theta _{11})\varphi _2+\theta _{21}.`$ (47)
Eqs. (43)-(45) possess a stationary solution in which the heap is separated into two regions where $`\theta _2(\varphi _2)<\theta <\theta _1(\varphi _2)`$ and $`\theta <\theta _2(\varphi _2)<\theta _1(\varphi _2)`$. This solution corresponds to small grains localized near the top and small grains near the bottom with a continuous transition between the two regions. However, Makse *et al.* (1997b) showed that this stationary solution is unstable if $`\delta =\theta _{22}\theta _{11}>0`$ and gives rise to the stratification pattern.
Similar effect of stratification patterns was observed experimentally in a thin slowly rotating drum which is more than half filled with a similar binary mixture Gray and Hutter (1997a), see Fig. 36. Periodic avalanches, occurring in the drum, lead to formation of strata by the same mechanism described above.
### VII.2 Axial segregation in rotating drums
The most common system in which granular segregation is studied is a rotating drum, or a partially filled cylinder rotating around its horizontal axis (see Section VI.4). When a polydisperse mixture of grains is rotated in a drum, strong radial segregation usually occurs within just a few revolutions. Small and rough particles aggregate to the center (core) of the drum, large and smooth particles rotate around the core (see Figs. 10 and 9). Since there is almost no shear flow in the bulk, the segregation predominantly occurs within a thin fluidized near-surface layer. For long narrow drums with the length much exceeding the radius, radial segregation is often followed by the axial segregation occurring at later stages (after several hundred revolutions) when the angle of repose of small particles exceeds that of large particles. As a result of axial segregation, a pattern of well segregated bands is formed Zik *et al.* (1994); Hill (1997) (see, e.g., Fig. 11) which slowly merge and coarsen. Depending on the rotation speed, coarsening can either saturate at a certain finite bandwidth at low rotation speeds when discrete avalanches provide granular transport Frette and Stavans (1997) or at higher rotation rates in a continuous flow regime it can lead to a final state in which all sand is separated in two bands Zik *et al.* (1994); Fiodor and Ottino (2003).
The axial segregation has been well known in the engineering community, it was apparently first observed by Oyama in 1939 Oyama (1939). The mechanism of axial segregation is apparently related to the different friction properties of grains which lead to different dynamical angles of repose. The latter are defined as the angle of the slope in the drum corresponding to continuous flow regime, however in real drums the free surface has a more complicated S-shape Zik *et al.* (1994); Elperin and Vikhansky (1998); Makse (1999); Orpe and Khakhar (2001). According to Zik *et al.* (1994) (see also Levine (1999)), if there is a local increase in concentration of particles with higher dynamic repose angle, the local slope there will be higher, and that will lead to a local bump near the top of the free surface and a dip near the bottom. As the particles tend to slide along the steepest descent path, more particles with higher repose angle will accumulate in this location, and the instability will develop. Zik *et al.* (1994) proposed a quantitative continuum model of axial segregation based on the equation for the conservation equation for the relative concentration of the two components (“glass” and “sand”), $`c(z,t)=(\rho _A\rho _B)/(\rho _A+\rho _B)`$,
$$_tc=\frac{C}{\rho _T}(\mathrm{tan}\theta _A\mathrm{tan}\theta _B)_z(1c^2)(1+y_x^2)\frac{y_z}{y_x}.$$
(48)
Here $`x`$ and $`z`$ are Cartesian horizontal coordinates across and along the axis of the drum, $`y(x,z,t)`$ describes the instantaneous free surface inside the drum, $`\rho _T=\rho _A+\rho _B`$, $`C`$ is a constant related to gravity and effective viscosity of granular material in the flowing layer. The term in angular brackets denotes the axial flux of the glass beads averaged over the cross-section of the drum. The profile of the free surface in turn should depend on $`c(z,t)`$. If $`(1+y_x^2)y_c/y_x<0`$, linearization of Eq.(48) leads to the diffusion equation with negative diffusion coefficient which exhibits segregation instability with growth rate proportional to the square of the wavenumber. It is easy to see that the term in angular brackets vanishes for a straight profile $`y_x=const(x)`$. However, for the experimentally observed S-shaped profile of the free surface Zik *et al.* (1994) calculated that the instability condition is satisfied when the drum is more than half full. While experiments show that axial segregation in fact observed even for less than 50% filling ratio, the model gives a good intuitive picture for the mechanism of the instability.
Recent experiments Hill (1997); Choo *et al.* (1997, 1998); Fiodor and Ottino (2003) have revealed interesting new features of axial segregation. Hill (1997) performed magnetic resonance imaging studies Hill (1997) which demonstrated that in fact the bands of larger particles usually have a core of smaller particles. More recent experiments by Fiodor and Ottino (2003) showed that small particles formed a shish kebab-like structure with bands connected by a rod-like core, while large particles formed disconnected rings. Choo *et al.* (1997, 1998) found that at early stages, the small-scale perturbations propagate across the drum in both directions (this was clearly evidenced by the experiments on the dynamics of pre-segregated mixtures Choo *et al.* (1997)), while at later times more long-scale static perturbations take over and lead to the emergence of quasi-stationary bands of separated grains (see Fig. 37). The slow coarsening process can be accelerated in a drum of a helical shape Zik *et al.* (1994). Alternatively, the bands can be locked in an axisymmetrical drum with the radius modulated along the axis Zik *et al.* (1994).
In order to account for the oscillatory behavior of axial segregation at the initial stage, Aranson *et al.* (1999b); Aranson and Tsimring (1999) generalized the model of Zik *et al.* (1994). The key assumption was that besides the concentration difference, there is an additional slow variable which is involved in the dynamics. Aranson *et al.* (1999b); Aranson and Tsimring (1999) conjectured that this variable is the instantaneous slope of the granular material (dynamic angle of repose) which unlike Eq.(48) is not slaved to the relative concentration $`c`$, but obeys its own dynamics. The equations of the model read
$`_tc=_z(D_zc+g(c)_z\theta ),`$ (49)
$`_t\theta `$ $`=`$ $`\alpha (\mathrm{\Omega }\theta +f(c))+D_\theta _{zz}\theta +\gamma _{zz}c.`$ (50)
The first term in the r.h.s. of Eq.(49) describes diffusion flux (mixing), and the second term describes differential flux of particles due to the gradient of the dynamic repose angle. This term is equivalent to the r.h.s. of Eq.(48) with a particular function $`g(c)=G_0(1c^2)`$. For simplicity, the constant $`G_0`$ can be eliminated by rescaling of distance $`xx/\sqrt{G_0}`$. The sign $`+`$ before this term means that the particles with the larger static repose angle are driven towards greater dynamic repose angle. This differential flux gives rise to the segregation instability. Since this segregation flux vanishes with $`g(c)`$ $`|c|1`$ (which correspond to pure $`A`$ or $`B`$ states), it provides a natural saturation mechanism for the segregation instability.
Parameter $`\mathrm{\Omega }`$ in the second equation is the normalized angular velocity of the drum rotation, and $`f(c)`$ is the static angle of repose which is an increasing function of the relative concentration Koeppe *et al.* (1998) (for simplicity it can be assumed linear, $`f(c)=F+f_0c`$). The constant $`F`$ can be eliminated by the substitution $`\theta \theta F`$. The first term in the r.h.s. of Eq.(50) describes the local dynamics of the repose angle ($`\mathrm{\Omega }`$ increases the angle, and $`\theta +f(c)`$ describes the equilibrating effect of the surface flow), and the term $`D_\theta _{xx}\theta `$ describes axial diffusive relaxation. The last term, $`\gamma _{xx}c`$, represents the lowest-order non-local contribution from an inhomogeneous distribution of $`c`$ (the first derivative $`_xc`$ cannot be present due to reflection symmetry $`xx`$). This term gives rise to the transient oscillatory dynamics of the binary mixture.
Linear stability analysis of a homogeneous state $`c=c_0;\theta _0=\mathrm{\Omega }+f_0c_0`$ reveals that for $`g_0f_0>\alpha D`$ long-wave perturbations are unstable, and if $`g_0\gamma >(D_\theta D)^2/4`$, short-wave perturbations oscillate and decay (two eigenvalues $`\lambda _{1,2}`$ are complex conjugate with negative real part), see Fig. 38. This agrees with the general phenomenology observed by Choo *et al.* (1997) both qualitatively and even quantitatively (Fig. 38b). The results of direct numerical solution of the full model (49),(50) are illustrated by Fig. 39. It shows that short-wave initial perturbations decay and give rise to more long-wave non-oscillatory modulation of concentration which eventually leads to well-separated bands. At long times (Fig. 39b) bands exhibit slow coarsening with the number of bands decreasing logarithmically with time (see also Frette and Stavans (1997); Levitan (1998); Fiodor and Ottino (2003)). This scaling follows from the exponentially weak interaction between interfaces separating different bands Aranson and Tsimring (1999); Fraerman *et al.* (1997).
While these continuum models of axial segregation showed a good qualitative agreement with the data, recent experimental observations demonstrate that the theoretical understanding of axial segregation is far from complete Ottino and Khakhar (2000). The interpretation of the second slow variable as the local dynamic angle of repose implies that in the unstable mode the slope and concentration modulation should be in phase, whereas in the decaying oscillatory mode, these two fields have to be shifted in phase. Further experiments Khan *et al.* (2004) showed that while the in-phase relationship in the asymptotic regime holds true, the quadrature phase shift in the transient oscillatory regime is not observed. That lead Khan *et al.* (2004) to hypothesize that some other slow variable other than the angle of repose (possibly related to the core dynamics) may be involved in the transient dynamics. However, so far experiments failed to identify which second dynamical field is necessary for oscillatory transient dynamics, so it remains an open problem. Another recent experimental observation by Khan and Morris (2005) suggested that instead of normal diffusion assumed in Eqs.(49),(50), a slower subdiffusion of particles in the core takes place, $`rt^\gamma `$ with the scaling exponent $`\gamma `$ close to 0.3. The most plausible explanation is that the apparent subdiffusive behavior is in fact a manifestation of nonlinear concentration diffusion which can be described by equation
$$_tc=_zD(c)_zc.$$
(51)
For example, for the generic concentration-dependent diffusion coefficient $`Dc`$, the asymptotic scaling behavior of the concentration $`c(z,t)`$ is given by the self-similar function $`cF(z/t^\alpha )/t^\alpha `$ for $`t\mathrm{}`$ with the scaling exponent $`\alpha =1/3`$ close to 0.3 observed experimentally. Experimentally observed scaling function $`F(x/t^\alpha )`$ appears to be consistent with that of Eq. (51) except for the tails of the distribution where $`c0`$ and the assumption $`Dc`$ is possibly violated. Normal diffusion behavior corresponding to $`D=const`$ and $`\alpha =1/2`$ is in strong disagreement with the experiment.
Newey *et al.* (2004) conducted studies of axial segregation in ternary mixtures of granular materials. It was found that for certain conditions bands of ternary mixtures oscillate axially. In contrast to the experiments of Choo *et al.* (1997, 1998), the oscillations of bands appear spontaneously from initially mixed state, which strongly indicates the supercritical oscillatory instability. While in binary mixtures the oscillations have the form of periodic mixing/demixing of bands, in the ternary mixtures the oscillations are in the form of periodic band displacements. It is likely that the mechanism of band oscillations in ternary mixtures is very different from that of binary mixtures. One of possible explanations could be that the third mixture component provides an additional degree of freedom necessary for oscillations. To demonstrate that we write phenomenological equations for the concentration differences $`C_A=c_1c_2`$ and $`C_B=c_2c_3`$, where $`c_{1,2,3}`$ are the individual concentrations. By analogy with Eq. (49) we write the system of coupled equations for the concentration differences $`C_{A,B}`$ linearized near the fully mixed state:
$`_tC_A`$ $`=`$ $`D_A_z^2C_A+\mu _A_z^2C_B,`$
$`_tC_B`$ $`=`$ $`D_B_z^2C_B+\mu _B_z^2C_A.`$ (52)
If the cross-diffusion terms have opposite signs, i.e. $`\mu _A\mu _B<0`$, the concentrations $`C_{A,B}`$ will exhibit oscillations in time and in space. Obviously this mechanism is intrinsic to ternary systems and has no counterpart in binary mixtures.
Parallel to the theoretical studies, molecular dynamics simulations have been performed Shoichi (1998); Rapaport (2002); Taberlet *et al.* (2004). Simulations allowed researchers to probe the role of material parameters which would be difficult to access in laboratory experiments. In particular, Rapaport (2002) addressed the role of particle-particle and wall-particle friction coefficients separately. It was found that the main role is played by the friction coefficients between the particles and the cylinder walls: if the friction coefficient between large particles and the wall is greater than that for smaller particles, the axial segregation always occur irrespective of the ratio of particle-particle friction coefficients. However, if the particle-wall coefficients are equal, the segregation may still occur if the friction among large particles is greater than among small particles. Taberlet *et al.* (2004) studied axial segregation in a system of grains made of identical material differing only by size. The simulations revealed rapid oscillatory motion of bands, which is not necessarily related to the slow band appearence/disappearence observed in experiments of Choo *et al.* (1997, 1998); Fiodor and Ottino (2003).
A different type of discrete element modelling of axial segregation was proposed by Yanagita (1999). This model builds upon the lattice-based sandpile model and replaces a rotating drum by a three-dimensional square lattice. Drum rotation is modelled by correlated displacement of particles on the lattice: particles in the back are shifted upward by one position, and the particles at the bottom are shifted to fill the voids. This displacement steepens the slope of the free surface, and once it reaches a critical value, particles slide down according to the rules similar to the sandpile model of Bak *et al.* (1987) but taking into account different critical slopes for different particles. This model despite its simplicity reproduced both radial and axial segregation patterns and therefore elucidated the critical components needed for adequate description of the phenomenon.
### VII.3 Other examples of granular segregation
As we have seen in the previous Section, granular segregation occurs in near-surface shear granular flows, such as in silos, hoppers, and rotating drums. However, other types of shear granular flows may also lead to segregation. For example, Taylor-Couette flow of granular mixtures between two rotating cylinders leads to formation of Taylor vortices and then in turn to segregation patterns Shinbrot (2004), see Fig. 40.
Pouliquen *et al.* (1997) observed granular segregation in a thin granular flow on an inclined plane. In this case, segregation apparently occurs as a result of an instability in which concentration mode is coupled with hydrodynamic mode. As a result, segregation occurs simultaneously with a fingering instability of the avalanche front (Fig. 7). As an implicit evidence of this relation between segregation and fingering instability, Pouliquen *et al.* (1997) found that mono-disperse granular material does not exhibit fingering instability. However, other experiments Shen (2002) indicate that in other conditions (more rapid flows), fingering instability may occur even in flows of mono-disperse granular materials. Thus, the segregation is likely a consequence rather than the primary cause of the fingering instability.
An interesting recent example of pattern formation caused by granular segregation in a horizontally shaken layer of binary granular mixture was presented by Mullin (2000, 2002); Reis and Mullin (2002). After several minutes of horizontal shaking with frequency 12.5 Hz and displacement amplitude 1 mm (which corresponds to the acceleration amplitude normalized by gravity $`\mathrm{\Gamma }=0.66`$), stripes were formed orthogonal to the direction of shaking. The width of the stripes was growing continuously with time as $`dt^{0.25}`$, thus indicating slow coarsening (Fig. 5). This power law is consistent with the diffusion-mediated mechanism of stripe merging. Reis and Mullin (2002) argued on the basis of experimental results on patterned segregation in horizontally shaken layers that the segregation bears features of the second-order phase transition. Critical slow-down was observed near the onset of segregation. The order parameter is associated with the combined filling fraction $`C`$, or the layer compacity,
$$C=\frac{N_sA_s+N_lA_l}{S}$$
(53)
where $`N_{s,l}`$ are numbers of particles in each species, $`A_{s,l}`$ are projected two-dimensional areas of the respective individual particles, and $`S`$ is the tray area. Ehrhard *et al.* (2005) proposed a simple numerical model to describe this phenomenon of segregation in horizontally vibrated layers. The model is based on a two-dimensional system of hard disks of mass $`m_\alpha `$ and radius $`R_\alpha `$ ($`\alpha =1,2`$ denote the species)
$$m_\alpha \dot{𝐯}_{\alpha i}=\gamma _i\left(𝐯_{\alpha i}𝐯_{tray}(t)\right)+\zeta _{\alpha i}(t)$$
(54)
where $`𝐯_i`$ is the particles velocity $`𝐯_{tray}(t)=A_0\mathrm{sin}(\omega t)`$ is oscillating tray velocity, $`\gamma `$ provides linear damping, and $`\zeta _{\alpha i}`$ is Gaussian white noise acting independently on each disk. The model reproduced segregation instability and subsequent coarsening of stripes. More realistic discrete element simulations were recently performed by Ciamarra *et al.* (2005). In these simulations a binary mixture of round disks of identical sizes but two different frictions with the bottom plate (in fact, velocity-dependent viscous drag was assumed), separated in alternating bands perpendicular to the oscillation direction irrespectively on initial conditions: both random mixed state and separated along the direction of oscillations state were used. Using particles of the same size eliminated the thermodynamic “excluded volume” mechanism for segregation, and the authors argued that the mechanism at work is related to the dynamical shear instability similar to the Kelvin-Helmholtz instability in ordinary fluids. It was confirmed by a numerical observation of the interfacial instability when two monolayers of grains with different friction constant were placed in contact along a flat interface parallel to the direction of horizontal oscillations. Similar instability is apparently responsible for ripple formation Scherer *et al.* (1999); Stegner and Wesfreid (1999).
Pooley and Yeomans (2004) proposed theoretical description of this experiment based on continuum model for periodically-driven two isothermal ideal gases which interact through frictional force. It was shown analytically that segregated stripes form spontaneously above critical forcing amplitude. While the model reproduces the segregation instability, apparently it does not exhibit coarsening of stripes observed in the experiment. Moreover, applicability of the isothermal ideal gas model to this experiment where the particles are almost at rest is an open question.
Similar coarsening effect in granular segregation in a particularly simple geometry was studied by Aumaître *et al.* (2001). They investigated the dynamics of a monolayer of grains of two different sizes in a dish shaken in a horizontal “swirling” motion. They observed that large particles tend to aggregate near the center of the cavity surrounded by small particles. The qualitative explanation of this effect follows from simple thermodynamic considerations (see above). Indeed, direct tracing of particle motion showed that the pressure in the area near the large particles is smaller than outside. But small particles do not follow the gradient of pressure and assemble near the center of the cavity because this gradient is counterbalanced by the force from large particles. The inverse of force acting on large particles leads to their aggregation near the center of the cavity. Aumaître *et al.* (2001) proposed a more quantitative model of segregation based on the kinetic gas theory and found satisfactory agreement with experimental data.
Burtally *et al.* (2002) studied spontaneous separation of vertically vibrated mixtures of particles of similar sizes but different densities (bronze and glass spheres). At low frequencies and at sufficient vibrational amplitudes, a sharp boundary between the lower layer of glass beads and the upper layer of the heavier bronze spheres was observed. At higher frequencies, the bronze particles emerge as a middle layer separating upper and lower glass bead layers. The authors argue that the effect of air on the granular motion is a relevant mechanism of particle separation. A somewhat similar conclusion was achieved by Möbius *et al.* (2001) in experiments with vertically-vibrated column of grains containing a large “intruder” particle.
Fiodor and Ottino (2003); Arndt *et al.* (2005) performed detailed experiments on axial segregation in slurries, or bi-disperse grain-water mixtures. A mixture of two types of spherical glass beads of two sizes were placed in a water-filled tube at the volume ratio 1:2. Authors found that both rotation rate and filling fraction play an important role in band formation. Namely, bands are less likely to form at lower fill levels (20-30%) and slower rotation rates (5-10 rpm). They mostly appear near the ends of the drum. At higher fill levels and rotation rates, bands form faster, and there are more of them throughout the drum. Fiodor and Ottino (2003); Arndt *et al.* (2005) also studied the relation between the bands visible on the surface, and the core of small beads, and found that for certain fill levels and rotation speeds, the core remains prominent at all times, while in other cases the core disappears completely between bands of small particles. They also observed an interesting oscillatory instability of interfaces between bands at high rotation speeds. All these phenomena still await theoretical modelling.
## VIII Granular materials with complex interactions
### VIII.1 Patterns in solid-fluid mixtures
Presence of interstitial fluid significantly complicates the dynamics of granular materials. Hydrodynamic flows lead to the viscous drag and anisotropic long-range interaction between particles. Even small amounts of liquid leads to cohesion among the particles which can have a profound effect on macroscopic properties of granular assemblies such as angle of repose, avalanching, ability to segregate, etc. (see for example Sec. VII.2 and Samadani and Kudrolli (2000, 2001); Tegzes *et al.* (2002); Li and McCarthy (2005)).
In this Section we will discuss the case when the volume fraction of fluid in the two-phase system is large, and the grains are completely immersed in fluid. This is relevant for many industrial applications, as well as for geophysical problems such as sedimentation, erosion, dune migration, etc.
One of the most technologically important examples of particle-laden flows is a fluidized bed. Fluidized beds have been widely used since German engineer Fritz Winkler invented the first fluidized bed for coal gasification in 1921. Typically, a vertical column containing granular matter is energized by a flow of gas or liquid. Fluidization occurs when the drag force exerted by the fluid on the granulate exceeds gravity. A uniform fluidization, the most desirable regime for most industrial applications, turns out to be prone to bubbling instability: bubbles of clear fluid are created spontaneously at the bottom, traverse the granular layer and destroy the uniform state Jackson (2000). Instabilities in fluidized beds is an active area of research in the engineering community, see Jackson (2000); Gidaspow (1994); Kunii and Levenspiel (1991). A shallow fluidized bed shows many similarities with mechanically vibrated layers, see Section V. In particular, modulations of airflow studied by Li *et al.* (2003) result in formation of subharmonic square and stripe patterns (see Fig. 23) similar to those in mechanically-vibrated systems Melo *et al.* (1994, 1995); Umbanhowar *et al.* (1996).
Wind and water driven granular flows play important roles in geophysical processes. Wind-blown sand forms dunes and beaches. The first systematic study of airborne (or aeolian) sand transport was conducted by R. Bagnold during Wold War II, see Bagnold (1954) who identified two primary mechanisms of sand transport: saltation and creep, and proposed the first empiric relation for the sand flux $`q`$ driven by wind shear stress $`\tau `$:
$$q=C_B\frac{\nu _a}{g}\sqrt{\frac{d}{D}}u_{}^3$$
(55)
where $`C_B=const`$, $`\nu _a`$ is air density, $`d`$ is the grain diameter, $`D=0.25`$ mm is a reference grain size, and $`u_{}=\sqrt{\tau /\nu _a}`$ is wind friction velocity. Later many refinements of Eq.(55) were proposed, see e.g. Pye and Tsoar (1991).
Nishimori and Ouchi (1993) proposed a simple theory which describes formation of ripples as well as dunes. The theory is based on a lattice model which incorporates separately saltation and creep processes. The model operates with the height of sand at each lattice side at discrete time $`n`$, $`h_n(x,y)`$. The full time step includes two substeps. The saltation substep is described as
$`\overline{h}_n(x,y)`$ $`=`$ $`h_n(x,y)q`$ (56)
$`\overline{h}_n(x+L(h(x,y)),y)`$ $`=`$ $`h_n(x+L(h_n(x,y)),y)+q`$
where $`q`$ is the height of grains being transferred from one (coarse grained) position $`(x,y)`$ to the other position $`(x+L,y)`$ on the lee side (wind is assumed blowing in the positive $`x`$ direction), $`L`$ is the flight length in one saltation which characterizes the wind strength. It is assumed that $`q`$ is conserved. Since the saltation length $`L`$ depends on multiple factors, the following simple approximation is accepted
$$L=L_0+bh_n(x,y)$$
(57)
with $`L_0`$ measuring wind velocity and $`b=const`$. The creep substep involves spatial averaging over neighboring sites in order to describe the surface relaxation due to gravity,
$`h_{n+1}(x,y)=\overline{h}_n(x,y)+`$ (58)
$`D\left[{\displaystyle \frac{1}{6}}{\displaystyle \underset{NN}{}}\overline{h}(x,y)+{\displaystyle \frac{1}{12}}{\displaystyle \underset{NNN}{}}\overline{h}(x,y)\overline{h}(x,y)\right],`$
where $`_{NN}`$ and $`_{NNN}`$ denote summation over the nearest neighbors and next nearest neighbors correspondingly, and $`D=const`$ is the surface relaxation rate. Despite its simplicity, simulation of the model reproduced formation of ripples and consequently arrays of barchan (crescent shaped) dunes, see Fig. 41. Nishimori and Ouchi (1993) found that above certain threshold an almost linear relation holds between the selected wavelength of the dune pattern and the “wind strength” $`L`$.
In the long-wave limit Eqs. (VIII.1)-(VIII.1) can be reduced to more traditional continuum models considered below.
In the continuum description of the evolution of the sand surface, the profile $`h`$ is governed by the mass conservation equation
$$\nu _s_th=𝐪,$$
(59)
where $`\nu _s`$ is the density of sand and $`𝐪`$ is the sand flux. In order to close Eqs. (55),(59), several authors proposed different phenomenological relations between shear stress at the bed surface $`\tau `$ and the height $`h`$, see e.g. Nishimori and Ouchi (1993); Kroy *et al.* (2002a, b); Hersen *et al.* (2004); Prigozhin (1999); Andreotti *et al.* (2002).
There are many theories generalizing Nishimori and Ouchi (1993) approach, see e.g. Caps and Vaanderwalle (2001). Prigozhin (1999) described the evolution of dunes by a system of two equations similar to the BCRE model discussed earlier in Sec. VI Bouchaud *et al.* (1994, 1995). One equation describes the evolution of the local height $`h`$ while another equation describes the density $`R`$ of particles rolling above the stationary sand bed profile (reptating particles),
$`_th`$ $`=`$ $`\mathrm{\Gamma }(h,R)f`$ (60)
$`_tR`$ $`=`$ $`𝐉+𝐐\mathrm{\Gamma }(h,R)`$ (61)
where $`\mathrm{\Gamma }`$ is the rolling-to-steady sand transition rate, $`𝐉`$ is the horizontal projection of the flux of rolling particles, $`𝐐`$ accounts for the influx of falling reptating grains, and $`f`$ is the erosion rate. With an appropriate choice of rate functions $`\mathrm{\Gamma },f,Q`$ and $`J`$, Eqs. (61) can reproduce many observed features of dune formation, such as initial instability of flat state, asymmetry of the dune profiles, coarsening and interaction of dunes, etc., see Fig. 42.
Thus, simplified models such as Nishimori and Ouchi (1993); Prigozhin (1999); Kroy *et al.* (2002a) have been successful in explaining many features of individual dune growth and evolution, see Fig. 43. However we should note that up to date none of the dune models have been able to address satisfactorily the wavelength selection in large-scale dune fields Hersen *et al.* (2004).
The phenomenon qualitatively similar to the dune formation occurs in an oscillatory fluid flow above a granular layer: sufficiently strong flow oscillations produce so-called vortex ripples on the surface of the underlying granular layer. These ripples are familiar to any beachgoer. Vortex ripple formation was first studied by Ayrton (1910); Bagnold (1956), and recently by Stegner and Wesfreid (1999); Scherer *et al.* (1999) and others. It was found that ripples emerge via a hysteretic transition, and are characterized by a near-triangular shape with slope angles close to the repose angle. The characteristic size of the ripples $`\lambda `$ is directly proportional to the displacement amplitude of the fluid flow $`a`$ (with a proportionality constant $`1.3`$) and is roughly independent on the frequency.
Andersen *et al.* (2001) introduced order parameter models for describing the dynamics of sand ripple patterns under oscillatory flow based on the phenomenological mass transport law between adjacent ripples. The models predict the existence of a stable band of wave numbers limited by secondary instabilities and coarsening of small ripples, in agreement with experimental observations.
Langlois and Valance (2005) studied underwater ripple formation on a two-dimensional sand bed sheared by viscous fluid. The sand transport is described by generalization of Eq. (55) taking into account both the local bed shear stress and the local bed slope. Linear stability analysis revealed that ripple formation is attributed to a growing longitudinal mode. The weakly nonlinear analysis taking into account resonance interaction of only three unstable modes revealed a variety of steady two-dimensional ripple patterns drifting along the flow at some speed.
Experiments in dune formation have been recently performed in water Betat *et al.* (1999). While water-driven and wind driven dunes and ripples have similar shape, the underlying physical processes are likely not the same due to a different balance between gravity and viscous drag in air and water.
Spectacular erosion patterns in sediment granular layers were observed in experiments with underwater flows Daerr *et al.* (2003); Malloggi *et al.* (2005a). In particular, a fingering instability of flat avalanche fronts was observed, see Fig. 44. These patterns are remarkably similar to those in thin films on inclined surfaces, both with clear and particle-laden fluids Troian *et al.* (1989); Zhou *et al.* (2005). In the framework of lubrication approximation the evolution of fluid film thickness $`h`$ is described by the following dimensionless equation following from the mass conservation law:
$$_th+\left\{\left[h^3^2h\right]\overline{D}h^3h\right\}+_xh^3=0$$
(62)
where dimensionless parameter $`\overline{D}`$ is inversely proportional to water surface tension. The instability occurs for small $`\overline{D}`$ values, i.e. in the large surface tension limit. However, despite visual similarity the physical mechanism leading to this fingering instability is not obvious: in fluid films the instability is driven (and stabilized) by the surface tension, whereas in the underwater granular flow fluid surface tension plays no role.
Duong *et al.* (2004) studied formation of periodic arrays of knolls in a slowly rotating horizontal cylinder filled with granular suspension, see Fig. 45. The solidified sediment knolls co-exist with freely circulating fluid. The authors applied variable viscosity fluid which formally allows simultaneous treatment of solid and liquid phase. In this model the effective flow viscosity $`\mu _s`$ diverges at the solid packing fraction $`\varphi _{rcp}`$,
$$\mu _s=\frac{\mu _0}{(1\varphi /\varphi _{rcp})^b}$$
(63)
where $`\mu _0`$ is the clear fluid viscosity and $`b`$ is an empirical coefficient. The model qualitatively reproduces the experiment, see Fig. 45. An interesting question in this context is whether there is a connection to the experiment by Shen (2002) where somewhat similar structures were obtained for the flow of “dry” particles in a horizontally rotating cylinder.
As it was mentioned in Sec. VII.3, Conway *et al.* (2004) reported that an air-fluidized vertical column of bi-disperse granular media sheared between counter-rotating cylinders exhibits formation of nontrivial vortex structure strongly reminiscent of Taylor vortices in conventional fluid, see e.g. Andereck *et al.* (1986). Authors argue that vortices in fluidized granular media, unlike Taylor vortices in fluid, are accompanied by the novel segregation-mixing mechanism specific for granular systems, see Fig. 40. Interestingly, no vortices were observed in a similar experiment in Couette geometry with monodisperse glass beads Losert *et al.* (2000).
Ivanova *et al.* (1996) studied patterns in a horizontal cylinder filled with sand/liquid mixture and subject to horizontal vibration. For certain vibration parameters standing wave patterns were observed at the sand/liquid interface. Authors argue that these wave patterns are similar to the Faraday ripples found at liquid/liquid interface under vertical vibration.
### VIII.2 Vortices in vibrated rods
In Section V we reviewed instabilities and collective motion in mechanically vibrated layers. In most experiments the particle shape was not important. However, strong particles anisotropy may give rise to non-trivial effects. Villarruel *et al.* (2000) observed onset of nematic order in packing of long rods in a narrow vertical tube subjected to vertical tapping. The rods initially compactify into a disordered state with predominantly horizontal orientation, but at later times (after thousands of taps) they align vertically, first along the walls, and then throughout the volume of the pipe. The nematic ordering can be understood in terms of the excluded volume argument put forward by Onsager (1949).
Blair *et al.* (2003a) studied the dynamics of vibrated rods in a shallow large aspect ratio system. Surprisingly, they found that vertical alignment of rods at large enough filling fraction $`n_f`$ and the amplitude of vertical acceleration ($`\mathrm{\Gamma }>2.2`$) can occur in the bulk, and it does not require side walls. Eventually, most of the rods align themselves vertically in a monolayer synchronously jumping on the plate, and engage in a correlated horizontal motion in the form of propagating domains of tilted rods, multiple rotating vortices etc, see Fig. 12 and Fig. 46. The vortices exhibit almost rigid body rotation near the core, and then the azimuthal velocity falls off, Fig. 47. The vortices merge in the course of their motion, and eventually a single vortex is formed in the cell.
Experiments showed that the rod motion occurs when the rods are tilted from the vertical, and it always occurs in the direction of tilt. In subsequent work Volfson *et al.* (2004) experimentally demonstrated that the correlated transport of bouncing rods is also found in quasi-one-dimensional geometry, and explained this effect using molecular dynamics simulations and a detailed description of inelastic frictional contacts between the rods and the vibrated plate. Effectively, bouncing rods become self-propelled objects similar to other self-propelled systems, for which large-scale coherent motion is often observed (bird flocks, fish schools, chemotactic microorganism aggregation, etc., see e.g. Grégoire and Chaté (2004); Helbing *et al.* (2000); Helbing (2001); Toner and Tu (1995)).
Aranson and Tsimring (2003) developed a phenomenological continuum theory describing coarsening and vortex formation in the ensemble of interacting rods. Assuming that the motion of rods is overdamped due to the bottom friction, the local horizontal velocity $`𝐯=(v_x,v_y)`$ of rods is of the form
$$𝐯=\left(p\alpha 𝐧f_0(n)\nu \right)/\zeta \nu ,$$
(64)
where $`\nu `$ is the density, $`p`$ is the hydrodynamic pressure, the tilt vector $`𝐧=(n_x,n_y)`$ is the projection of the rod director on the $`(x,y)`$ plane normalized by the rod length, i.e $`n=|𝐧|`$, and $`\zeta `$ is friction coefficient. According to Blair *et al.* (2003a); Volfson *et al.* (2004), the rods drift is determined by the average tilt of neighboring rods, thus the term $`\alpha 𝐧f_0(n)\nu `$ accounts for the average driving force from the vibrating bottom on the tilted rod. Eq. (64) combined with the mass conservation law yields
$$_t\nu =\mathrm{div}(𝐯\nu )=\zeta ^1\mathrm{div}\left(p\alpha 𝐧f_0(n)\nu \right).$$
(65)
To account for the experimentally observed phase separation and coarsening Aranson and Tsimring (2003) employed the Cahn-Hilliard approach (see Bray (1994) for review) by assuming that pressure $`p`$ can be obtained from the variation of a generic bistable “free energy” functional $`F`$ with respect to the density field $`\nu `$, $`p=\delta F/\delta \nu `$.
To close the description the equation for the evolution of tilt $`𝐧`$ is added on generic symmetry arguments:
$`_t𝐧`$ $`=`$ $`f_1(\nu )𝐧|𝐧|^2𝐧+`$ (66)
$`+`$ $`f_2(\nu )\left(\xi _1^2𝐧+\xi _2\mathrm{div}𝐧\right)+\beta \nu .`$
Here $`f_{1,2}`$ are certain functions of $`\nu `$, $`\xi _{1,2}`$ characterize diffusion coupling between the neighboring rods. Since the tilt field is not divergence-free, from the general symmetry considerations both $`\xi _{1,2}0`$ <sup>3</sup><sup>3</sup>3These constants are analogous to the first and second viscosity in ordinary fluids, see e.g. Landau and Lifshits (1959).
Numerical and analytic studies of Eqs. (65),(66) revealed phase coexistence, nucleation and coalescence of vortices in accord with the experiment, see Fig. 48.
An interesting experiment with anisotropic chiral particles was performed by Tsai *et al.* (2005). The role of particles was played by bend-wire objects which rotated in a preferred direction under vertical vibration. The experiments demonstrated that individual angular rotation of the particles was converted into the collective angular momentum of the granular gas of these chiral objects. The theoretical description of this system was formulated in the framework of two phenomenological equations for the density $`\nu `$ and center-of-mass momentum density $`\nu 𝐯`$ and the spin angular momentum density $`l=\mathrm{\Omega }`$ arising from the ration of particles around their center of mass, $`\mathrm{\Omega }`$ is the particle’s rotation frequency. Whereas the equations for density and velocity are somewhat similar to those for the vibrated rod system, the equation for the spin momentum clearly has no counterpart in the vibrated rod system and was postulated in the following form:
$$_tl+𝐯l=\tau \mathrm{\Gamma }^\mathrm{\Omega }\mathrm{\Gamma }(\mathrm{\Omega }\omega )+D_\mathrm{\Omega }^2\mathrm{\Omega }$$
(67)
where $`\tau `$ is the source of the angular rotation (due to chirality of particles), $`\omega `$ is coarse-grained or collective angular velocity, $`\mathrm{\Gamma }^\mathrm{\Omega }`$ and $`\mathrm{\Gamma }`$ are dissipative coefficients due to friction and $`D_\mathrm{\Omega }`$ is the angular momentum diffusion. Eq. (67) predicts, in agreement with the experiment, the onset of collective rotation of the gas of particles. Possibly, it also exhibits non-trivial spatio-temporal dynamics similar to those in the system of vibrated rods. However, due to the small number of particles (about 350) in the experiment the nontrivial collective regimes were not reported.
### VIII.3 Electrostatically driven granular media
Large ensembles of small particles display fascinating collective behavior when they acquire an electric charge and respond to competing long-range electromagnetic and short-range contact forces. Many industrial technologies face the challenge of assembling and separating such single- or multi-component micro and nano-size ensembles. Traditional methods, such as mechanical vibration and shear, are infective for very fine powders due to agglomeration, charging, etc. Electrostatic effects often change statistical properties of granular matter such as energy dissipation rate Sheffler and Wolf (2002), velocity distributions in granular gases Aranson and Olafsen (2002); Kohlstedt *et al.* (2005), agglomeration rates in suspensions Dammer and Wolf (2004), etc.
Aranson *et al.* (2000, 2002); Sapozhnikov *et al.* (2003a, 2004) studied electrostatically driven granular matter. This method relies on the collective interactions between particles due to a competition between short range collisions and long-range electromagnetic forces. Direct electrostatic excitation of small particles offers unique new opportunities compared to traditional techniques of mechanical excitation. It enables one to deal with extremely fine nonmagnetic and magnetic powders which are not easily controlled by other means.
In most experimental realizations, several grams of mono-dispersed conducting micro-particles were placed into a 1.5 mm gap between two horizontal $`30\times 30`$ cm<sup>2</sup> glass plates covered by transparent conducting layers of indium tin-dioxide. Typically $`45`$ $`\mu m`$ Copper or $`120`$ $`\mu m`$ Bronze spheres were used. Experiments were also performed with much smaller 1 $`\mu m`$ particles, Sapozhnikov *et al.* (2004). An electric field perpendicular to the plates was created by a high voltage source (0-3 kV) connected to the inner surface of each plate. Experiments were performed in air, vacuum, or in the cell filled with non-polar weakly-conducting liquid.
The basic principle of the electro-cell operation is as follows. A particle acquires an electric charge when it is in contact with the bottom conducting plate. It then experiences a force from the electric field between the plates. If the upward force induced by the electric field exceeds gravity, the particle travels to the upper plate, reverses charge upon contact, and is repelled down to the bottom plate. This process repeats in a cyclical fashion. In an air-filled or evacuated cell, the particle remains immobile at the bottom plate if the electric field $`E`$ is smaller than the first critical field $`E_1`$. For $`E>E_1`$ an isolated particle leaves the plate and starts to bounce. However, if several particles are in contact on the plate, screening of the electric field reduces the force on individual particles, and they remain immobile. A simple calculation shows that for the same value of the applied electric field the force acting on isolated particles exceeds by a factor of two the force acting on the particle inside the dense monolayer. However, if the field is larger than a second critical field value, $`E_2>E_1`$, all particles leave the plate, and the system of particles transforms into an uniform gas-like phase. When the field is decreased below $`E_2`$ ($`E_1<E<E_2`$), in air-filled or evacuated cells localized clusters of immobile particles spontaneously nucleate to form a static clusters (precipitate) on the bottom plate Aranson *et al.* (2000). The clusters exhibit the Ostwald-type ripening Meerson (1996); Bray (1994), see also Subsec. IV.3.
#### VIII.3.1 Coarsening of clusters
Results for the electrostatically driven system yielded the following asymptotic scaling law, see Fig. 50:
$$N\frac{1}{t}$$
(68)
where $`N`$ is the number of clusters and $`t`$ is time. Accordingly, the average cluster area $`A`$ increases with time as $`At`$. This behavior is consistent with the interface-controlled Ostwald ripening Meerson (1996).
A theoretical description of coarsening in an electrostatically driven granular system was developed by Aranson *et al.* (2000), Sapozhnikov *et al.* (2003). The theory was formulated in terms of the Ginzburg-Landau-type equation for the number density of immobile particles (precipitate or solid) $`n`$
$$_tn=^2n+\varphi (n,n_g)$$
(69)
where $`n_g`$ is the number density of bouncing particles (gas) $`n_g`$, and $`\varphi (n,n_g)`$ is a function characterizing a solid/gas conversion rate. The effectiveness of the solid/gas transitions is controlled by the local gas concentration $`n_g`$. It was assumed that the gas concentration is almost constant because the particle’s mean free pass in the gas state is very large. The gas concentration $`n_g`$ is coupled to $`n`$ due to total mass conservation constraint
$$Sn_g+n(x,y)𝑑x𝑑y=M,$$
(70)
where $`S`$ is the area of domain of integration, and $`M`$ is the total number of particles. Function $`\varphi (n,n_g)`$ is chosen in such a way as to provide bistable local dynamics of concentration corresponding to the hysteresis of the gas/solid transition. The above description yields a very similar temporal evolution of clusters (see Fig. 49) and produces a correct scaling for the number of clusters Eq. (68).
In the so-called sharp interface limit when the size of clusters is much larger than the width of interfaces between clusters and granular gas, Eq. (69) can be reduced to equations for the cluster radii $`R_i`$ (assuming that clusters have circular form):
$$\frac{dR_i}{dt}=\kappa \left(\frac{1}{R_c(t)}\frac{1}{R_i}\right),$$
(71)
where $`R_c`$ is critical cluster size, $`\kappa `$ is effective surface tension (experimental measurements of cluster surface tension were conducted by Sapozhnikov *et al.* (2003)). The critical radius $`R_c`$ is a certain function of the granular gas concentration $`n_g`$ that enters Eqs. (71) through the the conservation law Eq. (70) which in two dimensions reads
$$n_gS+\pi \underset{i=1}{\overset{N}{}}R_i^2=M.$$
(72)
The statistical properties of Ostwald ripening can be understood in terms of the probability distribution function $`f(R,t)`$ of cluster sizes. Following Lifshitz and Slyozov (1958, 1961); Wagner (1961) and neglecting cluster merger, one obtains in the limit $`N\mathrm{}`$ that the probability distribution $`f(R,t)`$ satisfies the continuity equation
$$_tf+_R\left(\dot{R}f\right)=0.$$
(73)
From the mass conservation in the limit of small gas concentration Eq. (72) one obtains an additional constraint:
$$\pi _0^{\mathrm{}}R^2f(R,t)𝑑R=M$$
(74)
Eqs. (73),(74) have a self-similar solution in the form
$$f(R,t)=\frac{1}{t^{3/2}}F\left(\frac{R}{\sqrt{t}}\right)$$
(75)
For the total number of clusters $`N=_0^{\mathrm{}}f𝑑R`$ the scaling Eq.(75) yields $`N1/t`$, which appears to be in a good agreement with the experiment, see Fig. 50. However, the cluster size distribution function appears to be in a strong disagreement, see Fig. 51. In particular, Lifshitz and Slyozov (1958, 1961); Wagner (1961) theory predicts the distribution with a cut-off (dotted line) whereas the experiment yields the function with an exponential tail. A much better agreement with the experiment was obtained when binary coalescence of clusters was incorporated in the Lifshitz-Slyozov-Wagner theory Sapozhnikov *et al.* (2005); Conti *et al.* (2002).The coalescence events become important for a finite area fraction of the clusters.
Ben-Naim and Krapivsky (2003) applied an exchange growth model to describe coarsening in granular media. In this theory the cluster growth rates are controlled only by the cluster area ignoring shape effects. Assuming that the number of particles in a cluster evolves via uncorrelated exchange of single particles with an other cluster the following equation for the density of clusters containing $`k`$ particles can be derived:
$$\frac{dA_k}{dt}=\underset{i,j}{}A_iA_jK_{ij}\left(\delta _{k,i+1}+\delta _{k,i1}2\delta _{k,i}\right)$$
(76)
where $`A_k`$ is the probability to find a cluster containing $`k`$ particles, $`K_{ij}`$ the exchange kernel and $`\delta _{k,i}`$ is the Kronecker symbol. For the choice of homogeneous kernel $`K_{ij}=(ij)^\lambda `$ with $`\lambda =1`$ this theory predicts correct scaling of the cluster size with time $`R\sqrt{t}`$ and exponential decay of the cluster size distribution function, as in the experiment. The choice of $`\lambda =1`$ is equivalent to the assumption that the exchange rate is determined by the size of the cluster. In the theory by Sapozhnikov *et al.* (2005) the cluster evolution is governed by the evaporation/deposition of particles at the interface of the cluster and controlled by the overall pressure of the granular gas. Thus, both theories predict the same scaling behavior, however the underlying assumptions are very different. A possible explanation for this may be that while the exchange growth model ignores the curvature of the cluster interface and the dependence on exchange rate on the pressure of granular gas, the agreement is obtained by tuning the adjustable parameter $`\lambda `$.
#### VIII.3.2 Dynamics of patterns in a fluid-filled cell
Sapozhnikov *et al.* (2003a) performed experiments with electrostatically driven granular media immersed in a weakly conducting non-polar fluid (toluene-ethanol mixture). Depending on the applied electric field and the ethanol concentration (which controls the conductivity of the fluid), a plethora of static and dynamic patterns were discovered, see Fig. 15. For relatively low concentrations of ethanol (below 3%), the qualitative behavior of the liquid-filled cell is not very different from that of the air-filled cell: clustering of immobile particles and coarsening were observed between two critical field values $`E_{1,2}`$ with the clusters being qualitatively similar to that of the air cell. However, when the ethanol concentration is increased, the phase diagram becomes asymmetric with respect to the direction of the electric field. Critical field magnitudes, $`E_{1,2}`$, are larger when the electric field is directed downward (“$`+`$” on the upper plate) and smaller when the field is directed upward (“$``$” on the upper plate). This difference increases with ethanol concentration. The observed asymmetry of the critical fields is apparently due to an excess negative charge in the bulk of the liquid.
The situation changes dramatically for higher ethanol concentrations: increasing the applied voltage leads to the formation of two new immobile phases: honeycomb (Fig. 15b) for the downward direction of the applied electric field, and two-dimensional crystal-type states for the upward direction.
A further increase of ethanol concentration leads to the appearance of a novel dynamic phase - condensate (Fig. 15c,d) where almost all particles are engaged in a circular vortex motion in the vertical plane, resembling Rayleigh-Bénard convection. The condensate co-exists with the dilute granular gas. The direction of rotation is determined by the polarity of the applied voltage: particles stream towards the center of the condensate near the top plate for the upward field direction and vice versa. The evolution of the condensate depends on the electric field direction. For the downward field, large structures become unstable due to the spontaneous formation of voids (Fig. 15d). These voids exhibit complex intermittent dynamics. In contrast, for the upward field, large vortices merge into one, forming an asymmetric object which often performs composite rotation in the horizontal plane. The pattern formation in this system is most likely caused by self-induced electro-hydrodynamic micro-vortices created by the particles in weakly-conducting fluids. These micro-vortices create long-range hydrodynamic vortex flows which often overwhelm electrostatic repulsion between likely-charged particles and introduce attractive dipole-like hydrodynamic interactions. Somewhat similar micro-vortices are known in driven colloidal systems, see e.g. Yeh *et al.* (1997).
Aranson and Sapozhnikov (2004) developed a phenomenological continuum theory of pattern formation for metallic micro-particles in a weakly conducting liquid subject to an electric field. Based on the analogy with the previously developed theory of coarsening in air-field cell Aranson *et al.* (2002), the model is formulated in terms of conservation laws for the number densities of immobile particles (precipitate) $`n_p`$ and bouncing particles (gas) $`n_g`$ averaged over the thickness of the cell:
$`_tn_p=𝐉_p+f,_tn_g=𝐉_gf.`$ (77)
Here $`J_{p,g}`$ are the mass fluxes of precipitate and gas respectively and the function $`f`$ describes gas/precipitate conversion which depends on $`n_{p,g}`$, electric field $`E`$ and local ionic concentration $`c`$. The fluxes are written as:
$$𝐉_{p,g}=D_{p,g}n_{p,g}+\alpha _{p,g}(E)𝐯_{}n_{p,g}(1\beta (E)n_{p,g}),$$
(78)
where $`v_{}`$ is horizontal hydrodynamic velocity, $`D_{p,g}`$ are precipitate/gas diffusivities. The last term, describing particles advection by fluid, is reminiscent of the Richardson-Zaki relation for a drag force frequently used in the engineering literature Richardson and Zaki (1954). The factor $`(1\beta (E)n_{p,g})`$ describes the saturation of flux at large particle densities $`n1/\beta `$ due to the decrease of void fraction. Terms $`\alpha _{p,g}`$ describe advection of particles by the fluid. Interestingly, in the limit of very large gas diffusion $`D_gD_p`$ and without advection terms ($`\alpha _{p,g}=0`$) the model reduces to Eqs. (69) and (70) applied for air-filled cell Aranson *et al.* (2002).
Eqs. (77) are coupled to the cross section averaged Navier-Stokes equation for vertical velocity $`v_z`$:
$$n_0(_tv_z+𝐯v_z)=\mu ^2v_z_zp+E_zq$$
(79)
where $`n_0`$ is the density of liquid (we set $`n_0=1`$), $`\mu `$ is the viscosity, $`p`$ is the pressure, and $`q`$ is the charge density. The last term describes the electric force acting on charged liquid. Horizontal velocity $`v_{}`$ is obtained from $`v_z`$ using the incompressibility condition $`_zv_z+_{}v_{}=0`$ in the approximation that vertical vorticity $`\mathrm{\Omega }_z=_xv_y_yv_x`$ is small compared to in-plane vorticity. This assumption allows one to find the horizontal velocity as a gradient of quasi-potential $`\varphi `$: $`𝐯_{}=_{}\varphi `$.
For an appropriate choice of the parameters the model Eqs. (77),(79) yields qualitatively correct phase diagram and the patterns observed in the experiment, see Figs. 15 and 52.
### VIII.4 Magnetic particles
Electric and magnetic interactions allow introduction of controlled long-range forces in granular systems. Blair *et al.* (2003a); Blair and Kudrolli (2003b); Stambaugh *et al.* (2004a, b) performed experimental studies with vibrofluidized magnetic particles. Several interesting phase transitions were reported, in particular, the formation of dense two-dimensional clusters and loose quasi-one-dimensional chains and rings. Blair *et al.* (2003a) considered pattern formation in a mixture of magnetic and non-magnetic (glass) particles of equal mass. The glass particles played the role of “phonons”, their concentration allowed an adjustment of the typical fluctuation velocity of the magnetic subsystem. The phase diagram delineating various regimes in this system is shown in Fig. 53. While the phase diagram shows some similarity with equilibrium dipolar fluids (such as phase coexistence), most likely there are differences due to the non-equilibrium character of granular systems.
Stambaugh *et al.* (2004a) performed experiments with relatively large particles (about 1.7 cm), and near the closed-packed density. It was found that particles form hexagonal closed-packed clusters in which the magnetic dipoles lay in the plane and assume circulating vortical patterns. For lower density ring patterns were observed. Experiments with mixture of particles with two different magnetic moments revealed segregation effects Stambaugh *et al.* (2004b). The authors argue that the static configurational magnetic energy is the primary factor in pattern selection.
Experiments by Blair and Kudrolli (2003b); Stambaugh *et al.* (2004a, b) were limited to a small number (about 10<sup>3</sup>) of large particles due to the intrinsic limitation of the mechanical vibrofluidization technique. Snezhko *et al.* (2005) performed experimental studies of 90 $`\mu m`$ Nickel micro-particles subjected to electrostatic excitation, see also Subsec. VIII.3. The electrostatic system allowed researchers to perform experiments with a very large number of particles (of the order of $`10^6`$) and a large aspect ratio of the experimental cell. Thus the transition between small chains and large networks (Fig. 13) was addressed in detail. An abrupt divergence of the chain length was found when the frequency of field oscillations decreased, resulting in the formation of a giant interconnected network.
Studies of the collective dynamics and pattern formation of magnetic particles are still in the early phases. While it is natural to assume that magnetic interaction plays a dominant role in pattern selection, further computational and theoretical studies of pattern formation in systems of driven dipolar particles are necessary. Besides a direct relevance for the physics of granular media, studies of magnetic granular media may provide an additional insight into the behavior of dipolar hard sphere fluids where the nature of solid/liquid transitions is still debated de Gennes and Pincus (1970); Levin (1999). Vibration or electrostatically fluidized magnetic particles can also be viewed as a macroscopic model of a ferrofluid, where similar experiments are technically difficult to perform.
## IX Overview and Perspectives
Studies of granular materials are intrinsically interdisciplinary and they borrow ideas and methods from other fields of physics such as statistical physics, mechanics, fluid dynamics, and the theory of plasticity. On the flip side, progress in understanding granular matter can be often applied to seemingly unrelated physical systems, such as ultra-thin liquid films, foams, colloids, emulsions, suspensions, and other soft condensed matter systems. The common feature shared by these systems is the discrete microstructure directly influencing macroscopic behavior. For example, the order parameter description similar to that of Sec. VI.1.1 was applied to stick-slip friction in ultra-thin films, Israelachvili *et al.* (1988); Urbach *et al.* (2004); Carlson and Batista (1996); Aranson *et al.* (2002c).
Lemaitre (2002); Lemaitre and Carlson (2004) applied the idea of shear-transformation zone (STZ) pioneered by Falk and Langer (1998) for amorphous solids both to granular matter and to the boundary lubrication problem in confined fluid. In this theory the plastic deformation is represented by a population of mesoscopic regions which may undergo non-affine deformations in response to stress. Concentration of STZs in amorphous material is somewhat similar to the order parameter (relative concentration of defects) introduced by Aranson *et al.* (2002c). A conceptually similar approach was proposed by Staron *et al.* (2002) who described the onset of fluidization as a percolation of the contact network with fully mobilized friction. Whereas derivation of the constitutive relations from first-principle microscopic rules is still a formidable challenge, these approaches are promising for understanding of not only the boundary lubrication problem, but also onset of motion in dense granular matter.
Flowing liquid foams and emulsions share many similarities with granular matter: they have internal discrete structure (bubbles and drops play the role of grains), and two different mechanisms are responsible for the transmission of stresses: elastic for small stress and visco-plastic above certain yield stress. However, there are additional complications: bubbles are highly deformable and, unlike granular matter, a number of particles may change due to the coalescence of bubbles.
Foams and granular materials often exhibit similar behavior, such as non-trivial stress relaxation and power-law distribution of rearrangement events Dennin and Knobler (1997). Stick-slip behavior was reported both for sheared foams Lauridsen *et al.* (2002) and granular materials Nasuno *et al.* (1997). Remarkably, recent experiments with two dimensional foams Lauridsen *et al.* (2004) and three dimensional emulsions Coussot *et al.* (2002a, b); DaCruz (2002) strongly suggest the coexistence between flowing (liquid) and jammed (solid) states reminiscent of that in granular matter. Furthermore, avalanche behavior reminiscent of granular flows down an inclined plane Daerr and Douady (1999) was reported by Coussot *et al.* (2002a) for clay suspensions, see Fig. 54. There are many approaches treating foams, gel and suspensions as complex fluids with specific stress-strain constitutive relation. For example, Fuchs and Cates (2002) used the analogy between glasses and dense colloidal suspensions and applied the mode coupling approach to understand the nonlinear rheology and yielding. Similar approaches can be possibly useful for granular materials Schofield and Oppenheim (1994).
Liu and Nagel (1998) suggested that a broad class of athermal soft matter systems (glasses, suspensions, granular materials) shows a universal critical behavior in the vicinity of solid-fluid or jamming transition, see Fig. 55. Whether jammed systems indeed have common features that can be described by a universal phase diagram is an open issue. An interesting question in this context is a possibility of thermodynamic description of driven, macroscopic, athermal systems like granular materials and foams in terms of some kind of effective temperature. Studies of interacting particles under shear Ono *et al.* (2002); O’Hern *et al.* (2004); Makse and Kurchan (2002); Xu and O’Hern (2005); Corvin *et al.* (2005) indicate that indeed under certain conditions it is possible to define an effective temperature (for example, from the equivalent of the Einstein-Stokes relation) for a broad class of athermal systems from comparison of the mechanical linear response with the corresponding time-dependent fluctuation-dissipation relation. However, the possibility of developing nonequilibrium thermodynamics of the basis of the effective temperature is under debate.
Granular systems exhibit many similarities with traffic flows and collective motion of self-propelled particles such as swimming bacteria, fish schools, bird flocks, etc., see for review Helbing (2001). In particular, jamming transition in granular media and traffic jams show similar features, such as hysteresis, and clusters formation. Moreover, continuum models of traffic flows are often cast in the form of modified Navier-Stokes equation with density-dependent viscosity, similar to granular hydrodynamics.
Let us discuss briefly some open questions in the physics of granular matter.
* Static vs. dynamic description. Commonly accepted models of rapid granular flows (granular hydrodynamics) and quasi-static dense flows (elastic and visco-plastic models) are very different, see e.g. Goldenberg and Goldhirsch (2002). However, near the fluidization transition, and in dense partially-fluidized flows, the differences between these two regimes become less obvious. The fluidization of sheared granular materials has many features of a first-order phase transition. The phenomenological partial fluidization theory in principle can be a bridge between the static and dynamic descriptions. The order parameter related to the local coordination number appears to be one of the hidden fields required for a consistent description of granular flows. One important question in this regard is the universality of the fluidization transition in different granular systems and geometries. On the opposite side of the fluidization transition, the static state of the granular matter can be described by the order parameter related to the percentage of static contacts with fully activated dry friction (critical contacts) Staron *et al.* (2002). It was shown that once these contacts form a percolation cluster, the granular pack slips and fluidization occurs. It is of obvious interest to relate this “static” order parameter and the “dynamics” order parameter discussed above. We see one of the main future challenges in the systematic derivation of the continuum theory valid both for flowing and static granular matter.
* Statistical mechanics of dense granular systems. Clearly, discrete grain structure plays a major role in the dynamics and inherent stochasticity of granular response. The number of particles in a typical granular assembly is large (10<sup>6</sup> or more) but it is much smaller than the Avogadro number. Traditional tools of statistical physics do not apply to dense granular systems since grains do not exhibit thermal Brownian motion. One of the alternative ways of describing statistics of granular media was suggested by Edwards and Grinev (1998) in which they proposed that volume rather than energy serves as the extensive variable in a static granular system, so that the role of temperature is played by the compactivity which is the derivative of the volume with respect to the usual entropy. Recent experiments Makse and Kurchan (2002); Schröter *et al.* (2005) aim to test this theory experimentally. Connecting Edwards theory with granular hydrodynamics will be an interesting challenge for future studies.
* Realistic simulations of three-dimensional granular flows. Even the most advanced simulations of granular flows in three dimensions Silbert *et al.* (2003); Silbert (2005) are limited to relatively small samples (e.g. $`100\times 40\times 40`$ particles box) and are very time consuming. The granular problems are inherently very stiff: while the collisions between particles are very short ($`O(10^4sec`$), the collective processes of interest may take many seconds or minutes. As a result, to the time step limitations a simulation of realistic hard particles is not feasible: the “simulations” particles have elastic moduli several orders of magnitude smaller than sand or glass. The particle softness may introduce unphysical artifacts in the overall picture of the motion. Different approaches to handling this problem will be necessary to advance the state of the art in simulations. New opportunity can be offered by the equation-free simulation method proposed by Kevrekidis *et al.* (2004). Another area of simulations which needs further refinements is an accurate account of dry friction. In the absence of a better solution current methods (see for review Luding (2004)) employ various approximate techniques to simulation dry friction, and accuracy of these methods can be questionable.
* Complex interactions. Understanding of dynamics of granular systems with complex interactions is certainly an intriguing and rapidly developing field. While interaction of grains with intersticial fluid is a traditional part of engineering research, effects of particle anisotropy, long-range electromagnetic interactions mediating collisions, adhesion, agglomeration and many others constitute a formidable challenge for theorists and a fertile field of future research.
* Granular physics on a nano-scale. There is a persistent trend in the industry such as powder metallurgy, pharmaceutical and various chemical technologies towards operating with smaller and smaller particles. Moreover, it was recognized recently that micro- and nano-particles can be useful for fabrication of desired ordered structures and templates for a broad range of nanotechnological applications through self-assembly processes. Self-assembly, the spontaneous organization of materials into complex architectures, constitutes one of the greatest hopes of realizing the challenge to create ever smaller nanostructures. It is a particulary attractive alternative to traditional approaches such as lithography and electron beam writing. Reduction of the particle size to micro- and nano -scales shifts the balance between forces controlling particle interaction because the dominant interactions depend on the particle size. While for macroscopic grains the dynamics are governed mostly by the gravity, collisional and frictional forces, for micro- and nano-particles the dominant interactions include long-range electromagnetic forces, short- range van der Waals interactions, etc. Nevertheless, some concepts and ideas developed in the “traditional” granular physics were successfully applied to understand dynamic self-assembly of microparticles Sapozhnikov *et al.* (2003a, 2004) and even biological microtubules Aranson and Tsimring (2005). We expect to see more and more efforts in this direction.
## Acknowledgments
The authors thank Dmitrii Volfson, Alexey Snezhko, Maksim Saposhnikov, Jie Li, Adrian Daerr, Bob Behringer, Jerry Gollub, Thomas Halsey, Denis Ertas, Harry Swinney, Jeff Olafsen, Eli Ben-Naim, Valerii Vinokur, Wai Kwok, George Crabtree, Paul Umbanhowar, Francisco Melo, Eric Clement, Jacques Prost, Philippe Claudin, Julio Ottino, Devang Khakhar, Jean-Philippe Bouchaud, Olivier Pouliquen, Jacques Duran, Anaël Lemaître, Evelyne Kolb, Hugues Chaté, Gary Grest, Arshad Kudrolli, Douglas Durian, Peter Schiffer, Leo Silbert, Wolfgang Losert, Daniel Blair, Paul Chaikin, Henrich Jaeger, Sid Nagel, Leo Kadanoff, Thomas Witten, Sue Coppersmith, Baruch Meerson, Ray Goldstein, Chay Goldenberg, Isaak Goldhirsch, Robert Ecke, Thorsten Pöschel, Alexandre Valance, James Dufty, James Jenkins, Dietrich Wolf, Haye Hinrichsen, Lorenz Kramer, Len Pismen, Martin van Hecke, Wim van Saarloos, Guenter Ahlers, Jacob Israelachvili, James Langer, Pierre-Gilles de Gennes and many others for useful discussions. This work was supported by the Office of the Basic Energy Sciences at the United States Department of Energy, grants W-31-109-ENG-38, and DE-FG02-04ER46135.
The review was partly written when one of us (I.A) was attending Granular Session in Institute Henry Poincaré, Paris, and Granular Physics Program, Kavli Institute for Theoretical Physics in Santa Barbara. |
warning/0507/gr-qc0507003.html | ar5iv | text | # Mass and angular momentum of asymptotically AdS or flat solutions in the topologically massive gravity
## I Introduction
In the presence of a cosmological constant, the source-free field equations of the (2+1 dimensional) topologically massive gravity (TMG) theory read
$$R_{\mu \nu }\frac{1}{2}g_{\mu \nu }R+\mathrm{\Lambda }g_{\mu \nu }+\frac{1}{\mu }C_{\mu \nu }=0,C^{\mu \nu }\frac{1}{\sqrt{g}}ϵ^{\mu \alpha \beta }_\alpha \left(R_\beta ^\nu \frac{1}{4}\delta _\beta ^\nu R\right),$$
where the Cotton tensor $`C_{\mu \nu }`$ is the three dimensional analogue of the Weyl tensor and is symmetric, traceless and identically conserved; the parameter $`\mu `$ is the coupling constant for the gravitational Chern-Simons term in the action and corresponds to the mass of the linearized TMG excitations at $`\mathrm{\Lambda }=0`$ (see des for details and for $`\mathrm{\Lambda }0`$ see dba1 ). A minimally supersymmetric extension of this theory was also constructed long time ago deskay .
The well known BTZ metric btz , as well as the Anti-de Sitter (AdS) and the Schwarzschild-dS spacetimes, are solutions to TMG theory with a cosmological constant in a ‘trivial’ manner since their Cotton tensors vanish identically. There are, however, other known ‘nontrivial’ solutions; i.e. those spacetimes that obey the ‘full’ TMG equations, not just their Einstein part (or the cosmological Einstein part in the relevant cases) alone. The first example that we know is Deser’s gravitational anyons which are only solutions of the linearized TMG equations any . The ‘fully’ nonlinear, ‘nontrivial’ solutions include the Vuorio solution vuo and its generalization to solutions with a constant twist per . There are also exact static/stationary solutions for spinning point sources for which the spin and the mass of the sources obey a certain relation ort , par . A class of cosmological-type solutions is given by finite action exact solutions of TMG nub that are also useful for a classification of homogeneous solutions. There also exists a two parameter solution to TMG theory with a cosmological constant which seems to have properties similar to the BTZ solution nut , gur . This solution is not asymptotically AdS in its original form, but then its asymptotically AdS form, which is obtained by imposing a certain relation between $`\mathrm{\Lambda }`$ and $`\mu `$, is equivalent to the BTZ metric. Another class of solutions which asymptotically approach extremal BTZ black holes, but are geodesically complete with no event horizons, were also given cle . Finally, the first nontrivial example of a solution to TMG that preserves half of supersymmetry was found in ds ; moreover, the solutions in cle seem to be related to the supersymmetric solutions by a certain choice of parameters and a coordinate transformation. Another two-parameter family of black hole solutions that are obtained by an analytical continuation of the Vuorio solution, but fail to be asymptotically AdS like their ancestor, was recently given in dal .
In a recent work db1 , a concrete and rigorous definition of conserved gravitational charges (particularly energy and angular momentum) were given in a ‘surface’ integral form about their flat or asymptotically AdS backgrounds in TMG theory. It is only natural to consider the exact solutions listed in the previous paragraph as explicit examples whose gravitational charges can be calculated a la db1 . We do this for the BTZ and the only nontrivial supersymmetric solutions of TMG that are asymptotically AdS or flat in this paper. This should also help to better understand the physical properties of these examples and to clarify the physical meanings of some of the parameters that explicitly show up in them.
## II The conserved gravitational charges of the TMG theory
Let us start by giving a brief outline of how gravitational charges are defined in TMG. (We refer the reader to db1 and db2 for details.) Assume that the deviation, $`h_{\mu \nu }`$, of the actual spacetime metric $`g_{\mu \nu }=\overline{g}_{\mu \nu }+h_{\mu \nu }`$ from an asymptotically AdS metric (or the background) $`\overline{g}_{\mu \nu }`$, which obeys
$$\overline{R}_{\mu \alpha \nu \beta }=\mathrm{\Lambda }(\overline{g}_{\mu \nu }\overline{g}_{\alpha \beta }\overline{g}_{\mu \beta }\overline{g}_{\alpha \nu }),\overline{R}_{\mu \nu }=2\mathrm{\Lambda }\overline{g}_{\mu \nu },\overline{R}=6\mathrm{\Lambda },$$
is employed for constructing “linearized gravity” in the usual sense with the usual assumptions db2 . Then the ‘linearized’ part of the Ricci tensor <sup>1</sup><sup>1</sup>1Here $`hh_{\mu \nu }\overline{g}^{\mu \nu }`$, all indices are raised and lowered with the background metric $`\overline{g}_{\mu \nu }`$ and also all covariant differentiations are carried with respect to $`\overline{g}_{\mu \nu }`$.
$$R_{\mu \nu }^L=\frac{1}{2}(\overline{\text{ }\text{ }\text{ }\text{ }}h_{\mu \nu }\overline{}_\mu \overline{}_\nu h+\overline{}^\sigma \overline{}_\nu h_{\sigma \mu }+\overline{}^\sigma \overline{}_\mu h_{\sigma \nu }),$$
and the linearized Ricci scalar
$$R^L(R_{\mu \nu }g^{\mu \nu })^L=R_{\mu \nu }^L\overline{g}^{\mu \nu }2\mathrm{\Lambda }h=\overline{\text{ }\text{ }\text{ }\text{ }}h+\overline{}_\mu \overline{}_\nu \overline{h}^{\mu \nu }2\mathrm{\Lambda }h,$$
can be used for finding the linearized cosmological Einstein and the Cotton tensors as
$`𝒢_{\mu \nu }`$ $``$ $`(G_{\mu \nu }+\mathrm{\Lambda }g_{\mu \nu })^L=R_{\mu \nu }^L{\displaystyle \frac{1}{2}}\overline{g}_{\mu \nu }R^L2\mathrm{\Lambda }h_{\mu \nu },`$
$`C_L^{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{\overline{g}}}}ϵ^{\mu \alpha \beta }\overline{g}_{\beta \sigma }\overline{}_\alpha \left(R_L^{\sigma \nu }2\mathrm{\Lambda }h^{\sigma \nu }{\displaystyle \frac{1}{4}}\overline{g}^{\sigma \nu }R_L\right).`$
Now one can find a background conserved and gauge invariant charge (corresponding to each background Killing vector $`\overline{\xi }^\mu `$) which is given as the sum of the following three terms:
$$Q^\mu (\overline{\xi })=\frac{1}{8\pi G}_{}𝑑S_i\left(q_E^{\mu i}(\overline{\xi })+\frac{1}{2\mu }q_E^{\mu i}(\overline{\mathrm{\Xi }})+\frac{1}{2\mu }q_C^{\mu i}(\overline{\xi })\right),$$
(1)
where <sup>2</sup><sup>2</sup>2Here one in fact has $`𝒢^{\mu \nu }(G^{\mu \nu }+\mathrm{\Lambda }g^{\mu \nu })^L`$ but then since $`\overline{G}_{\mu \nu }+\mathrm{\Lambda }\overline{g}_{\mu \nu }=0`$, moving the indices of the linearized cosmological Einstein tensor $`𝒢`$ can equivalently be carried out with the background metric $`\overline{g}_{\mu \nu }`$.
$`q_E^{\mu i}(\overline{\xi })`$ $``$ $`\sqrt{\overline{g}}(\overline{\xi }_\nu \overline{}^\mu h^{i\nu }\overline{\xi }_\nu \overline{}^ih^{\mu \nu }+\overline{\xi }^\mu \overline{}^ih\overline{\xi }^i\overline{}^\mu h`$ (2)
$`+h^{\mu \nu }\overline{}^i\overline{\xi }_\nu h^{i\nu }\overline{}^\mu \overline{\xi }_\nu +\overline{\xi }^i\overline{}_\nu h^{\mu \nu }\overline{\xi }^\mu \overline{}_\nu h^{i\nu }+h\overline{}^\mu \overline{\xi }^i),`$
$`q_C^{\mu i}(\overline{\xi })`$ $``$ $`ϵ^{\mu i\beta }𝒢_{\nu \beta }\overline{\xi }^\nu +ϵ^{\nu i\beta }𝒢_\beta ^\mu \overline{\xi }_\nu +ϵ^{\mu \nu \beta }𝒢_\beta ^i\overline{\xi }_\nu ,`$ (3)
and $`\overline{\mathrm{\Xi }}^\beta ϵ^{\alpha \nu \beta }\overline{}_\alpha \overline{\xi }_\nu /\sqrt{\overline{g}}`$ is another background Killing vector constructed out of $`\overline{\xi }`$. Here $``$ is a spatial 2-dimensional hypersurface, $``$ is its 1-dimensional boundary and $`i`$ denotes the space direction orthogonal to the boundary $``$ with the corresponding line element $`dS_i`$. $`G`$ denotes the 3-dimensional Newton’s constant and the charge has been normalized by the overall factor $`8\pi G`$ in (1).
## III The BTZ black hole
To set the stage properly, let us take the BTZ solution btz
$$ds^2=\left(M\frac{r^2}{\mathrm{}^2}\right)dt^2Jdtd\varphi +r^2d\varphi ^2+\frac{dr^2}{M+\frac{r^2}{\mathrm{}^2}+\frac{J^2}{4r^2}}$$
(4)
as a first example. The correct black hole vacuum background is found by setting $`M=0`$, $`J=0`$ in (4) (see btz for a discussion on this) which is clearly locally AdS:
$$ds^2=\frac{r^2}{\mathrm{}^2}dt^2+\frac{\mathrm{}^2}{r^2}dr^2+r^2d\varphi ^2.$$
The timelike $`\overline{\xi }^\mu =(/t)^\mu `$ and the spacelike $`\overline{\zeta }^\mu =(/\varphi )^\mu `$ Killing vectors can be used in finding the conserved energy and the angular momentum, respectively. The surface integral (1) at some finite distance $`r`$ from the origin yields the following non gauge-invariant quantities, which afterwards give the ‘true’ energy and angular momentum that are only to be measured at infinity <sup>3</sup><sup>3</sup>3Throughout we have chosen the Newton constant $`G`$ in (1) such that one finds the usual ADM pair $`(M,J)`$ in the limit $`\mu \mathrm{}`$.:
$`E(r)`$ $`=`$ $`{\displaystyle \frac{4r^4(J\mu M\mathrm{}^2)+J\mathrm{}^2r^2(\mu J4M)+J^3\mathrm{}^2}{4\mu \mathrm{}^2r^4+4\mu M\mathrm{}^4r^2\mu J^2\mathrm{}^4}},`$
$`L(r)`$ $`=`$ $`{\displaystyle \frac{8r^5(\mu JM)+Jr^3(J8\mu M\mathrm{}^2)+J^2\mathrm{}^2r(2\mu JM)}{2\mu r(4r^44M\mathrm{}^2r^2+J^2\mathrm{}^2)}}.`$
As a result, one obtains the energy and the angular momentum in the limit as $`r\mathrm{}`$ to be
$$E=M\frac{J}{\mu \mathrm{}^2}\text{and}L=J\frac{M}{\mu }.$$
These quantities dal ; dkt are obviously different from the ADM charges of the BTZ black hole btz ; the Cotton part clearly has a nontrivial contribution to the conserved charges. Amusingly enough, the angular momentum vanishes when the two parameters $`M`$ and $`J`$ are related by $`M=\mu J`$, in which case $`E=M(11/(\mu ^2\mathrm{}^2))`$. Thus, if furthermore $`\mu ^2\mathrm{}^2=1`$, then the BTZ black hole is left with no ‘energy’ and ‘angular momentum’ in the TMG context!
As a brief remark on the charged version of the BTZ solution mtz , we note that since the “electric potential” rises logarithmically in $`D=3`$, even a cursory look suggests that a single charged black hole will have divergent energy. In fact the authors of mtz define the energy of their charged rotating solution only upto an infinite constant factor. In this respect, the gauge invariant energy in the sense of db1 is naturally found to be divergent for this case.
## IV The supersymmetric solution
The half supersymmetry preserving solution given in ds is described by the metric
$$ds^2=f^2(\rho )dt^2+d\rho ^2+h^2(\rho )\left[d\varphi +a(\rho )dt\right]^2$$
(5)
and depending on whether the cosmological constant $`\mathrm{\Lambda }=1/\mathrm{}^2<0`$ is present or not, the metric functions are given by either <sup>4</sup><sup>4</sup>4The constants $`\beta _0`$ and $`\beta _3`$ were set equal to 1 in ds . Here we keep them for later convenience.
i) nonvanishing cosmological constant:
$`f(\rho )`$ $`=`$ $`f_0e^{2\rho /\mathrm{}}X^{1/2},h(\rho )=h_0X^{1/2},a(\rho )=a_0+k{\displaystyle \frac{f_0}{h_0}}e^{2\rho /\mathrm{}}X^1,`$
$`X(\rho )`$ $``$ $`\beta _0+\beta _1e^{2\rho /\mathrm{}}+\beta _2e^{(1/\mathrm{}\mu k)\rho },`$ (6)
or
ii) vanishing cosmological constant:
$`f(\rho )`$ $`=`$ $`f_0Y^{1/2},h(\rho )=h_0Y^{1/2},a(\rho )=a_0+k{\displaystyle \frac{f_0}{h_0}}Y^1,`$
$`Y(\rho )`$ $``$ $`\beta _3e^{\mu k\rho }\mu \beta _4(\omega _0+k\rho ).`$ (7)
Here $`f_0`$, $`h_0`$, $`a_0`$, $`\beta _i`$ ($`i=0,1,\mathrm{},4`$) and $`\omega _0`$ are all real constants that arise from the integration of the field equations whereas $`k=\pm 1`$ is a free parameter that comes from the solution of the Killing spinor equation on the supersymmetry side.
In ds , it was impossible to explicitly invert the functional relation $`r=h(\rho )`$ for the case of the nonvanishing cosmological constant so that the metric could be brought to the well-studied BTZ form btz , and the vast literature on that metric could be suitably adopted for an analysis of the physical meanings of the integration constants above. Instead, a much more complicated analysis was carried out by studying the quasilocal mass and the quasilocal angular momentum which was developed in do in an AdS background. We refer the reader to ds for the details of this.
Here a brief remark stating the differences between these quasilocal charges and the gravitational charges in the sense of db1 ; db2 are in order perhaps: The quasilocal energy in a spatially bounded region (such as an asymptotically AdS background for our case) is defined as minus the ‘time’ rate of change of the classical gravitational action. An analogous definition exists also for the quasilocal angular momentum. (Please see byork and the references therein for the attempts to define “quasilocal gravitational charges”.) The definition of gravitational charges in TMG, however, are much more natural since these gauge invariant conserved (global) charges are forged into being by the Gauss law and the presence of asymptotic Killing symmetries db1 ; db2 .
For the time being, let us concentrate on the case of nonvanishing cosmological constant $`\mathrm{\Lambda }=1/\mathrm{}^20`$. By substituting the metric functions (6) in the metric (5), one obtains
$$ds^2=d\rho ^2\frac{f_0^2e^{4\rho /\mathrm{}}}{X(\rho )}dt^2+\frac{f_0^2e^{4\rho /\mathrm{}}}{X(\rho )}\left(dt+k\frac{h_0}{f_0}e^{2\rho /\mathrm{}}X(\rho )(d\varphi a_0dt)\right)^2$$
(8)
after some simplifications.
In ds , it was found that the quasilocal mass was $`a_0`$ times the quasilocal angular momentum (see (39) of ds ) and the asymptotic behavior of the metric was examined through the metric function $`a(r)`$ and hence $`a_0`$. It was shown that for $`a`$ to vanish asymptotically as $`r\mathrm{}`$, $`a_0`$ had to be chosen either as 0 or as $`kf_0/(h_0\beta _1)`$. Whether $`a_0`$ has a physical meaning or not (and whether it can be set equal to zero or not), one should be able to make the simple change of variable $`d\theta =d\varphi a_0dt`$ in the metric (8). The outcome is simply
$$ds^2=d\rho ^2+2kf_0h_0e^{2\rho /\mathrm{}}dtd\theta +h_0^2X(\rho )d\theta ^2.$$
Another simple redefinition of the coordinates as $`u=kf_0t`$ and $`v=h_0\theta `$ can always be made at this stage and one arrives at the final form
$$ds^2=d\rho ^2+2e^{2\rho /\mathrm{}}dudv+\left(\beta _0+\beta _1e^{2\rho /\mathrm{}}+\beta _2e^{(1/\mathrm{}\mu k)\rho }\right)dv^2.$$
(9)
It is obvious that one of the integration constants in (9) can be set to 1 by simple coordinate rescalings. The curvature invariants of this metric can be calculated easily: the Ricci scalar $`R=6/\mathrm{}^2`$ and $`R_{\mu \nu }R^{\mu \nu }=12/\mathrm{}^4`$, moreover this solution is asymptotically AdS (for $`1/\mathrm{}\mu k<0`$) with no curvature singularities. When $`\beta _0=\beta _2=0`$, the metric is the AdS metric in the Poincaré coordinates. For this case, even when one starts with $`\beta _1=0`$, one can still introduce it back by a simple coordinate redefinition as $`\stackrel{~}{u}=u\beta _1v/2`$.
There is yet another alternative way to understand the emergence of the constants $`\beta _0`$ and $`\beta _1`$ in the expression for $`X(\rho )`$. These two terms can be thought of as describing a gravitational wave in AdS. One can use the technique developed by Garfinkle and Vachaspati gar which permits the addition of a wave to an already existing solution when there is a null Killing vector present. A detailed discussion of this method, its extension to various supergravity theories and to theories that include nontrivial matter couplings can be found in myers . One can easily follow the footsteps of myers and conclude without any difficulties that the Garfinkle-Vachaspati method is also applicable to the TMG theory. A brief outline of this technique can also be found in the appendix of sad and here we will use the notation outlined there. In our case, one starts from (9) with $`X(\rho )=0`$, i.e. the AdS metric
$$ds^2=d\rho ^2+2e^{2\rho /\mathrm{}}dudv$$
(10)
which satisfies $`R_{\mu \nu }=(2/\mathrm{}^2)g_{\mu \nu }`$. Using the null Killing vector $`k^\mu =(/v)^\mu `$, the scalar $`\mathrm{\Omega }`$ (of sad ) is calculated easily as $`\mathrm{\Omega }=\mathrm{\Omega }_0e^{2\rho /\mathrm{}}`$, where $`\mathrm{\Omega }_0`$ is an arbitrary constant. Following relevant steps, one obtains $`\mathrm{\Phi }(\rho )=\mathrm{\Phi }_1\mathrm{}\mathrm{\Phi }_0e^{2\rho /\mathrm{}}`$, where $`\mathrm{\Phi }_0`$ and $`\mathrm{\Phi }_1`$ are arbitrary real constants. Using these, defining $`\mathrm{\Omega }_0\mathrm{}\mathrm{\Phi }_0\beta _0`$ and $`\mathrm{\Omega }_0\mathrm{\Phi }_1\beta _1`$, one obtains (9) with $`\beta _2`$ set equal to zero in $`X(\rho )`$.
In fact the metric with $`\beta _2=0`$ has showed up earlier in different contexts: It corresponds to a generalized Kaigorodov metric kai . It is obtainable from the AdS metric by an $`SL(2,R)`$ transformation brec and its equivalence to the extremal limit of the BTZ black hole btz can be shown brec ; pope . A crucial point is that the boundaries of the AdS and the extremal BTZ metric are different (see pope for details). When $`\beta _20`$, one can again remove the constant $`\beta _0`$ by a shift in the $`\rho `$ coordinate.
Another observation that needs to be stated is that in fact the constants $`\beta _0`$, $`\beta _1`$ and $`\beta _2`$ can be taken as arbitrary functions of $`v`$ and the metric (9) is also a solution to the TMG equations with a cosmological constant when $`X(\rho )`$ is replaced by
$$X(\rho ,v)\beta _0(v)+\beta _1(v)e^{2\rho /\mathrm{}}+\beta _2(v)e^{(1/\mathrm{}\mu k)\rho }.$$
These arbitrary functions can be thought of as describing the profile of the gravitational wave then. However when they are left arbitrary, it is highly probable that the supersymmetry is completely broken.
Let us now calculate the gravitational charges associated with the metric (9) using the procedure outlined in section II. It is clear that the metric (10) can be used as the background with its timelike $`\overline{\xi }^\mu =(/u+/v)^\mu `$ and spacelike $`\overline{\zeta }^\mu =(/u+/v)^\mu `$ Killing vectors yielding the energy and the angular momentum, respectively. After a tedious calculation, one finds that the integrand (that is, the terms inside the parentheses) in (1) is given by
$$E(\rho )=\delta _u^\mu \delta _\rho ^i\frac{1}{2\mu \mathrm{}^2}\left\{4\beta _0(1+\mu \mathrm{})+\beta _2e^{(1/\mathrm{}\mu k)\rho }(1+\mu k\mathrm{})\left[1+(k+2)\mu \mathrm{}\right]\right\},$$
and obviously depending on the sign of $`1/\mathrm{}\mu k`$, one finds that the energy at the boundary of AdS $`(\rho \mathrm{})`$ is <sup>5</sup><sup>5</sup>5Here we have chosen the Newton constant $`G`$ in (1) accordingly.
$$E=\{\begin{array}{cc}2(k/\mathrm{})(1+k)(\beta _0+\beta _2),& 1/\mathrm{}\mu k=0\hfill \\ 2\beta _0(1+\mu \mathrm{})/(\mu \mathrm{}^2),& 1/\mathrm{}\mu k<0\hfill \end{array}.$$
As for the angular momentum, one finds that it is equal to the energy: $`L=E`$.
The steps that have been taken up until this point can also be repeated in an analogous fashion for the case of vanishing cosmological constant. The metric that corresponds to (8) in such a process is simply found as
$$ds^2=d\rho ^2+2dtd\theta +\left(\beta _3e^{\mu k\rho }\mu \beta _4(\omega _0+k\rho )\right)d\theta ^2.$$
(11)
If one starts from the flat metric $`ds^2=d\rho ^2+2dtd\theta `$ and applies the Garfinkle-Vachaspati method to add a gravitational wave to this spacetime, one readily finds $`\mathrm{\Omega }=\mathrm{\Omega }_0`$ and $`\mathrm{\Phi }(\rho )=\mathrm{\Phi }_1+\mathrm{\Phi }_0\rho `$, where $`\mathrm{\Phi }_0`$, $`\mathrm{\Phi }_1`$ and $`\mathrm{\Omega }_0`$ are arbitrary real constants, and these, with the definitions $`\mathrm{\Omega }_0\mathrm{\Phi }_1\mu \beta _4\omega _0`$ and $`\mathrm{\Omega }_0\mathrm{\Phi }_0\mu k\beta _4`$, lead to the metric (11) with $`\beta _3`$ set equal to zero in $`Y(\rho )`$.
Once again the constants $`\beta _4`$, $`\omega _0`$ and $`\beta _3`$ can be taken as arbitrary functions of $`\theta `$ and the metric (11) is also a solution to the TMG equations when $`Y(\rho )`$ is replaced by
$$Y(\rho ,\theta )\beta _3(\theta )e^{\mu k\rho }\mu \beta _4(\theta )(\omega _0(\theta )+k\rho ).$$
These arbitrary functions can be thought of as describing the wave profile again, but with these functions in place, it may be that there is no supersymmetry left to preserve then.
As for the gravitational charges related to (11), the background to work with is simply the flat metric $`ds^2=d\rho ^2+2dtd\theta `$ and the Killing vectors needed are just the timelike $`\overline{\xi }^\mu =(/t+/\theta )^\mu `$ and the spacelike $`\overline{\zeta }^\mu =(/t+/\theta )^\mu `$ vectors. The steps leading to (1) can easily be repeated by setting $`\mathrm{\Lambda }=0`$ in the relevant places and replacing the covariant derivatives with respect to $`\overline{g}_{\mu \nu }`$ with ordinary derivatives. One then finds that the terms inside the parentheses in (1) is given by
$$E(\rho )=\delta _t^\mu \delta _\rho ^i\left(\mu k\beta _4+\frac{1}{2}\mu \beta _3(2k1)e^{\mu k\rho }\right),$$
and depending on the sign of $`\mu k`$, one finds the energy as $`\rho \mathrm{}`$ to be
$$E=\{\begin{array}{cc}0,& \mu =0\hfill \\ \mu k\beta _4,& \mu k>0\hfill \end{array}.$$
The angular momentum is again given by $`L=E`$.
## V Conclusions
In this paper, we showed how the physical properties of a known supersymmetric solution of the full cosmological TMG theory can be better understood by the Garfinkle-Vachaspati method and then determined its conserved charges for bulk asymptotically flat and constant curvature backgrounds. Even though the question of how the supersymmetric version of the TMG theory can be obtained from any compactification of M-theory or, for that matter, any higher dimensional supergravity theory remains open, provided that an exact form of the CFT dual of TMG can be formulated on the boundary of AdS, this particular supersymmetric solution should be suitable for understanding the AdS/CFT duality in the infinite momentum frame brec , pope . Another open question that deserves attention is the problem of finding a supersymmetric matter coupled extension of the TMG theory. Looking for the charged versions of the metrics studied here within this model and their relations with the ones presented in ds2 would be certainly worth the effort.
## VI Acknowledgments
We thank G. Clément for useful discussions. This work is partially supported by the Scientific and Technical Research Council of Turkey (TÜBİTAK); the work of B.T. is also supported by the “Young Investigator Fellowship” of the Turkish Academy of Sciences (TÜBA) and by a TÜBİTAK Kariyer Grant. |
warning/0507/astro-ph0507147.html | ar5iv | text | # Late-time Radio Observations of 68 Type Ibc Supernovae: Strong Constraints on Off-Axis Gamma-ray Bursts
## 1. Introduction
It is now generally accepted that long duration gamma-ray bursts (GRBs) give rise to engine-driven relativistic jets as well as non-relativistic spherical supernova (SN) explosions. The first example of this GRB-SN connection came with the discovery of the Type Ic supernova, SN 1998bw, associated with GRB 980425 ($`d36`$ Mpc; Galama et al. 1998; Pian et al. 2000). The unusually fast photospheric velocities and exceptionally bright radio emission of SN 1998bw indicated $`10^{52}`$ erg of kinetic energy and mildly relativistic ejecta (bulk Lorentz factor, $`\mathrm{\Gamma }3`$; Kulkarni et al. 1998; Iwamoto et al. 1998; Li & Chevalier 1999; Woosley, Eastman & Schmidt 1999). In comparison with other core-collapse events ($`E_{KE}10^{51}`$ erg and ejecta speeds, $`v0.1c`$), SN 1998bw was considered a hyper-energetic supernova (“hypernova”; Iwamoto et al. 1998). Broad optical absorption lines were also observed in the Type Ic SNe 2003dh and 2003lw, associated with the cosmological GRBs 030329 and 031203, indicative of comparably large photospheric velocities (Matheson et al. 2003b; Malesani et al. 2004). Together, these observations appear to suggest that broad spectral features are characteristic of GRB-associated SNe.
In addition to events with prompt gamma-ray emission, the GRB-SN connection also implies the existence of “orphan” supernovae whose relativistic jets are initially beamed away from our line of sight (Rhoads 1999; Paczynski 2001). Since the discovery of SN 1998bw, several broad-lined SNe have been identified locally ($`d100`$ Mpc) and are currently estimated to represent $`5\%`$ of the Type Ibc supernova (SNe Ibc) population (Podsiadlowski et al. 2004). Given their spectral similarity to the GRB-associated SNe, it has been argued that local broad-lined supernovae can be used as signposts for GRBs.
Thus, associations with poorly-localized BATSE bursts have been invoked for the broad-lined SNe 1997cy, 1997ef and 1999E<sup>1</sup><sup>1</sup>1We note that SNe 1997cy and 1999E were initially classified as Type IIn supernovae while Hamuy et al. (2003) later showed convincing evidence that they are hydrogen-rich Type Ia events similar to SN 2002ic. (Germany et al. 2000; Turatto et al. 2000; Wang & Wheeler 1998; Mazzali, Iwamoto & Nomoto 2000; Rigon et al. 2003). In addition, association with off-axis GRBs have also been claimed. In the case of SN 2002ap, broad optical absorption lines and evidence for mildly asymmetric ejecta (based on spectropolarimetry measurements) were interpreted to support an off-axis GRB jet (Kawabata et al. 2002; Totani 2003, but see Leonard et al. 2002).
More recently, an off-axis GRB model has been proposed for SN 2003jd, for which photospheric velocities upward of 40,000 $`\mathrm{km}\mathrm{s}^1`$ were measured at early time (Filippenko, Foley & Swift 2003; Matheson et al. 2003a). More intriguingly, late-time ($`t400`$ days) spectra showed double-peaked emission lines of light-elements, attributed to an asymmetric explosion (Kawabata et al. 2004). Mazzali et al. (2005) argue that these observations can be understood if SN 2003jd was accompanied by a highly collimated GRB jet initially directed $`70`$ degrees away from our line-of-sight.
Regardless of viewing angle, however, strong afterglow emission eventually becomes visible as the decelerating GRB jets spread laterally and the emission becomes effectively isotropic. As the jets spread into our line-of-sight, a rapid increase of broadband synchrotron emission is observed on a timescale of a few weeks to several years. This late-time emission is most easily detected at long wavelengths (Perna & Loeb 1998; Levinson et al. 2002; Waxman 2004). Targeting local Type Ibc supernovae with late-time radio observations has thus become the preferred method to search for evidence of off-axis GRBs (Stockdale et al. 2003; Soderberg, Frail & Wieringa 2004).
Using early radio observations ($`t100`$ days) we have already limited the fraction of SNe Ibc harboring on-axis (or mildly off-axis) GRBs to be $`3\%`$ (Berger et al. 2003a). In this paper, we present late-time ($`t0.5`$ to 20 yr) radio observations for 68 local Type Ibc supernovae, including SN 2003jd and five additional broad-lined events, making this the most comprehensive study of late-time radio emission from SNe Ibc. We use these data to constrain the SN fraction associated with GRB jets regardless of viewing angle assumptions, constraining even those initially beamed perpendicular to our line-of-sight.
## 2. Radio Observations
### 2.1. Type Ic SN 2003jd
SN 2003jd was discovered on 2003 October 25.2 UT within host galaxy MCG -01-59-021 ($`d_L81`$ Mpc; Burket et al. 2003). In Table 1 we summarize our radio observations for SN 2003jd, spanning $`8569`$ days after the explosion<sup>2</sup><sup>2</sup>2Here we assume an approximate explosion date of 2003 October 21 UT, based on pre-explosion images (Burket et al. 2003).. All observations were conducted with the Very Large Array<sup>3</sup><sup>3</sup>3The Very Large Array is a facility of the National Science Foundation operated under cooperative agreement by Associated Universities, Inc. (VLA) in the standard continuum mode with a bandwidth of $`2\times 50`$ MHz centered at 4.86, 8.46 or 22.5 GHz. We used 3C48 and 3C147 (J0137+331 and J0542+498) for flux calibration, while J2323-032 was used to monitor the phase. Data were reduced using standard packages within the Astronomical Image Processing System (AIPS).
No radio emission was detected at the optical SN position during our early observations. Our radio limits imply that SN 2003jd was a factor of $`100`$ less luminous than SN 1998bw on a comparable timescale. We conclude that SN 2003jd, like the majority of SNe Ibc, did not produce relativistic ejecta along our line-of-sight.
We re-observed SN 2003jd at $`t1.6`$ yrs to search for radio emission from an off-axis GRB jet. No emission was detected, implying a limit of $`F_\nu <45\mu `$Jy ($`3\sigma `$) at 8.46 GHz.
### 2.2. Late-time data on Local Type Ibc Supernovae
We supplement these data with late-time ($`t0.520`$ year) radio observations for 67 local ($`d_L200`$ Mpc) SNe Ibc, summarized in Table 2. Eleven objects were observed at moderately late-time as part of our on-going VLA program to characterize the early ($`t100`$ days) radio emission from SNe Ibc (Soderberg et al., in prep). The remaining 54 objects were observed on a later timescale ($`t1`$ year) and were taken from the VLA archive<sup>4</sup><sup>4</sup>4http://e2e.nrao.edu/archive/. We note that five of these supernovae (SNe 1997dq, 1997ef, 1998ey, 2002ap, 2002bl) were spectroscopically observed to have broad optical absorption lines, similar to SN 1998bw.
All VLA observations were conducted at 8.46 GHz (except for SN 1991D at 4.86 GHz) in the standard continuum mode with a bandwidth of $`2\times 50`$ MHz. Data were reduced using AIPS, and the resulting flux density measurements for this sample of SNe Ibc is given in Table 2. With the exception of SN 2001em, from which radio emission from the non-relativistic, spherical supernova ejecta is still detected at late-time (Stockdale et al. 2005; Bietenholz & Bartel 2005, but see Granot & Ramirez-Ruiz 2004), none of the SNe Ibc show radio emission above our average detection limit of $`0.15`$ mJy ($`3\sigma `$). In comparison with SN 1998bw, only SN 2001em shows a comparable radio luminosity on this timescale. These results are consistent with the earlier report by Stockdale et al. (2003).
In Figure 1 we plot the radio observations for this sample of SNe Ibc, in addition to late-time radio data for SN1954A (Eck, Cowan & Branch 2002) and SN 1984L (Soderberg, Frail & Wieringa 2004).
## 3. Off-Axis Models for Gamma-ray Bursts
### 3.1. An Analytic Approach
Waxman (2004) present an analytic model for the late-time radio emission from a typical GRB viewed significantly away from the collimation axis. In this model, the GRB jet is initially characterized by a narrow opening angle, $`\theta _j`$ few degrees, while the viewing angle is assumed to be large, $`\theta _{\mathrm{obs}}1`$ radian. As the jet sweeps up circumstellar material (CSM) and decelerates, it eventually undergoes a dynamical transition to sub-relativistic expansion (Frail, Waxman & Kulkarni 2000). The timescale for this non-relativistic transition is estimated at $`t_{NR}0.2(E_{51}/n_0)^{1/3}`$ yr ($`0.3E_{51}/A_{}`$ yr) in the case of a homogeneous (wind-stratified) medium (Waxman 2004). Here, $`E_{51}`$ is the beaming-corrected ejecta energy normalized to $`10^{51}`$ erg and $`n_0`$ is the circumstellar density of the homogeneous medium (interstellar medium; ISM) normalized to 1 particle cm<sup>-3</sup>. For a wind-stratified medium, $`A_{}`$ defines the circumstellar density in terms of the progenitor mass loss rate, $`\dot{M}`$, and wind velocity, $`v_w`$, such that $`\dot{M}/4\pi v_w=5\times 10^{11}A_{}\mathrm{g}\mathrm{cm}^1`$, and thus $`A_{}=1`$ for $`\dot{M}=10^5M_{}\mathrm{yr}^1`$ and $`v_w=10^3\mathrm{km}\mathrm{s}^1`$ (Li & Chevalier 1999).
Once sub-relativistic, the jets spread sideways, rapidly intersecting our line-of-sight as the ejecta approach spherical symmetry. At this point the afterglow emission is effectively isotropic and appears similar to both on-axis and off-axis observers. The broadband emission observed from the sub-relativistic ejecta is described by a standard synchrotron spectrum, characterized by three break frequencies: the synchrotron self-absorption frequency, $`\nu _a`$, the characteristic synchrotron frequency, $`\nu _m`$, and the synchrotron cooling frequency, $`\nu _c`$. On timescales comparable to the non-relativistic transition, $`\nu _a`$ and $`\nu _m`$ are typically below the radio band while $`\nu _c`$ is generally near the optical (Frail, Waxman & Kulkarni 2000; Berger, Kulkarni & Frail 2004; Frail et al. 2005). Making the usual assumption that the kinetic energy is partitioned between relativistic electrons and magnetic fields ($`ϵ_e`$ and $`ϵ_B`$, respectively), and that these fractions are constant throughout the evolution of the jet, Waxman (2004) estimate the radio luminosity of the sub-relativistic, isotropic emission to be
$`L_\nu `$ $``$ $`8.0\times 10^{29}\left({\displaystyle \frac{ϵ_e}{0.1}}\right)\left({\displaystyle \frac{ϵ_B}{0.1}}\right)^{3/4}n_0^{3/4}E_{51}`$
$`\times \left({\displaystyle \frac{\nu }{10\mathrm{G}\mathrm{H}\mathrm{z}}}\right)^{1/2}\left({\displaystyle \frac{t}{t_{\mathrm{NR}}}}\right)^{9/10}\mathrm{erg}\mathrm{s}^1\mathrm{Hz}^1`$
for the ISM case, while for a wind-stratified medium
$`L_\nu `$ $``$ $`2.1\times 10^{29}\left({\displaystyle \frac{ϵ_e}{0.1}}\right)\left({\displaystyle \frac{ϵ_B}{0.1}}\right)^{3/4}A_{}^{9/4}E_{51}^{1/2}`$
$`\times \left({\displaystyle \frac{\nu }{10\mathrm{G}\mathrm{H}\mathrm{z}}}\right)^{(p1)/2}\left({\displaystyle \frac{t}{t_{\mathrm{NR}}}}\right)^{3/2}\mathrm{erg}\mathrm{s}^1\mathrm{Hz}^1.`$
Here it is assumed that the electrons are accelerated into a power-law distribution, $`N(\gamma )\gamma ^p`$ with $`p=2.0`$. These equations reveal that the strength of the non-relativistic emission is strongly dependent on the density of the circumstellar medium (especially in the case of a wind) and is best probed at low frequencies.
While this analytic model provides robust predictions for the afterglow emission at $`t>t_{\mathrm{NR}}`$, it does not describe the early evolution or the transition from relativistic to sub-relativistic expansion. At early time, the observed emission from an off-axis GRB is strongly dependent on the viewing angle and dynamics of the jet. To investigate this early afterglow evolution and the transition to sub-relativistic expansion, we developed a detailed semi-analytic model, described below.
### 3.2. A Semi-analytic Model
In modeling the afterglow emission from an off-axis GRB jet, we adopt the standard framework for a adiabatic blastwave expanding into either a uniform or wind stratified medium (Sari 1997; Granot & Sari 2002). We assume a uniform, sharp-edged jet such that Lorentz factor and energy are constant over the jet surface. The hydrodynamic evolution of the jet is fully described in Oren, Nakar & Piran (2004). As the bulk Lorentz factor of the ejecta approaches $`\mathrm{\Gamma }1`$, the jets begin to spread laterally at the sound speed <sup>5</sup><sup>5</sup>5Since the spreading behavior of relativistic GRB jets is poorly constrained by observations, we assume negligible spreading during this phase. We adopt this conservative assumption since it produces the faintest off-axis light-curves.. Our off-axis light-curves are obtained by integrating the afterglow emission over equal arrival time surface. We note that these resulting light-curves are in broad agreement with Model 2 of Granot et al. (2002) and are consistent with Waxman’s analytic model (§3.1) on timescales, $`tt_{\mathrm{NR}}`$.
Over-plotted in Figure 1 are our off-axis models calculated for both wind-stratified and homogeneous media at an observing frequency of $`\nu _{\mathrm{obs}}=8.46`$ GHz. We assume standard GRB parameters of $`E_{51}=A_{}=n=1`$, $`ϵ_B=ϵ_e=0.1`$, $`p=2.2`$ and $`\theta _j=5^\mathrm{o}`$, consistent with the typical values inferred from broadband modeling of GRBs (Panaitescu & Kumar 2002; Yost et al. 2003; Chevalier, Li & Fransson 2004). We compute model light-curves for off-axis viewing angles between 30 and 90 degrees. As clearly shown in the figure, the majority of our late-time SNe Ibc limits are significantly fainter than all of the model light-curves, constraining even the extreme case where $`\theta _{\mathrm{obs}}=90^\mathrm{o}`$.
## 4. SN 2003jd: Constraints on the off-axis jet
Based on the double-peaked profiles observed for the nebular lines of neutral oxygen and magnesium, Mazzali et al. (2005) argue that SN 2003jd was an aspherical, axisymmetric explosion viewed near the equatorial plane. They suggest that this asymmetry may be explained if the SN explosion was accompanied by a tightly collimated and relativistic GRB jet, initially directed $`70`$ degrees from our line-of-sight. This hypothesis is consistent with the observed lack of prompt gamma-ray emission (Hurley et al. 2003) as well as the absence of strong radio and X-ray emission at early time (Soderberg, Kulkarni & Frail 2003; Watson et al. 2003).
Our radio observation of SN 2003jd at $`t1.6`$ years imposes strong constraints on the putative off-axis GRB jet. While the early data constrain only mildly off-axis jets ($`\theta _{\mathrm{obs}}30^\mathrm{o}`$), our late-time epoch constrains even those jets initially directed perpendicular to our line-of-sight. As shown in Figure 1, our radio limit is a factor of $`200`$ ($`20`$) fainter than that predicted for a typical GRB expanding into a homogeneous (wind-stratified) medium, even in the extreme case where $`\theta _{\mathrm{obs}}90^\mathrm{o}`$. Given the assumption of typical GRB parameters, we conclude that our late-time radio limit is inconsistent with the presence of an off-axis GRB jet. We note that the model assumptions and physical parameters of our off-axis afterglow light-curves are identical to those adopted by Mazzali et al. (2005).
We next explore the range of parameters ruled out by our deep radio limits. As shown in Equations 1 and 2, the luminosity of the late-time emission is a function of the ejecta energy, the density of the circumstellar medium and the equipartition fractions. To investigate the effect of energy and density on the late-time radio luminosity, we fix the equipartition fractions to $`ϵ_e=ϵ_B=0.1`$, chosen to be consistent with the values typically inferred from afterglow modeling of cosmological GRBs (Panaitescu & Kumar 2002; Yost et al. 2003).
In Figure 2, we illustrate how each radio epoch for SN 2003jd maps to a curve within the two-dimensional parameter space of kinetic energy and circumstellar density for an off-axis GRB. Here we adopt our semi-analytic model (§3.2) for a wind-stratified medium, along with a typical electron index of $`p=2.2`$ and a viewing angle of $`\theta _{\mathrm{obs}}=90^\mathrm{o}`$; the faintest model for a given set of equipartition fractions. By comparing the luminosity limit for SN 2003jd at a particular epoch with the off-axis model prediction for that time, we exclude the region of parameter space rightward of the curve since this region produces a jet which is brighter than the observed limit. The union of these regions represents the total parameter space ruled out for an associated GRB. As shown in this figure, the total excluded parameter space extends from $`A_{}0.03`$ and $`E10^{47}`$ to $`10^{52}`$.
We compare these constraints with the beaming-corrected kinetic energies and CSM densities for 18 cosmological GRBs (Table 3). Here we make the rough approximation that $`A_{}n_0`$; a reasonable assumption for circumstellar radii near $`10^{18}`$ cm. As shown in Figure 2, these GRBs span the region of parameter space roughly bracketed by $`A_{}0.002`$ to 100 and $`E2\times 10^{49}`$ to $`4\times 10^{51}`$. The majority of the bursts (13 out of 18) fall within the excluded region of parameter space for SN 2003jd. We conclude that SN 2003jd was not likely associated with a typical GRB at a confidence level of $`72\%`$.
## 5. Local Type Ibc Supernovae: Further Constraints
While physical parameters atypical of the cosmological GRB population can be invoked to hide an off-axis GRB for SN 2003jd, it is exceedingly unlikely for atypical parameters to dominate a large statistical sample of SNe Ibc. Motivated thus, we searched for off-axis GRBs in the 67 local Type Ibc SNe for which we have compiled late-time ($`t0.530`$ yr) radio observations. Applying the method described in §4 we produce exclusion regions in the $`E_{51}A_{}`$ parameter space for each SN. Figure 3 shows the resulting contours for all 68 SNe, including SN 2003jd and five broad-lined events. For the twenty SNe with early radio limits (Berger et al. 2003a; Berger, Kulkarni & Chevalier 2002) we combine late- and early-time data to provide further constraints.
In Figure 4 we compile all 68 exclusion regions to quantify the $`E_{51}A_{}`$ parameter space constrained by this statistical sample. Contours map the regions excluded by incremental fractions of our sample. As in the case of SN 2003jd, all curves rule out bursts with $`A_{}1`$ and $`E10^{50}`$ erg. Moreover, 50% exclude $`A_{}0.1`$ and $`E10^{49}`$ erg. For comparison, the mean ejecta energy and CSM density values for cosmological GRBs are $`E4.4\times 10^{50}`$ erg and $`A_{}=n_01.2`$.
Focusing on the subsample of broad-lined SNe, we emphasize that our deep limits rule out both putative GRB jets directed along our line-of-sight (e.g. SN 1997ef) as well as those which are initially beamed off-axis (e.g. SN 2002ap and SN 2003jd). In particular, the large exclusion region for SN 2002ap (see Figure 3) implies that an extremely low CSM density, less than $`A_{}3\times 10^3`$, is needed to suppress the emission from an associated GRB. This is a factor of $`10`$ below the density inferred from modeling of the early radio emission (Berger, Kulkarni & Chevalier 2002) and we therefore conclude that an off-axis GRB model is inconsistent with our late-time observations of SN 2002ap. In Figure 4 we show that this entire sample of six broad-lined SNe rule out bursts with energies $`E10^{49}`$ erg, and 50% even rule out $`E10^{47}`$ erg (all assuming a typical $`A_{}=1`$).
We next address the limits on an association with GRBs defined by the cosmological sample (Table 3). For each SN in our sample we calculate the fraction of observed GRBs that lie in its exclusion region. We then determine the probability of finding null-detections for our entire sample by calculating the product of the individual probabilities. We find that the probability that every Type Ibc supernova has an associated GRB is $`1.1\times 10^{10}`$. We further rule out a scenario in which one in ten SNe Ibc is associated with a GRB at a confidence level of $`90\%`$. For the broad-lined events alone we rule out the scenario that every event is associated with a GRB at a confidence level of $`84\%`$. Confidence levels are shown as a function of GRB/SN fraction in Figure 5.
## 6. Discussion and Conclusions
We present late-time radio observations for 68 local Type Ibc supernovae, including six broad-lined SNe (“hypernovae”), making this the most comprehensive study of late-time radio emission from SNe Ibc. None of these objects show evidence for bright, late-time radio emission that could be attributed to off-axis jets coming into our line-of-sight. Comparison with our most conservative off-axis GRB afterglow models reveals the following conclusions:
(1) Less than $`10\%`$ of Type Ibc supernovae are associated with GRBs. These data impose an empirical constraint on the GRB beaming factor, $`f_b^1`$, where $`f_b=(1\mathrm{cos}\theta _j)`$. Assuming a local GRB rate of $`0.5\mathrm{Gpc}^3\mathrm{yr}^1`$ (Schmidt 2001; Perna, Sari & Frail 2003; Guetta, Piran & Waxman 2005) and an observed SNe Ibc rate of $`4.8\times 10^4\mathrm{Gpc}^3\mathrm{yr}^1`$ (Marzke et al. 1998; Cappellaro, Evans & Turatto 1999; Folkes et al. 1999), we constrain the GRB beaming factor to be $`f_b^1\times 10^4`$. Adopting a lower limit of $`f_b^1>13`$ (Levinson et al. 2002), the beaming factor is now observationally bound by $`f_b^1[1310^4]`$, consistent with the observed distribution of jet opening angles (Frail et al. 2001; Guetta, Piran & Waxman 2005).
(2) Despite predictions that most or all broad-lined SNe Ibc harbor GRB jets (Podsiadlowski et al. 2004), our radio observations for six broad-lined events (SNe 1997dq, 1997ef, 1998ey, 2002ap, 2002bl and 2003jd) reveal no evidence for association with typical (or even sub-energetic) GRBs. While unusual physical parameters can suppress the radio emission from off-axis jets in any one SN, it is unlikely that all six broad-lined events host atypical GRBs. We observationally rule out the scenario in which every broad-lined SN harbors GRB jets with a confidence level of $`84\%`$.
(3) While low CSM densities (e.g. $`A_{}0.1`$) can suppress the emission from off-axis GRB jets, such values are inconsistent with the mass loss rates measured from local Wolf-Rayet stars ($`0.69.5\times 10^5\mathrm{M}_{}\mathrm{yr}^1`$; Cappa, Goss & van der Hucht 2004), thought to be the progenitors of long-duration gamma-ray bursts.
(4) While we have so far considered only the signature from a highly collimated GRB jet, these late-time radio data also impose constraints on the presence of broader jets and/or jet cocoons. As demonstrated by GRBs 980425 and 030329, the fraction of energy coupled to mildly relativistic and mildly collimated ejecta can dominate the total relativistic energy budget (Kulkarni et al. 1998; Berger et al. 2003b). Less sensitive to to the effects of beaming and viewing geometry, broad jets are more easily probed at early time ($`t100`$ days) when the emission is brightest. Still, we note that the majority of our late-time radio limits are significantly fainter than GRBs 980425 and 030329 on a comparable timescale, thus constraining even mildly relativistic ejecta.
These conclusions, taken together with the broad spectral features observed for GRB-associated SNe 1998bw, 2003dh and 2003lw, motivate the question: what is the connection between GRBs and local Type Ibc supernovae? While current optical data suggest that all GRB-SNe are broad-lined, our late-time radio observations clearly show that the inverse is not true: broad optical absorption lines do not serve as a reliable proxy for relativistic ejecta. This suggests that their observed large photospheric velocities and asymmetric ejecta are often merely characteristics of the non-relativistic SN explosion and thus manifestations of the diversity within SNe Ibc.
The authors thank Doug Leonard, Paolo Mazzali, Dale Frail, Brian Schmidt and Avishay Gal-Yam for helpful discussions. As always, the authors thank Jochen Greiner for maintaining his GRB page. A.M.S. is supported by the NASA Graduate Student Research Program. E.B. is supported by NASA through Hubble Fellowship grant HST-HF-01171.01 awarded by the STScI, which is operated by the Association of Universities for Research in Astronomy, Inc., for NASA, under contract NAS 5-26555. |
warning/0507/astro-ph0507330.html | ar5iv | text | # Model Simulations of a Shock-Cloud Interaction in the Cygnus Loop
## 1. Introduction
The interaction of shock waves with the interstellar medium (ISM) such as those associated with supernovae, stellar winds, bipolar flows, H II regions, or spiral density waves is a fundamental process in interstellar gas dynamics and is key to understanding the evolution and structure of the ISM. The highly nonlinear interaction between supernova generated shocks and interstellar clouds is often not suited to analytic approaches but requires a multidimensional hydrodynamics study of the shock-cloud problem using high resolution methods.
A hydrodynamical study of a shocked ISM cloud was made by Klein, McKee, & Colella (1994, hereafter KMC94), who found that the cloud may be destroyed by a series of instabilities associated with the post-shock flow of inter-cloud gas past the cloud. Earlier work on this problem includes that of Woodward (1976) and Nittman, Falle, & Gaskell (1982). More recently, Poludnenko, Frank, & Blackman (2002) studied the role that internal cloud structure plays in the destruction of the cloud.
For investigating shocked ISM cloud physics, the interaction between a supernova (SN) shock and low density diffuse ISM clouds is of particular interest. Supernova remnants (SNRs) shape and enrich the chemical and dynamical structure of the ISM which, in turn, affect the evolution of subsequent SNRs. The details of just how SN generated shock waves interact with interstellar clouds are not well understood.
There are several limiting factors in attempting to compare model simulations to observed SNR shock cloud interactions. While models can be viewed edge on and rotated in two or three dimensions, shocked interstellar clouds are viewed only in projection, which leads to a complex appearing shock structure due to multiple and overlapping shocks. In addition, one observes only a single epoch, i.e., a ‘snapshot’, of the interaction. These factors make it difficult to understand and model the time dependent kinematics and detailed dynamical processes of the interaction. Also, unlike how they are often modeled, real interstellar clouds are neither cylindrical or spherical in shape nor sharp edged, with interiors very likely non-uniform in density. Furthermore, many shocked interstellar clouds are dense enough so that radiative losses, which can alter the overall dynamics of the shock-cloud interaction, are important (Mellema et al., 2002; Fragile et al., 2004). Finally, the inclusion of an embedded magnetic field can drastically alter the dynamics of the interaction. For instance, a strong, ordered magnetic field can suppress dynamical instability growth predicted by fluid dynamical simulations (Mac Low et al., 1994; Fragile et al., 2005).
In looking for an ‘ideal’ shock-cloud interaction, one would like to avoid many of the aforementioned effects and the Cygnus Loop supernova remnant affords several distinct advantages. Because of the remnant’s large angular size (2.8$`{}_{}{}^{}\times `$ 3.5), low foreground extinction ($`E[BV]=0.08`$ mag; Parker 1967; Fesen, Blair, & Kirshner 1982), and wide range of shock conditions, the Cygnus Loop is one of the better locations for studying the ISM shock physics of middle-aged remnants. At a distance of 550$`{}_{}{}^{+110}{}_{80}{}^{}`$ pc (Blair et al., 2005), it has a physical size of 27 $`\times `$ 33.5 pc. Located 8.5 below the galactic plane, the remnant lies in a multi-phase medium containing large ISM clouds with a hydrogen density of $`n=510`$ cm<sup>-3</sup>, surrounded by a lower density inter-cloud component of $`n0.10.2`$ cm<sup>-3</sup> (DeNoyer, 1975).
Recently, Patnaude et al. (2002) studied a small, isolated cloud along the southwest limb of the Cygnus Loop which met many of the desired cloud properties for investigating shock-cloud interactions. This cloud is relatively small ($``$ 2$`\mathrm{}`$ in radius; $`0.32`$ pc at $`550`$ pc; Blair et al., 2005), exhibits a fairly uncomplicated, line-of-sight internal structure, and lies isolated from other shocked ISM clouds. Moreover, the shock-cloud interaction is dominated by non-radiative, or ‘Balmer-dominated’ filaments, indicating that the cloud-shock dynamics is not significantly affected by post-shock radiative losses.
Here we present a new analysis of this small shocked cloud. Proper motion measurements and inferred shock velocities of individual filaments in and surrounding the cloud are presented. These results were used to estimate the initial conditions for hydrodynamical model simulations of a shock interaction with an unmagnetized, lumpy cloud. In $`\mathrm{\S }2`$, these new observations are presented as well as the technique used to measure the filament proper motions. Model parameter estimates are then discussed in $`\mathrm{\S }3`$. Our hydrodynamical models are presented in $`\mathrm{\S }4`$, where proper motion and density estimates are implemented in the model initial conditions. Model results are presented in $`\mathrm{\S }5`$, and they are compared to the southwest cloud in $`\mathrm{\S }6`$ with our conclusions in $`\mathrm{\S }7`$.
## 2. Observations
Narrow passband H$`\alpha `$ images of the southwest region of the Cygnus Loop were obtained on 7 July 1992 and 29 August 2003 using the MDM 2.4 m Hiltner telescope. For the July 1992 images, four 600 s H$`\alpha `$ filter (FWHM = 80 Å) exposures were acquired with a Loral 2048 $`\times `$ 2048 front side illuminated CCD yielding a spatial resolution of $`0\stackrel{}{\mathrm{.}}343`$ pixel<sup>-1</sup>. Details of the 1992 observations and subsequent data reduction can be found in Patnaude et al. (2002).
Two 1000 s H$`\alpha `$ filter (FWHM = 15 Å) exposures were taken in August 2003 with a SITe 2048 $`\times `$ 2048 back side illuminated CCD with a resolution of $`0\stackrel{}{\mathrm{.}}275`$ pixel<sup>-1</sup>. Using IRAF<sup>1</sup><sup>1</sup>1IRAF is distributed by the National Optical Astronomy Observatories (NOAO), which is operated by the Association of Universities for Research in Astronomy, Inc. (AURA) under cooperative agreement with the National Science Foundation., the data were bias-subtracted, flat-fielded, and cosmic-ray hits were removed. The resulting 2003 epoch image is shown in Figure 1. Globally, the cloud’s morphology is nearly identical to that seen in the 1992 images (Figs. 2 & 3 Patnaude et al., 2002), but close inspection between the two epoch images showed measurable proper motions for both internal cloud structures and the surrounding thin shock front filaments.
## 3. Analysis and Results
Following the procedure described by Thorstensen, Fesen, & van den Bergh (2001), the coordinate systems of the two H$`\alpha `$ images were aligned using DAOPHOT in conjunction with the USNO-A2.0 catalog. The datasets were then rebinned to an effective image scale of $`0\stackrel{}{\mathrm{.}}1`$ pix<sup>-1</sup>. This rebinning introduced a small global offset of $`0\stackrel{}{\mathrm{.}}07`$ pix<sup>-1</sup> between the two images, uniform across the entire field of view.
### 3.1. Proper Motion Measurements
Individual filament regions for the proper motion analysis were selected based on their projection onto the plane of the sky, the complexity of the filament and surrounding regions, and the brightness of the filament feature. Based on these criteria, 14 filaments within the cloud, including both Balmer dominated and radiative filaments, and 21 regions from the surrounding Balmer dominated shock front were chosen (see Figure 1).
One-dimensional intensity profiles were extracted for each region and the pixel shift in each shock filament was computed using the IRAF task xcsao, which is based on the software of Tonry & Davis (1979). While this task was written to compute relative radial velocities via the cross-correlation function between two spectra, the cross-correlation function yields accurate filament motions in terms pixel shifts between two images. For thin Balmer-dominated filaments and bright and sharp cloud shock features, the cross-correlation analysis was able to match the shock fronts between the two epochs and measure the pixel shift between the two data sets to an accuracy of 10% – 15%. The results from this analysis are listed in terms of proper motion (mas yr<sup>-1</sup>) and transverse velocity (km s<sup>-1</sup>) in Table 1. The quoted velocities assume a distance of 550$`{}_{}{}^{+110}{}_{80}{}^{}`$ pc (Blair et al., 2005).
Example data and cross-correlation functions for two regions are shown in Figure 2. These one-dimensional filaments and cross-correlation functions are representative of the data for the non-radiative filaments (Fig. 2, left) and internal cloud filaments (Fig. 2, right), where there is often a plateau of emission (from shocked cloud material) downstream from the shock front and then a steep rise in emission at the cloud shock front.
Quoted errors in Table 1 include the signal to noise in the filament region, the curvature of the filament region, the profile of the filament when convolved with the image PSF (FWHM<sub>1992</sub> = $`1\stackrel{}{\mathrm{.}}0`$; FWHM<sub>2003</sub> = $`0\stackrel{}{\mathrm{.}}7`$), and the contribution from the background, local nebulosity, and adjacent filaments. For well resolved filaments, the cumulative effect of these errors is $``$ 0.6 pixel, or about 10 km s<sup>-1</sup>.
### 3.2. Estimate of Cloud Parameters
As discussed in Patnaude et al. (2002), this shock-cloud interaction is nearly tangent along the line of sight to the southwestern limb of the southern region of the Cygnus Loop. The cloud is being run over by the remnant’s shock front which is moving roughly east-west. We have divided the cloud into two regions: the ISM shock front and the cloud-shock region. Based on the filament velocities listed in Table 1, we estimated an interstellar shock velocity of 250 km s<sup>-1</sup> associated with the Balmer-dominated filaments. The wide range of Balmer filament velocities observed ($`140260`$ km s<sup>-1</sup>) may be due in part to density fluctuations around the cloud and the fact that only one component of the filament velocity is measured. That is, for filaments which are highly curved, the space velocity of the shock might be 200 km s<sup>-1</sup>, but the local measured velocity might be in a direction other than perpendicular to our line of sight. Furthermore, though filaments were chosen based upon selective criteria, factors such as low signal to noise as well as adjacent, overlapping filaments contributed in some cases to a poorer cross-correlation between the two images. Nonetheless, our estimated shock velocity $``$250 km s<sup>-1</sup> is consistent with the X-ray shock velocity of $``$ 300 km s<sup>-1</sup> inferred from the ROSAT PSPC measurements (Patnaude et al., 2002).
The shocked cloud can be further divided into regions where the cloud-shock is interacting with the cloud, and where it is interacting with cloud clumps or “cloudlets”. In general, the inferred shock velocities vary widely (65–140 km s<sup>-1</sup>). This suggests that the density structure of the cloud is fairly complex, as the cloud-shock appears to have been slowed less in certain areas relative to others.
Based on these measurements, we adopt a shock velocity range inside the cloud of $`60100`$ km s<sup>-1</sup>. These estimated cloud-shock velocities in turn imply a range of density contrasts in the cloud. Assuming ram pressure equilibrium, ($`\rho _av_s^2\rho _cv_{cs}^2`$) the density contrast between the cloud and the ambient medium, $`\chi `$ $``$ $`\rho _c/\rho _a`$, is $``$ 4–17, with higher values representing areas populated with cloudlets, and lower values representing regions of low density within the cloud.
The low density nature of this cloud permits us to view its internal shock dynamics. We have used the structure and spacing of the internal shock fronts to estimate the clumpiness of the cloud. The easiest place to do this is at the western nose of the cloud (Regions C8–C10). Measurements suggest that the post-shock spacing of clumps in the cloud is $``$ 10$`\mathrm{}`$ – 30$`\mathrm{}`$. The upper limits corresponds to shocks which are more highly evolved, while the lower limit corresponds to ‘small’ shocks. Based on the size of the cloud (2$`\mathrm{}`$ – 4$`\mathrm{}`$ diameter), the cloudlets likely account for 30% of the total cloud volume. Furthermore, based on a maximum compression of $``$ 4 for cloud clumps, we estimate the spacing between cloudlets to be $``$ 4–5 cloudlet radii ($`a_{cloudlet}`$).
## 4. Hydrodynamical Models
For modeling this shock-cloud interaction, we have used the numerical hydrodynamics software VH-1 (Blondin & Knerr, 1992; Stevens et al., 1992), which implements the piecewise parabolic method (PPM) to solve the equations of gas dynamics (Colella, & Woodward, 1984). The PPM approach incorporates a fixed computational grid to evolve the standard conservation equations of mass, momentum, and energy. We assume an ideal, inviscid fluid with a ratio of specific heats, $`\gamma `$, equal to $`5/3`$. VH-1 does not explicitly treat the collisionless shock physics associated with Balmer-dominated shocks. However, the goal of this study is to understand how a blast wave interacts with a diffuse ISM cloud. For these purposes, VH-1 serves as an excellent tool for tracing the motion and dynamics of this interaction.
Simulations were performed on a 1440 $`\times `$ 1440 Cartesian grid. The fiducial physical size of the square grid is 1 $`\times `$ 1, with dx = dy = 6.9 $`\times `$ 10<sup>-4</sup>, The size of the cloud sets the scale of the models. The cloud is $``$ 2$`\mathrm{}`$ in radius ($``$ 10<sup>18</sup> cm $`\times `$ d$`_{550\mathrm{p}c}`$). On average, the cloud radius is 30% of the computational grid, or $``$ 450 cells per cloud radius. Therefore, the physical length scale of the grid $`\mathrm{\Delta }`$x $``$ 2 $`\times `$ 10<sup>15</sup> cm cell<sup>-1</sup>.
We estimate the importance of radiative losses by calculating the cooling time scale and comparing it to dynamically relevant time scales (mainly, the cloud crushing time and the pressure variation time scale). In general, radiative losses will be considered important if the cooling time is comparable to or shorter than the cloud crushing time. We estimate the cooling time using the approximation of Kahn (1976), $`t_{cool}`$ = $`Cv^3/\rho `$, where $`v`$ is the cloud shock velocity in units of km s<sup>-1</sup>, $`\rho `$ is the cloud density in units of gm cm<sup>-3</sup>, and $`C`$ is a constant $`=`$ 6.0 $`\times `$ 10<sup>-35</sup>. Assuming a cloud shock velocity of 140 km s<sup>-1</sup> and a cloud density of 10<sup>-24</sup> gm cm<sup>-3</sup>, we estimate a cooling time $`t_{cool}`$ $`>`$ 5000 yr. In contrast, the cloud crushing time, $`t_{cc}`$ $``$ $`\chi ^{1/2}a_0/v_b`$, is $``$ 3800 yr, for a blast wave velocity of 250 km s<sup>-1</sup> and an initial cloud radius of $`a_0`$ = 0.3 pc. Furthermore, the pressure variation time scale is $``$ 0.1$`t_{cc}`$ (KMC94), which is $`<`$ $`t_{cool}`$. Thus, it appears that neglecting the effects of cooling in our models will not have a significant impact on our results. This is supported by Patnaude et al. (2002) who showed that this cloud is only weakly radiative.
Under the assumptions that radiative losses are not dynamically important and that magnetic fields are not present (or that the cloud is only weakly magnetized), the shock-cloud interaction can be wholly defined by two parameters (KMC94), the shock Mach number, and $`\chi `$, the density contrast of the cloud. Assuming an ISM sound speed of $``$ 10–15 km s<sup>-1</sup> and a blast wave velocity of 250 km s<sup>-1</sup>, we estimate a shock Mach number of $`M`$ $``$ 20. Furthermore, as pointed out in $`\mathrm{\S }3`$, we estimate a cloud to ISM density contrast of $`\chi `$ $``$ 6.
To further simplify the problem, we chose a set of non-dimensional variables such that the ISM density $`\rho _a`$ is set to the ratio of specific heats, $`\gamma `$ = $`5/3`$, and the ISM pressure, $`P_a`$, is set to unity. The ISM sound speed, $`c_a`$, is thus set to 1 ($`c_a`$ = $`(\gamma P_a/\rho _a)^{1/2}`$), and the shock velocity $`v_b`$ is just the shock Mach number.
Model results by KMC94 suggested that a cloud radius should be at least of order 120 cells. For our models here, we chose the main cloud to have a radius of 300-500 cells ($``$ 10<sup>17-18</sup> cm). The internal cloudlets have sizes which are 10–20% of the cloud radius. Therefore, the cloudlets are only 33% the suggested size. This small cloudlet size limits our ability to resolve instability growth along their boundaries. However, the goal here is to understand the global, internal morphology of the cloud, rather than small scale mixing within the cloud.
We broke the cloud’s density structure into two parts: A background density profile, and a clumping or density perturbation distribution. The background density distribution is chiefly responsible for the large-scale shock features, such as how the shock drapes over the cloud edges, while also defining the initial internal cloud shock velocity. In contrast, the internal cloud density perturbations lead to the formation of small scale shock structures within the cloud and have little effect on the cloud–ambient medium boundary layer.
Based on the cloud’s emission features (Fig. 1), we assume that the large scale density structure of the cloud is smoothly varying. The interface between the blast wave and the cloud shock, seen along the northern and southern periphery of the cloud, appears smooth. This suggests that the cloud is surrounded by a low density envelope. There is no evidence to suggest that the central density is sharply peaked, so we assume that at some inner radius the density profile turns over and becomes relatively constant throughout. Therefore, we assume that the cloud consists of a smoothly varying core surrounded by a low density envelope. A function which fits this description is a truncated Lorentzian coupled to a power law core:
$$\chi (r)=\{\begin{array}{cc}\frac{\chi _{max}}{1+r_i^2/r_0^2}\left(1+\mathrm{\Delta }\left(\frac{r_ir}{r_i}\right)^\alpha \right)\hfill & 0rr_i\hfill \\ \frac{\chi _{max}}{1+r^2/r_0^2}\hfill & r_i<ra_0\hfill \end{array}$$
(1)
where $`r_0^2=a_0^2/(\chi _{max}1)`$, $`\mathrm{\Delta }`$ $`=`$ 0.11 sets the maximum core density, and $`\alpha `$ is the power law index $`=`$ 0.5 which ensures continuity across the core-envelope interface.
While the fine-scale structure of density perturbations within the cloud is not known, the lack of observed dynamical instabilities (at the resolution of the observations), suggests that such perturbations are probably smoothly varying. Therefore, for models using the above cloud density distribution, we chose to model the cloudlets as Gaussians. The spacing of the Gaussians is such that $`\mathrm{\Delta }r`$ $``$ 4$`\sigma `$ between the cloudlet cores, which is consistent with the optimal spacing of $`\mathrm{\Delta }r`$ $``$ $`4.2a_{cloudlet}`$ suggested by Poludnenko, Frank, & Blackman (2002). Individual model parameters are listed in Table 2. For comparison, we include models of cylindrical clouds with similar density distributions.
## 5. Results
Our basic shock-cloud interaction, Model 1, is shown in the top panel of Figure 3. This model is of a Mach 20 shock interacting with a cloud of radius $`a_0`$ = 0.3 and constant density contrast $`\chi `$ = 6. Figure 3 shows the model at $`t`$ = 1.4, 2.6, 3.6, and 4.7 $`\times `$ 10<sup>-2</sup>. Panel c ($`t=3.6\times 10^2`$), shows the model at approximately one cloud crushing time (the cloud crushing time, given by Equation 2.3 of KMC94, is $`t_{cc}`$ $``$ 3.7 $`\times `$ 10<sup>-2</sup>). According to KMC94, the growth time for Kelvin-Helmholtz (K-H) instabilities is $`t_{KH}`$ $``$ $`\chi ^{1/2}/(kv_{rel})`$, where $`v_{rel}`$ is the relative velocity between the shocked cloud and the shocked ambient medium.
The relative velocity between the post cloud-shock material and the post shock ambient material is given to first order by the relative jump conditions between the cloud and the ambient medium. Since $`M`$ $``$ $`\chi ^{1/2}`$, $`t_{KH}\chi ^{1/2}t_{cc}`$ (for $`ka_0`$). Higher wavenumber perturbations will form on a shorter time-scale. In Panel c of Figure 3, there is clear evidence for K-H growth along the backside of the cloud. In fact, it is evident that K-H growth occurs much earlier (top, Panel b, Fig. 3).
Model 2 is shown in the bottom panel of Figure 3. This model has the density distribution described in Equation 1, with an inner radius $`r_i`$ = $`0.6a_0`$. The evolution of this model is markedly different than that of Model 1 (Fig 3). The $`\chi `$ = 10 listed in Table 2 is what the density would be at the center of the Lorentzian. However, since the Lorentzian is truncated at $`r_i`$, the effective $`\chi `$ is much lower, by an amount $`1/(1+r_i^2/r_0^2)`$. Therefore, the maximum $`\chi `$ in the cloud is $``$ 3, or half the $`\chi `$ of Model 1. More importantly than the lower $`\chi `$ between the two models, Model 2 does not show signs of the instability growth seen in Model 1. This interesting result lends weight to the notion of a smooth boundary between the ISM and an embedded ISM cloud.
While Models 1 and 2 accentuate the differences which arise between a smoothly varying density distribution and that of a sharp edged cylinder, the remaining models (3–6) simulate how the internal density structure affects the cloud shock. Models 3 and 4 (Fig. 4) are of cylindrical clouds of $`\chi `$ = 6 and 3. Both clouds contain perturbations with a $`\chi `$ of 10 above the cloud density (or 30 and 60 times the ambient density). In Model 3 ($`\chi `$ = 6), the flow around the cloud is still strongly influenced by the higher cloud density. The inclusion of cloudlets results in the formation of shock structure within the cloud (not seen in Model 1). However, in the late time (Panel d), the morphology of the cloud is still similar to the late time morphology of Model 1.
The evolution of Model 4, on the other hand, is more strongly influenced by the cloudlets. This is because the density contrast between the cloud and the ambient medium here is only 3, and thus the cloud shock velocity is not significantly different than the blast wave velocity ($`\sqrt{3}`$ lower). Instead, the blast wave is more influenced by the high density cloudlets (relative to the cloud), as seen in Figure 4 (bottom). While this model reproduces the observed shock diffraction (c.f. Fig. 1, the presence of instability growth along the (albeit low density) model cloud boundary is not something observed in the observations.
Models 5 and 6 represent our best approximations to this diffuse ISM cloud. Physically, the model distributions represent cool, low density clouds surrounded by warm, lower density envelopes. Within the clouds, cold dense cloudlets are interspersed here on a regular grid. Cloudlet formation is beyond the scope of this paper but is likely a thermal, rather than gravitational condensation.
Model 5, shown in Figure 5 (top), is of a Mach 10 shock over-running a cloud with a density distribution given by Equation 1. The cloudlets have a $`\chi `$ of 10 and a maximum extent of $`a_{cloudlet}`$ = $`0.05`$. Furthermore, the inner radius of the cloud core is $`0.6a_0`$, which results in some of the cloudlets being outside of the cloud core.
As seen in Figure 5, the blast wave shock is hardly slowed by it’s interaction with the cloud, similar to Model 4. However, the high density cloudlets do significantly alter the cloud shock structure. Compared to Model 2 (Fig. 3), the cloudlets appear to play a significant role in slowing the cloud shock.
Model 6 differs from Model 5 in three ways: First, the $`\chi `$ of the cloud is 8, rather than 10; secondly, the $`\chi `$ of the cloudlets is increased to 15, and thirdly, the radius of the cloudlets is $`0.03`$. Model 6 is shown in Figure 5 (bottom). Here one sees that the $`\chi `$ of the cloud is so low that it barely slows the shock. Moreover, and probably more importantly, the spacing of the cloudlets is such that they do not feel the effects of their neighbors (i.e. $`\mathrm{\Delta }r`$ $`>`$ $`4.2a_{cloudlet}`$ and the diffracted shocks are not significantly curved).
Several features appear in the model simulations which are not observed. Prominent in all the models is the formation of a bow shock behind the cloud. A bow shock is not seen in Figure 1 simply because it is moving back into previously ionized material, so that there is no neutral population to excite. In models containing cloud clumps (Models 3–6), fingers and mushroom heads are prominent in the post shock flow. Three-dimensional models for shock cloud interactions show that many of these features are unstable and will not persist in three dimensions due to turbulent effects in the post shock flow (Stone & Norman, 1992).
## 6. Discussion
As seen in Figure 1, the southwest cloud of the Cygnus Loop represents a fairly uncomplicated case for investigating many of the basic phenomena of a shock-cloud interaction. The low density of the cloud implies that the cloud shocks will be largely non-radiative in nature. Compared to other regions of the Cygnus Loop (Levenson & Graham, 2001, and references therein), the low density nature of this small cloud allows us to view internal cloud shocks. Furthermore, while the east-west extent of the cloud is not known, the location of the forward shock not interacting with the cloud suggests that this shock-cloud interaction is relatively young. Thus, this cloud presents a good test case to model the interaction between a SNR shock and an ISM cloud.
### 6.1. Comparisons to Other Shock Models
There have been several previous studies on shock-cloud interactions. Perhaps the best current model for comparison is that of KMC94. Though there is not a one-to-one comparison between our Model 1 and their models due to the differing initial condition, many of their conclusions are observed in Model 1 such as the formation of K-H instabilities on the order of $`t_{cc}`$ (top, Panel c, Figure 3), like that found in KMC94.
On the other hand, there have been few published studies concerning the interaction between a shock and a cloud with a smoothly varying density. Our models seen in Models 2, 5, and 6 suggest that much of the instability growth observed in previous studies is related to the chosen geometry. ISM clouds are often modeled as cylinders or spheres with sharp, well defined boundaries. Yet, the real boundary between the ISM and embedded diffuse clouds is likely to be less distinct. However, models where the density varies over a certain distance such as that described by a hyperbolic tangent (Poludnenko, Frank, & Blackman, 2002) can sometimes lead to spurious instability growth like that seen in the sharp edged cylindrical case.
In regard to the internal cloud density structure (‘cloudlets’), Poludnenko, Frank, & Blackman (2002) found that the principle parameter is the spacing between the cloudlets. Their models suggested that there exists a critical separation between cloudlets of $``$ 4.2$`a_{cloudlet}`$, and not surprisingly, the cloudlet spacing in Model 5 is about this value. Furthermore, as pointed out by Poludnenko, Frank, & Blackman (2002), a larger $`\chi _{cloudlet}`$ combined with a larger cloudlet spacing does not result in dynamics which are similar to the case of a lower $`\chi _{cloudlet}`$ combined with a smaller cloudlet spacing. Instead, as evidenced by Model 6, the larger separation, regardless of $`\chi _{cloudlet}`$, results in what are essentially multiple, independent interactions between the cloudlets and the cloud-shock.
### 6.2. The Southwest Cloud
While the southwest cloud represents a valuable laboratory for investigating the shock-cloud interaction, as evidenced in Figure 1, it is still highly complex on small scales. Hence, the models presented here only approximate its global properties. Based on Figure 1, the cloud has a radius (N-S direction) of 1–2$`\mathrm{}`$. At the assumed distance of 550 pc, this corresponds to 0.16–0.33 pc, or $``$ 0.5–1 $`\times `$ 10<sup>18</sup> cm. In Model 5, the fiducial radius of the cloud is $`0.35`$. Using Model 5 as our potential model for the southwest cloud, the length scale of the model is thus 1.5 $`\times `$ 10<sup>15</sup> cm cell<sup>-1</sup>.
The density of the ISM in this region is estimated to be $``$ 0.1 – 0.3 cm<sup>-3</sup>. The maximum $`\chi `$ for Model 5 is 10, but in reality the density profile is truncated at an inner radius $`r_i`$ = 0.21$`a_0`$. At $`r_i`$ = 0.21$`a_0`$, $`\chi `$ $``$ 4 for Model 5. This agrees with our lower density estimate of $`\chi `$ $``$ 4.5 from ram pressure arguments. Thus, the cloud particle density is $``$ 0.4 – 1.2 cm<sup>-3</sup> with a $`\chi `$ = 10 for the individual clumps in Model 5 ($`n`$ $``$ 1.0 – 3.0 cm<sup>-3</sup>). The lower shock velocities seen in the cloud suggest cloudlet $`\chi `$’s as high as 17, but the difference between a Gaussian profile with a peak density of 10 and one of 17 is minimal.
Based on the size of the grid cell and the shock velocity, the ambient shock traverses one cell in $``$ 10<sup>8</sup> s $``$ 3 yr. The time difference in the proper motion analysis is about 10 years; that is, the ambient shock has traveled 3–5 cells. In Figure 5, the top panels show the density at $`t`$ = 2.2 – 7.3 $`\times `$ 10<sup>-2</sup>. The simulation begins at $`t`$ = 0. and the shock first hits the cloud at $`t`$ = 0.005. The radius of the cloud is $``$ 500 cells. Therefore, the ambient shock has been traveling for $``$ 2000 yr when it reaches the cloud midpoint.
From the X-ray derived shock velocity of $``$ 300 km s<sup>-1</sup>, Patnaude et al. (2002) estimated the age of the interaction to be $``$ 1200 years, so the modeled cloud size and the shock velocity appear reasonable. At the current epoch, the forward shock is 1$`\mathrm{}`$ – 2$`\mathrm{}`$ ahead of the cloud shock. From Figure 5, the cloud shock lags behind the blast wave by 10% (bottom, Fig. 5, Panel b). This corresponds to a physical distance of 1.9 $`\times `$ 10<sup>17</sup> cm, or $`0\stackrel{}{\mathrm{.}}5`$ at a distance of 550 pc. By Panel c of Model 5, the blast wave is 20% farther along than the cloud shock. Here, the morphology of Model 5 closely matches that of the southwest cloud (Fig. 1).
The observed internal cloud structure in the H$`\alpha `$ image (Fig. 1) is not that unlike the modeled shock cloud internal structure seen in Figure 5. In general, the cloud-shocks seen in the H$`\alpha `$ image are $``$ $`0\stackrel{}{\mathrm{.}}5`$ tip to tip. This scale is consistent with the approximate size of the internal shocks seen in our Model 5 (Fig. 5). The cloud shocks have survived the 10 years between observations. The models, however, show that internal shocks are straightened out over a course of a few hundred time steps ($``$ 200 yr). However, over the short time we are concerned with here, the shock structure of the cloud shock looks remarkably similar to that of the southwest cloud.
## 7. Conclusions
A relatively isolated, low-density ISM cloud situated along the southwest limb of the Cygnus Loop provides a particularly clear view of the early stages of a SNR shock – ISM cloud interaction. The combination of multi-epoch observations and high resolution numerical modeling of this cloud has provided some new insights regarding how shocks overrun ISM clouds. The southwest cloud’s isolation and low-density has also allowed us to view its internal density structure and make inferences concerning the cloud’s initial density distribution.
Our specific findings are:
1) Using multi-epoch H$`\alpha `$ observations of a small, isolated ISM cloud in the southwest portion of the Cygnus Loop, we measured proper motions of Balmer-dominated shock filaments which wrap around the cloud, as well as the proper motion of several internal cloud shocks. The Balmer-dominated filaments have transverse velocities of $``$ 200–250 km s<sup>-1</sup>, while the shock filaments internal to the cloud have transverse velocities of 65 – 140 km s<sup>-1</sup>.
2) The shocked cloud’s morphology does not show many of the dynamical instabilities predicted by previous shock-cloud models. This suggests that there is no abrupt boundary or edge for diffuse ISM clouds. A sharp density rise between the cloud and the ISM would lead to steep velocity gradients at the shocked cloud – shocked ISM interface. These steep gradients would in turn lead to the onset of Kelvin-Helmholtz instabilities, which are not observed. This conclusion contrasts with the shock-cloud interaction seen in the southeast of the Cygnus Loop, where the blast wave is thought to be interacting with a large, dense cloud, and instability growth is clearly seen along the cloud-shock boundary.
3) Our model hydrodynamic simulations suggest that ISM clouds are best modeled as a constant or smoothly varying core density embedded in lower density envelope which tapers to the surrounding ISM. Ram pressure equilibrium arguments suggest a cloud–ISM density contrast for this cloud of $`\chi `$ = 5 – 17, with the lower $`\chi `$’s corresponding to the diffuse regions of the cloud and the upper limit of 17 corresponding to the dense cloud clumps.
4) A definite spacing of dense, small “cloudlets” inside the cloud is needed to generate the cloud’s internal morphology as seen in the H$`\alpha `$ image. As pointed out by Poludnenko, Frank, & Blackman (2002), clumps spaced too closely together interact with the shock as if they were one large clump, while those spaced too far apart behave as a set of individual clouds. Our models are consistent with the optimal spacing of $`d_{crit}`$ $``$ 4$`a_{cloudlet}`$ (Poludnenko, Frank, & Blackman, 2002). The observed internal cloud shock diffraction caused by these cloudlets is a short lived phenomena. As the cloud shock interacts with the cloudlets, the diffracted shocks re-order themselves on a time scale of order a few cloudlet crossing times.
We wish to thank John Blondin for both making the VH-1 code available, and answering several questions regarding its use, and John Raymond for useful suggestions regarding our results. We also thank the anonymous referee for several helpful comments during the preparation of this manuscript. |
warning/0507/nucl-ex0507010.html | ar5iv | text | # Accuracy of traditional Legendre estimators of quadrupole ratios for the 𝑁→Δ transition
(July 6, 2005)
## Abstract
We evaluate the accuracy of traditional estimators often used to extract $`N\mathrm{\Delta }`$ quadrupole ratios from cross section angular distributions for pion electroproduction. We find that neither $`M_{1+}`$ dominance nor $`\mathrm{}1`$ truncation is sufficiently accurate for this purpose. Truncation errors are especially large for $`R_{EM}`$, for which it is also essential to perform Rosenbluth separation. The accuracy of similar truncated Legendre analyses for $`E_{0+}`$, $`S_{0+}`$, and especially $`M_1`$ is even worse.
preprint: Internal report to e91011 collaboration
Historically the most important indications of deformation of low-lying baryons have been the quadrupole ratios for electromagnetic excitation of the $`N\mathrm{\Delta }(1232)`$ transition. Magnetic dipole excitation dominates and is represented by the $`M_{1+}`$ multipole amplitude while nonzero values for the electric and scalar (longitudinal) multipoles, $`E_{1+}`$ and $`S_{1+}`$, arise either from nonspherical contributions to the wave functions or from higher-order dynamical contributions to the electromagnetic transition. The quadrupole ratios are defined as
$`R_{EM}=\mathrm{Re}{\displaystyle \frac{E_{1+}}{M_{1+}}}`$ (1a)
$`R_{SM}=\mathrm{Re}{\displaystyle \frac{S_{1+}}{M_{1+}}}`$ (1b)
evaluated at the physical mass of the resonance, $`W=M_\mathrm{\Delta }1.232`$ GeV. Determination of quadrupole ratios for isospin-3/2 amplitudes requires measurements of two charge states, such as $`p\pi ^0`$ and $`n\pi ^+`$. Complex multipole amplitudes have been deduced for $`Q^2=0`$ using polarization data for pion photoproduction Blanpied01 , but few experiments for $`Q^2>0`$ have provided sufficient information to perform an actual multipole analysis. Instead, most experimental determinations of $`N\mathrm{\Delta }`$ transition form factors Kalleicher97 ; Frolov99 ; Joo02 ; Sparveris05 rely upon estimators derived from multipole expansions for the angular dependence of unpolarized cross sections using two simplifying assumptions: 1) only multipoles with $`\mathrm{}1`$ contribute, which is described as $`sp`$ truncation; and 2) only terms involving $`M_{1+}`$ are retained, which is described as $`M_{1+}`$ dominance. Although the reliability of these assumptions has been questioned before, the improved kinematic completeness and statistical precision of modern experiments warrants re-examination of their accuracy. In this Brief Report, we consider the accuracy of traditional quadrupole estimators for $`Q^21`$ (GeV/$`c`$)<sup>2</sup> where a nearly model-independent multipole analysis of recoil-polarization response functions for the $`p(\stackrel{}{e},e^{}\stackrel{}{p})\pi ^0`$ reaction disagrees appreciably with the traditional Legendre analysis of the cross section data Kelly05c ; Kelly05d .
The unpolarized differential cross section for $`\gamma _vNN\pi `$ in the $`\pi N`$ center of momentum frame takes the form
$$\frac{d\sigma }{d\mathrm{\Omega }_\pi }=\nu _0\left(ϵ_SR_L+R_T+\sqrt{2ϵ_S(1+ϵ)}R_{LT}\mathrm{sin}\theta \mathrm{cos}\varphi +ϵR_{TT}\mathrm{sin}^2\theta \mathrm{cos}2\varphi \right)$$
(2)
where $`\nu _0`$ is a phase-space factor, $`ϵ`$ is the transverse polarization of the virtual photon, $`ϵ_S=ϵQ^2/q^2`$, and $`(\theta ,\varphi )`$ are polar and azimuthal pion angles relative to the $`\stackrel{}{q}`$ vector and the electron scattering plane. The response functions can be expanded in Legendre series
$$R_\lambda =\underset{n=0}{\overset{\mathrm{}}{}}A_n^\lambda P_n(\mathrm{cos}\theta )$$
(3)
where $`\lambda \{L,T,LT,TT\}`$. The expansion coefficients, $`A_n^\lambda `$, are functions of $`(W,Q^2)`$ that can be fit to the angular distribution of the differential cross section. Each of those coefficients can in turn be expressed as a multipole expansion containing terms of the form $`\mathrm{Re}B_\mathrm{}\pm C_\mathrm{}\pm ^{}`$ where $`B,C\{M,E,S\}`$ are magnetic, electric, or scalar multipole amplitudes for specified $`\mathrm{}`$ and $`j=\mathrm{}\pm 1/2`$. In principle, experimentally determined Legendre coefficients include contributions from arbitrarily large $`\mathrm{}`$ and are not limited either by $`sp`$ truncation or by $`M_{1+}`$ dominance.
Truncated multipole expansions of the Legendre coefficients used in quadrupole estimators are given in Eq. (4) where the contributions that satisfy $`M_{1+}`$ dominance are listed first and where the remaining terms include some of the lowest multipolarity contributions of other types but are not necessarily arranged in order of numerical importance.
$`A_0^L`$ $`=`$ $`|S_{0+}|^2+8|S_{1+}|^2+|S_1|^2+8|S_2|^2+27|S_{2+}|^2+\mathrm{}`$ (4a)
$`A_0^T`$ $`=`$ $`2|M_{1+}|^2+|E_{0+}|^2+|M_1|^2+6|E_{1+}|^2]+6|M_2|^2+2|E_2|^2+9|M_{2+}|^2+18|E_{2+}|^2+\mathrm{}`$ (4b)
$`A_0^{TT}`$ $`=`$ $`{\displaystyle \frac{3}{2}}|M_{1+}|^2\mathrm{Re}[M_{1+}^{}(3E_{1+}+3M_1+12M_3+3E_3+2M_{3+}+10E_{3+})]`$
$`+`$ $`{\displaystyle \frac{9}{2}}|E_{1+}|^2+{\displaystyle \frac{3}{2}}|E_2|^2+24|E_{2+}|^2{\displaystyle \frac{9}{2}}|M_2|^212|M_{2+}|^2`$
$`+`$ $`\mathrm{Re}[3E_{0+}^{}(E_2+M_2M_{2+}+E_{2+})+E_{1+}^{}(3M_121E_312M_3+12M_{3+})`$
$`+`$ $`M_1^{}(3E_3+3M_3+10E_{3+}10M_{3+})+\mathrm{}]`$
$`A_1^{LT}`$ $`=`$ $`3\mathrm{R}\mathrm{e}[M_{1+}^{}(2S_{1+}3S_3+4S_{3+})+S_{0+}^{}(E_2M_2+M_{2+}4E_{2+})`$
$``$ $`E_{0+}^{}(2S_23S_{2+})2E_{1+}^{}(S_{1+}+S_1)2M_1S_{1+}^{}+\mathrm{}]`$
$`A_2^L`$ $`=`$ $`8|S_{1+}|^2+8|S_2|^2+{\displaystyle \frac{216}{7}}|S_{2+}|^2+\mathrm{Re}[S_{0+}^{}(8S_2+18S_{2+})+8S_{1+}S_1^{}+\mathrm{}]`$ (4e)
$`A_2^T`$ $`=`$ $`|M_{1+}|^2+\mathrm{Re}[M_{1+}^{}(6E_{1+}2M_1+{\displaystyle \frac{24}{7}}M_3+6E_3+{\displaystyle \frac{144}{7}}M_{3+})]`$
$`+`$ $`3|E_{1+}|^2|E_2|^2+{\displaystyle \frac{108}{7}}|E_{2+}|^2+3|M_2|^2+{\displaystyle \frac{36}{7}}|M_{2+}|^2`$
$`+`$ $`\mathrm{Re}[E_{0+}^{}(2E_26M_2+6M_{2+}+12E_{2+})6M_1E_{1+}^{}+\mathrm{}]`$
Table 1 shows that the number of independent terms in the multipole expansions of these Legendre coefficients increases very rapidly with the maximum $`\mathrm{}`$ permitted. Complete expansions for $`\mathrm{}_{\text{max}}6`$ can be found in Ref. e91011\_mpquad but as $`\mathrm{}_{\text{max}}`$ increases they quickly become too unwieldy to display here or to use in practical applications. The Legendre coefficients are usually obtained by numerical integration of response functions against Legendre functions instead of by these algebraic formulas, but both methods do agree.
The assumption of $`M_{1+}`$ dominance omits any terms that do not involve $`M_{1+}`$, which strongly inhibits the proliferation of terms but is not sufficient in itself to extract quadrupole ratios from cross section data. Combined with $`sp`$ truncation, these expansions reduce to
$`A_0^L`$ $``$ $`0`$ (5a)
$`A_0^T`$ $``$ $`2|M_{1+}|^2`$ (5b)
$`A_0^{TT}`$ $``$ $`\mathrm{Re}[({\displaystyle \frac{3}{2}}M_{1+}+3E_{1+}+3M_1)M_{1+}^{}]`$ (5c)
$`A_1^{LT}`$ $``$ $`6\mathrm{R}\mathrm{e}[S_{1+}M_{1+}^{}]`$ (5d)
$`A_2^L`$ $``$ $`0`$ (5e)
$`A_2^T`$ $``$ $`\mathrm{Re}[(M_{1+}+6E_{1+}2M_1)M_{1+}^{}]`$ (5f)
Thus, one obtains the traditional quadrupole estimators
$`\stackrel{~}{R}_{EM}`$ $`=`$ $`{\displaystyle \frac{3(A_2^T+ϵA_2^L)2A_0^{TT}}{12(A_0^T+ϵA_0^L)}}`$ (6a)
$`\stackrel{~}{R}_{SM}`$ $`=`$ $`{\displaystyle \frac{A_1^{LT}}{3(A_0^T+ϵA_0^L)}}`$ (6b)
where $`A_n^L`$ is included because most experiments have not used Rosenbluth separation to isolate $`A_n^T`$. When Rosenbluth separation is available, one can simply use $`ϵ0`$ in Eq. (6). Therefore, it is useful to define the accuracy parameters
$`f_{EM}(\mathrm{}_{\text{max}},ϵ)`$ $`=`$ $`{\displaystyle \frac{1}{R_{EM}}}\left({\displaystyle \frac{3(A_2^T+ϵA_2^L)2A_0^{TT}}{12(A_0^T+ϵA_0^L)}}\right)_{\mathrm{}\mathrm{}_{\text{max}}}{\displaystyle \frac{\stackrel{~}{R}_{EM}}{R_{EM}}}`$ (7a)
$`f_{SM}(\mathrm{}_{\text{max}},ϵ)`$ $`=`$ $`{\displaystyle \frac{1}{R_{SM}}}\left({\displaystyle \frac{A_1^{LT}}{3(A_0^T+ϵA_0^L)}}\right)_{\mathrm{}\mathrm{}_{\text{max}}}{\displaystyle \frac{\stackrel{~}{R}_{SM}}{R_{SM}}}`$ (7b)
where asymptotic equality refers to the limit $`\mathrm{}_{\text{max}}\mathrm{}`$. Despite their appealing simplicity, it is clear that many contributions are omitted and the accuracy of the traditional estimators is a numerical issue that can be addressed either theoretically using model calculations or experimentally using additional polarization measurements to extract complex multipole amplitudes directly.
The convergence of these expansions is evaluated in Table 2 using $`p\pi ^0`$ multipole amplitudes for $`W=1.232`$ GeV and $`Q^2=1.0`$ (GeV/$`c`$)<sup>2</sup> from MAID2003 Drechsel99 ; MAID2003 . First, we observe that $`A_n^L`$ contributions are not negligible: the contribution of $`A_0^L`$ to the denominators of Eq. (6) reduces the estimated quadrupole ratios by about $`4\%`$ without Rosenbluth separation when $`ϵ1`$. (Note that $`ϵ=0.949`$ at $`W=1.232`$ GeV in Ref. Kelly05c .) Even though $`A_2^L`$ is much smaller, its effect upon $`f_{EM}`$ is even larger because the strong cancellation between $`A_0^{TT}`$ and $`A_2^T+ϵA_2^L`$ amplifies the dependence on $`ϵ`$. Therefore, the assumption of $`M_{1+}`$ dominance is not sufficiently accurate to measure $`R_{EM}`$ without Rosenbluth separation. Even with Rosenbluth separation, one should not expect better than about $`15\%`$ accuracy for either quadrupole ratio using the traditional Legendre analysis (see the bottom of last two columns of Table 2). Second, it is clear that $`sp`$ truncation is not valid either because contributions with $`\mathrm{}>1`$ are not negligible. Cancellation between contributions to the numerator of $`\stackrel{~}{R}_{EM}`$ also amplifies truncation errors and convergence is not necessarily monotonic as $`\mathrm{}_{\text{max}}`$ increases. Contributions to $`A_0^L`$ and $`A_0^T`$ are nonnegative, but the signs for other Legendre coefficients are mixed. While the magnitudes of multipole amplitudes for $`\mathrm{}>1`$ do tend to decrease, their coefficients in Eq. (4) tend to increase with $`\mathrm{}`$. Thus, convergence becomes a delicate numerical issue.
Under the present conditions, we find that $`\mathrm{Re}(M_1E_{1+}^{})`$ is the most important contribution to $`\stackrel{~}{R}_{EM}`$ neglected by $`M_{1+}`$ dominance and is approximately $`40\%`$ of the leading term. Thus, $`M_{1+}`$ dominance is not very accurate. The fact that $`f_{EM}`$ approaches $`0.88`$ for $`ϵ=0`$ is actually nothing more than a lucky conspiracy among the magnitudes and signs for a very large number of smaller terms, many of which are not especially small individually. However, most experiments omit Rosenbluth separation. Similarly, the second most important contribution to $`\stackrel{~}{R}_{SM}`$ is $`\mathrm{Re}(S_{0+}E_2^{})`$ but is only about $`6\%`$ of the leading term; hence, $`f_{SM}`$ converges more rapidly. The details of this analysis are obviously model dependent, but qualitatively similar results are obtained for other models as well. Although $`f_{EM}`$ for $`ϵ=0`$ is slightly closer to unity than $`f_{SM}`$ for the present analysis, the greater susceptibility of $`\stackrel{~}{R}_{EM}`$ to truncation errors through its reliance upon delicate cancellations suggests that the traditional Legendre analysis is intrinsically less reliable for $`R_{EM}`$ than for $`R_{SM}`$, with or without Rosenbluth separation.
It is often argued that the traditional Legendre analysis should be more accurate for the isospin-3/2 channel than for the $`p\pi ^0`$ reaction because the resonant multipoles should share a common phase and become pure imaginary at the physical mass, thereby suppressing background contributions. Leaving aside the propagation of errors involved in extracting isospin-3/2 amplitudes by combining two independent experiments, we can address the intrinsic accuracy of this analysis method using model calculations also. The convergence of the accuracy parameters for isospin-3/2 quadrupole ratios is examined in Table 3. Again we find that Rosenbluth separation is required for $`\stackrel{~}{R}_{EM}`$. Interestingly, $`f_{EM}`$ deteriorates as $`\mathrm{}_{\text{max}}`$ increases and the final accuracy of $`\stackrel{~}{R}_{EM}`$ is worse for isospin-3/2 than for $`p\pi ^0`$ even with $`ϵ=0`$. The cancellations are severe, the method is unstable, and calculations for $`\stackrel{~}{R}_{EM}`$ are highly model-dependent.
Figure 1 compares traditional quadrupole estimators with $`R_{EM}^{(p\pi ^0)}`$ and $`R_{SM}^{(p\pi ^0)}`$ for MAID2003 at $`Q^2=1.0`$ (GeV/$`c`$)<sup>2</sup>. Ideally the estimators would be most accurate in the immediate vicinity of the physical mass, $`W=M_\mathrm{\Delta }1.232`$ GeV, but neither actually has that property. Rosenbluth separation is most important for $`R_{EM}`$, but even with separation the residual error at $`M_\mathrm{\Delta }`$ is significant at the level of experimental precision that is now possible.
Similarly, Fig. 2 shows the $`Q^2`$ dependence for the accuracies of the traditional quadrupole estimators using MAID2003 $`p\pi ^0`$ multipole amplitudes for $`W=1.232`$ GeV. Solid curves use $`ϵ=0`$, corresponding to Rosenbluth separation, while dashed curves use $`ϵ=0.9`$, typical of many experiments. Even with Rosenbluth separation, neither quadrupole estimator can be trusted to better than about $`20\%`$ and their accuracy deteriorates at larger $`Q^2`$ as $`M_{1+}`$ dominance breaks down. Therefore, truncation errors can seriously affect the $`Q^2`$ dependence of quadrupole amplitudes deduced from Legendre coefficients. The most extensive recent study of $`p\pi ^0`$ quadrupole ratios for $`0.4Q^21.8`$ (GeV/$`c`$)<sup>2</sup> used Eq. (6) without Rosenbluth separation Joo02 . We do not advocate adjustment of such results using Fig. 2, at least at this time, because the shapes of $`f_{EM}`$ and $`f_{SM}`$ are model dependent and MAID2003 does not describe all of the low-lying multipole amplitudes at $`Q^2=1`$ (GeV/$`c`$)<sup>2</sup> from Ref. Kelly05d sufficiently well to be confident of its predictions for the $`Q^2`$ dependencies of these ratios. Instead, we claim that accurate measurements of the quadrupole ratios require multipole analysis of both polarization and cross section data.
Finally, other simple estimators
$`\mathrm{Re}E_{0+}M_{1+}^{}`$ $``$ $`A_1^T/2`$ (8a)
$`\mathrm{Re}S_{0+}M_{1+}^{}`$ $``$ $`A_0^{LT}`$ (8b)
$`\mathrm{Re}M_1M_{1+}^{}`$ $``$ $`(2A_0^T+2A_0^{TT}+A_2^T)/8`$ (8c)
based upon $`M_{1+}`$ dominance and $`sp`$ truncation are sometimes quoted Joo02 ; Burkert04 . Note that Rosenbluth separation is required. However, using MAID2003 $`p\pi ^0`$ multipoles at $`(W,Q^2)=(1.232,1.0)`$ with $`\mathrm{}5`$, the ratios between the right- and left-hand sides of Eq. (8) are $`1.74`$, $`0.77`$, and $`9.75`$. Most notably, the numerical contribution of $`\mathrm{Re}M_1M_{1+}^{}`$ is only the fifth largest term in the multipole expansion of the specified combination of Legendre coefficients. Therefore, these estimators are worthless under these conditions.
In summary, we have performed a detailed numerical analysis of truncation errors in quadrupole ratios deduced from Legendre coefficients fit to cross section angular distributions. We find that neither $`M_{1+}`$ dominance nor $`sp`$ truncation is reliable and that one cannot expect better than $`20\%`$ accuracy from this method. Truncation errors are especially important for $`R_{EM}`$. Furthermore, accurate results for $`R_{EM}`$ also require Rosenbluth separation, which was not performed in recent studies of the $`Q^2`$ dependence of the quadrupole ratios. The accuracy of truncated Legendre analyses of $`E_{0+}`$, $`S_{0+}`$, and especially $`M_1`$ is even worse. Polarization data for pion electroproduction are needed to perform nearly model-independent multipole analyses that provide complex amplitudes and do not depend upon unjustifiable truncation schemes.
###### Acknowledgements.
The support of the U.S. National Science Foundation under grant PHY-0140010 is gratefully acknowledged. |
warning/0507/hep-ph0507229.html | ar5iv | text | # Gamma Rays from Heavy Neutralino Dark Matter
## I Introduction
One of the most favoured dark matter candidates since 20 years or so, is the lightest supersymmetric particle goldberg ; ellisetal . Most of the expected phenomenology has been worked out, and is implemented in freely available, extensive computer packages like DarkSusy darksusy and micrOMEGAs micromegas . The vast majority of the analyses has been performed in various constrained versions of the minimal supersymmetric extension of the Standard Model (MSSM) (for reviews, see e.g. neut ), where either radiative breaking of supersymmetry or the GUT condition on gaugino masses (or both) have been imposed. It is clear from all these studies that a suitable neutralino candidate for dark matter is indeed available, with relic density as measured by WMAP ($`\mathrm{\Omega }_{\text{CDM}}h^2=0.113\pm 0.009`$, where $`\mathrm{\Omega }_{\text{CDM}}`$ denotes the ratio of cold dark matter to critical density and $`h`$ is the Hubble constant in units of $`100\text{ km}\text{ s}^1\text{ Mpc}^1`$) wmap , and with great prospects of being detected either at the CERN LHC, or at the various direct and indirect detection experiments of halo dark matter that are currently operational or being constructed.
However, one may also ask the unpleasant but not unrealistic question what happens if supersymmetry is not found at the LHC. If supersymmetry is still present, that would probably mean that the mass scale of the lightest supersymmetric particles is beyond the kinematical reach of the accelerator. For neutralino masses at the TeV scale, also the scattering cross section for direct detection would necessarily be small, as would many of the rates (antiprotons, positrons) for indirect detection. An exception seems to be gamma-ray detection, firstly because rates do not fall off as rapidly BUB and secondly due to eminent new gamma-ray telescopes, most clearly demonstrated by the recent spectacular performance of HESS, in particular as regards the multi-TeV signal that has been observed towards the galactic center hess . With such new astrophysical gamma-ray instruments of unprecedented size and energy resolution either in operation hess ; magic ; cangaroo or under construction veritas ; glast , it is appropriate to investigate possible levels of signals and spectral signatures of heavy dark matter particle annihilation.
Most of the previous calculations have been carried out at tree-level, with radiative corrections typically (and correctly) believed to be at the few percent level. In some cases, radiative corrections may on the contrary relieve the annihilation rate from inhibiting factors having to do with the Majorana fermion property of neutralinos and the fact that annihilation in the dark matter halo effectively takes place at rest goldberg . One example of this is the radiative “correction” of $`\chi \chi f\overline{f}`$ through the emission of a photon in the final state. Here, the first-order corrected cross section can be many orders of magnitude larger than the tree-level result lbe89 .
A second example of an unexpectedly large cross section is that of the second order, loop-induced $`\gamma \gamma `$ and $`Z\gamma `$ annihilation 2gamma ; zgamma ; BUB , where in the high mass, pure higgsino (or wino ullio ) limit the branching ratio normalized to the lowest order rate can reach percent level, despite the naive expectation of being 2 to 3 orders of magnitude smaller. The origin of this enhancement has only recently been fully understood, and is explained by nonperturbative, binding energy effects in the special situation of having very small (i.e. galactic) velocities and very large dark matter masses as well as small mass differences between the neutralino and the lightest chargino hisano .
In this Letter, we focus on gamma rays from neutralino annihilation into charged gauge boson pairs and show that there is yet another, previously neglected enhancement mechanism, appearing already at first order in $`\alpha _{\mathrm{em}}`$: radiative processes containing one photon in addition to the weak bosons in the final state, give a new source of photons which peaks near the highest possible energy (the mass of the neutralino). This turns out to be a very beneficial effect for the potential detection of neutralinos with masses $`m_\chi 1`$ TeV, as the normally soft spectrum of continuous photons coming from the fragmentation of $`W`$ or $`Z`$ bosons (see, e.g., BUB ; hisano ; fornengo ) gets a high-energy supplement with a clearly distinguishable signature. This is reminiscent of the case of Kaluza-Klein dark matter, where internal bremsstrahlung in annihilation processes with charged lepton final states dominates the gamma-ray spectrum at the highest energies kk (see also birkedal ).
## II Gamma rays from neutralino annihilations
In most models, the lightest stable supersymmetric particle is the lightest neutralino, henceforth just “the neutralino”, which is a linear combination of the superpartners of the gauge and Higgs fields,
$$\chi \stackrel{~}{\chi }_1^0=N_{11}\stackrel{~}{B}+N_{12}\stackrel{~}{W}^3+N_{13}\stackrel{~}{H}_1^0+N_{14}\stackrel{~}{H}_2^0.$$
(1)
In order not to overclose the universe, a TeV-scale neutralino must generally have a very large higgsino fraction, $`Z_h\left|N_{13}\right|^2+\left|N_{14}\right|^2`$, if the usual GUT condition $`M_1M_2/2`$ is imposed; otherwise a heavy wino would also be acceptable. In the following, we therefore focus on higgsino-like neutralinos, with $`Z_h1`$ and $`N_{13}\pm N_{14}`$ <sup>1</sup><sup>1</sup>1 If the neutralino is a pure wino, one gets identical results as for the (anti-) symmetric combination of higgsinos that we consider here, apart from an overall factor of 16 that multiplies all quoted cross sections. A pure bino state, on the other hand, does not couple to $`W`$ at lowest order at all.. For the high masses we are interested in, the annihilation rate into charged gauge bosons often dominates. Internal bremsstrahlung in these final states are therefore of great interest to investigate. Moreover, we note that all our results are almost independent of the relative velocity $`v`$ of the annihilating neutralino pair. Analytical expressions are therefore presented in the limit of vanishing velocity, but should be applicable both at the time of freeze-out ($`v/c1/6`$) and to annihilating neutralinos in the galactic halo today ($`v/c10^3`$).
For a pure higgsino, the only contribution to the lowest order annihilation cross section into charged gauge bosons comes from a $`t`$-channel exchange of a chargino; it is given by
$$(\sigma v)_{WW}=\frac{g^4}{32\pi }\frac{\left(m_\chi ^2m_W^2\right)\sqrt{1m_W^2/m_\chi ^2}}{\left(m_\chi ^2+m_{\chi _1^\pm }^{}{}_{}{}^{2}m_W^2\right)^2},$$
(2)
where $`m_{\chi _1^\pm }`$ and $`m_W`$ are the lightest chargino and $`W`$ masses, respectively.
Let us now consider radiative corrections with a photon in the final state in addition to the $`W`$ pair. Just as at lowest order, the potential $`s`$-channel exchanges of $`Z`$ and Higgs bosons vanish, and the only Feynman diagrams that contribute are shown in Fig. 1. To zeroth order in $`ϵm_W/m_\chi `$, and retaining a leading logarithmic term, the resulting photon multiplicity is given by
$`{\displaystyle \frac{\mathrm{d}N_\gamma ^W}{\mathrm{d}x}}{\displaystyle \frac{\text{d}(\sigma v)_{WW\gamma }/\text{d}x}{(\sigma v)_{WW}}}`$
$`{\displaystyle \frac{\alpha _{\mathrm{em}}}{\pi }}[{\displaystyle \frac{4(1x+x^2)^2\mathrm{ln}(2/ϵ)}{(1x)x}}`$
$`{\displaystyle \frac{2(412x+19x^222x^3+20x^410x^5+2x^6)}{(2x)^2(1x)x}}`$
$`+{\displaystyle \frac{2(824x+42x^237x^3+16x^43x^5)\mathrm{ln}(1x)}{(2x)^3(1x)x}}`$
$`+\delta ^2\left({\displaystyle \frac{2x(2(2x)x)}{(2x)^2(1x)}}+{\displaystyle \frac{8(1x)\mathrm{ln}(1x)}{(2x)^3}}\right)`$
$`+\delta ^4({\displaystyle \frac{x(x1)}{(2x)^2}}+{\displaystyle \frac{(x1)(22x+x^2)\mathrm{ln}(1x)}{(2x)^3}})],`$ (3)
where $`xE_\gamma /m_\chi `$ and $`\delta (m_{\chi _1^\pm }m_\chi )/m_W`$.
Several interesting features can be identified in this expression. For large mass shifts $`\delta `$ the last two terms dominate. They originate from longitudinally polarized charged gauge bosons in the final state, which are forbidden in the lowest order process because of the different CP properties of the initial and final state dreesnojiri . Remember that in the limit of vanishing relative velocity, the initial state must be an $`S`$-wave with pseudoscalar quantum numbers due to the Majorana nature of the neutralino. The emission of a photon, on the other hand, will open up this channel in the $`{}_{}{}^{1}S_{0}^{}`$ partial wave, potentially leading to very large cross sections <sup>2</sup><sup>2</sup>2 Unitarity is in general restored and therefore forbids too large cross sections, which can be understood from the equivalence theorem between longitudinal gauge bosons and would-be Goldstone modes of the Higgs sector dreesnojiri .. However, in supersymmetric scenarios with a heavy higgsino-like neutralino one usually expects a mass shift $`\delta <ϵ`$ dreesnojiri ; hisano , in which case the longitudinal part (last two terms) can be neglected.
For small mass shifts, the cross section is instead dominated by the production of transversely polarized gauge bosons. This results in a peak in the spectrum at high energies that becomes more and more pronounced for higher neutralino masses. The appearance of this peak can be understood by observing that for very heavy neutralino masses the transversely polarized $`W`$ bosons can be treated as light and thus behave in the same way as infrared photons radiated from the neutralino/chargino line in Fig. 1. The mechanism that takes place is, in other words, an amusing reflection of QED infrared behaviour also for $`W`$ bosons: The kinematical situation when the photon and one of the $`W`$s leave the annihilation point each with maximal energy, gets an enhancement, since it is automatically accompanied by a very soft $`W`$. This is also reflected in the symmetric appearance of the $`x0`$ and $`x1`$ poles in the first terms of Eq. (II).
As an illustrative example, we have chosen a typical higgsino-like MSSM model, fulfilling all experimental constraints, as specified in Table 1 (similar models are found in, e.g., the focus point region of mSUGRA). The resulting photon spectrum from internal bremsstrahlung of W pair final states is shown in Fig. 2. The symmetry around $`x0.5`$ in the spectrum indicates the related nature of the peak and the infrared divergence. For completeness, we have also included a very high mass (10 TeV) higgsino model which has received some attention recently hisano ; Boudjema:2005hb (even though thermal production of such a neutralino in general gives a too large $`\mathrm{\Omega }_{\text{CDM}}`$, unless one allows for finetuning of parameters like the psedudoscalar Higgs mass Profumo:2005xd ). In addition, the case of a hypothetical model with a very large mass shift is shown (where the contributions from longitudinal $`W`$ bosons dominate at high energies).
Let us now consider those contributions to the gamma-ray spectrum from the decay of heavy neutralinos that have been studied earlier. Secondary gamma rays are produced in the fragmentation of the $`W`$ pairs, mainly through the decay of neutral pions. In addition to the secondary spectrum, there are line signals from the direct annihilation of a neutralino pair into $`\gamma \gamma `$ 2gamma and $`Z\gamma `$ zgamma . Due to the high mass of the neutralino, these lines cannot be resolved but effectively add to each other at an energy equal to the neutralino mass.
For comparison, again using the model of Table 1, Fig. 3 shows the contributions from secondary photons hisano and the line signals, as well as the new source of photons from the internal bremsstrahlung diagrams of Fig. 1.
The practical importance of the latter contribution can be appreciated even more, when considering a finite detector resolution of 15 %, which is typical for atmospheric Cherenkov telescopes in that energy range; the result is a smeared spectrum as shown in Fig. 4. One can see that, although the strength of the $`\gamma \gamma `$ and $`Z\gamma `$ lines already are surprisingly large BUB , the contribution from the internal bremsstrahlung further enhances this peak by a factor of 2. The signal is also dramatically increased at lower energies, thereby filling out the “dip” just below the peak; this latter effect will of course become even more pronounced for better detector resolutions.
## III Conclusions and Discussion
In this Letter, we have presented important radiative corrections to the gamma-ray spectrum from heavy neutralino annihilations. They contribute a characteristic peak shape at the highest energies, competing with the $`\gamma \gamma `$ and $`Z\gamma `$ line signals in today’s detectors.
When it comes to predicting absolute gamma-ray fluxes, there is, unfortunately, a great uncertainty which is primarily due to the unknown dark matter clustering properties of the Milky Way halo. Particle physics alone predicts, in lowest order perturbation theory, a flux which falls roughly as $`m_\chi ^2`$. On the other hand, the possible nonperturbative enhancement discussed in hisano is important for heavy neutralinos, and may give a substantial boost to the gamma-ray signal. The process discussed here should benefit from a similar boost as the line signal treated in hisano , since the enhancement is due to binding effects in the initial state. Thus, even if the radiated $`W`$ boson is soft compared to the initial state total mass, it still causes a virtuality of one of the neutralinos which is much greater than the binding energy. Therefore, the factorization of the process into long distance and short distance kernels performed in hisano should be valid also here, and our curves for the ratios of cross sections ($`\mathrm{d}N_\gamma ^W/\mathrm{d}x`$) should not be affected much.
Since the absolute gamma-ray flux - although not unlikely considerable in size - is hard to predict, it is rather the spectral shape that eventually may separate a dark matter signal from the background. The HESS observations of the galactic center, e.g., show a power-law energy spectrum with a spectral index of about $`2`$ up to at least 10 TeV hess , and for Higgsinos as heavy as that, we expect a spectrum that is too hard to explain the full data set Profumo:2005xd . Nevertheless, even a lighter higgsino could still partly contribute; once one has access to better statistics, a characteristic distortion in the spectrum would then be distinguishable at the dark matter particle’s mass. One should furthermore bear in mind that the best prospects for detection might therefore not be found near the galactic center, but rather for sources with a low or at least well understood background. Examples for this could be nearby dark matter clumps or intermediate mass black holes Bertone:2005xz ; a thorough analysis of the detectional prospects for these candidates, however, is beyond the scope of this Letter.
Finally, we note that the radiative corrections presented in this Letter add to the total annihilation cross section, quite independently of the relative velocity of the annihilating neutralinos. This is relevant for any precise calculation of neutralino relic densities and is therefore of relevance for computer packages like DarkSusy and micrOMEGAs.
We thank J. Edsjö for useful discussions and support with the DarkSusy package. L.B. is grateful to the Swedish Science Research Council (VR) for support. |
warning/0507/cond-mat0507395.html | ar5iv | text | # Comparison study of DFA and DMA methods in analysis of autocorrelations in time series
## 1 Introduction
The main problem discussed in the context of stochastic time series in various physical, biological, financial and economical processes is the presence of autocorrelations in data. One of the technique to check whether such autocorrelations are present in time series is based on the investigation of the fractal structure in time series and is related to the scaling exponent H, sometimes denoted also as $`\alpha `$ and called Hurst exponent. It plays a significant role as the main concept upon which fluctuations of a time series around its local trend (drift) are formed and it may be considered as the one of the crucial points responsible for ’genetic code’ of time series of various origin. For the purpose of mentioned above fractal analysis one can introduce the scaling exponent $`\alpha `$ as follows.
Let $`x(t)`$ ($`t=1,\mathrm{},L`$) is the time series defined for discrete time points $`t`$. By rescaling time axis $`\gamma `$ times (e.g. enlarging it $`\times 10^n`$), one reveals the tiny structure of time series not visible for smaller resolution ($`\gamma 1`$). The fractal structure of the series comes from the relation:
$`x^{}(t^{})\mathrm{\Gamma }x(\gamma ^1t)x(t)`$ (1)
where $``$ means similarity correspondence.
The above formula indicates that the magnitude of rescaled time series $`x(\gamma ^1t)`$ should be simultaneously increased $`\mathrm{\Gamma }`$ times in order to satisfy full (local) equivalence of $`x(t)`$ and $`x^{}(t^{})`$ series.
It turns out that the scaling factor $`\mathrm{\Gamma }`$ can be expressed in terms of time rescaling factor $`\gamma `$ with the use of Hurst-Hausdorff $`\alpha `$ exponent ($`\alpha >0`$):
$`\mathrm{\Gamma }=\gamma ^\alpha `$ (2)
The commonly accepted methods to measure $`\alpha `$ exponent are Rescaled Range Analysis (R/S), spectral density analysis , and Detrended Fluctuation Analysis (DFA) . Recently, new method called Detrended Moving Average (DMA) has also been proposed . In this article we will focus on the latter two methods due to large uncertainties in spectral density analysis and problems with R/S predictions in nonstationary series. Searches for better understanding how results of these two methods relate to each other are in progress . A DFA method was first developed for biological purposes and then applied also to finances . It is a detrendisation technique basically measuring fluctuations of a given time series around its local trend as a function of the trend length. Let us recall the main steps of this method:
1. A given signal $`x(t)`$ ($`t=1,\mathrm{},L`$) of time series is divided into $`L/\tau `$ not overlapping boxes of length $`\tau `$ each.
2. A polynomial fit $`x_{\tau ,k}`$ is constructed in each box representing the local trend in that box, where $`k`$ is the order of polynomial fit.
3. A detrended signal $`X_{\tau ,k}(t)`$ is found:
$`X_{\tau ,k}(t)=x(t)x_{\tau ,k}(t)`$ (3)
and then its fluctuation (standard deviation)$`F_{DFA}(\tau ,k)`$ is calculated
$`F_{DFA}(\tau ,k)=\left({\displaystyle \frac{1}{L}}{\displaystyle \underset{t=1}{\overset{L}{}}}X_{\tau ,k}^2(t)\right)^{1/2}`$ (4)
4. From the basic differential stochastic equation of the time series $`x(t)`$ with a local drift $`\mu (t)`$ and a local dispersion $`\sigma (t)`$
$`dx(t)=\mu (t)dt+\sigma (t)dX(t)`$ (5)
one expects the power law behavior:
$`F_{DFA}(\tau ,k)\tau ^{\alpha (k)}`$ (6)
where $`\alpha (k)`$ is the searched Hurst exponent.
The last equation enables to calculate $`\alpha `$ exponent directly from log-log linear fit:
$`\mathrm{log}F_{DFA}(\tau ,k)\alpha (k)\mathrm{log}\tau `$ (7)
It can be proved that $`\alpha (k)`$ depends very weakly on $`k`$ so in most application one takes linear function $`(k=1)`$ as a good candidate for $`x_{\tau ,k}`$. This approach will also be used in our paper.
It turns out that the bigger $`\alpha `$ the more ’quiet’ time series is, i.e. a signal fluctuates in a more correlated way. In fact, for $`0<\alpha <1/2`$ we have negative autocorrelations (antipersistence) in time series. On the other hand, if $`1/2<\alpha 1`$, there are positive autocorrelations (persistence) in signal. The case $`\alpha =1/2`$ corresponds to completely uncorrelated signal, so called integer Brownian walk. An existing link between $`\alpha `$ exponent and the probability that a given trend will last in the immediate future if it did so in the immediate past gives an additional hint about trend changes forecast possibility .
A Detrended Moving Average (DMA) technique looks very similar to DFA. The main difference one meets here is that instead of linear or polynomial detrendisation procedure in equally sized boxes, one uses moving average of a given length $`\lambda `$. The basic steps of DMA analysis are then:
1. A simple moving average of length $`\lambda `$ ($`\lambda =1,\mathrm{},L`$) is constructed for $`x(t)`$ series $`(t\lambda )`$:
$`x(t)_\lambda ={\displaystyle \frac{1}{\lambda }}{\displaystyle \underset{k=0}{\overset{\lambda 1}{}}}x(tk)`$ (8)
2. A detrended signal is found similarly to Eq. (3):
$`X_\lambda (t)=x(t)x(t)_\lambda `$ (9)
and its fluctuation within a window of size $`\lambda `$ reads now:
$`F_{DMA}(\lambda )=\left({\displaystyle \frac{1}{L\lambda +1}}{\displaystyle \underset{t=\lambda }{\overset{L}{}}}X_\lambda ^2(t)\right)^{1/2}`$ (10)
3. Similarly to DFA a power law should be observed
$`\mathrm{log}F_{DMA}(\lambda )\alpha \mathrm{log}\lambda `$ (11)
where $`\alpha `$ is the searched scaling Hurst exponent.
The DMA technique is less complicated and seems to be faster in practical application than DFA algorithm. However, so far no final clear conclusion has been reached regarding mutual relationship between DFA and DMA results for the same series.This article contributes to the above area of interest.
## 2 DMA–DFA Comparison Study
Preliminary results obtained for some real financial series suggest that $`\alpha _{DMA}`$ values are lower than corresponding $`\alpha _{DFA}`$ results. It seems to be confirmed for the set of artificial time series of length $`L2^{18}`$ constructed with the use of Random Midpoint Displacement (RMD) algorithm where one finds $`\alpha _{DFA}\alpha _{DMA}+0.05`$ . This supports the existence of systematic displacement between DFA and DMA results, at least for longer series. In many practical applications however, the length of time series we deal with is shorter (e.g. finance, biology, genetics, medicine), especially if one looks at the local $`\alpha `$ exponent value rather than the global one .
To attack the problem of mutual dependence between DMA and DFA results for series of various length, let us first look at the set of artificial arithmetic integer Brownian time series of length $`L=3\times 10^4`$ with discrete time interval $`\mathrm{\Delta }t=1`$, i.e.:
$`x(L\mathrm{\Delta }t)=x_0+{\displaystyle \underset{k=1}{\overset{L}{}}}\mathrm{\Delta }x_k`$ (12)
where $`\mathrm{\Delta }x_k`$ ($`k=1,\mathrm{},L`$) are centered and normalized displacements generated by random number generator.
Two cases with opposite relation $`\alpha _{DMA}`$ vs $`\alpha _{DFA}`$ are shown in Fig. 1. In the first case $`\alpha _{DFA}>\alpha _{DMA}`$ and $`\alpha _{DFA}\alpha _{DMA}=0.02`$, in the other one $`\alpha _{DFA}<\alpha _{DMA}`$ and $`\alpha _{DMA}\alpha _{DFA}=0.04`$. Thus no systematic relationship is produced.
This induces to treat the problem statistically, i.e. one should find statistical distributions of Hurst exponents measured within two methods for artificial series of various length. It seems to be interesting to compare two statistics and to work out correlations between scaling exponents measured within DMA and DFA techniques for the same sample of time series.
For this purpose we took samples of arithmetic Brownian time series of length $`L`$ in the range $`10^210^5`$. Each sample contained $`N65000`$ series of fixed length. We tried to cover uniformly the whole range of $`L`$ in log-scale keeping $`LL_0q^n`$ with the approximate log step $`q7/4`$ to create variety of lengths.
For any sample of fixed length series the averaged scaling range $`\tau `$ or $`\lambda `$ has been calculated for defined number of candidates ($`30`$) and the corresponding standard deviation $`\sigma _\tau `$ ($`\sigma _\lambda `$). The scaling range was taken as the range of $`\tau `$ or $`\lambda `$ variables strictly obeying scaling laws of Eqs. (7),(11) and assumed to terminate respectively at $`\tau \sigma _\tau `$ for DFA and $`\lambda \sigma _\lambda `$ for DMA. Only series with regression statistical correlation coefficient $`R^2>0.98`$ were taken into account for $`\alpha `$ exponent extraction. For any sample of time series a statistical distribution of $`\alpha _{DFA}`$ and $`\alpha _{DMA}`$ frequencies has been built.
The full range of obtained distribution results is shown in Fig.2-9. The first observation one makes is that for any length $`L`$ both distributions fit very well normal distributions, but with different parameters for the gaussian curve. We made all plots also for centered and normalized $`\alpha `$ frequencies in semi-log scale (Fig.2(b,c)–9(b,c)). Only small deviations from the normal distribution are observed in tails - basically due to smaller statistics there. A good correspondence with gaussian curve is confirmed also in Kolmogorov and Anderson-Darling tests, whose results are displayed in Table 1 and shown for chosen lengths $`L`$ in Fig.10.
One may notice that the standard deviation $`\sigma _{DFA}`$ of $`\alpha _{DFA}`$ scaling parameters is always smaller than the corresponding standard deviation $`\sigma _{DMA}`$ of $`\alpha _{DMA}`$ exponents, and both standard deviations decrease when $`L`$ grows. This can be explained in terms of different sensitivity of DFA and DMA techniques to the presence of random autocorrelations in time series. Such autocorrelations are naturally randomly distributed in any sample of generated time series and hence a distribution of $`\alpha `$ exponent is normal. The probability of random autocorrelations is bigger for short time series, where all statistical fluctuations manifest in a more vivid way. When $`L`$ increases, their influence on the presumed global autocorrelation in series can be neglected. Therefore, both standard deviations $`\sigma _{DFA}`$ and $`\sigma _{DMA}`$ drop with increasing $`L`$. However, we always observe $`\sigma _{DFA}<\sigma _{DMA}`$, what indicates that DMA technique is more sensitive to such ”autocorrelation noise” than DFA one.
One may look at this problem also from another side - like in Fig. 11. Here we have drawn several plots of DFA and DMA analysis, i.e. $`\mathrm{ln}F`$ vs $`\mathrm{ln}\tau `$ or $`\mathrm{ln}\lambda `$ for several corresponding artificial Brownian series of length $`L=1000`$. It is seen that deviations from the strict power law behavior, if occur, are more drastic for DMA than for DFA case and the dispersion of produced slopes is also larger for DMA than for DFA, despite the fact that DMA plots are more smooth in comparison with DFA ones.
The next observation concerns the mean values. One gets $`\alpha _{DFA}_N<\alpha _{DMA}_N`$ for all $`L`$, where $`._N`$ is taken over a sample of $`N`$ time series. A clear shift of the central DMA values to the right with respect to DFA ones (see Figs. 2–9(a)) does not suggest however the presence of systematic relation between $`\alpha _{DMA}`$ and $`\alpha _{DFA}`$. Indeed evaluating the correlation coefficient (values are shown in the description of Fig. 2–9(a)):
$`corr(\alpha _{DFA},\alpha _{DMA})={\displaystyle \frac{\alpha _{DFA}\alpha _{DMA}_N\alpha _{DFA}_N\alpha _{DMA}_N}{\sigma _{DFA}\sigma _{DMA}}}`$ (13)
one finds it increasing with $`L`$, but it never indicates the full correlation. Its value is maximal for large $`L`$, where $`corr(\alpha _{DFA},\alpha _{DMA})0.8`$ for $`L10^410^5`$.
This situation is graphically illustrated in Fig. 12, where a correlation plot $`\alpha _{DFA}`$ vs $`\alpha _{DMA}`$ is shown for Hurst exponent values obtained for $`L=3000`$, $`L=10000`$ and $`L=30000`$ series. From the asymmetry of plots against diagonal one notices that DMA gives higher values than DFA method in most series. This result is independent on the length of time series. In fact the percentage excess of cases $`n_+`$, where $`\alpha _{DMA}>\alpha _{DFA}`$ over the cases where $`\alpha _{DMA}<\alpha _{DFA}`$ ($`n_{}`$), i.e.:
$`\delta _\pm ={\displaystyle \frac{n_+n_{}}{n_++n_{}}}`$ (14)
changes from $`20\%25\%`$ for series with $`L<10000`$ up to $`50\%`$ for longer series.
It is obvious therefore that the mean of difference $`\delta _{DFADMA}`$, where
$`\delta _{DFADMA}=\alpha _{DFA}\alpha _{DMA}_N`$ (15)
is not a good measure of ’distance’ between two investigated methods. It is more convenient to define this distance in a standard way, i.e.:
$`\mathrm{\Delta }_{DFADMA}=\left((\alpha _{DFA}\alpha _{DMA})^2_N\right)^{1/2}`$ (16)
The sufficient number of time series samples of various length has been worked out to find a relationship $`\mathrm{\Delta }_{DFADMA}(L)`$. The polynomial best fit for the collected data is drawn in Fig. 13 with error bars coming from the uncertainties in slope determination. This plot indicates that the average displacement between $`\alpha _{DFA}`$ and $`\alpha _{DMA}`$ exponents for a given time series ranges from $`15\%`$ for series with $`L10^3`$, down to $`2\%`$ for long series ($`L10^5`$). The latter value is much smaller than one reported in . The fastest drop in DFA-DMA distance is observed for medium length series, i.e. when $`L10^310^4`$. For such series $`\mathrm{\Delta }_{DFADMA}`$ makes on the average $`10\%`$ of $`\alpha _{DFA}`$ value.
This might be of interest if more detailed study of $`\alpha `$ exponent is required for more exact predictions to be made(e.g. heart diseases, finances, etc.). The plot in Fig. 13 may also suggests that $`\mathrm{\Delta }_{DFADMA}0`$ when $`L\mathrm{}`$. The latter case has not been explored in details.
## 3 Conclusions
We report from the analysis of artificial Brownian integer time series and from the collected data that, on the average, DMA method overestimates Hurst exponent values in comparison with DFA technique. This result contradicts to some previous hypothesis in literature. The DMA method seems to be also more sensitive to the presence of random fluctuations in autocorrelations in time series than DFA analysis does. In many practical situations, especially for shorter series, it might be a disadvantage leading to the false signal of not really existing, global autocorrelations in time series.
The mean distance between two methods, i.e. the mean difference between $`\alpha _{DFA}`$ and $`\alpha _{DMA}`$ exponents calculated for the time series of given length $`L`$ is a decreasing function of $`L`$. For shorter series ($`L6000)`$ this distance reaches $`15\%`$ what might be important in precise determination of $`\alpha `$ exponent for such series.
There are some open questions. It is not exactly clear where the scaling law exactly starts or terminates, so one needs a more strict requirements how the scaling range should be determined for DFA and DMA techniques and how uncertainties in the choice of scaling range are related to uncertainties in the scaling exponent $`\alpha `$. This work is now in progress . |
warning/0507/math0507235.html | ar5iv | text | # Analyticity of Entropy Rate of Hidden Markov Chains
## 1 Introduction
For $`m,n`$ with $`mn`$, we denote a sequence of symbols $`y_m,y_{m+1},\mathrm{},y_n`$ by $`y_m^n`$. Consider a stationary stochastic process $`Y`$ with a finite set of states $`\{1,2,\mathrm{},B\}`$ and distribution $`p(y_m^n)`$. Denote the conditional distributions by $`p(y_{n+1}|y_m^n)`$. The entropy rate of $`Y`$ is defined as
$$H(Y)=\underset{n\mathrm{}}{lim}E_p(\mathrm{log}(p(y_0|y_n^1))),$$
where $`E_p`$ denotes expectation with respect to the distribution $`p`$.
Let $`Y`$ be a stationary first order Markov chain with
$$\mathrm{\Delta }(i,j)=p(y_1=j|y_0=i).$$
It is well known that
$$H(Y)=\underset{i,j}{}p(y_0=i)\mathrm{\Delta }(i,j)\mathrm{log}\mathrm{\Delta }(i,j).$$
A *hidden Markov chain* $`Z`$ (or function of a Markov chain) is a process of the form $`Z=\mathrm{\Phi }(Y)`$, where $`\mathrm{\Phi }`$ is a function defined on $`\{1,2,\mathrm{},B\}`$ with values $`\{1,2,\mathrm{},A\}`$. Often a hidden Markov chain is defined as a Markov chain observed in noise. It is well known that the two definitions are equivalent (the equivalence is typified by Example 4.1).
For a hidden Markov chain, $`H(Z)`$ turns out (see Equation (2.4) below) to be the integral of a certain function defined on a simplex with respect to a measure due to Blackwell . However Blackwell’s measure is somewhat complicated and the integral formula appears to be difficult to evaluate in most cases.
Recently there has been a rebirth of interest in computing the entropy rate of a hidden Markov chain, and many approaches have been adopted to tackle this problem. For instance, some researchers have used Blackwell’s measure to bound the entropy rate and others introduced a variation on bounds due to .
In a new direction, have studied the variation of the entropy rate as parameters of the underlying Markov chain vary. These works motivated us to consider the general question of whether the entropy rate of a hidden Markov chain is smooth, or even analytic , as a function of the underlying parameters. Indeed, this is true under mild positivity assumptions:
###### Theorem 1.1.
Suppose that the entries of $`\mathrm{\Delta }`$ are analytically parameterized by a real variable vector $`\stackrel{}{\epsilon }`$. If at $`\stackrel{}{\epsilon }=\stackrel{}{\epsilon }_0`$,
1. For all $`a`$, there is at least one $`j`$ with $`\mathrm{\Phi }(j)=a`$ such that the $`j`$-th column of $`\mathrm{\Delta }`$ is strictly positive – and –
2. Every column of $`\mathrm{\Delta }`$ is either all zero or strictly positive,
then $`H(Z)`$ is a real analytic function of $`\stackrel{}{\epsilon }`$ at $`\stackrel{}{\epsilon }_0`$.
Note that this theorem holds if all the entries of $`\mathrm{\Delta }`$ are positive. The more general form of our hypotheses is very important (see Example 4.1).
Real analyticity at a point is important because it means that the function can be expressed as a convergent power series in a neighborhood of the point. The power series can be used to approximate or estimate the function. For convenience of the reader, we recall some basic concepts of analyticity in Section 3.
Several authors have observed that the entropy rate of a hidden Markov chain can be viewed as the top Lyapunov exponent of a random matrix product . Results in show that under certain conditions the top Lyapunov exponent of a random matrix product varies analytically as either the underlying Markov process varies analytically or as the matrix entries vary analytically, but not both. However, when regarding the entropy rate as a Lyapunov exponent of a random matrix product, the matrix entries depend on the underlying Markov process. So, the results from Lyapunov theory do not appear to apply directly. Nevertheless, much of the main idea of our proof of Theorem 1.1 is essentially contained in Peres . In contrast to Peres’ proof, we do not use the language of Lyapunov exponents and we use only basic complex analysis and no functional analysis. Also the hypotheses in do not carry over to our setting. To the best of our knowledge the statement and proof of Theorem 1.1 has not appeared in the literature. For analyticity of certain other statistical quantities, see also related work in the area of statistical physics in .
After discussing background in Sections 2 and 3, we prove Theorem 1.1 in Section 4. As an example, we show that the entropy rate of a hidden Markov chain obtained by observing a binary Markov chains in binary symmetric noise, with noise parameter $`\epsilon `$, is analytic at any $`\epsilon =\epsilon _00`$, provided that the Markov transition probabilities are all positive.
In Section 5, we infer from the proof of Theorem 1.1 a general principle to determine a domain of analyticity for the entropy rate. We apply this to the case of hidden Markov chains obtained from binary Markov chains in binary symmetric noise to find a lower bound on the radius of convergence of a power series in $`\epsilon `$ at $`\epsilon _0=0`$. Given the recent results of , which compute the derivatives of all orders at $`\epsilon _0=0`$, this gives an explicit power series for entropy rate near $`\epsilon _0=0`$.
In Section 6, we show how to relax the conditions of Theorem 1.1 and apply this to give more examples where the entropy rate is analytic.
The entropy rate can fail to be analytic. In Section 7 we give examples and then give a complete set of necessary and sufficient conditions for analyticity in the special case of binary hidden Markov chains with an unambiguous symbol, i.e., a symbol which can be produced by only one symbol of the Markov chain.
Finally in Section 8, we resort to more advanced techniques to prove a stronger version, Theorem 8.1, of Theorem 1.1. This result gives a sense in which the hidden Markov chain itself varies analytically with $`\stackrel{}{\epsilon }`$. The proof of this result requires some measure theory and functional analysis, along with ideas from equilibrium states , which are reviewed in Appendix C. Our first proof of Theorem 1.1 was derived as a consequence of Theorem 8.1. It also follows from Theorem 8.1 that, in principle, many statistical properties in addition to entropy rate vary analytically.
Most results of this paper were first announced in .
## 2 Iteration on the Simplex
Let $`W`$ be the simplex, comprising the vectors
$$\{w=(w_1,w_2,\mathrm{},w_B)^B:w_i0,\underset{i}{}w_i=1\},$$
and let $`W_a`$ be all $`wW`$ with $`w_i=0`$ for $`\mathrm{\Phi }(i)a`$. Let $`W^{}`$ denote the complex version of $`W`$, i.e., $`W^{}`$ denotes the complex simplex comprising the vectors
$$\{w=(w_1,w_2,\mathrm{},w_B)^B:\underset{i}{}w_i=1\},$$
and let $`W_a^{}`$ denote the complex version of $`W_a`$, i.e., $`W_a^{}`$ consists of all $`wW^{}`$ with $`w_i=0`$ for $`\mathrm{\Phi }(i)a`$. For $`aA`$, let $`\mathrm{\Delta }_a`$ denote the $`B\times B`$ matrix such that $`\mathrm{\Delta }_a(i,j)=\mathrm{\Delta }(i,j)`$ for $`j`$ with $`\mathrm{\Phi }(j)=a`$, and $`\mathrm{\Delta }_a(i,j)=0`$ otherwise. For $`aA`$, define the scalar-valued and vector-valued functions $`r_a`$ and $`f_a`$ on $`W`$ by
$$r_a(w)=w\mathrm{\Delta }_a\mathrm{𝟏},$$
and
$$f_a(w)=w\mathrm{\Delta }_a/r_a(w).$$
Note that $`f_a`$ defines the action of the matrix $`\mathrm{\Delta }_a`$ on the simplex $`W`$. For any fixed $`n`$ and $`z_n^0`$, define
$$x_i=x_i(z_n^i)=p(y_i=|z_i,z_{i1},\mathrm{},z_n),$$
(2.1)
(here $``$ represent the states of the Markov chain $`Y`$,) then from Blackwell , $`\{x_i\}`$ satisfies the random dynamical iteration
$$x_{i+1}=f_{z_{i+1}}(x_i),$$
(2.2)
starting with
$$x_{n1}=p(y_{n1}=).$$
(2.3)
We remark that Blackwell showed that
$$H(Z)=\underset{a}{}r_a(w)\mathrm{log}r_a(w)dQ(w),$$
(2.4)
where $`Q`$, known as Blackwell’s measure, is the limiting probability distribution, as $`n\mathrm{}`$, of $`\{x_0\}`$ on $`W`$. However, we do not use Blackwell’s measure explicitly in this paper.
Next, we consider two metrics on a compact subset $`S`$ of the interior of a subsimplex $`W^{}`$ of $`W`$. Without loss of generality, we assume that $`W^{}`$ consists of all points from $`W`$ with the last $`Bk`$ coordinates equal to $`0`$. The Euclidean metric $`d_\text{E}`$ on $`S`$ is defined as usual, namely for $`u,vS`$,
$$u=(u_1,u_2,\mathrm{},u_B),v=(v_1,v_2,\mathrm{},v_B)S,$$
we have
$$d_\text{E}(u,v)=\sqrt{(u_1v_1)^2+(u_2v_2)^2+\mathrm{}+(u_kv_k)^2}.$$
The Hilbert metric $`d_\text{B}`$ on $`S`$ is defined as follows:
$$d_\text{B}(u,v)=\underset{ijk}{\mathrm{max}}\mathrm{log}\left(\frac{u_i/u_j}{v_i/v_j}\right).$$
The following result is well known (for instance, see ). For completeness, we give a detailed proof in Appendix A.
###### Proposition 2.1.
$`d_\text{E}`$ and $`d_\text{B}`$ are equivalent (denoted by $`d_\text{E}d_\text{B}`$) on any compact subset $`S`$ of the interior of a subsimplex $`W^{}`$ of $`W`$, i.e., there are positive constants $`C_1<C_2`$ such that for any two points $`u,vS`$,
$$C_1d_\text{B}(u,v)<d_\text{E}(u,v)<C_2d_\text{B}(u,v).$$
###### Proposition 2.2.
Assume that at $`\stackrel{}{\epsilon }_0`$, $`\mathrm{\Delta }`$ satisfies conditions $`1`$ and $`2`$ of Theorem 1.1. Then for sufficiently large $`n`$ and all choices of $`a_1,\mathrm{},a_n`$ and $`b`$, the mapping $`f_{a_n}f_{a_{n1}}\mathrm{}f_{a_1}`$ is a contraction mapping under the Euclidean metric on $`W_b`$.
###### Proof.
$`\widehat{W}_b=f_b(W)`$ is a compact subset of the interior of some subsimplex $`W_b^{}`$ of $`W_b`$; this subsimplex corresponds to column indices $`j`$ such that $`\mathrm{\Phi }(j)=b`$ and the $`j`$-th column is strictly positive. Therefore one can define the Hilbert metric accordingly on $`\widehat{W}_b`$. Each $`f_a`$ is a contraction mapping on each $`\widehat{W}_b`$ under the Hilbert metric ; namely there exists $`0<\rho <1`$ such that for any $`a`$ and $`b`$, and for any two points $`u,v\widehat{W}_b`$,
$$d_\text{B}(f_a(u),f_a(v))<\rho d_\text{B}(u,v).$$
Thus, for any choices of $`a_2,a_3,\mathrm{},a_n`$, we have
$$d_\text{B}(f_{a_n}f_{a_{n1}}\mathrm{}f_{a_2}(u),f_{a_n}f_{a_{n1}}\mathrm{}f_{a_2}(v))<\rho ^{n1}d_\text{B}(u,v).$$
By Proposition 2.1, there exists a positive constant $`C`$ such that
$$d_\text{E}(f_{a_n}f_{a_{n1}}\mathrm{}f_{a_2}(u),f_{a_n}f_{a_{n1}}\mathrm{}f_{a_2}(v))<C\rho ^{n1}d_\text{E}(u,v).$$
Let $`L`$ be a universal Lipschitz constant for any $`f_c:W_bW_c^{}`$ with respect to the Euclidean metric. Choose $`n`$ large enough such that $`C\rho ^{n1}<1/L`$. So, for sufficiently large $`n`$, any composition of the form $`f_{a_n}\mathrm{}f_{a_1}`$ is a Euclidean contraction on $`W_b`$.
###### Remark 2.3.
Using a slightly modified proof, one can show that for sufficiently large $`n`$, any composition of the form $`f_{a_n}\mathrm{}f_{a_1}`$ is a Euclidean contraction on the whole simplex $`W`$.
## 3 Brief background on analyticity
In this section, we briefly review the basics in complex analysis for the purpose of this paper. For more details, we refer to .
A real (or complex) function of several variables is analytic at a given point if it admits a convergent Taylor series representation in a real (or complex) neighborhood of the given point. We say that it is real (or complex) analytic in a neighborhood if it is real (or complex) analytic at each point of the neighborhood.
The relationship between real and complex analytic functions is as follows: 1) Any real analytic function can be extended to a complex analytic function on some complex neighborhood; 2) Any real function obtained by restricting a complex analytic function from a complex neighborhood to a real neighborhood is a real analytic function.
The main fact regarding analytic functions used in this paper is that the uniform limit of a sequence of complex analytic functions on a fixed complex neighborhood is complex analytic. The analogous statement does not hold (in fact, fails dramatically!) for real analytic functions.
As an example of a real-valued parametrization of a matrix, consider:
$$\mathrm{\Delta }(\epsilon )=\left[\begin{array}{ccc}2\epsilon & \epsilon & 13\epsilon \\ \epsilon & 1\epsilon \mathrm{sin}(\epsilon )& \mathrm{sin}(\epsilon )\\ \epsilon ^2& \epsilon ^3& 1\epsilon ^2\epsilon ^3\end{array}\right].$$
Denote the states of $`\mathrm{\Delta }`$ by $`\{1,2,3\}`$ and let $`\mathrm{\Phi }(1)=\mathrm{\Phi }(2)=0,\mathrm{\Phi }(3)=1`$. Each entry of $`\mathrm{\Delta }`$ is a real analytic function of $`\epsilon `$ at any given point $`\epsilon =\epsilon _0`$. For $`\epsilon _0>0`$ and sufficiently small, $`\mathrm{\Delta }`$ is stochastic (i.e., each row sums to 1 and each entry is nonnegative) and in fact strictly positive (i.e., each entry is positive). According to Theorem 1.1, for such values of $`\epsilon _0`$, the entropy rate of the hidden Markov chain defined by $`\mathrm{\Delta }(\epsilon )`$ and $`\mathrm{\Phi }`$ is real analytic as a function of $`\epsilon `$ at $`\epsilon _0`$. .
While we typically think of analytic parametrizations as having the “look” of the preceding example, there is a conceptually simpler parametrization – namely, parameterize an $`n\times n`$ matrix $`\mathrm{\Delta }`$ by its entries themselves; if $`\mathrm{\Delta }`$ is required to be stochastic, we choose the parameters to be any set of $`n1`$ entries in each row (so, the real variable vector is an $`n(n1)`$-tuple). Clearly, for analyticity it does not matter which entries are chosen. We call this the natural parametrization.
Suppose that $`H(Z)`$ is analytic with respect to this parametrization. Then, $`H(Z)`$ viewed as a function of any other analytic parametrization of the entries of $`\mathrm{\Delta }`$ is the composition of two analytic functions and thus must be analytic. We thus have that the following two statements are equivalent.
* $`H(Z)`$ is analytic with respect to the natural parameterization.
* $`H(Z)`$ is analytic with respect to any analytic parameterization.
We shall use this implicitly through the paper.
## 4 Proof of Theorem 1.1
Notation: We rewrite $`\mathrm{\Delta }`$, $`Z`$, $`f_a(x)`$, $`p(z_0|z_{\mathrm{}}^1)`$ with parameter vector $`\stackrel{}{\epsilon }`$ as $`\mathrm{\Delta }^\stackrel{}{\epsilon }`$, $`Z^\stackrel{}{\epsilon }`$, $`f_a^\stackrel{}{\epsilon }(x)`$ and $`p^\stackrel{}{\epsilon }(z_0|z_{\mathrm{}}^1)`$, respectively. We use the notation $`\widehat{W}_a`$ to mean $`f_a^{\epsilon _0}(W)`$. Let $`\mathrm{\Omega }_{}=\mathrm{\Omega }_{}(r)`$ denote the set of points of distance at most $`r`$ from $`\stackrel{}{\epsilon }_0`$ in the complex parameter space $`^m`$. Let $`N_b=N_b(R)`$ denote the set of all points in $`W_b^{}`$ of distance at most $`R`$ from $`\widehat{W}_b`$.
We first prove that for some $`r>0`$, $`\mathrm{log}p^\stackrel{}{\epsilon }(z_0|z_{\mathrm{}}^1)`$ can be extended to a complex analytic function of $`\stackrel{}{\epsilon }\mathrm{\Omega }_{}(r)`$ and that $`|\mathrm{log}p^\stackrel{}{\epsilon }(z_0|z_{\mathrm{}}^1)\mathrm{log}p^\stackrel{}{\epsilon }(\widehat{z}_0|\widehat{z}_{\mathrm{}}^1)|`$ decays exponentially fast in $`n`$, when $`z_n^0=\widehat{z}_n^0`$, uniformly in $`\stackrel{}{\epsilon }\mathrm{\Omega }_{}(r)`$.
Note that for each $`a,b`$, $`f_a^\stackrel{}{\epsilon }(w)`$ is a rational function of the entries of $`\mathrm{\Delta }^\stackrel{}{\epsilon }`$ and $`w\widehat{W}_b`$. So, by viewing the real vector variables $`\stackrel{}{\epsilon }`$ and $`w`$ as complex vector variables, we can naturally extend $`f_a^\stackrel{}{\epsilon }(w)`$ to a complex-valued function of complex vector variables $`\stackrel{}{\epsilon }`$ and $`w`$. Since $`\mathrm{\Delta }`$ satisfies conditions $`1`$ and $`2`$ at $`\stackrel{}{\epsilon }_0`$, for sufficiently small $`r`$ and $`R`$, the denominator of $`f_a^\stackrel{}{\epsilon }(w)`$ is nonzero for $`\stackrel{}{\epsilon }`$ in $`\mathrm{\Omega }_{}(r)`$ and $`w`$ in $`N_b(R)`$. Thus, $`f_a^\stackrel{}{\epsilon }(w)`$ is a complex analytic function of $`(\stackrel{}{\epsilon },w)`$ in the neighborhood $`\mathrm{\Omega }_{}(r)\times N_b(R)`$.
Assuming conditions $`1`$ and $`2`$, we claim that $`\mathrm{\Delta }`$ has an isolated (in modulus) maximum eigenvalue $`1`$ at $`\stackrel{}{\epsilon }_0`$. To see this, we apply Perron-Frobenius theory as follows. By permuting the indices, we can express:
$$\mathrm{\Delta }=\left[\begin{array}{cc}U& 0\\ V& 0\end{array}\right]$$
where $`U`$ is the submatrix corresponding to indices with positive columns. The nonzero eigenvalues of $`\mathrm{\Delta }`$ are the same as the eigenvalues of $`U`$, which is a positive stochastic matrix. Such a matrix has isolated (in modulus) maximum eigenvalue $`1`$.
The stationary distribution $`p^\stackrel{}{\epsilon }(y=)`$ (the eigenvector corresponding to the maximum eigenvalue $`1`$) is a rational function of the entries of $`\mathrm{\Delta }^\stackrel{}{\epsilon }`$, since it is a solution of the equation $`v\mathrm{\Delta }^\stackrel{}{\epsilon }=v`$. So, in the same way as for $`f_a^\stackrel{}{\epsilon }(w)`$ we can naturally extend $`p^\stackrel{}{\epsilon }(y=)`$ to a complex analytic function $`p^\stackrel{}{\epsilon }(y=)`$ on $`\mathrm{\Omega }_{}`$.
Extending (2.1) for each $`i`$, we define
$$x_i^\stackrel{}{\epsilon }=x_i^\stackrel{}{\epsilon }(z_n^i)=p^\stackrel{}{\epsilon }(y_i=|z_n^i),$$
(4.5)
by iterating the following complexified random dynamical system (extending (2.2) and (2.3)):
$$x_{i+1}^\stackrel{}{\epsilon }=f_{z_{i+1}}^\stackrel{}{\epsilon }(x_i^\stackrel{}{\epsilon }),$$
(4.6)
starting with
$$x_{n1}^\stackrel{}{\epsilon }=p^\stackrel{}{\epsilon }(y_{n1}=).$$
(4.7)
By Proposition 2.2, for sufficiently large $`n`$, we can replace the set of mappings $`\{f_a^{\epsilon _0}\}`$ with the set $`\{f_{a_n}^{\epsilon _0}f_{a_{n1}}^{\epsilon _0}\mathrm{}f_{a_1}^{\epsilon _0}\}`$ and then assume that each $`f_a^{\epsilon _0}`$ is a Euclidean contraction on each $`W_b`$ with contraction coefficient $`\rho <1`$. Since $`\widehat{W}_b`$ is compact and the definition of $`\rho `$-contraction is given by strict inequality, we can choose $`r`$ and $`R`$ sufficiently small such that
$$f_a^\stackrel{}{\epsilon }\text{ is a Euclidean }\rho \text{contraction on each }N_b(R),\epsilon \mathrm{\Omega }_{}(r).$$
(4.8)
Further, we claim that by choosing $`r`$ still smaller, if necessary,
$$x_i^\stackrel{}{\epsilon }_bN_b(R),\text{ for all }i,n\text{ and all choices of }z_n^i,\epsilon \mathrm{\Omega }_{}(r).$$
(4.9)
To see this, fixing $`\rho `$ and $`R`$, choose $`r`$ so small that
$$|f_a^\stackrel{}{\epsilon }(x)f_a^{\stackrel{}{\epsilon }_0}(x)|R(1\rho ),x_b\widehat{W}_b,\epsilon \mathrm{\Omega }_{}(r)$$
(4.10)
and
$$|p^\stackrel{}{\epsilon }()p^{\stackrel{}{\epsilon }_0}()|R(1\rho ),\epsilon \mathrm{\Omega }_{}(r).$$
(4.11)
Now consider the difference
$$x_{i+1}^\stackrel{}{\epsilon }x_{i+1}^{\stackrel{}{\epsilon }_0}$$
$$=f_{z_{i+1}}^\stackrel{}{\epsilon }(x_i^\stackrel{}{\epsilon })f_{z_{i+1}}^{\stackrel{}{\epsilon }_0}(x_i^{\stackrel{}{\epsilon }_0})=f_{z_{i+1}}^\stackrel{}{\epsilon }(x_i^\stackrel{}{\epsilon })f_{z_{i+1}}^\stackrel{}{\epsilon }(x_i^{\stackrel{}{\epsilon }_0})+f_{z_{i+1}}^\stackrel{}{\epsilon }(x_i^{\stackrel{}{\epsilon }_0})f_{z_{i+1}}^{\stackrel{}{\epsilon }_0}(x_i^{\stackrel{}{\epsilon }_0}).$$
(4.12)
Then by (4.8) , (4.10) and (4.11), and (4.12), for $`i>n1`$, we have
$$|x_{i+1}^\stackrel{}{\epsilon }x_{i+1}^{\stackrel{}{\epsilon }_0}|\rho |x_i^\stackrel{}{\epsilon }x_i^{\stackrel{}{\epsilon }_0}|+R(1\rho ).$$
So,
$$|x_{i+1}^\stackrel{}{\epsilon }x_{i+1}^{\stackrel{}{\epsilon }_0}|R,$$
and thus for all $`i`$, we have $`x_{i+1}^\stackrel{}{\epsilon }_bN_b(R)`$, yielding (4.9). Each $`x_i^\epsilon `$ is the composition of analytic functions on $`\mathrm{\Omega }_{}(r)`$ and so is complex analytic on $`\mathrm{\Omega }_{}(r)`$.
For $`0n_1,n_2\mathrm{}`$, we say two sequences $`\{z_{n_1}^0\}`$ and $`\{\widehat{z}_{n_2}^0\}`$ have a common tail if there exists $`n0`$ with $`nn_1,n_2`$ such that $`z_i=\widehat{z}_i,ni0`$ (denoted by $`z_{n_1}^0\stackrel{n}{}\widehat{z}_{n_2}^0`$).
Let
$$x_i^\stackrel{}{\epsilon }=x_i^\stackrel{}{\epsilon }(z_{n_1}^i)=p^\stackrel{}{\epsilon }(y_i=|z_{n_1}^i),$$
$$\widehat{x}_i^\stackrel{}{\epsilon }=\widehat{x}_i^\stackrel{}{\epsilon }(\widehat{z}_{n_2}^i)=p^\stackrel{}{\epsilon }(y_i=|\widehat{z}_{n_2}^i).$$
Then we have
$$x_{i+1}^\stackrel{}{\epsilon }=f_{z_{i+1}}^\stackrel{}{\epsilon }(x_i^\stackrel{}{\epsilon }),\widehat{x}_{i+1}^\stackrel{}{\epsilon }=f_{z_{i+1}}^\stackrel{}{\epsilon }(\widehat{x}_i^\stackrel{}{\epsilon }).$$
From (4.8) and (4.9), it follows that there exists a positive constant $`L`$ independent of $`n_1`$ and $`n_2`$ such that
$$|x_0^\stackrel{}{\epsilon }\widehat{x}_0^\stackrel{}{\epsilon }|L\rho ^n.$$
(4.13)
Naturally
$$p^\stackrel{}{\epsilon }(z_0|z_n^1)=\underset{\{y_0:\mathrm{\Phi }(y_0)=z_0\}}{}\underset{y_1}{}\mathrm{\Delta }^\stackrel{}{\epsilon }(y_1,y_0)p^\stackrel{}{\epsilon }(y_1|z_n^1).$$
(4.14)
Then, there is a positive constant $`L^{}`$, independent of $`n_1,n_2`$, such that
$$|p^\stackrel{}{\epsilon }(z_0|z_{n_1}^1)p^\stackrel{}{\epsilon }(\widehat{z}_0|\widehat{z}_{n_2}^1)|L^{}\rho ^n.$$
(4.15)
Since $`\mathrm{\Delta }^{\stackrel{}{\epsilon }_0}`$ satisfies conditions $`1`$ and $`2`$, $`p^\stackrel{}{\epsilon }(z_0|z_n^1)`$ is bounded away from $`0`$, uniformly in $`\stackrel{}{\epsilon }\mathrm{\Omega }_{}`$, $`n`$ and choices of $`z_n^1`$; thus there is a positive constant $`L^{\prime \prime }`$, independent of $`n_1,n_2`$, such that
$$|\mathrm{log}p^\stackrel{}{\epsilon }(z_0|z_{n_1}^1)\mathrm{log}p^\stackrel{}{\epsilon }(\widehat{z}_0|\widehat{z}_{n_2}^1)|L^{\prime \prime }\rho ^n.$$
(4.16)
Since for each $`y\{1,\mathrm{},B\}`$, $`p^\stackrel{}{\epsilon }(y)`$ is analytic, from
$$p^\stackrel{}{\epsilon }(z)=\underset{\mathrm{\Phi }(y)=z}{}p^\stackrel{}{\epsilon }(y),$$
we deduce that $`p^\stackrel{}{\epsilon }(z)`$ is analytic. Furthermore since $`p^\stackrel{}{\epsilon }(z_0|z_n^1)`$ is analytic on $`\mathrm{\Omega }_{}`$, we conclude $`p^\stackrel{}{\epsilon }(z_n^0)`$ is analytic on $`\mathrm{\Omega }_{}`$.
Choose $`\sigma `$ so that
$$1<\sigma <1/\rho .$$
If $`r`$ and $`R`$ are chosen sufficiently small, then
$$\underset{z_0}{}|p^\stackrel{}{\epsilon }(z_0|z_n^1)|\sigma ,\epsilon \mathrm{\Omega }_{}(r)\text{ and all sequences }z$$
(4.17)
and
$$\underset{z_0}{}|p^\stackrel{}{\epsilon }(z_0)|\sigma ,\epsilon \mathrm{\Omega }_{}(r).$$
(4.18)
Then we have
$$\underset{z_{n1}^0}{}|p^\stackrel{}{\epsilon }(z_{n1}^0)|=\underset{z_{n1}^0}{}|p^\stackrel{}{\epsilon }(z_{n1}^1)p^\stackrel{}{\epsilon }(z_0|z_{n1}^1)|\underset{z_{n1}^1}{}|p^\stackrel{}{\epsilon }(z_{n1}^1)|\underset{z_0}{}|p^\stackrel{}{\epsilon }(z_0|z_{n1}^1)|\sigma \underset{z_n^0}{}|p^\stackrel{}{\epsilon }(z_n^0)|,$$
implying
$$\underset{z_{n1}^0}{}|p^\stackrel{}{\epsilon }(z_{n1}^0)|\sigma ^{n+2}.$$
(4.19)
Let
$$H_n^\stackrel{}{\epsilon }(Z)=\underset{z_n^0}{}p^\stackrel{}{\epsilon }(z_n^0)\mathrm{log}p^\stackrel{}{\epsilon }(z_0|z_n^1)$$
and
$$\rho _1=\rho \delta <1,$$
then we have
$$|H_{n+1}^\stackrel{}{\epsilon }(Z)H_n^\stackrel{}{\epsilon }(Z)|=|\underset{z_{n1}^0}{}p^\stackrel{}{\epsilon }(z_{n1}^0)\mathrm{log}p^\stackrel{}{\epsilon }(z_0|z_{n1}^1)\underset{z_n^0}{}p^\stackrel{}{\epsilon }(z_n^0)\mathrm{log}p^\stackrel{}{\epsilon }(z_0|z_n^1)|$$
$$=|\underset{z_{n1}^0}{}p^\stackrel{}{\epsilon }(z_{n1}^0)(\mathrm{log}p^\stackrel{}{\epsilon }(z_0|z_{n1}^1)\mathrm{log}p^\stackrel{}{\epsilon }(z_0|z_n^1))|\sigma ^2L^{\prime \prime }\rho _1^n;$$
here the latter inequality follows from (4.16) and (4.19). Thus, for $`m>n`$,
$$|H_m^\stackrel{}{\epsilon }(Z)H_n^\stackrel{}{\epsilon }(Z)|\sigma ^2L^{\prime \prime }(\rho _1^n+\mathrm{}+\rho _1^{m1})\frac{\sigma ^2L^{\prime \prime }\rho _1^n}{1\rho _1}.$$
This establishes the uniform convergence of $`H_n^\stackrel{}{\epsilon }(Z)`$ to a limit $`H_{\mathrm{}}^\stackrel{}{\epsilon }(Z)`$. By Theorem $`\mathrm{2.4.1}`$ of , the uniform limit of complex analytic functions on a fixed complex neighborhood is analytic on that neighborhood, and so $`H_{\mathrm{}}^\stackrel{}{\epsilon }(Z)`$ is analytic on $`\mathrm{\Omega }_{}`$.
For real $`\stackrel{}{\epsilon }`$, $`H_{\mathrm{}}^\stackrel{}{\epsilon }(Z)`$ coincides with the entropy rate function $`H(Z^\stackrel{}{\epsilon })`$, and so Theorem 1.1 follows.
###### Example 4.1.
Consider a binary symmetric channel with crossover probability $`\epsilon `$. Let $`\{Y_n\}`$ be the input Markov chain with the transition matrix
$$\mathrm{\Pi }=\left[\begin{array}{cc}\pi _{00}& \pi _{01}\\ \pi _{10}& \pi _{11}\end{array}\right].$$
(4.20)
At time $`n`$ the channel can be characterized by the following equation
$$Z_n=Y_nE_n,$$
where $``$ denotes binary addition, $`E_n`$ denotes the i.i.d. binary noise with $`p_E(0)=1\epsilon `$ and $`p_E(1)=\epsilon `$, and $`Z_n`$ denotes the corrupted output. Then $`(Y_n,E_n)`$ is jointly Markov, so $`\{Z_n\}`$ is a hidden Markov chain with the corresponding
$$\mathrm{\Delta }=\left[\begin{array}{cccc}\pi _{00}(1\epsilon )& \pi _{00}\epsilon & \pi _{01}(1\epsilon )& \pi _{01}\epsilon \\ \pi _{00}(1\epsilon )& \pi _{00}\epsilon & \pi _{01}(1\epsilon )& \pi _{01}\epsilon \\ \pi _{10}(1\epsilon )& \pi _{10}\epsilon & \pi _{11}(1\epsilon )& \pi _{11}\epsilon \\ \pi _{10}(1\epsilon )& \pi _{10}\epsilon & \pi _{11}(1\epsilon )& \pi _{11}\epsilon \end{array}\right];$$
here, $`\mathrm{\Phi }`$ maps states $`1`$ and $`4`$ to $`0`$ and maps states $`2`$ and $`3`$ to $`1`$. This class of hidden Markov chains has been studied extensively (e.g., ).
By Theorem 1.1, when $`\epsilon `$ and $`\pi _{ij}`$’s are positive, the entropy rate $`H(Z)`$ is analytic as a function of $`\epsilon `$ and $`\pi _{ij}`$’s. This still holds when $`\epsilon =0`$ and the $`\pi _{ij}`$’s are positive, because in this case, we have
$$\mathrm{\Delta }=\left[\begin{array}{cccc}\pi _{00}& 0& \pi _{01}& 0\\ \pi _{00}& 0& \pi _{01}& 0\\ \pi _{10}& 0& \pi _{11}& 0\\ \pi _{10}& 0& \pi _{11}& 0\end{array}\right].$$
## 5 Domain of Analyticity
Suppose $`\mathrm{\Delta }`$ is analytically parameterized by a vector variable $`\stackrel{}{\epsilon }`$, and Conditions 1 and 2 in Theorem 1.1 are satisfied at $`\stackrel{}{\epsilon }=\stackrel{}{\epsilon }_0`$. In principle, the proof of Theorem 1.1 determines a neighborhood $`\mathrm{\Omega }_{}(r)`$ of $`\stackrel{}{\epsilon }_0`$ on which the entropy rate is analytic. Specifically, if one can find $`\rho ,r`$ and $`R`$ such that all of the following hold, then the entropy rate is analytic on $`\mathrm{\Omega }_{}(r)`$.
1. Find $`\rho `$ such that each $`f_a^{\epsilon _0}`$ is a Euclidean $`\rho `$-contraction on each $`W_b`$. Then choose positive $`r,R`$ such that for all $`\stackrel{}{\epsilon }\mathrm{\Omega }_{}(r)`$, each $`f_a^\stackrel{}{\epsilon }`$ is a Euclidean $`\rho `$-contraction on each $`N_b(R)`$ (see (4.8)).
2. Next find $`r`$ smaller (if necessary) such that for all $`\stackrel{}{\epsilon }\mathrm{\Omega }_{}(r)`$, the image of the stationary vector of $`\mathrm{\Delta }^\stackrel{}{\epsilon }`$, under any composition of the mappings $`\{f_a^\epsilon \}`$, stays within $`_bN_b(R)`$ (see (4.9)). Note that the argument in the proof shows that this holds if (4.10) and (4.11) hold.
3. Finally, find $`r,R`$ such that the sum of the absolute values of the complexified conditional probabilities, conditioned on any given past symbol sequence, is $`<1/\rho `$, and similarly for the sum of the absolute values of the complexified stationary probabilities (see (4.17) and (4.18)).
In fact, the proof shows that one can always find such $`\rho ,r,R`$, but in condition $`1`$ above one may need to replace $`f_a`$’s by all $`n`$-fold compositions of the $`f_a`$’s, for some $`n`$.
Recall from Example 4.1 the family of hidden Markov chains $`Z^\epsilon `$ determined by passing a binary Markov chain through a binary symmetric channel with cross-over probability $`\epsilon `$. Recall that $`H(Z^\epsilon )`$ is an analytic function of $`\epsilon `$ at $`\epsilon =0`$ when the Markov transition probabilities are all positive. We shall determine a complex neighborhood of $`0`$ such that the entropy rate, as a function of $`\epsilon `$, is analytic on this neighborhood.
Let $`u_n=p(y_n=0|z_1^n)`$ and $`v_n=p(y_n=1|z_1^n)`$. For $`z_{n+1}=1`$ we have
$$u_{n+1}=\frac{\epsilon (\pi _{00}u_n+\pi _{10}v_n)}{\epsilon (\pi _{00}u_n+\pi _{10}v_n)+(1\epsilon )(\pi _{01}u_n+\pi _{11}v_n)},$$
$$v_{n+1}=\frac{(1\epsilon )(\pi _{01}u_n+\pi _{11}v_n)}{\epsilon (\pi _{00}u_n+\pi _{10}v_n)+(1\epsilon )(\pi _{01}u_n+\pi _{11}v_n)}.$$
Since $`u_n+v_n=1`$, $`u_{n+1}`$ is a function of $`u_n`$; let $`g_1`$ denote this function.
For $`z_{n+1}=0`$ we have
$$u_{n+1}=\frac{(1\epsilon )(\pi _{00}u_n+\pi _{10}v_n)}{(1\epsilon )(\pi _{00}u_n+\pi _{10}v_n)+\epsilon (\pi _{01}u_n+\pi _{11}v_n)},$$
$$v_{n+1}=\frac{\epsilon (\pi _{01}u_n+\pi _{11}v_n)}{(1\epsilon )(\pi _{00}u_n+\pi _{10}v_n)+\epsilon (\pi _{01}u_n+\pi _{11}v_n)}.$$
Again, $`u_{n+1}`$ is a function of $`u_n`$; let $`g_0`$ denote this function.
And for the conditional probability, we have
$$p(z_n=0|z_1^{n1})=((1\epsilon )\pi _{00}+\epsilon \pi _{01})u_n+((1\epsilon )\pi _{10}+\epsilon \pi _{11})v_n.$$
Since $`u_n+v_n=1`$, $`p(z_n=0|z_1^{n1})`$ is a function of $`u_n`$; let $`r_0`$ denote this function. And
$$p(z_n=1|z_1^{n1})=(\epsilon \pi _{00}+(1\epsilon )\pi _{01})u_n+(\epsilon \pi _{10}+(1\epsilon )\pi _{11})v_n.$$
Again, $`p(z_n=1|z_1^{n1})`$ is a function of $`u_n`$; let $`r_1`$ denote this function.
Note that $`g_0,g_1,r_0,r_1`$ are all implicitly parameterized by $`\epsilon `$. The stationary vector $`(\pi _0,\pi _1)`$ of $`Y`$, which doesn’t depend on $`\epsilon `$, is equal to $`(\pi _{10}/(\pi _{10}+\pi _{01}),\pi _{01}/(\pi _{10}+\pi _{01}))`$.
We shall choose $`\rho `$ with $`0<\rho <1`$, $`r>0`$ and $`R>0`$ such that for all $`\epsilon `$ with $`|\epsilon |<r`$
1. $`g_0`$ and $`g_1`$ are $`\rho `$-contraction mappings on $`R`$-neighborhoods of $`0`$ and $`1`$ in the complex plane,
2. the set of all $`\{g_{a_n}g_{a_{n1}}\mathrm{}g_{a_1}(\pi _0)\}`$) are within the $`R`$-neighborhoods of $`0`$ and $`1`$,
3. and $`|r_0(u)|+|r_1(u)|<1/\rho `$ for $`u`$ in $`R`$-neighborhoods of $`0`$ and $`1`$ in the complex plane.
By the general principle above, the entropy rate should be analytic on $`|\epsilon |<r`$.
More concretely, condition $`1`$, $`2`$ and $`3`$ translate to (here $`\rho <1`$):
1. $`|g_0^{}(u)|<\rho `$, $`|g_1^{}(u)|<\rho `$ on ($`|\epsilon |<r`$ and $`|u|<R`$) and ($`|\epsilon |<r`$ and $`|1u|<R`$),
2. $`\mathrm{max}\{|g_0(0)1|,|g_0(1)1|,|g_1(0)|,|g_1(1)|\}<R(1\rho )`$ on $`|\epsilon |<r`$ (this follows from (4.10); (4.11) is trivial since the stationary vector of $`Y`$ doesn’t depend on $`\epsilon `$),
3. $`|r_0(u)|+|r_1(u)|<1/\rho `$ on ($`|\epsilon |<r`$ and $`|u|<R`$) and ($`|\epsilon |<r`$ and $`|1u|<R`$).
A straightforward computation shows that the following conditions guarantee conditions $`1`$, $`2`$, $`3`$:
$$0<\frac{\sqrt{r(|\pi _{00}\pi _{11}+\pi _{10}\pi _{11}+\pi _{10}\pi _{01}\pi _{10}\pi _{11}|r+|(\pi _{00}\pi _{11}+\pi _{10}\pi _{01})|)}}{\pi _{11}|\pi _{10}\pi _{11}|r(|\pi _{00}\pi _{10}\pi _{01}+\pi _{11}|r+|\pi _{01}\pi _{11}|)R}<\sqrt{\rho },$$
$$0<\frac{\sqrt{r(|\pi _{00}\pi _{11}+\pi _{10}\pi _{11}+\pi _{10}\pi _{01}\pi _{10}\pi _{11}|r+|(\pi _{00}\pi _{11}+\pi _{10}\pi _{01})|)}}{\pi _{01}|\pi _{00}\pi _{01}|r(|\pi _{00}\pi _{10}\pi _{01}+\pi _{11}|r+|\pi _{01}\pi _{11}|)R}<\sqrt{\rho },$$
$$0<\frac{\sqrt{r(|\pi _{11}\pi _{00}+\pi _{01}\pi _{00}+\pi _{01}\pi _{10}\pi _{01}\pi _{00}|r+|\pi _{11}\pi _{00}\pi _{01}\pi _{10}|)}}{\pi _{00}|\pi _{01}\pi _{00}|r(|\pi _{00}\pi _{10}+\pi _{11}\pi _{01}|r+|\pi _{10}\pi _{00}|)R}<\sqrt{\rho },$$
$$0<\frac{\sqrt{r(|\pi _{11}\pi _{00}+\pi _{01}\pi _{00}+\pi _{01}\pi _{10}\pi _{01}\pi _{00}|r+|\pi _{11}\pi _{00}\pi _{01}\pi _{10}|)}}{\pi _{10}|\pi _{11}\pi _{10}|r(|\pi _{00}\pi _{10}+\pi _{11}\pi _{01}|r+|\pi _{10}\pi _{00}|)R}<\sqrt{\rho },$$
$$0<\frac{r\pi _{00}}{\pi _{01}|\pi _{00}\pi _{01}|r}<R(1\rho ),0<\frac{r\pi _{10}}{\pi _{11}|\pi _{10}\pi _{11}|r}<R(1\rho ),$$
$$0<\frac{r\pi _{11}}{\pi _{10}|\pi _{11}\pi _{10}|r}<R(1\rho ),0<\frac{r\pi _{01}}{\pi _{00}|\pi _{01}\pi _{00}|r}<R(1\rho ),$$
$$(|\pi _{00}\pi _{01}\pi _{10}+\pi _{11}|r+|\pi _{01}\pi _{11}|)R+|\pi _{10}\pi _{11}|r+\pi _{11},$$
$$+(|\pi _{01}\pi _{00}+\pi _{10}\pi _{11}|r+|\pi _{00}\pi _{10}|)R+|\pi _{11}\pi _{10}|r+\pi _{10}<1/\rho ,$$
$$(|\pi _{10}\pi _{11}\pi _{00}+\pi _{01}|r+|\pi _{11}\pi _{01}|)R+|\pi _{00}\pi _{01}|r+\pi _{01}$$
$$+(|\pi _{11}\pi _{10}+\pi _{00}\pi _{01}|r+|\pi _{10}\pi _{00}|)R+|\pi _{01}\pi _{00}|r+\pi _{00}<1/\rho .$$
In other words, for given $`\rho `$ with $`0<\rho <1`$, choose $`r`$ and $`R`$ to satisfy all the constraints above. Then the entropy rate is an analytic function of $`\epsilon `$ on $`|\epsilon |<r`$.
## 6 Relaxed Conditions
We do not know a complete set of necessary and sufficient conditions on $`\mathrm{\Delta }`$ and $`\mathrm{\Phi }`$ that guarantee analyticity of entropy rate. However, in this section, we show how the hypotheses in Theorem 1.1 can be relaxed and still guarantee analyticity. We then give several examples. In Section 7, we do give a a complete set of necessary and sufficient conditions for a very special class of hidden Markov chains.
In this section, we assume that $`\mathrm{\Delta }`$ has a simple maximum eigenvalue $`1`$; this implies that $`\mathrm{\Delta }`$ has a unique stationary vector $`\stackrel{}{s}`$.
For a mapping $`f`$ from $`W_b`$ to $`W`$ and $`wW_b`$. Let $`f^{}`$ denote the first derivative of $`f`$ at $`w`$ restricted to the subspace spanned by directions parallel to the simplex $`W_b`$ and let $``$ denote the Euclidean norm of a linear mapping. We say that $`\{f_a:aA\}`$ is *eventually contracting* at $`wW_b`$ if there exists $`n`$ such that for any $`a_0,a_1,\mathrm{},a_nA`$, $`(f_{a_n}f_{a_{n1}}\mathrm{}f_{a_0})^{}(w)`$ is strictly less than $`1`$. We say that $`\{f_a:aA\}`$ is *contracting* at $`wW_b`$ if it is *eventually contracting* at $`w`$ with $`n=0`$. Using the mean value theorem, one can show that if $`\{f_a:aA\}`$ is *contracting* at each $`w`$ in a compact convex subset $`K`$ of $`W_b`$ then each $`f_a`$ is a contraction mapping on $`K`$.
Let $`L`$ denote the limit set of $`\{(f_{a_n}f_{a_{n1}}\mathrm{}f_{a_0})(\stackrel{}{s})\}`$.
###### Theorem 6.1.
If at $`\mathrm{\Delta }=\widehat{\mathrm{\Delta }}`$,
1. $`1`$ is a simple eigenvalue for $`\widehat{\mathrm{\Delta }}`$,
2. For every $`a`$ and all $`w`$ in $`L`$, $`r_a(w)>0`$,
3. For every $`b`$, $`\{f_a:aA\}`$ is eventually contracting at all $`w`$ in the convex hull of the intersection of $`L`$ and $`W_b`$,
then $`H(Z)`$ is analytic at $`\mathrm{\Delta }=\widehat{\mathrm{\Delta }}`$.
###### Proof.
Let $`𝒳`$ denote the right infinite shift space $`\{a_0^{\mathrm{}}:a_iA\}`$. Let $`L_\delta `$ be the set of all points in $`W`$ of distance at most $`\delta `$ from $`L`$. Choose $`\delta `$ so small that
* For every $`aA`$ and $`w`$ in $`L_\delta `$, $`r_a(w)>0`$ – and –
* For every $`b`$, $`\{f_a:aA\}`$ is eventually contracting at all $`w`$ in the convex hull of the intersection of $`L_\delta `$ and $`W_b`$.
Since the convex hull $`K_\delta `$ of the intersection of $`L_\delta `$ and $`W_b`$ is compact, there exists $`n`$ such that for any $`a_0,a_1,\mathrm{},a_nA`$ and any $`wK_\delta `$, $`(f_{a_n}f_{a_{n1}}\mathrm{}f_{a_0})^{}(w)`$ is strictly less than $`1`$. For simplicity, we may assume that for each $`a`$, $`\{f_a\}`$ is contracting on $`K_\delta `$, and so each $`f_a`$ is a contraction mapping on $`K_\delta `$. Since $`L_\delta K_\delta `$, it follows that $`f_a(L_\delta )L_\delta `$ and so each $`f_a`$ is a contraction mapping on $`L_\delta `$.
For any $`c_0^{\mathrm{}}𝒳`$, there exists $`n`$ such that $`\{(f_{c_n}f_{c_{n1}}\mathrm{}f_{c_0})(\stackrel{}{s})\}L_\delta `$. Let $`𝒳_{c_0^{\mathrm{}}}^n`$ denote the cylinder set $`\{a_0^{\mathrm{}}:a_0=c_0,a_1=c_1,\mathrm{},a_n=c_n\}`$. Since $`\{f_a:aA\}`$ is a contraction mapping on $`L_\delta `$, we conclude that for any $`a_0^{\mathrm{}}𝒳_{c_0^{\mathrm{}}}^n`$ and all $`mn`$, $`\{(f_{a_m}f_{a_{m1}}\mathrm{}f_{a_0})(\stackrel{}{s})\}L_\delta `$. By the compactness of $`𝒳`$, we can find finitely many such cylinder sets to cover $`𝒳`$. Consequently we can find $`n`$ such that for any $`a_0^{\mathrm{}}𝒳`$ and any $`mn`$ , we have $`\{(f_{a_m}f_{a_{m1}}\mathrm{}f_{a_0})(\stackrel{}{s})\}L_\delta `$. We can now apply the proof of Theorem 1.1 – namely, we can use the contraction (along any symbolic sequence $`z_n^0`$) to extend $`H_n(Z)=H(Z_0|Z_n^1)`$ from real to complex and prove the uniform convergence of $`H_n(Z)`$ to $`H(Z)`$ in complex parameter space. ∎
###### Remark 6.2.
(1) If $`\widehat{\mathrm{\Delta }}`$ has a strictly positive column (or more generally, there is a $`j`$ such that for all $`i`$, there exists $`n`$ such that $`\widehat{\mathrm{\Delta }}_{ij}^n>0`$), then condition $`1`$ of Theorem 6.1 holds by Perron-Frobenius theory.
(2) If for each symbol $`a`$, $`\widehat{\mathrm{\Delta }}_a`$ is row allowable (i.e., no row is all zero), then $`r_a(w)>0`$ for all $`wW`$ and so condition $`2`$ of Theorem 6.1 holds.
Theorem 6.1 relaxes the positivity assumptions of Theorem 1.1. Indeed given conditions $`1`$ and $`2`$ of Theorem 1.1, by Remark 6.2, conditions 1 and 2 of Theorem 6.1 hold. For condition $`3`$ of Theorem 6.1, first observe that $`L`$ is contained in $`_bf_b(W)`$. Using the equivalence of the Euclidean metric and the Hilbert metric, Proposition 2.2 shows that for every $`b`$, $`\{f_a:aA\}`$ is eventually contracting on $`f_b(W)`$, which is a convex set containing the intersection of $`L`$ and $`W_b`$.
Theorem 6.1 also applies to many cases not covered by Theorem 1.1. For instance, suppose that some column of $`\widehat{\mathrm{\Delta }}`$ is strictly positive and each $`\widehat{\mathrm{\Delta }}_a`$ is row allowable. By Remark 6.2, Theorem 6.1 applies whenever we can guarantee condition 3. For this, it is sufficient to check that for each $`a,b`$, $`f_a`$ is a contraction, with respect to the Euclidean metric, on the convex hull of the intersection of $`L`$ with each $`W_b`$. This can be done by explicitly computing derivatives.
###### Example 6.3.
Consider a hidden Markov chain $`Z`$ defined by :
$$\widehat{\mathrm{\Delta }}=\left[\begin{array}{cccc}a_{11}& a_{12}& a_{13}& a_{14}\\ a_{21}& a_{22}& a_{23}& a_{24}\\ a_{31}& a_{32}& a_{33}& a_{34}\\ a_{41}& a_{42}& a_{43}& a_{44}\end{array}\right],$$
with $`\mathrm{\Phi }(1)=\mathrm{\Phi }(2)=0`$ and $`\mathrm{\Phi }(3)=\mathrm{\Phi }(4)=1`$. We assume that some column of $`\widehat{\mathrm{\Delta }}`$ is strictly positive and both $`\widehat{\mathrm{\Delta }}_0`$ and $`\widehat{\mathrm{\Delta }}_1`$ are row allowable.
Parameterize $`W_0`$ by $`(y,1y,0,0)`$ and parameterize $`W_1`$ by $`(0,0,y,1y)`$ (with $`y[0,1]`$). We can explicitly compute the derivatives of $`f_0`$ and $`f_1`$ with respect to $`y`$:
$$f_0^{}|_{(y,1y,0,0)}=\frac{a_{11}a_{22}a_{12}a_{21}}{((a_{11}+a_{12}a_{21}a_{22})y+a_{21}+a_{22})^2},$$
$$f_0^{}|_{(0,0,y,1y)}=\frac{a_{31}a_{42}a_{32}a_{41}}{((a_{31}+a_{32}a_{41}a_{42})y+a_{41}+a_{42})^2},$$
$$f_1^{}|_{(y,1y,0,0)}=\frac{a_{13}a_{24}a_{14}a_{23}}{((a_{13}+a_{14}a_{23}a_{24})y+a_{23}+a_{24})^2},$$
$$f_1^{}|_{(0,0,y,1y)}=\frac{a_{33}a_{44}a_{34}a_{43}}{((a_{33}+a_{34}a_{43}a_{44})y+a_{43}+a_{44})^2},$$
Note that the row allowability condition guarantees that the denominators in these expressions never vanish.
Choose $`a_{ij}`$’s such that each of these derivatives is less than 1; then we conclude that the entropy rate is analytic at $`\widehat{\mathrm{\Delta }}`$. One way to do this is to make each of the $`2\times 2`$ upper/lower left/right matrices singular.
Or choose the $`a_{ij}`$’s such that
$$\widehat{\mathrm{\Delta }}=\left[\begin{array}{cccc}\alpha _1& & \beta _1& 0\\ 0& \alpha _2& 0& \beta _2\\ \lambda _1& & \eta _1& 0\\ 0& \lambda _2& 0& \eta _2\end{array}\right]$$
where $`0<\alpha _1<\alpha _2`$, $`0<\beta _1<\beta _2`$, $`0<\lambda _1<\lambda _2`$, $`0<\eta _1<\eta _2`$ and $``$ denote a real positive number. Let $`(s_2,s_4)`$ be the Perron eigenvalue of the stochastic matrix:
$$\left[\begin{array}{cc}\alpha _2& \beta _2\\ \lambda _2& \eta _2\end{array}\right].$$
Then $`\stackrel{}{s}=(0,s_2,0,s_4)`$ is the stationary vector of $`\mathrm{\Delta }`$ corresponding to the simple eigenvalue $`1`$. Let $`w_0=(0,1,0,0)`$ and $`w_1=(0,0,0,1)`$. One checks that for $`n0`$, $`f_{a_n}f_{a_{n1}}\mathrm{}f_{a_0}(\stackrel{}{s})=w_{a_n}`$. Therefore $`L`$ consists of $`\{w_0,w_1\}`$. Using the expressions above, we see that
$$f_0^{}|_{w_0}=\alpha _1/\alpha _2<1,f_0^{}|_{w_1}=\lambda _1/\lambda _2<1,$$
$$f_1^{}|_{w_0}=\beta _1/\beta _2<1,f_1^{}|_{w_1}=\eta _1/\eta _2<1.$$
So, $`f_0`$ and $`f_1`$ are contraction mappings at $`\{w_0,w_1\}`$, and so condition 3 holds. Thus, the entropy rate $`H(Z)`$ is analytic at $`\widehat{\mathrm{\Delta }}`$.
## 7 Hidden Markov Chains with Unambiguous Symbol
###### Definition 7.1.
A symbol $`a`$ is called unambiguous if $`\mathrm{\Phi }^1(a)`$ contains only one element.
###### Remark 7.2.
Note that unambiguous symbol is referred to as “singleton clump” in some ergodic theory work, such as .
When an unambiguous symbol is present, the entropy rate can be expressed in a simple way: letting $`a_1`$ be an unambiguous symbol,
$$H(Z)=\underset{a_{i_j}a_1}{}p(a_{i_n}a_{i_{n1}}\mathrm{}a_{i_2}a_1)H(z|a_{i_n}a_{i_{n1}}\mathrm{}a_{i_2}a_1).$$
(7.21)
In this section, we focus on the case of a binary hidden Markov chain, in which 0 is unambiguous. Then, we can rewrite (7.21) as
$$H(Z^\stackrel{}{\epsilon })=p^\stackrel{}{\epsilon }(0)H^\stackrel{}{\epsilon }(z|0)+p^\stackrel{}{\epsilon }(10)H^\stackrel{}{\epsilon }(z|10)+\mathrm{}+p^\stackrel{}{\epsilon }(1^{(n)}0)H^\stackrel{}{\epsilon }(z|1^{(n)}0)+\mathrm{},$$
(7.22)
where $`1^{(n)}`$ denotes the sequence of $`n`$ 1’s and
$$H^\stackrel{}{\epsilon }(z|1^{(n)}0)=p^\stackrel{}{\epsilon }(0|1^{(n)}0)\mathrm{log}p^\stackrel{}{\epsilon }(0|1^{(n)}0)p^\stackrel{}{\epsilon }(1|1^{(n)}0)\mathrm{log}p^\stackrel{}{\epsilon }(1|1^{(n)}0).$$
###### Example 7.3.
Fix $`a,b,\mathrm{},h>0`$ and for $`\epsilon 0`$ let
$$\mathrm{\Delta }(\epsilon )=\left[\begin{array}{ccc}\epsilon & a\epsilon & b\\ g& c& d\\ h& e& f\end{array}\right].$$
Assume $`a,b,\mathrm{},h>0`$ are chosen such that $`\mathrm{\Delta }(\epsilon )`$ is stochastic. The symbols of the Markov chain are the matrix indices $`\{1,2,3\}`$. Let $`Z^\epsilon `$ be the binary hidden Markov chain defined by: $`\mathrm{\Phi }(1)=0`$ and $`\mathrm{\Phi }(2)=\mathrm{\Phi }(3)=1`$. We claim that $`H(Z^\epsilon )`$ is not analytic at $`\epsilon =0`$.
Let $`\pi (\epsilon )`$ be the stationary vector of $`\mathrm{\Delta }(\epsilon )`$ (which is unique since $`\mathrm{\Delta }(\epsilon )`$ is irreducible). Observe that
$$p^\epsilon (0)=\pi _1(\epsilon ),p^\epsilon (00)=\pi _1(\epsilon )\epsilon ,$$
and for $`n1`$.
$$p^\epsilon (1^{(n)}0)=\pi _1(\epsilon )(a\epsilon ,b)\left[\begin{array}{cc}c& d\\ e& f\end{array}\right]^{n1}\left(\begin{array}{c}1\\ 1\end{array}\right).$$
Since $`\mathrm{\Delta }(\epsilon )`$ is irreducible, $`\pi (\epsilon )`$ is analytic in $`\epsilon `$ and positive. Now,
$$p^\epsilon (0)H^\epsilon (z|0)=p^\epsilon (00)\mathrm{log}p^\epsilon (0|0)p^\epsilon (10)\mathrm{log}p^\epsilon (1|0).$$
(7.23)
The first term in (7.23) is
$$p^\epsilon (00)\mathrm{log}p^\epsilon (0|0)=\pi _1(\epsilon )\epsilon \mathrm{log}\epsilon ,$$
which is not analytic (or even differentiable at $`\epsilon =0`$). The second term in (7.23) is
$$p^\epsilon (10)\mathrm{log}p^\epsilon (1|0)=\pi _1(\epsilon )(a\epsilon +b)\mathrm{log}(\pi _1(\epsilon )(a\epsilon +b)),$$
which is analytic at $`\epsilon =0`$. Thus, $`H^\epsilon (z|0)`$ is not analytic at $`\epsilon =0`$. Similarly it can be shown that all of the terms of (7.22), other than $`H^\epsilon (z|0)`$, are analytic at $`\epsilon =0`$. Since the matrix
$$\left[\begin{array}{cc}c& d\\ e& f\end{array}\right]$$
has spectral radius $`<1`$, the terms of (7.22) decay exponentially; it follows that the infinite sum of these terms is analytic. Thus, $`H(Z^\epsilon )`$ is the sum of two functions of $`\epsilon `$, one of which is analytic and the other is not analytic at $`\epsilon =0`$. Thus, $`H(Z^\epsilon )`$ is not analytic at $`\epsilon =0`$.
###### Example 7.4.
Fix $`a,b,\mathrm{},g>0`$ and consider the stochastic matrix
$$\mathrm{\Delta }(\epsilon )=\left[\begin{array}{ccc}e& a& b\\ f\epsilon & c& \epsilon \\ g& 0& d\end{array}\right].$$
The symbols of the Markov chain are the matrix indices $`\{1,2,3\}`$. Again let $`Z^\epsilon `$ be the binary hidden Markov chain defined by $`\mathrm{\Phi }(1)=0`$ and $`\mathrm{\Phi }(2)=\mathrm{\Phi }(3)=1`$. We show that $`H(Z^\epsilon )`$ is analytic at $`\epsilon =0`$ when $`cd`$, and not analytic when $`c=d`$. Note that
$$p^\epsilon (0)=\pi _1(\epsilon ),$$
and for $`n1`$.
$$p^\epsilon (1^{(n)}0)=\pi _1(\epsilon )(a,b)\left[\begin{array}{cc}c& \epsilon \\ 0& d\end{array}\right]^{n1}\left(\begin{array}{c}1\\ 1\end{array}\right).$$
When $`cd`$, we assume $`c>d`$, then
$$\left[\begin{array}{cc}c& \epsilon \\ 0& d\end{array}\right]^n=\left[\begin{array}{cc}c^n& \epsilon c^{n1}\frac{1(d/c)^n}{1d/c}\\ 0& d^n\end{array}\right].$$
Since $`\mathrm{\Delta }(\epsilon )`$ is irreducible, $`\pi (\epsilon )`$ is analytic in $`\epsilon `$ and positive. Simple computation leads to:
$$p^\epsilon (1|1^{(n)}0)=(ac^n+a\epsilon c^{n1}\frac{1(d/c)^n}{1d/c}+bd^n)/(ac^{n1}+a\epsilon c^{n2}\frac{1(d/c)^{n1}}{1d/c}+bd^{n1})$$
$$=(ac^2+a\epsilon c\frac{1(d/c)^n}{1d/c}+bd^2(d/c)^{n2})/(ac+\epsilon \frac{1(d/c)^{n1}}{1d/c}+bd(d/c)^{n2}),$$
and
$$p^\epsilon (0|1^{(n)}0)=((f\epsilon )ac^{n1}+g(a\epsilon c^{n2}\frac{1(d/c)^{n1}}{1d/c}+bd^{n1}))/(ac^{n1}+a\epsilon c^{n2}\frac{1(d/c)^{n1}}{1d/c}+bd^{n1})$$
$$=((f\epsilon )ac+g(a\epsilon \frac{1(d/c)^{n1}}{1d/c}+bd(d/c)^{n2}))/(ac+\epsilon \frac{1(d/c)^{n1}}{1d/c}+bd(d/c)^{n2}).$$
In this case all terms are analytic. Again since
$$\left[\begin{array}{cc}c& \epsilon \\ 0& d\end{array}\right]$$
has spectral radius $`<1`$, the term $`p^\epsilon (1^{(n)}0)H^\epsilon (z|1^{(n)}0)`$ is exponentially decaying with respect to $`n`$. Therefore the infinite sum of these terms is also analytic, and so the entropy rate is a real analytic function of $`\epsilon `$.
When $`c=d`$, we have
$$p^\epsilon (1|1^{(n)}0)=(ac^{n+1}+a\epsilon (n+1)c^n+bc^{n+1})/(ac^n+a\epsilon nc^{n1}+bc^n)$$
$$=(ac^2+a\epsilon (n+1)c+bc^2)/(ac+a\epsilon n+bc),$$
and
$$p^\epsilon (0|1^{(n)}0)=((f\epsilon )ac^n+ga\epsilon nc^{n1}+gbc^n)/(ac^n+a\epsilon nc^{n1}+bc^n)$$
$$=((f\epsilon )ac+ga\epsilon n+gbc)/(ac+a\epsilon n+bc).$$
For any $`n`$, consider a small neighborhood $`N_n`$ of $`(a+b)c/an`$ in $``$ such that $`(a+b)c/ajN_n`$ only holds for $`j=n`$. When $`\epsilon (a+b)c/an`$, the complexified term $`p^\epsilon (1^{(n)}0)H^\epsilon (z|1^{(n)}0)\mathrm{}`$. Meanwhile, the sum of all the other terms can be analytically extended to $`N_n`$ (from any path $`I`$ from a positive $`\epsilon `$ to $`(a+b)c/an`$ with $`(a+b)c/ajI`$ for $`jn`$). Thus, by the uniqueness of analytic continuation of $`H(Z^\epsilon )`$, we conclude that $`H(Z^\epsilon )`$ blows up when one approaches $`(a+b)c/an`$ and therefore is not analytic at $`\epsilon =0`$ (although it is smooth from the right at $`\epsilon =0`$).
The two examples above show that under certain conditions the entropy rate of a binary hidden Markov chain with unambiguous symbol can fail to be analytic at the boundary. We now show that these examples typify all the types of failures of analyticity at the boundary (in the case of a binary hidden Markov chains with an unambiguous symbol).
We will need the following result.
###### Lemma 7.5.
Let $`A(\stackrel{}{\epsilon })`$ be an analytic parameterization of complex matrices. Let $`\lambda `$ be the spectral radius of $`A(\stackrel{}{\epsilon }_0)`$. Then for any $`\eta >0`$, there exists a complex neighborhood $`\mathrm{\Omega }`$ of $`\stackrel{}{\epsilon }_0`$ and positive constant $`C`$ such that for all $`\stackrel{}{\epsilon }\mathrm{\Omega }`$ and all $`i,j,k`$
$$|A_{ij}^k(\stackrel{}{\epsilon })|C(\lambda +\eta )^k.$$
###### Proof.
Following , we consider
$$(IzA)^1=I+zA+z^2A^2+\mathrm{}.$$
And
$$(IzA)^1=\frac{Adj(IzA)}{det(IzA)}=\frac{Adj(IzA)}{(1\lambda _1z)(1\lambda _2z)\mathrm{}(1\lambda _nz)},$$
where $`\lambda _1,\mathrm{},\lambda _n`$ are the eigenvalues of $`A`$. So every entry of $`(IzA)^1`$ takes the form:
$$(p_0+p_1z+\mathrm{}+p_mz^m)\underset{j=1}{\overset{n}{}}\underset{i=0}{\overset{\mathrm{}}{}}\lambda _j^iz^i$$
$$=\underset{k=0}{\overset{\mathrm{}}{}}\underset{u=0}{\overset{m}{}}p_u\underset{i_1+i_2+\mathrm{}+i_n=ku}{}\lambda _1^{i_1}\lambda _2^{i_2}\mathrm{}\lambda _n^{i_n}z^k.$$
Since the eigenvalues of a complex matrix vary continuously with entries, the lemma follows. ∎
Now let $`S(n)`$ denote the set of all the $`n\times n`$ complex matrices with isolated (in modulus) maximum eigenvalue.
###### Lemma 7.6.
$`S(n)`$ is connected.
###### Proof.
let $`A,BS(n)`$, then we consider their Jordan forms:
$$A=U\mathrm{diag}(\lambda _1,C)U^1,B=V\mathrm{diag}(\eta _1,D)V^1,$$
here $`\lambda _1,\eta _1`$ are maximum eigenvalues for $`A,B`$, respectively, $`C,D`$ correspond to other Jordan blocks, and $`U,VGL(n,)`$ (here $`GL(n,)`$ denotes the set of all the $`n\times n`$ nonsingular complex matrices). Since $`GL(n,)`$ is connected , it suffices to prove that there is a path in $`S(n)`$ from $`\mathrm{diag}(\lambda _1,C)`$ to $`\mathrm{diag}(\eta _1,D)`$. This is straightforward: first connect $`\mathrm{diag}(\lambda _1,C)`$ to $`\mathrm{diag}(\eta _1,\eta _1/\lambda _1C)`$ by a continuous rescaling; then connect $`\eta _1/\lambda _1C`$ to $`D`$ by the path $`t\eta _1/\lambda _1C+(1t)D`$ (the path $`\mathrm{diag}(\eta _1,t\eta _1/\lambda _1C+(1t)D)`$ stays within $`S(n)`$ since the matrices along this path are upper triangular with all diagonal entries, except $`\eta _1`$, of modulus less than $`|\eta _1|`$). ∎
For a complex analytic function $`f(z_1,z_2,\mathrm{},z_n)`$, let $`V(f)`$ denote the “hypersurface” defined by $`f`$, namely
$$V(f)=\{(z_1,z_2,\mathrm{},z_n)^n:f(z_1,z_2,\mathrm{},z_n)=0\}.$$
Now let $`\mathrm{\Omega }`$ denote a connected open set in $`^n`$. It is well known that the following Lemma holds (for completeness, we include a brief proof).
###### Lemma 7.7.
$`\mathrm{\Omega }\backslash V(f)`$ is connected.
###### Proof.
For simplicity, we first assume $`\mathrm{\Omega }`$ is a ball $`B_r(z_0)`$ (here $`z_0^n`$ is the center of the ball and $`r`$ is the radius, i.e., $`B_r(z_0)=\{z^n:|zz_0|<r\}`$) in $`^n`$. For any two distinct point $`P,Q\mathrm{\Omega }\backslash V(f)`$, consider the “complex line”
$$L_{}^{PQ}=\{zP+(1z)Q:z\}.$$
$`L_{}^{PQ}V(f)\mathrm{\Omega }`$ consists of only isolated points (A non-constant one variable complex analytic function must have isolated zeros in the complex plane ). It then follows that for the compact real line segment:
$$L_{}^{PQ}=\{tP+(1t)Q:t[0,1]\},$$
$`L_{}^{PQ}V(f)\mathrm{\Omega }`$ consists of only finitely many points. Certainly one can choose an arc in $`L_{}^{PQ}\mathrm{\Omega }`$ to avoid these points and connect $`P`$ and $`Q`$. This implies that $`\mathrm{\Omega }\backslash V(f)`$ is connected.
In the general case, $`\mathrm{\Omega }`$ is a connected open set in $`^n`$. Let $`I`$ be an arc in $`\mathrm{\Omega }`$ connecting $`P`$ and $`Q`$, and let $`\{B_{r_j}(z_j)\}`$ be a collection of balls covering $`I`$ such that each $`B_{r_j}(z_j)B_{r_{j+1}}(z_{j+1})\varphi `$. Pick a point $`P_j`$ in $`B_{r_j}(z_j)B_{r_{j+1}}(z_{j+1})`$ such that $`P_j\mathrm{\Omega }\backslash V(f)`$. Applying the same argument as above to every ball $`B_{r_j}(z_j)`$, we see that $`P`$ is connected to $`Q`$ in $`\mathrm{\Omega }\backslash V(f)`$ through the points $`P_j`$’s. Thus we prove the lemma. ∎
###### Theorem 7.8.
Let $`\mathrm{\Delta }`$ be an irreducible stochastic $`d\times d`$ matrix. Write $`\mathrm{\Delta }`$ in the form:
$$\mathrm{\Delta }=\left[\begin{array}{cc}a& r\\ c& B\end{array}\right]$$
(7.24)
where $`a`$ is a scalar and $`B`$ is a $`(d1)\times (d1)`$ matrix. Let $`\mathrm{\Phi }`$ be the function defined by $`\mathrm{\Phi }(1)=0`$, and $`\mathrm{\Phi }(2)=\mathrm{}=\mathrm{\Phi }(d)=1`$. Then for any parametrization $`\mathrm{\Delta }(\stackrel{}{\epsilon })`$ such that $`\mathrm{\Delta }(\stackrel{}{\epsilon }_0)=\mathrm{\Delta }`$, letting $`Z^\stackrel{}{\epsilon }`$ denote the hidden Markov chain defined by $`\mathrm{\Delta }(\stackrel{}{\epsilon })`$ and $`\mathrm{\Phi }`$, $`H(Z^\stackrel{}{\epsilon })`$ is analytic at $`\stackrel{}{\epsilon }_0`$ if and only if
1. $`a>0`$, and $`rB^jc>0`$ for $`j=0,1,\mathrm{}`$.
2. The maximum eigenvalue of $`B`$ is simple and strictly greater in absolute value than the other eigenvalues of $`B`$.
###### Proof.
Proof of sufficiency.
We write
$$\mathrm{\Delta }(\stackrel{}{\epsilon })=\left[\begin{array}{cc}a(\stackrel{}{\epsilon })& r(\stackrel{}{\epsilon })\\ c(\stackrel{}{\epsilon })& B(\stackrel{}{\epsilon })\end{array}\right],$$
(7.25)
where $`a(\stackrel{}{\epsilon })`$ is a scalar and $`B(\stackrel{}{\epsilon })`$ is a $`(d1)\times (d1)`$ matrix.
Since $`\mathrm{\Delta }(\stackrel{}{\epsilon }_0)`$ is stochastic and irreducible, its spectral radius is 1, and 1 is a simple eigenvalue of $`\mathrm{\Delta }`$. Thus, if $`\mathrm{\Omega }`$ is sufficiently small, for all $`\stackrel{}{\epsilon }\mathrm{\Omega }`$, any fixed row $`\pi (\stackrel{}{\epsilon })=(\pi _1(\stackrel{}{\epsilon }),\pi _2(\stackrel{}{\epsilon }),\mathrm{},\pi _d(\stackrel{}{\epsilon }))`$ of $`Adj(I\mathrm{\Delta }(\stackrel{}{\epsilon }))`$ is a left eigenvector of $`\mathrm{\Delta }(\stackrel{}{\epsilon })`$ associated with eigenvalue $`1`$ and is an analytic function of $`\stackrel{}{\epsilon }`$. Normalizing, we can assume that $`\pi (\stackrel{}{\epsilon })\mathrm{𝟏}=\mathrm{𝟏}`$, $`\pi (\stackrel{}{\epsilon })`$ is analytic in $`\stackrel{}{\epsilon }`$, and $`\pi (\stackrel{}{\epsilon }_0)>0`$.
The entries of $`r(\stackrel{}{\epsilon }),B(\stackrel{}{\epsilon }),`$ and $`c(\stackrel{}{\epsilon })`$ are real analytic in $`\stackrel{}{\epsilon }`$ and can be extended to complex analytic functions in a complex neighborhood $`\mathrm{\Omega }`$ of $`\stackrel{}{\epsilon }_0`$. Thus, for all $`n`$, $`\pi _1(\stackrel{}{\epsilon })r(\stackrel{}{\epsilon })B(\stackrel{}{\epsilon })^{n1}\mathrm{𝟏}`$ and $`\pi _1(\stackrel{}{\epsilon })r(\stackrel{}{\epsilon })B(\stackrel{}{\epsilon })^{n1}c(\stackrel{}{\epsilon })`$ can be extended to complex analytic functions on $`\mathrm{\Omega }`$ (in fact, each of these functions is a polynomial in $`\stackrel{}{\epsilon }`$).
Since $`B(\stackrel{}{\epsilon }_0)`$ is a proper sub-matrix of the irreducible stochastic matrix $`\mathrm{\Delta }(\stackrel{}{\epsilon }_0)`$, its spectral radius is strictly less than 1. Thus, by Lemma 7.5, there exists $`0<\lambda ^{}<1`$ and a constant $`C_1>0`$, such that for some complex neighborhood $`\mathrm{\Omega }`$ of $`\stackrel{}{\epsilon }_0`$, all $`\stackrel{}{\epsilon }\mathrm{\Omega }`$, and all n,
$$|B_{ij}^n(\stackrel{}{\epsilon })|<C_1(\lambda ^{})^n.$$
Since $`\pi _1(\stackrel{}{\epsilon })`$, $`r(\stackrel{}{\epsilon })`$ and $`c(\stackrel{}{\epsilon })`$ are continuous in $`\stackrel{}{\epsilon }`$, there is a constant $`C_2>0`$ such that for all $`\stackrel{}{\epsilon }\mathrm{\Omega }`$ and all $`n`$:
$$|\pi _1(\stackrel{}{\epsilon })r(\stackrel{}{\epsilon })B(\stackrel{}{\epsilon })^n\mathrm{𝟏}|<C_2(\lambda ^{})^n.$$
(7.26)
We will need the following result, proven in Appendix B.
###### Lemma 7.9.
Let
$$a(\stackrel{}{\epsilon },n)\frac{\pi _1(\stackrel{}{\epsilon })r(\stackrel{}{\epsilon })B(\stackrel{}{\epsilon })^n\mathrm{𝟏}}{\pi _1(\stackrel{}{\epsilon })r(\stackrel{}{\epsilon })B(\stackrel{}{\epsilon })^{n1}\mathrm{𝟏}}$$
and
$$b(\stackrel{}{\epsilon },n)\frac{\pi _1(\stackrel{}{\epsilon })r(\stackrel{}{\epsilon })B(\stackrel{}{\epsilon })^{n1}c(\stackrel{}{\epsilon })}{\pi _1(\stackrel{}{\epsilon })r(\stackrel{}{\epsilon })B(\stackrel{}{\epsilon })^{n1}\mathrm{𝟏}}.$$
For a sufficiently small neighborhood $`\mathrm{\Omega }`$ of $`\stackrel{}{\epsilon }_0`$, both $`a(\stackrel{}{\epsilon },n)`$ and $`b(\stackrel{}{\epsilon },n)`$ are bounded from above and away from zero, uniformly in $`\stackrel{}{\epsilon }\mathrm{\Omega }`$ and $`n`$.
Define
$$H_n^\stackrel{}{\epsilon }=a(\stackrel{}{\epsilon },n)\mathrm{log}a(\stackrel{}{\epsilon },n)b(\stackrel{}{\epsilon },n)\mathrm{log}b(\stackrel{}{\epsilon },n),$$
where $`a(\stackrel{}{\epsilon },n)`$ and $`b(\stackrel{}{\epsilon },n)`$ are as in Lemma 7.9. Choosing $`\mathrm{\Omega }`$ to be a smaller neighborhood of $`\stackrel{}{\epsilon }_0`$, if necessary, $`a(\stackrel{}{\epsilon },n)`$ and $`b(\stackrel{}{\epsilon },n)`$ are constrained to lie in a closed disk not containing $`0`$. Thus for all $`n`$, $`H_n^\stackrel{}{\epsilon }`$ is an analytic function of $`\stackrel{}{\epsilon }`$, with $`|H_n^\stackrel{}{\epsilon }|`$ bounded uniformly in $`\stackrel{}{\epsilon }\mathrm{\Omega }`$ and $`n`$. Since $`\pi _1(\stackrel{}{\epsilon })r(\stackrel{}{\epsilon })B(\stackrel{}{\epsilon })^{n1}\mathrm{𝟏}`$ is analytic on $`\mathrm{\Omega }`$ and exponentially decaying (by (7.26)), the infinite series
$$H^\stackrel{}{\epsilon }(Z)=\pi _1(\stackrel{}{\epsilon })H_0^\stackrel{}{\epsilon }+\pi _1(\stackrel{}{\epsilon })r(\stackrel{}{\epsilon })\mathrm{𝟏}H_1^\stackrel{}{\epsilon }+\mathrm{}+\pi _1(\stackrel{}{\epsilon })r(\stackrel{}{\epsilon })B(\stackrel{}{\epsilon })^{n1}\mathrm{𝟏}H_n^\stackrel{}{\epsilon }+\mathrm{}$$
(7.27)
converges uniformly on $`\mathrm{\Omega }`$ and thus defines an analytic function on $`\mathrm{\Omega }`$.
Note that for $`\stackrel{}{\epsilon }0`$,
$$p^\stackrel{}{\epsilon }(1^{(n)}0)=\pi _1(\stackrel{}{\epsilon })r(\stackrel{}{\epsilon })B(\stackrel{}{\epsilon })^{n1}\mathrm{𝟏}$$
(7.28)
and
$$p^\stackrel{}{\epsilon }(01^{(n)}0)=\pi _1(\stackrel{}{\epsilon })r(\stackrel{}{\epsilon })B(\stackrel{}{\epsilon })^{n1}c(\stackrel{}{\epsilon }).$$
(7.29)
By (7.28), (7.29), and the expression for entropy rate in the case of an unambiguous symbol (given at the beginning of this section), $`H^\stackrel{}{\epsilon }(Z)`$ agrees with the entropy rate when $`\mathrm{\Delta }(\stackrel{}{\epsilon })0`$, as desired.
###### Remark 7.10.
We show how sufficiency relates to Theorem 6.1. Namely, the assumptions in Theorem 7.8 imply those of Theorem 6.1. Condition $`1`$ of Theorem 6.1 follows from the fact that $`\mathrm{\Delta }`$ is assumed irreducible. For conditions 2 and 3 of Theorem 6.1, one first notes that the image of $`f_0`$ is a single point $`W_0`$, and the $`f_1`$-orbit of $`W_0`$ and $`f_1`$-orbit of $`\stackrel{}{s}`$ converge to a point $`p_1`$. It follows that $`L`$ is the union of $`W_0`$, the $`f_1`$-orbit of $`W_0`$ and $`p_1`$. The assumptions in Theorem 7.8. imply that $`r_a>0`$ on $`L`$ (i.e., condition 2 of Theorem 6.1 holds) and that for sufficiently large $`n`$, the $`n`$-fold composition of $`f_1`$ is contracting on the convex hull of the intersection of $`L`$ and $`W_1`$ (so condition 3 of Theorem 6.1 holds). To see the latter, one uses the ideas in the proof of sufficiency.
Proof of necessity
We first consider condition $`2`$. We shall use the natural parameterization and view $`H(Z)`$ as a function of $`\mathrm{\Delta }`$, or more precisely of $`(B,r)`$. Note that there is a one-to-one correspondence between $`\mathrm{\Delta }`$ and $`(B,r)`$; we shall use this correspondence throughout the proof.
Suppose $`\mathrm{\Delta }`$ doesn’t satisfy condition $`2`$, however $`H(Z)`$ is analytic at $`\mathrm{\Delta }`$ with respect to the natural parameterization. In other words, suppose there exists a complex neighborhood $`N_\mathrm{\Delta }`$ of $`\mathrm{\Delta }`$ (here $`N_\mathrm{\Delta }`$ corresponds to $`N_B\times N_r`$ where $`N_B`$ is neighborhood of $`B`$ and $`N_r`$ is neighborhood of $`r`$) such that $`H(Z)`$ can be analytically extended to $`N_\mathrm{\Delta }`$, while the corresponding $`B`$ doesn’t have isolated (in modulus) maximum eigenvalue.
We first claim there exists $`\stackrel{~}{\mathrm{\Delta }}N_\mathrm{\Delta }`$ with $`\stackrel{~}{r}\stackrel{~}{B}^k\mathrm{𝟏}=0`$, here $`\stackrel{~}{r}`$ and $`\stackrel{~}{B}`$ correspond to $`\stackrel{~}{\mathrm{\Delta }}`$ and $`\stackrel{~}{B}`$ has distinct eigenvalues (in modulus). Indeed we can first (for simplicity) perturb $`\mathrm{\Delta }`$ to $`\stackrel{~}{\mathrm{\Delta }}`$ such that the corresponding $`\stackrel{~}{B}`$ has distinct eigenvalues in modulus. Then
$$\stackrel{~}{B}=\stackrel{~}{U}\mathrm{diag}(\stackrel{~}{\lambda }_1,\stackrel{~}{\lambda }_2,\mathrm{},\stackrel{~}{\lambda }_{d1})\stackrel{~}{U}^1$$
$$=(\stackrel{~}{v}_1,\stackrel{~}{v}_2,\mathrm{},\stackrel{~}{v}_{d1})\mathrm{diag}(\stackrel{~}{\lambda }_1,\stackrel{~}{\lambda }_2,\mathrm{},\stackrel{~}{\lambda }_{d1})(\stackrel{~}{w}_1^t,\stackrel{~}{w}_2^t,\mathrm{},\stackrel{~}{w}_{d1}^t)^t$$
where $`|\stackrel{~}{\lambda }_1|>|\stackrel{~}{\lambda }_2|>\mathrm{}>|\stackrel{~}{\lambda }_{d1}|`$, and $`\stackrel{~}{v}_i,\stackrel{~}{w}_i`$’s are appropriately scaled right and left eigenvectors of $`\stackrel{~}{B}`$, respectively. Then we have
$$r\stackrel{~}{B}^k\mathrm{𝟏}=r\stackrel{~}{v}_1\stackrel{~}{w}_1\mathrm{𝟏}\stackrel{~}{\lambda }_1^k+r\stackrel{~}{v}_2\stackrel{~}{w}_2\mathrm{𝟏}\stackrel{~}{\lambda }_2^k+\mathrm{}+r\stackrel{~}{v}_{d1}\stackrel{~}{w}_{d1}\mathrm{𝟏}\stackrel{~}{\lambda }_{d1}^k.$$
Further consider a perturbation of $`B`$ from
$$\stackrel{~}{B}=\stackrel{~}{U}\mathrm{diag}(\stackrel{~}{\lambda }_1,\stackrel{~}{\lambda }_2,\mathrm{},\stackrel{~}{\lambda }_{d1})\stackrel{~}{U}^1$$
to
$$\stackrel{~}{B}=V\stackrel{~}{U}\mathrm{diag}(\stackrel{~}{\lambda }_1,\stackrel{~}{\lambda }_2,\mathrm{},\stackrel{~}{\lambda }_{d1})\stackrel{~}{U}^1V^1,$$
where $`V`$ is a complex matrix close to the $`(d1)\times (d1)`$ identity matrix $`I_{d1}`$. So we can pick $`V`$ such that $`\stackrel{~}{v}_1\stackrel{~}{w}_1V^1\mathrm{𝟏}0`$, $`\stackrel{~}{v}_1\stackrel{~}{w}_1V^1\stackrel{~}{c}0`$, $`\stackrel{~}{v}_2\stackrel{~}{w}_2V^1\mathrm{𝟏}0`$. Clearly $`\stackrel{~}{v}_1\stackrel{~}{w}_1V^1\mathrm{𝟏}`$ is not proportional to $`\stackrel{~}{v}_2\stackrel{~}{w}_2V^1\mathrm{𝟏}`$. Then by a further perturbation of $`r`$ to $`\stackrel{~}{r}`$, we can simultaneously require that $`\stackrel{~}{r}\stackrel{~}{v}_1\stackrel{~}{w}_1\mathrm{𝟏}0`$, $`\stackrel{~}{r}\stackrel{~}{v}_1\stackrel{~}{w}_1\stackrel{~}{c}0`$, $`\stackrel{~}{r}\stackrel{~}{v}_2\stackrel{~}{w}_2\mathrm{𝟏}0`$, $`|\stackrel{~}{r}\stackrel{~}{v}_1\stackrel{~}{w}_1\mathrm{𝟏}||\stackrel{~}{r}\stackrel{~}{v}_2\stackrel{~}{w}_2\mathrm{𝟏}|`$, where we redefine $`\stackrel{~}{v}_i=V\stackrel{~}{v}_i`$ and $`\stackrel{~}{w}_i=\stackrel{~}{w}_iV^1`$. For any $`\theta `$ and $`\eta >0`$, it can be checked that
$$\underset{k=0}{\overset{\mathrm{}}{}}\{z^k:|ze^{i\theta }|<\eta \}=\backslash \{0\}.$$
Since $`\stackrel{~}{\lambda }_2`$ is a perturbation of $`\stackrel{~}{\lambda }_1`$, it follows that for large enough $`k`$, one can perturb $`\stackrel{~}{\lambda }_2`$ to satisfy the equation
$$\left(\stackrel{~}{\lambda }_2/\stackrel{~}{\lambda }_1\right)^k=\frac{\stackrel{~}{r}\stackrel{~}{v}_1\stackrel{~}{w}_1\mathrm{𝟏}\stackrel{~}{r}\stackrel{~}{v}_3\stackrel{~}{w}_3\mathrm{𝟏}(\stackrel{~}{\lambda }_3/\stackrel{~}{\lambda }_1)^k\mathrm{}\stackrel{~}{r}\stackrel{~}{v}_{d1}\stackrel{~}{w}_{d1}\mathrm{𝟏}(\stackrel{~}{\lambda }_{d1}/\stackrel{~}{\lambda }_1)^k}{\stackrel{~}{r}\stackrel{~}{v}_2\stackrel{~}{w}_2\mathrm{𝟏}},$$
with $`|\stackrel{~}{\lambda }_2||\stackrel{~}{\lambda }_1|`$ and $`|\stackrel{~}{\lambda }_2|`$ strictly greater than $`|\stackrel{~}{\lambda }_j|`$ for $`j3`$. Thus we prove the claim.
We now pick a positive matrix $`\widehat{\mathrm{\Delta }}N_\mathrm{\Delta }`$ with corresponding $`\widehat{r}`$ and $`\widehat{B}`$. We then pick $`\stackrel{~}{\mathrm{\Delta }}N_\mathrm{\Delta }`$ with corresponding $`\stackrel{~}{r}`$ and $`\stackrel{~}{B}`$ (with distinct eigenvalues in modulus) such that $`\stackrel{~}{r}\stackrel{~}{B}^{k_1}\mathrm{𝟏}=0`$ for some $`k_1`$, and we can further require that $`\stackrel{~}{r}\stackrel{~}{v}_1\stackrel{~}{w}_1\mathrm{𝟏}0`$, $`\stackrel{~}{r}\stackrel{~}{v}_1\stackrel{~}{w}_1\stackrel{~}{c}0`$ (see the proof for the previous claim), where as before, $`\stackrel{~}{v}_1,\stackrel{~}{w}_1`$ are eigenvectors corresponding to the largest eigenvalue of $`\stackrel{~}{B}`$. According to Lemma 7.6, there is an arc $`I_1S(d1)`$ connecting $`\widehat{B}`$ to $`\stackrel{~}{B}`$; we then connect $`\widehat{r}`$ and $`\stackrel{~}{r}`$ using an arc $`I_2`$ in $`^{d1}`$. According to Lemma 7.7, we can choose the arc $`I=(I_1,I_2)`$ to avoid the hypersurface $`V((rv_1w_1\mathrm{𝟏})(rv_1w_1c))^{(d1)^2}\times ^{d1}`$; in other words, we can assume that along the path $`I`$, $`rv_1w_1\mathrm{𝟏}0`$ and $`rv_1w_1c0`$; here $`v_1,w_1,c`$ are determined by the variable matrix $`B`$ along the path $`I_1`$ and $`r`$ is the variable point along path $`I_2`$ (we remind the reader that the coordinates of $`v_1`$ and $`w_1`$ are all analytic functions of the entries of $`B`$). We then claim that there is a neighborhood $`N_I`$ of $`I`$ such that $`V_kN_I\varphi `$ and $`W_kN_I\varphi `$ hold for only finitely many $`k`$, where $`V_k=\{(B,r):rB^k\mathrm{𝟏}=0\}`$ and $`W_k=\{(B,r):rB^kc=0\}`$. Indeed for any $`\mathrm{\Delta }I`$ with corresponding $`BS(d1)`$, by the Jordan form we have
$$rB^k\mathrm{𝟏}=rv_1w_1\mathrm{𝟏}\lambda _1^k+o(\lambda _1^k),$$
where $`\lambda _1`$ is the isolated maximum eigenvalue and $`v_1,w_1`$ are appropriately scaled right and left eigenvectors of $`B`$, respectively. Since $`rv_1w_1\mathrm{𝟏}0`$ on $`I`$, there exists a complex connected neighborhood $`N_I`$ of $`I`$ such that $`rv_1w_1\mathrm{𝟏}0`$ on $`N_I`$ and $`rv_1w_1\mathrm{𝟏}\lambda _1^k`$ dominates uniformly on $`N_I`$ (see Lemma 7.5). Consequently, $`|rB^k\mathrm{𝟏}|>0`$ on $`N_I`$ for large enough $`k`$. In other words, $`V_kN_I\varphi `$ holds for only finitely many $`k`$. Similarly since $`rv_1w_1c0`$ on $`I`$, there exists a complex neighborhood $`N_I`$ of $`I`$ (here we use the same notation for a possibly different neighborhood) such that $`W_kN_I\varphi `$ holds only for finitely many $`k`$. From now on, we assume such $`k`$’s are less than some $`K`$, which depends on $`N_I`$.
We claim that we can further choose $`I`$ and find a new neighborhood $`N_I`$ in $`^{d1}\times S(d1)`$ of $`I`$ such that $`V_kN_I\varphi `$ holds only for $`k=k_1`$ and $`W_kN_I=\varphi `$ for all $`k`$. Consider $`\stackrel{~}{\mathrm{\Delta }}`$ with corresponding $`\stackrel{~}{B}`$, let $`F_i=F_i(\stackrel{~}{B})=\{r:r\stackrel{~}{B}^i\mathrm{𝟏}=0\}`$, which is a hyperplane orthogonal to the vector $`\stackrel{~}{B}^i\mathrm{𝟏}`$ in $`^{d1}`$. Similarly we define $`G_i=G_i(\stackrel{~}{B})=\{r:r\stackrel{~}{B}^i\stackrel{~}{c}=0\}`$. Recall that $`\stackrel{~}{B}=\stackrel{~}{U}\mathrm{diag}(\lambda _1,\lambda _2,\mathrm{},\lambda _{d1})\stackrel{~}{U}^1`$; we can require that $`\stackrel{~}{U}^1\mathrm{𝟏}`$ has no zero coordinates by a small perturbation of $`\stackrel{~}{U}`$ if necessary. We then show that $`F_i`$’s and $`G_j`$’s define different hyperplanes in $`^{d1}`$. Indeed suppose $`F_i=F_j`$. It follows that $`\stackrel{~}{U}\mathrm{diag}(\stackrel{~}{\lambda }_1^i,\stackrel{~}{\lambda }_2^i,\mathrm{},\stackrel{~}{\lambda }_{d1}^i)\stackrel{~}{U}^1\mathrm{𝟏}`$ is proportional to $`\stackrel{~}{U}\mathrm{diag}(\stackrel{~}{\lambda }_1^j,\stackrel{~}{\lambda }_2^j,\mathrm{},\stackrel{~}{\lambda }_{d1}^j)\stackrel{~}{U}^1\mathrm{𝟏}`$. It then follows that $`(\stackrel{~}{\lambda }_1^i,\stackrel{~}{\lambda }_2^i,\mathrm{},\stackrel{~}{\lambda }_{d1}^i)`$ is proportional to $`(\stackrel{~}{\lambda }_1^j,\stackrel{~}{\lambda }_2^j,\mathrm{},\stackrel{~}{\lambda }_{d1}^j)`$. However since not all eigenvalues have the same modulus, this implies that $`i=j`$. With a perturbation of $`\stackrel{~}{c}`$ (equivalently a perturbation of row sums of $`\stackrel{~}{B}`$), if necessary, we conclude that the $`F_i`$’s and $`G_i`$’s determine different hyperplanes, i.e., $`F_iF_j`$, $`G_iG_j`$ for $`ijK`$, and $`F_iG_j`$ for all $`i,j`$. Thus, with a perturbation of $`\stackrel{~}{r}`$ if necessary, we can choose a new $`\stackrel{~}{\mathrm{\Delta }}`$ contained in $`V_{k_1}`$, but not contained in any $`V_k`$ with $`kk_1`$ or $`W_k`$ for all $`k`$. Again by Lemma 7.7, one can choose a new $`I`$ inside original $`N_I`$, connecting $`\widehat{\mathrm{\Delta }}`$ and $`\stackrel{~}{\mathrm{\Delta }}`$, to avoid all $`V_k`$’s and $`W_k`$’s except $`V_{k_1}`$, then choose a smaller new neighborhood $`N_I`$ of the new $`I`$ to make sure that $`V_kN_I\varphi `$ only holds for $`k=k_1`$ and $`W_kN_I=\varphi `$ for all $`k`$.
Since the perturbed complex matrix $`B`$ still has spectral radius strictly less than $`1`$, all the complexified terms in the entropy rate formula (see (7.27)) with $`kk_1`$ are exponentially decaying and thus sum up to an analytic function on $`N_I`$.(i.e., the sum of these terms can be analytically continued to $`N_I`$), while the unique analytic extension of the $`k_1`$-th term on $`N_I`$ blows up as one approaches $`V_{k_1}N_I`$ from $`\widehat{\mathrm{\Delta }}`$. Again by the uniqueness of analytic extension of $`H(Z)`$ on $`N_I`$, this would be a contradiction to the assumption that $`H(Z)`$ is analytic at $`\mathrm{\Delta }`$ (here we are applying the uniqueness theorem of analytic continuation of a function of several complex variables, see page $`21`$ in ). Thus we prove the necessity of condition $`2`$.
We now consider condition $`1`$. Suppose $`\mathrm{\Delta }`$ doesn’t satisfies condition $`1`$, namely $`a=0`$ or $`rB^kc=0`$ for some $`k`$, however $`H(Z)`$ is analytic at $`\mathrm{\Delta }`$. With the proof above for the necessity of condition $`2`$, we can now assume the corresponding $`BS(d1)`$.
If $`a=0`$, consider any perturbation of $`\mathrm{\Delta }`$ to $`\mathrm{\Delta }_1`$ such that $`\stackrel{~}{B}S(d1)`$, $`\stackrel{~}{r}\stackrel{~}{v}_1\stackrel{~}{w}_1\mathrm{𝟏}0`$, $`\stackrel{~}{r}\stackrel{~}{v}_1\stackrel{~}{w}_1\stackrel{~}{c}0`$, $`\stackrel{~}{r}\stackrel{~}{B}^k\mathrm{𝟏}0`$ and $`\stackrel{~}{r}\stackrel{~}{B}^k\stackrel{~}{c}0`$ for all $`k`$ (here we follow the notation as in the proof of necessity of condition $`2`$). Then using similar arguments, we can prove the sum of all the terms except the first term in the entropy rate formula (see (7.27)) can be analytically extended to $`\stackrel{~}{\mathrm{\Delta }}`$. However this implies that $`a\mathrm{log}a`$ is a well-defined analytic function on some neighborhood of $`0`$ in $``$, which is a contradiction. Similar arguments can be applied to the case that $`rB^kc=0`$ for some $`k`$’s. Thus we prove the necessity of condition $`1`$. ∎
## 8 Analyticity of a Hidden Markov Chain in a Strong Sense
In this section, we show that if $`\mathrm{\Delta }`$ is analytically parameterized by a real variable vector $`\stackrel{}{\epsilon }`$, and at $`\stackrel{}{\epsilon }_0`$, $`\mathrm{\Delta }`$ satisfies conditions $`1`$ and $`2`$ of Theorem 1.1, then the hidden Markov chain itself is a real analytic function of $`\stackrel{}{\epsilon }`$ at $`\stackrel{}{\epsilon }_0`$ in a strong sense. We assume (for this section only) that the reader is familiar with the basics of measure theory and functional analysis . Our approach uses a connection between the entropy rate of a hidden Markov chain and symbolic dynamics explored in .
Let $`𝒳`$ denote the set of left infinite sequences with finite alphabet. A cylinder set is a set of the form: $`(\{x_{\mathrm{}}^0:x_0=z_0,\mathrm{},x_n=z_n\})`$. The Borel sigma-algebra is the smallest sigma-algebra containing the cylinder sets. A Borel probability measure (BPM) $`\nu `$ on $`𝒳`$ is a measure on the Borel measurable sets of $`𝒳`$ such that $`\nu (𝒳)=1`$. Such a measure is uniquely determined by its values on the cylinder sets.
For real $`\stackrel{}{\epsilon }`$, consider the measure $`\nu ^\stackrel{}{\epsilon }`$ on $`𝒳`$ defined by:
$$\nu ^\stackrel{}{\epsilon }(\{x_{\mathrm{}}^0:x_0=z_0,\mathrm{},x_n=z_n\})=p^\stackrel{}{\epsilon }(z_n^0).$$
(8.30)
Note that $`H(Z)`$ can be rewritten as
$$H^\stackrel{}{\epsilon }(Z)=\mathrm{log}p^\stackrel{}{\epsilon }(z_0|z_{\mathrm{}}^1)d\nu ^\stackrel{}{\epsilon }.$$
(8.31)
Usually, the Borel sigma-algebra is defined to be the smallest sigma-algebra containing the open sets; in this case, the open sets are defined by the metric: for any two elements $`\xi `$ and $`\eta `$ in $`𝒳`$, define $`d(\xi ,\eta )=2^k`$ where $`k=inf\{|i|:\xi _i\eta _i\}`$. The metric space $`(𝒳,d)`$ is compact.
Let $`C(𝒳)`$ be the space of real-valued continuous functions on $`𝒳`$. Then $`C(𝒳)`$ is a Banach space (i.e., complete normed linear space) with the sup norm $`||f||_{\mathrm{}}=sup\{|f(x)|:x𝒳\}`$. Then any BPM $`\nu `$ acts as a bounded linear functional on $`C(𝒳)`$, namely $`\nu (f)=f𝑑\nu `$. As such, the set of BPM’s is a subset of the dual space, $`C(𝒳)^{}`$, which is itself a Banach space; the norm of a BPM $`\nu `$ is defined: $`\nu =sup_{\{fC(𝒳):f_{\mathrm{}}=1\}}f𝑑\nu `$. In fact, since $`𝒳`$ is compact, $`C(𝒳)^{}`$ is the linear span of the BPM’s.
It makes sense to ask if $`\stackrel{}{\epsilon }\nu ^\stackrel{}{\epsilon }`$ is analytic as a mapping from the parameter space to $`C(𝒳)^{}`$; by definition, this would mean that $`\nu ^\stackrel{}{\epsilon }`$ can be expressed as a power series in the coordinates of $`\stackrel{}{\epsilon }`$. However, as the following example shows, this mapping is not even continuous.
Let $`𝒳`$ be the set of binary left infinite sequences. Let $`\nu _p`$ denote the i.i.d. $`(p,1p)`$ measure, with $`0<p<1`$. Let
$$S_p=\{xX:\underset{n\mathrm{}}{lim}(1/n)(\mathrm{log}p_{x_1}+\mathrm{}+\mathrm{log}p_{x_n})=p\mathrm{log}p(1p)\mathrm{log}(1p)\}.$$
Note that $`S_p`$ is a Borel measurable set. By the strong law of large numbers, $`\nu _p(S_p)=1`$. Clearly, for distinct $`p`$, $`S_p`$ are disjoint. Thus, for $`qp`$, $`\nu _q(S_p)=0`$.
Any Borel measurable set $`S`$ can be approximated by a finite union of cylinder sets in the following sense: given $`\delta >0`$ and $`p(0,1)`$, there is a finite union $`C`$ of cylinder sets such that $`|\nu _q(S)\nu _q(C)|<\delta `$ for all $`q`$ in a neighborhood of $`p`$. Applying this fact to $`S=S_p`$, and denoting $`C_{(p,\delta )}=C`$, we obtain
$$1=\nu _p(S_p)\nu _q(S_p)|\nu _p(S_p)\nu _p(C_{(p,\delta )})|+|\nu _p(C_{(p,\delta )})\nu _q(C_{(p,\delta )})|+|\nu _q(C_{(p,\delta )})\nu _q(S_p)|$$
$$2\delta +|\nu _p(C_{(p,\delta )})\nu _q(C_{(p,\delta )})|.$$
If $`\delta <1/2`$, then $`\nu _q(C_{(p,\delta )})`$ cannot converge to $`\nu _p(C_{(p,\delta )})`$ as $`qp`$. Since the characteristic function of a finite union of cylinder sets is continuous, this shows that the map $`p\nu _p`$ from $``$ to $`C(X)^{}`$ is discontinuous.
On the other hand, using the work of Ruelle , we now show that $`\stackrel{}{\epsilon }\nu ^\stackrel{}{\epsilon }`$ is analytic as a mapping from the parameter space to another natural space.
For $`fC(𝒳)`$, define $`var_n(f)=sup\{|f(\xi )f(\xi ^{})|:\xi _i=\xi _i^{}\text{for}in\}`$. We denote by $`F^\theta `$ the subset of $`fC(𝒳)`$ such that
$$f_\theta \underset{n0}{sup}(\theta ^nvar_n(f))<+\mathrm{}.$$
$`F^\theta `$ is a Banach space with the norm $`f=\mathrm{max}(|f|_{\mathrm{}},f_\theta )`$. Using complex functions instead of real functions, one defines $`F_{}^\theta `$ similarly.
In the following theorem, we prove the analyticity of a hidden Markov chain in a strong sense.
###### Theorem 8.1.
Suppose that the entries of $`\mathrm{\Delta }`$ are analytically parameterized by a real variable vector $`\stackrel{}{\epsilon }`$. If at $`\stackrel{}{\epsilon }=\stackrel{}{\epsilon }_0`$, $`\mathrm{\Delta }`$ satisfies conditions $`1`$ and $`2`$ in Theorem 1.1, then the mapping $`\stackrel{}{\epsilon }\mathrm{log}p^\stackrel{}{\epsilon }(z_0|z_{\mathrm{}}^1)`$ is analytic at $`\stackrel{}{\epsilon }_0`$ from the real parameter space to $`F^\rho `$ (here $`\rho `$ is the contraction constant in the proof of Theorem 1.1). Moreover the mapping $`\stackrel{}{\epsilon }\nu ^\stackrel{}{\epsilon }`$ is analytic at $`\stackrel{}{\epsilon }_0`$ from the real parameter space to $`(F^\rho )^{}`$, the dual space (i.e., bounded linear functionals) on $`F^\rho `$.
###### Proof.
For complex $`\stackrel{}{\epsilon }`$, by (4.16), one shows that $`\mathrm{log}p^\stackrel{}{\epsilon }(z_0|z_{\mathrm{}}^1)`$ can be defined on $`\mathrm{\Omega }_{}`$ as the uniform (in $`\stackrel{}{\epsilon }`$ and $`z𝒳`$) limit of $`\mathrm{log}p^\stackrel{}{\epsilon }(z_0|z_n^1)`$ as $`n\mathrm{}`$, and $`\mathrm{log}p(z_0|z_{\mathrm{}}^1)`$ belongs to $`F^\rho `$. By (4.5), (4.6), (4.7) and (4.14) it follows that $`p^\stackrel{}{\epsilon }(z_0|z_n^1)`$ is analytic on $`\mathrm{\Omega }_{}`$. As a result of (4.16), if $`\mathrm{\Delta }`$ satisfies conditions $`1`$ and $`2`$, for fixed $`z𝒳`$, $`\mathrm{log}p^\stackrel{}{\epsilon }(z_0|z_{\mathrm{}}^1)`$ is the uniform limit of analytic functions and hence is analytic on $`\mathrm{\Omega }_{}`$ (see Theorem $`\mathrm{2.4.1}`$ of ).
Using (4.16) and the Cauchy integral formula in several variables (which expresses the derivative of an analytic function at a point as an integral of a closed curve around the point), we obtain the following. There is a positive constant $`C^{}`$ such that whenever $`z_{\mathrm{}}^0\stackrel{n}{}\widehat{z}_{\mathrm{}}^0`$, for all $`\stackrel{}{\epsilon }\mathrm{\Omega }_{}`$
$$|D_\stackrel{}{\epsilon }(\mathrm{log}p^\stackrel{}{\epsilon }(z_0|z_{n_1}^1))D_\stackrel{}{\epsilon }(\mathrm{log}p^\stackrel{}{\epsilon }(\widehat{z}_0|\widehat{z}_{n_2}^1))|C^{}\rho ^n.$$
(8.32)
Therefore for arbitrary yet fixed $`z_{\mathrm{}}^0`$, the components of the derivatives of $`\mathrm{log}p^\stackrel{}{\epsilon }(z_0|z_{\mathrm{}}^1)`$ with respect to $`\stackrel{}{\epsilon }`$ are also in $`F_{}^\rho `$.
Furthermore, we prove that the mapping $`\stackrel{}{\epsilon }\mathrm{log}p^\stackrel{}{\epsilon }(z_0|z_{\mathrm{}}^1)`$ is complex differentiable (therefore analytic) from $`\mathrm{\Omega }_{}`$ to $`F_{}^\theta `$. Let $`f(\stackrel{}{\epsilon };)=\mathrm{log}p^\stackrel{}{\epsilon }()`$. It suffices to prove that
$$f(\stackrel{}{\epsilon }+\stackrel{}{h};)f(\stackrel{}{\epsilon };)D_\stackrel{}{\epsilon }f|_\stackrel{}{\epsilon }(\stackrel{}{h};)_{\mathrm{}}o(\stackrel{}{h}).$$
(8.33)
and
$$f(\stackrel{}{\epsilon }+\stackrel{}{h};)f(\stackrel{}{\epsilon };)D_\stackrel{}{\epsilon }f|_\stackrel{}{\epsilon }(\stackrel{}{h};)_\theta o(\stackrel{}{h}).$$
(8.34)
Again applying the Cauchy integral formula in several variables, it follows that there exists a positive constant $`C^{\prime \prime }`$ such that for all $`\stackrel{}{\epsilon }\mathrm{\Omega }_{}`$ we have
$$|D_\stackrel{}{\epsilon }^2f|_\stackrel{}{\epsilon }(\stackrel{}{h},\stackrel{}{h};z)|C^{\prime \prime }|\stackrel{}{h}|^2$$
(8.35)
and whenever $`z_{\mathrm{}}^0\stackrel{n}{}\widehat{z}_{\mathrm{}}^0`$,
$$_0^1(1t)|(D_\stackrel{}{\epsilon }^2f|_\stackrel{}{\epsilon }(\stackrel{}{h},\stackrel{}{h};z)D_\stackrel{}{\epsilon }^2f|_\stackrel{}{\epsilon }(\stackrel{}{h},\stackrel{}{h};\widehat{z}))|dtC^{\prime \prime }|\stackrel{}{h}|^2\rho ^n,$$
(8.36)
From the Taylor formula with integral remainder, we have:
$$f(\stackrel{}{\epsilon }+\stackrel{}{h};z)f(\stackrel{}{\epsilon };z)D_\stackrel{}{\epsilon }f|_\stackrel{}{\epsilon }(\stackrel{}{h};z)=_0^1(1t)D_\stackrel{}{\epsilon }^2f|_{\stackrel{}{\epsilon }+t\stackrel{}{h}}(\stackrel{}{h},\stackrel{}{h};z)dt.$$
(8.37)
To prove (8.33), use (8.35) and (8.37). To prove (8.34), use (8.36) and (8.37). Therefore $`\stackrel{}{\epsilon }\mathrm{log}p^\stackrel{}{\epsilon }()`$ is analytic as a mapping from $`\mathrm{\Omega }_{}`$ to $`F_{}^\rho `$. Restricting the mapping $`\stackrel{}{\epsilon }\mathrm{log}p^\stackrel{}{\epsilon }(z_0|z_{\mathrm{}}^1)`$ to the real parameter space, we conclude that it is real analytic (as a mapping into $`F^\rho `$). Using this and the theory of equilibrium states ), the “Moreover” is proven in Appendix C. ∎
###### Corollary 8.2.
Suppose that at $`\stackrel{}{\epsilon }_0`$, $`\mathrm{\Delta }`$ satisfies conditions $`1`$ and $`2`$ in Theorem 1.1, and $`\stackrel{}{\epsilon }f^\stackrel{}{\epsilon }F^\rho `$ be analytic at $`\stackrel{}{\epsilon }_0`$, then $`\stackrel{}{\epsilon }\nu ^\stackrel{}{\epsilon }(f^\stackrel{}{\epsilon })`$ is analytic at $`\stackrel{}{\epsilon }_0`$. In particular, we recover Theorem 1.1: $`\stackrel{}{\epsilon }H^\stackrel{}{\epsilon }(Z)`$ is analytic at $`\stackrel{}{\epsilon }_0`$.
###### Proof.
The map
$$\mathrm{\Omega }F^\rho \times (F^\rho )^{}$$
$$\stackrel{}{\epsilon }(f^\stackrel{}{\epsilon },\nu ^\stackrel{}{\epsilon })\nu ^\stackrel{}{\epsilon }(f^\stackrel{}{\epsilon })$$
is analytic at $`\stackrel{}{\epsilon }_0`$, as desired. ∎
Acknowledgements: We are grateful to Wael Bahsoun, Joel Feldman, Robert Israel, Izabella Laba, Erik Ordentlich, Yuval Peres, Gadiel Seroussi, Wojciech Szpankowski and Tsachy Weissman for helpful discussions.
## Appendices
## Appendix A Proof of Proposition 2.1
###### Proof.
Without loss of generality, we assume $`S`$ is convex (otherwise consider the convex hull of $`S`$). It follows from standard arguments that max norm and sum norm are equivalent. More specifically, for another metric $`d_1`$ defined by
$$d_1(u,v)=\sqrt{\underset{ijk}{}\mathrm{log}^2(\frac{u_i/u_j}{v_i/v_j}.)},$$
we have $`d_\text{B}d_1`$. For metric $`d_2`$ defined by
$$d_2(u,v)=\sqrt{\underset{ijk}{}(u_i/u_jv_i/v_j)^2}.$$
Applying mean value theorem to $`\mathrm{log}`$ function, one concludes that $`d_1d_2`$. Note that
$`u_iv_i`$ $`=`$ $`{\displaystyle \frac{u_i}{u_1+u_2+\mathrm{}+u_k}}{\displaystyle \frac{v_i}{v_1+v_2+\mathrm{}+v_k}}`$
$`=`$ $`{\displaystyle \frac{1}{u_1/u_i+u_2/u_i+\mathrm{}+u_k/u_i}}{\displaystyle \frac{1}{v_1/v_i+v_2/v_i+\mathrm{}+v_k/v_i}}`$
Applying the mean value theorem to function $`f`$, defined as
$$f(x_1,x_2,\mathrm{},x_B)=\frac{1}{x_1+x_2+\mathrm{}+x_k},$$
we conclude that there exists $`\xi S`$ such that
$$u_iv_i=f|_\xi (u_1/u_iv_1/v_i,\mathrm{},u_k/u_iv_k/v_i).$$
It follows from Cauchy inequality that there exists a positive constant $`D_1`$ such that
$$d_𝐄(u,v)<D_1d_2(u,v).$$
Similarly consider $`u_i/u_jv_i/v_j`$, and apply mean value theorem to function $`g`$, defined as $`g(x,y)=x/y`$, we show that there exists a positive constant $`D_2`$ such that
$$d_2(u,v)<D_2d_𝐄(u,v).$$
Namely $`d_2d_E`$. Thus the claim in this Proposition follows, namely there exist two positive constant $`C_1<C_2`$ such that for any two points $`u,vS`$,
$$C_1d_\text{B}(u,v)<d_\text{E}(u,v)<C_2d_\text{B}(u,v).$$
## Appendix B Proof of Lemma 7.9:
Recall that for a non-negative matrix $`B`$, the *canonical form* of $`B`$ is:
$$B=\left[\begin{array}{cccc}B_{11}& B_{12}& \mathrm{}& B_{1n}\\ 0& B_{22}& \mathrm{}& B_{2n}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& B_{nn}\end{array}\right],$$
where $`B_{ii}`$ is either an irreducible matrix (called irreducible components) or a $`1\times 1`$ zero matrix.
Condition 2 in Theorem 7.8 is equivalent to the statement that $`B=B(\stackrel{}{\epsilon }_0)`$ has a unique irreducible component of maximal spectral radius and that this component is primitive. Let $`C`$ denote the square matrix obtained by restricting $`B`$ to this component and let $`S_C`$ denote the set of indices corresponding to this component. Let $`\lambda _1`$ denote the spectral radius of $`B`$, equivalently the spectral radius of $`C`$.
Let $`\lambda _1(\stackrel{}{\epsilon })`$ denote the largest, in modulus, eigenvalue of $`B(\stackrel{}{\epsilon })`$. Since the entries of $`B(\stackrel{}{\epsilon })`$ are analytic in $`\stackrel{}{\epsilon }`$ and $`\lambda _1`$ is simple, it follows that if the complex neighborhood $`\mathrm{\Omega }`$ is chosen sufficiently small, then $`\lambda _1(\stackrel{}{\epsilon })`$ is analytic function of $`\stackrel{}{\epsilon }\mathrm{\Omega }`$.
The columns (resp., rows) of $`Adj(\lambda _1(\stackrel{}{\epsilon })IB(\stackrel{}{\epsilon }))`$ are right (resp., left) eigenvectors of $`B(\stackrel{}{\epsilon })`$ corresponding to $`\lambda _1(\stackrel{}{\epsilon })`$. By choosing $`x(\stackrel{}{\epsilon })`$ (resp. $`y(\stackrel{}{\epsilon })`$) to be a fixed column (resp. row) of $`Adj(\lambda _1(\stackrel{}{\epsilon })IB(\stackrel{}{\epsilon }))`$ and then replacing $`x(\stackrel{}{\epsilon })`$ and $`y(\stackrel{}{\epsilon })`$ by appropriately rescaled versions, we may assume that:
* $`x(\stackrel{}{\epsilon }_0),y(\stackrel{}{\epsilon }_0)0`$, and they are positive on $`S_C`$
* $`y(\stackrel{}{\epsilon })x(\stackrel{}{\epsilon })=1`$
* $`x(\stackrel{}{\epsilon })`$ and $`y(\stackrel{}{\epsilon })`$ are analytic in $`\stackrel{}{\epsilon }\mathrm{\Omega }`$
Let
$$V(\stackrel{}{\epsilon })=\lambda _1(\stackrel{}{\epsilon })x(\stackrel{}{\epsilon })y(\stackrel{}{\epsilon })$$
and
$$U(\stackrel{}{\epsilon })=B(\stackrel{}{\epsilon })V(\stackrel{}{\epsilon }).$$
Then $`V(\stackrel{}{\epsilon })`$ is the restriction of $`B(\stackrel{}{\epsilon })`$ to the subspace corresponding to $`\lambda _1(\stackrel{}{\epsilon })`$ and $`U(\stackrel{}{\epsilon })`$ is the restriction to the subspace corresponding to the remainder of the spectrum of $`B(\stackrel{}{\epsilon })`$. It follows that
$$U(\stackrel{}{\epsilon })V(\stackrel{}{\epsilon })=0=V(\stackrel{}{\epsilon })U(\stackrel{}{\epsilon }).$$
Let $`\mu (\stackrel{}{\epsilon })`$ denote the spectral radius of $`U(\stackrel{}{\epsilon })`$. By condition $`2`$, $`\mu (\stackrel{}{\epsilon }_0)<\lambda _1(\stackrel{}{\epsilon }_0)`$. Thus, there is a constant $`\nu >0`$ such that if the neigbourhood $`\mathrm{\Omega }`$ is sufficiently small, then for all $`\stackrel{}{\epsilon }\mathrm{\Omega }`$
$$\mu (\stackrel{}{\epsilon })<\nu <|\lambda _1(\stackrel{}{\epsilon })|.$$
Thus, by Lemma 7.5, and making still $`\mathrm{\Omega }`$ smaller if necessary, there is a constant $`K_1>0`$ such that for all $`i,j`$, all $`n`$ and all $`\stackrel{}{\epsilon }\mathrm{\Omega }`$,
$$|U_{ij}^n(\stackrel{}{\epsilon })|<K_1\nu ^n.$$
(B.38)
Let $`r=r(\stackrel{}{\epsilon }_0)`$, $`c=c(\stackrel{}{\epsilon }_0)`$, $`x=x(\stackrel{}{\epsilon }_0)`$ and $`y=y(\stackrel{}{\epsilon }_0)`$. In the following we will show that the irreducibility of $`\mathrm{\Delta }`$ will rule out the possibility that $`c`$ is non-zero only in non-maximal spectral radius irreducible components of $`B`$, and so we can extend $`a(\stackrel{}{\epsilon },n)`$ and $`b(\stackrel{}{\epsilon },n)`$ from real to complex.
Let $`s_0S_C`$. Since $`\mathrm{\Delta }(\stackrel{}{\epsilon }_0)`$ is irreducible and $`r`$ is nonnegative, but not the zero vector, for some $`j_0`$, $`(rB^{j_0})_{s_0}>0`$. Similarly, for any index $`s_1`$ other than 1 of the underlying Markov chain, there exists $`j_1`$ such that $`B_{s_0s_1}^{j_1}>0`$. Choose $`s_1`$ to be any index such that $`c_{s_1}>0`$. Since $`C`$ is primitive, it then follows that there is a constant $`K_2`$ such that for sufficiently large $`n`$,
$$rxyc\lambda _1^n+rU^nc=rV^nc+rU^nc=rB^nc>K_2\lambda _1^n,$$
which by (B.38) implies that $`rxyc>0`$. Therefore if $`\mathrm{\Omega }`$ is sufficiently small, there exists a positive constant $`K_4`$ such that
$$|r(\stackrel{}{\epsilon })x(\stackrel{}{\epsilon })y(\stackrel{}{\epsilon })c(\stackrel{}{\epsilon })|>K_4,$$
for $`\stackrel{}{\epsilon }\mathrm{\Omega }`$.
Let $`K_3`$ be an upper bound on the entries of $`|x(\stackrel{}{\epsilon })|,|y(\stackrel{}{\epsilon })|,|r(\stackrel{}{\epsilon })|`$ and $`|c(\stackrel{}{\epsilon })|`$.
Thus, for all $`n`$ and all $`\stackrel{}{\epsilon }\mathrm{\Omega }`$, we have
$$|r(\stackrel{}{\epsilon })B^n(\stackrel{}{\epsilon })c(\stackrel{}{\epsilon })||r(\stackrel{}{\epsilon })U^n(\stackrel{}{\epsilon })c(\stackrel{}{\epsilon })|+|r(\stackrel{}{\epsilon })V^n(\stackrel{}{\epsilon })c(\stackrel{}{\epsilon })|||^2K_3^2K_1\nu ^n+||^2K_3^4|\lambda _1(\stackrel{}{\epsilon })|^n$$
and
$$|r(\stackrel{}{\epsilon })B^n(\stackrel{}{\epsilon })c(\stackrel{}{\epsilon })||r(\stackrel{}{\epsilon })V^n(\stackrel{}{\epsilon })c(\stackrel{}{\epsilon })||r(\stackrel{}{\epsilon })U^n(\stackrel{}{\epsilon })c(\stackrel{}{\epsilon })|K_4|\lambda _1(\stackrel{}{\epsilon })|^n||^2K_3^2K_1\nu ^n.$$
With similar upper and lower bounds for $`|(r(\stackrel{}{\epsilon })B^n(\stackrel{}{\epsilon })\mathrm{𝟏}|`$, it follows that for sufficiently large $`n`$ and all $`\stackrel{}{\epsilon }\mathrm{\Omega }`$,
$$\frac{\pi _1(\stackrel{}{\epsilon })r(\stackrel{}{\epsilon })B(\stackrel{}{\epsilon })^n\mathrm{𝟏}}{\pi _1(\stackrel{}{\epsilon })r(\stackrel{}{\epsilon })B(\stackrel{}{\epsilon })^{n1}\mathrm{𝟏}}$$
and
$$\frac{\pi _1(\stackrel{}{\epsilon })r(\stackrel{}{\epsilon })B(\stackrel{}{\epsilon })^{n1}c(\stackrel{}{\epsilon })}{\pi _1(\stackrel{}{\epsilon })r(\stackrel{}{\epsilon })B(\stackrel{}{\epsilon })^{n1}\mathrm{𝟏}}$$
are uniformly bounded from above and away from zero. By condition $`1`$, for any finite collection of $`n`$, there is a (possibly smaller) neighborhood $`\mathrm{\Omega }`$ of $`\stackrel{}{\epsilon }_0`$, such that for all $`\stackrel{}{\epsilon }\mathrm{\Omega }`$, these quantities are uniformly bounded from above and away from zero. This completes the proof of Lemma 7.9 ( and therefore the proof of sufficiency for Theorem 7.8.)
## Appendix C $`\stackrel{}{\epsilon }\nu ^\stackrel{}{\epsilon }`$ is analytic
In this appendix, we follow the notation in Section 8. Let $`\tau :𝒳𝒳`$ be the right shift operator, which is a continuous mapping on $`𝒳`$ under the topology induced by the metric $`d`$. For $`fC(𝒳)`$, one defines the pressure via a variational principle :
$$P(f)=\underset{\mu M(𝒳,\tau )}{sup}\left(H_\mu (\tau )+f𝑑\mu \right),$$
where $`M(𝒳,\tau )`$ denotes the set of $`\tau `$-invariant probability measures on $`𝒳`$ and $`H_\mu (\tau )`$ denotes measure-theoretic entropy. A member $`\mu `$ of $`M(𝒳,\tau )`$ is called an *equilibrium state* for $`f`$ if $`P(f)=H_\mu (T)+f𝑑\mu `$.
For $`fC(𝒳)`$ the Ruelle operator $`_f:C(𝒳)C(𝒳)`$ is defined by
$$(_fh)(x)=\underset{y\tau ^1x}{}e^{f(y)}h(y).$$
The connection between pressure and the Ruelle operator is as follows . When $`fF^\theta `$, $`P(f)`$ is $`\mathrm{log}\lambda `$, where $`\lambda `$ is the spectral radius of $`_f`$. The restriction of $`_f`$ to $`F^\theta `$ still has spectral radius $`\lambda `$, and $`\lambda `$ is isolated from all other eigenvalues of the restricted operator. Using this, Ruelle applied standard perturbation theory for linear operators to conclude that pressure $`P(f)`$ is real analytic on $`F^\theta `$. Moreover, he showed that each $`fF^\theta `$ has a unique equilibrium state $`\mu _f`$ and the first order derivative of $`fP(f)`$ on $`F^\theta `$ is $`\mu _f`$, viewed as a linear functional on $`F^\theta `$. So, the analyticity of $`P(f)`$ implies that the equilibrium state $`\mu _f`$ is also analytic in $`fF^\theta `$.
We first claim that for $`f(\stackrel{}{\epsilon },z)=\mathrm{log}p^\stackrel{}{\epsilon }(z_0|z_{\mathrm{}}^1)`$, we have $`\mu _{f(\stackrel{}{\epsilon },)}=\nu ^\stackrel{}{\epsilon }`$ as in (8.30).
To see this, first observe that the spectral radius $`\lambda `$ of $`=_{f(\stackrel{}{\epsilon },)}`$ is 1; this follows from the observations:
* the function $`\overline{1}`$ which is identically 1 on $`𝒳`$ is a fixed point of $``$ – and –
* (see Proposition 5.16 of ) $`^n(\overline{1})/\lambda ^n`$ converges to a strictly positive function.
Thus $`P(f(\stackrel{}{\epsilon },))=0`$. So, for $`\mu ^\stackrel{}{\epsilon }=\mu _{f(\stackrel{}{\epsilon },)}`$, we have
$$h_{\mu ^\stackrel{}{\epsilon }}(\tau )+f(\stackrel{}{\epsilon },)𝑑\mu ^\stackrel{}{\epsilon }=0.$$
But from (8.31), we have
$$h_\nu (\tau )+f(\stackrel{}{\epsilon },)𝑑\nu ^\stackrel{}{\epsilon }=0.$$
By uniqueness of the equilibrium state, we thus obtain $`\mu _{f(\stackrel{}{\epsilon },)}=\nu ^\stackrel{}{\epsilon }`$ as claimed.
Since $`\stackrel{}{\epsilon }f(\stackrel{}{\epsilon },)`$ is analytic, it then follows that $`\stackrel{}{\epsilon }\nu ^\stackrel{}{\epsilon }`$ is analytic, thereby completing the proof of Theorem 8.1. |
warning/0507/physics0507061.html | ar5iv | text | # Quantum control of optical four-wave mixing with femtosecond 𝜔-3𝜔 laser pulses: coherent ac Stark nonlinear spectroscopy
## Abstract
The four-wave mixing produced with two ultrashort phase-locked $`\omega `$-$`3\omega `$ laser pulses propagating coherently in a two-level system in the infrared spectral region is shown to depend on the pulses relative phase. The Maxwell-Bloch equations are solved beyond the rotating-wave approximation to account for field frequencies which are largely detuned from the atomic resonance. The relative phase dominating the efficiency of the coupling to the $`5\omega `$ anti-Stokes Raman component is determined by sign of the total ac Stark shift induced in the system, in such a way that the phase influence disappears precisely where the ac Stark effect due to both pulses is compensated. This fundamental quantum interference effect can be the basis for nonlinear ultrafast optical spectroscopy techniques.
Phenomena arising from the coherent control of nonlinear interactions are of importance in fields as diverse as optoelectronics and materials research, in high harmonic generation, in photoionization or molecular dissociation, and in biological applications such as spectroscopy and imaging, among others shapiro ; photodiss ; Watanabe ; photoion ; Bandrauk ; Brown ; Xu . Optical quantum coherent control is based on the fact that the phases of interfering transition amplitudes in light-matter interactions can be controlled through the optical phase of coherent light sources that drive the interaction, in such a way that the transition rates to final states and the dynamics at various stages of the process can be modified shapiro .
In a recent paper pra1 , a theoretical investigation on the quantum coherent control of the optical transient four-wave mixing of two intense phase-locked femtosecond laser pulses of central angular frequencies $`\omega `$ and $`3\omega `$ propagating in a two-level atom (TLA) was reported. It was shown how the nonlinear ($`\chi ^{(3)}`$) coupling to the anti-Stokes Raman field at frequency $`5\omega `$ depends critically on the initial relative phase $`\varphi `$ of the propagating pulses. In Ref. pra1 , the study was centered to intense pulses in the visible and ultraviolet spectral regions, with frequencies at resonance or lower than the atomic transition. The phenomena observed in pra1 , however, can be scaled to various laser and material parameters. In the infrared spectral range, for instance, experiments on ultrafast molecular dynamics are frequently performed by help of two-color pump probe nonlinear spectroscopy techniques. In this type of studies, due to the high intensities inherent to ultrashort (subpicosecond) pulses, nonlinear effects such as stimulated Raman processes may become important and are often utilized as a complementary tool to gain information IRspectroscopy .
In this Letter, we address the conditions for quantum coherent control of transient four-wave mixing interactions with two phase-locked $`\omega `$-$`3\omega `$ femtosecond laser pulses in the mid-IR spectral region. Our purpose is to reveal the basic physics for a strict two-level medium and to this end the Maxwell-Bloch TLA will be considered. We will study pulses with a duration of 300 fs (with a spectral width of $`35`$ cm<sup>-1</sup>) and peak intensities as $``$10<sup>8</sup> W/cm<sup>2</sup>, which are typically used in infrared nonlinear spectroscopy experiments IRspectroscopy . We will examine the influence of frequency detuning by considering the field at $`3\omega `$ above resonance with respect to the atomic transition, something that was not considered in pra1 . Under these conditions, we will show that the quantum interferences leading to a dominating efficiency of the nonlinear coupling to the anti-Stokes Raman component are governed by the ac Stark shift induced in the system. Furthermore, we will demonstrate that the relative spectral amplitude of the anti-Stokes fields produced by phase-locked pulses cancels for frequency detunings that compensate the ac Stark effect. This fundamental interference effect might be the basis for a novel ultrafast spectroscopy tool based on coherent control, which we name coherent ac Stark nonlinear spectroscopy (CSNS) for use later below.
The pulse propagation is modeled by means of the Maxwell-Bloch equations beyond the rotating-wave approximation, allowing the resonant as well as the non-resonant regimes of the system to be described shapiro . The equations are written as
$`{\displaystyle \frac{H}{t}}`$ $`=`$ $`{\displaystyle \frac{1}{\mu _0}}{\displaystyle \frac{E}{z}},`$
$`{\displaystyle \frac{E}{t}}`$ $`=`$ $`{\displaystyle \frac{1}{ϵ_0}}{\displaystyle \frac{H}{z}}{\displaystyle \frac{N_{at}\mu }{ϵ_0T_2}}(\rho _1T_2\omega _{12}\rho _2),`$
$`{\displaystyle \frac{\rho _1}{t}}`$ $`=`$ $`{\displaystyle \frac{1}{T_2}}\rho _1+\omega _{12}\rho _2,`$ (1)
$`{\displaystyle \frac{\rho _2}{t}}`$ $`=`$ $`{\displaystyle \frac{1}{T_2}}\rho _2+{\displaystyle \frac{2\mu }{\mathrm{}}}E\rho _3\omega _{12}\rho _1,`$
$`{\displaystyle \frac{\rho _3}{t}}`$ $`=`$ $`{\displaystyle \frac{1}{T_1}}(\rho _3\rho _{30}){\displaystyle \frac{2\mu }{\mathrm{}}}E\rho _2,`$
where $`H(z,t)`$ and $`E(z,t)`$ represent the magnetic and electric fields propagating along the $`z`$ direction, respectively, $`\mu _0`$ and $`ϵ_0`$ are the magnetic permeability and electric permittivity of free space, respectively, $`N_{at}=2\times 10^{24}`$ m<sup>-3</sup> is the density of polarizable atoms, $`\mu =4.2\times 10^{29}`$ Cm is the effective dipole coupling coefficient, $`T_1=T_2=1`$ ps are the excited-state lifetime and dephasing time, respectively, $`\rho _1`$ and $`\rho _2`$ are the real and imaginary components of the polarization, and $`\omega _{12}`$ is the transition resonance angular frequency of the two level medium, considered in the present simulations in the mid-IR region at 3000 nm (see Fig. 1). The population difference is $`\rho _3`$, and $`\rho _{30}`$ represents its initial value. An hyperbolic secant two-color pulse that can be expressed as
$`E(t)=E_\omega (t)+E_{3\omega }(t)=E_0sech((tt_0)/t_p)\times `$
$`\left[cos(\omega (tt_0))+cos(3\omega (tt_0)+\varphi )\right]`$ (2)
is externally injected to the system. The peak input electric field amplitude $`E_0`$ is chosen the same for both pulses and results in an intensity of $`4.0\times 10^8`$ W/cm<sup>2</sup>. The duration of the pulses is given by $`t_p=\tau _p/1.763`$, with $`\tau _p=300`$ fs being the full width at half maximum (FWHM) of the pulse intensity envelope. $`t_0`$ gives the offset position of the pulse center at $`t=0`$, and it is the reference value for the phase of the pulses. The central angular frequencies of the pulses are $`\omega `$ and $`3\omega `$, and $`\varphi `$ is the relative phase. The propagating system has been resolved numerically by means of a standard finite difference time domain method described elsewhere pra1 .
Figure 2 shows the field spectra at different propagation lengths in the case that the central angular frequency of the field $`E_{3\omega }`$ is at resonance with the atomic transition ($`3\omega =\omega _{12}`$). The spectrum on the top is for the initial pulses. The succeeding plots show the evolution of the spectrum as the pulses propagate through the medium. In this case, we can observe the effect of the absorption of the pulse at $`3\omega `$, together with the appearance of other spectral contributions at $`5\omega `$, $`7\omega `$, and $`9\omega `$, which are produced as a result of the coupling of the fields through the third order nonlinearity ($`\chi ^{(3)}`$) of the medium. It is clear that the conversion to the anti-Stokes Raman component ($`5\omega `$) depends on the relative phase between the pulses. For the parameter values of the results shown in Fig. 2, the coupling to the fifth harmonic ($`5\omega `$) is more efficient for $`\varphi =0`$ than for $`\varphi =\pi `$, an scenario that was already reported in Ref. pra1 for a transition in the visible region. Hence here we confirm that observed in Ref. pra1 for a different system, particularly, considering a mid-IR atomic transition. It is important to note that this phase dependence effect involves the anti-Stokes Raman component ($`5\omega `$) only, not the $`7\omega `$ nor the $`9\omega `$ spectral components, which remain insensitive to the initial relative phase of the pulses.
We now turn to the study of the frequency detuning of the fields with respect to the atomic transition. We will show that there is a central pulse frequency $`\omega `$ at which the relative phase dependence of the coupling to the anti-Stokes Raman component, which has been previously discussed in Fig. 2, disappears. We will observe that this effect can occur because the ac Stark shifts produced in the medium by the fields $`E_\omega `$ and $`E_{3\omega }`$ can be compensated when $`\omega <\omega _{12}<3\omega `$. We will then conclude that the ac Stark shift in the medium governs the relative phase dependence of the coherent four-wave coupling.
Indeed, the ac Stark frequency shift $`\mathrm{\Delta }\omega `$ produced by a field of frequency $`\omega `$ which is not near resonance in a transition of frequency $`\omega _{12}`$ can be expressed as acStark :
$`\mathrm{\Delta }\omega `$ $`=`$ $`{\displaystyle \frac{\omega _{12}\omega }{2}}`$
$`\pm {\displaystyle \frac{1}{2}}\left[(\omega _{12}\omega )^2+4\beta ^2(1+{\displaystyle \frac{\omega _{12}\omega }{\omega _{12}+\omega }})\right]^{1/2},`$
where $`+`$ is for $`\omega <\omega _{12}`$ and $``$ is for $`\omega \omega _{12}`$, $`\beta =\mu |E|/(2\mathrm{})`$, with $`|E|`$ being the amplitude of the field, $`\omega _{12}`$ the frequency of the atomic transition, and $`\omega `$ the field angular frequency. The conditions $`\beta /\omega <<1`$ and $`\beta /\omega _{12}<<1`$ must be met for Eq. Quantum control of optical four-wave mixing with femtosecond $`\omega `$-$`3\omega `$ laser pulses: coherent ac Stark nonlinear spectroscopy to be valid acStark . Far from resonance, where
$`{\displaystyle \frac{(\omega _{12}\omega )^2}{\beta ^2}}`$ $`>>`$ $`{\displaystyle \frac{8\omega _{12}}{\omega _{12}+\omega }},`$ (4)
Eq. (Quantum control of optical four-wave mixing with femtosecond $`\omega `$-$`3\omega `$ laser pulses: coherent ac Stark nonlinear spectroscopy) becomes
$`\mathrm{\Delta }\omega \beta ^2\left[{\displaystyle \frac{1}{\omega _{12}\omega }}+{\displaystyle \frac{1}{\omega _{12}+\omega }}\right].`$ (5)
Note that the last term inside the brackets in Eq. (5) is important far from resonance, where the rotating wave approximation does not apply. Clearly from Eq. (5), when $`\omega >\omega _{12}`$ we have $`\mathrm{\Delta }\omega <0`$, and the separation in frequency of the states ($`\omega _{12}+\mathrm{\Delta }\omega `$) appears to be less than in the absence of the field. Contrarily, for $`\omega <\omega _{12}`$, the ac Stark frequency shift is positive.
In the case of the present investigation, the main contributions to the ac Stark effect come from the two components of the propagating pulses $`E_\omega `$ and $`E_{3\omega }`$, which have central angular frequencies $`\omega `$ and $`3\omega `$, respectively, and the conditions for Eq. (5) are clearly met for the parameter values of our study. Therefore, requiring that the combined ac Stark effect is null
$`0={\displaystyle \frac{1}{\omega _{12}\omega }}+{\displaystyle \frac{1}{\omega _{12}+\omega }}+{\displaystyle \frac{1}{\omega _{12}3\omega }}+{\displaystyle \frac{1}{\omega _{12}+3\omega }},`$ (6)
we obtain $`\omega =\omega _{12}/\sqrt{5}`$.
We will next show that the ac Stark cancellation frequency $`\omega =\omega _{12}/\sqrt{5}`$ is indeed observed with accuracy from numerical simulations. In Figure 3, the spectral amplitude of the anti-Stokes $`5\omega `$ component is shown for different values of the detuning between the fields and the atomic transition at a propagation length such as $`z=25\mu `$m. It is useful to define the parameter $`\eta =(\omega _{12}/\omega )^2`$, which sets the value of the detuning of the $`E_\omega `$ and $`E_{3\omega }`$ fields. The spectra corresponding to $`\varphi =0`$ is shown by solid lines in Fig. 3, while the spectra for $`\varphi =\pi `$ is represented by the dashed lines. Clearly, there is a switch in the tendency to dominate the conversion to the $`5\omega `$ anti-Stokes component. For $`\eta <5`$, the ac Stark shift due to $`E_\omega `$ dominates over the shift induced by $`E_{3\omega }`$, and therefore the resulting separation of the states due to the combined Stark effect appears to be larger than in the absence of fields. In this situation, the coupling to the $`5\omega `$ anti-Stokes field is enhanced for $`\varphi =\pi `$, as it can be observed in Fig. 3, left plots. Contrarily, for $`\eta >5`$, when the ac Stark shift due to $`E_{3\omega }`$ dominates over the shift induced by $`E_\omega `$, the coupling to the $`5\omega `$ anti-Stokes Raman component is enhanced for $`\varphi =0`$ (right plots in Fig. 3). In this last case the resulting separation in frequency of the states due to the total Stark effect is less than in the absence of fields. Furthermore, as shown in Fig. 3, the switching occurs at $`\eta =5`$, where the total ac Stark shift is cancelled, as expected from Eq. 6. Note that for the numerical simulations to agree with the prediction of the analytical theory, the expression for the Stark effect beyond the rotating wave approximation needs to be considered (see Eq. (33) in Ref. acStark ).
We now look at the relative spectral amplitude of the anti-stokes Raman fields, which we define as
$`\mathrm{\Delta }E_{5\omega }(z)`$ $`=`$ $`{\displaystyle \frac{E_{5\omega }^0(z)E_{5\omega }^\pi (z)}{E_{5\omega }^0(z)+E_{5\omega }^\pi (z)}},`$ (7)
with $`E_{5\omega }^0(z)`$ being the spectral amplitude at the generated anti-Stokes ($`5\omega `$) frequency in the case that the initial relative phase is $`\varphi =0`$ (see Fig. 3), and $`E_{5\omega }^\pi (z)`$ being the spectral amplitude at $`5\omega `$ in the case that the initial relative phase between the pulses is $`\varphi =\pi `$. Figure 4 shows the results obtained from the numerical simulations for $`z=25\mu `$m. Relative differences in amplitudes as $`|\mathrm{\Delta }E_{5\omega }|0.5`$ can readily be produced in the cases considered in our simulations. Moreover, we have checked that the relative phase dependence at $`\eta =5.0`$ remains null for propagation distances as long as $`z=100\mu `$m.
We have therefore demonstrated that considering the coherent propagation of two-color phase-locked femtosecond pulses in a two-level medium, with central angular frequencies $`\omega `$ and $`3\omega `$, one can find an angular frequency $`\omega `$ at which the ac Stark effect produced by the propagating pulses is cancelled. At this frequency value, which requires the pulse with central frequency $`3\omega `$ to be above resonance with respect to the atomic transition, the phase dependence of the transient four-wave coupling through the $`\chi ^{(3)}`$ nonlinear susceptibility of the medium disappears. This is a fundamental quantum interference effect that to the best of our knowledge has not been reported before.
In the present Letter, we have considered atomic frequencies which lie in the mid-IR region, which are of interest for applications e.g. in infrared nonlinear spectroscopic techniques. It has to be stressed however that the phenomena that we report can be scaled to several material and pulse parameters. We can hence imagine straightforward applications for nonlinear ultrafast spectroscopy techniques based on the coherent control of subpicosecond two-color $`\omega `$-$`3\omega `$ propagating laser pulses. Indeed, the production of phase related $`\omega `$-$`3\omega `$ pulses is frequently accomplished by some frequency trippling mechanism with the subsequent variation of the phase of one of the pulses in order to obtain experimental control over the relative phase. Although our analysis has obviously been simplified by considering two well isolated levels as a first approach, the switching effect discussed here is of a fundamental level, and in that sense it should be observed experimentally in particular media where the two-level approximation is met. For instance, some gaseous atoms have well isolated resonances (as e.g. rubidium in the visible region). Also, the two-level approximation can be used for studying coherent effects in materials with a broad distribution of transitions, such as inhomogeneously broadened resonance lines in gases and in condensed matter, and even for inhomogeneous quasicontinuous energy bands as in semiconductors twolevelatoms ; twolevelexperiments . CSNS can hence provide information on the transitions being probed by measuring the relative spectral amplitude of the anti-Stokes Raman fields with a simple scan of the laser frequency. The ac Stark mediated coherent control scenario reported in this Letter has therefore a broad interest and may be the basis for future studies on more complex systems.
Support from the Programa Ramón y Cajal of the Spanish Ministry of Science and Technology and from project FIS2004-02587 is acknowledged. |
warning/0507/cond-mat0507048.html | ar5iv | text | # Thermodynamics of An Ideal Generalized Gas: II Means of Order 𝛼
## I Mathematical versus physical inequalities
Since the harmonic-arithmetic-geometric mean inequalities are particular manifestations of the property that power means are increasing functions of their order, it was thought that the second law inequality could be derived from this property. This would avoid having to resort to experiment to determine the sign of the entropic change in processes involving non-quasi static changes.
Apart from an exercise that can be found in Sommerfeld’s book *Thermodynamics and Statistical Physics* Sommerfeld (1956), where he shows that the increase in entropy of an ideal classical gas (ICG) which has come to equilibrium is “a generalization of the inequality between arithmetical and geometrical means”, the real impetus began with a series of papers in the early ’80’s by Landsberg and co-worker PTL-bis (1987) to generalize the arithmetic-geometric mean inequality to cases of negative heat capacities. Sidhu Sidhu (1980) generalized their results to arbitrary power means of order greater than one, and excluded negative powers on the basis that they would be in conflict with the third law.
Yet all these results could be found in an earlier paper by Cashwell and Everett Cashwell (1967) who, because they dealt with temperature dependent heat capacities, were dealing with an ideal generalized gas (IGG). Their work was, however, limited to processes of pure thermal conduction in which a system was divided into a number of cells whose initial temperatures were all not equal and subsequently allowed to interact thermally by replacing the adiabatic partitions by diathermal ones. The mass fractions played the role of a complete probability distribution. Probabilistic notions naturally arise when the uncontrollable processes concerning heat transfer among the cells occur. The final equilibrium state was characterized by a common mean temperature determined from the conservation of the internal energy. The second law followed from the property that the final common temperature was greater than that which would have been obtained in an entropy-conserving equilibration, or that the power means are monotonically increasing functions of their order.
It was realized that in an entropy-conserving equilibration Landsberg (1978) a final common temperature would be reached that would be lower than an energy-conserving equilibration, implying that energy has been extracted from the $`n`$-body system for the performance of external work. This would then provoke a negative change in the internal energy making it comparable to an entropy evolution criterion. It would thus appear that the first and second laws have exchanged roles. The problem inherent to such a formulation, apart from processes involving pure thermal conduction, is that the power means determined from an energy- or entropy-conserving equilibration are not comparable: the temperature dependencies are different but the volume dependencies are not, since both internal energy and entropy are first-order homogeneous functions of the volume. Thus, for processes other than pure heat conduction, the two laws are not comparable. Moreover, in the case of pure deformations, there would be no evolution criterion at all because both potentials are first-order homogeneous functions of the volume.
For processes involving pure thermal conduction, a formulation different from the Cashwell-Everett one was given by Abriata Abriata (1979). In that formulation, an increase in entropy, on the average, is associated with going from a less uniform to a more uniform temperature distribution, or becoming ‘less spread out’. As such it can be formulated as a problem of majorization Marshall *et. al* (1979), with the internal energy, or the $`L`$-potential, playing the role of a Schur convex function of the emprical temperature. Although the entropy need not be a Schur concave function of the empirical temperature, it must necessarily be one of the energy. In this way a continuous temperature distribution involving the Schur convex function of the energy is contrasted to a possibly discrete energy distribution related to the Schur concave function of the entropy. Majorization is considered in the last section.
## II Comparable means and the laws of thermodynamics
The conventional forms of the first and second laws are only comparable for processes involving pure heat conduction. Thus, Cashwell and Everett’s results Cashwell (1967) constitute a particular case where the internal energy, $`E`$, and the $`L`$-potential differ by a constant factor since the volume is held constant. Their results cannot be generalized to processes involving work.
Means are said to be comparable if there exists an inequality Hardy *et. al.* (1952)
$`𝔐_f(z)=f^1\left({\displaystyle \underset{i=1}{\overset{n}{}}}p_if(z_i)\right)`$ (1)
$`g^1\left({\displaystyle \underset{i=1}{\overset{n}{}}}p_ig(z_i)\right)=𝔐_g(z)`$
between them, where $`\{p_i\}`$ is a complete probability distribution, and $`f^1(z)`$ the inverse function. A well-known necessary and sufficient condition for inequality (1) to hold is that the composite function $`gf^1`$ be convex on the interval $`[z_c,z_h]`$, if $`g`$ is increasing Hardy *et. al.* (1952). Inequality (1) is satisfied by power means, where the generators are $`g=z^\alpha `$ and $`f=z^{\alpha 1}`$.
Since the concept of absolute entropy has no meaning in thermodynamics, one cannot distinguish between the Nernst and Planck formulations of the third law Einbinder (1948). This is translated into the property of equivalent means Hardy *et. al.* (1952): in order for
$$𝔐_F(z)=𝔐_f(z),$$
it is both necessary and sufficient that
$$F=af+b,$$
where $`a`$ and $`b`$ are constants, and $`a0`$.
Consider a system comprised of $`n`$ cells, adiabiatically isolated from the environment. Initially the walls of the cells are rigid and adiabatic. When the walls are replaced by deformable, diathermal ones, the initial probability that the $`i`$th cell will have a linear dimension, $`R_i=V_i^{1/(qr)}`$, and an empirical temperature, $`t_i=T^{1/r}`$, is $`p_i`$. Probabilities enter naturally when dealing with processes of heat exchange: Heat is the uncontrollable form of work Thomson (1888), and temperature is its measure.
If $`\stackrel{}{z}`$ is an $`n`$-tuple of initial values of the adiabatic variables, $`z_i=\left(t_iR_i\right)^r`$, and $`\stackrel{}{\zeta }`$ is an $`n`$-tuple of their final, equilibrium, values then the average changes in $`L`$ and $`S`$ when the subsystems have been placed into thermal and mechanical contact are
$$\overline{\mathrm{\Delta }L}=\underset{i=1}{\overset{n}{}}p_i_{z_i}^{\zeta _i}𝑑L(z)=𝔐_\alpha ^\alpha (\zeta )𝔐_\alpha ^\alpha (z),$$
(2)
and
$`\overline{\mathrm{\Delta }S}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}p_i{\displaystyle _{z_i}^{\zeta _i}}{\displaystyle \frac{dL(z)}{z}}`$ (3)
$`=`$ $`{\displaystyle \frac{\alpha }{\alpha 1}}\left\{𝔐_{\alpha 1}^{\alpha 1}(\zeta )𝔐_{\alpha 1}^{\alpha 1}(z)\right\}.`$
The vanishing of either relation, (2) or (3) would determine a uniform mean value of $`\zeta `$. From an $`L`$-equilibration it would be $`𝔐_\alpha (z)`$, while from an $`S`$-equilibration the final mean value would be $`𝔐_{\alpha 1}(z)`$. On the strength of (1), the former would be larger than the latter.
Consider an $`L`$-equilibrating transition where (2) vanishes. Then, if we divide the cells into two groups those for which $`z_i𝔐_\alpha (z)`$, marked by a “$``$ ” on the upper limit of the summation sign, and those for which $`z_i>𝔐_\alpha (z)`$, indicated by a “$`>`$ ” on the summation sign, we find
$`{\displaystyle \underset{i=1}{\overset{}{}}}p_i{\displaystyle _{z_i}^{𝔐_\alpha (z)}}{\displaystyle \frac{dL(z)}{z}}>{\displaystyle \frac{1}{𝔐_\alpha (z)}}{\displaystyle \underset{i=1}{\overset{}{}}}p_i{\displaystyle _{z_i}^{𝔐_\alpha (z)}}𝑑L(z)`$
$`=`$ $`{\displaystyle \frac{1}{𝔐_\alpha (z)}}{\displaystyle \underset{i=1}{\overset{>}{}}}p_i{\displaystyle _{𝔐_\alpha (z)}^{z_i}}𝑑L(z)>{\displaystyle \underset{i=1}{\overset{>}{}}}p_i{\displaystyle _{𝔐_\alpha (z)}^{z_i}}{\displaystyle \frac{dL(z)}{z}}.`$
Since $`\overline{\mathrm{\Delta }S}>0`$ for $`\overline{\mathrm{\Delta }L}=0`$, no matter what the initial $`n`$-tuple of $`z`$ values are, the inequality Cashwell *et. al.* (1967)
$$\underset{i=1}{\overset{n}{}}p_i_{z_i}^{\zeta _i}\frac{dL(z)}{z}\underset{i=1}{\overset{n}{}}p_i_{z_i}^{𝔐_\alpha (z)}\frac{dL(z)}{z}$$
(4)
is a consequence of the fact that the last term in
$`{\displaystyle \underset{i=1}{\overset{n}{}}}p_i{\displaystyle _{z_i}^{\zeta _i}}{\displaystyle \frac{dL(z)}{z}}`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}p_i{\displaystyle _{z_i}^{𝔐_\alpha (z)}}{\displaystyle \frac{dL(z)}{z}}{\displaystyle \underset{i=1}{\overset{n}{}}}p_i{\displaystyle _{\zeta _i}^{𝔐_\alpha (z)}}{\displaystyle \frac{dL(z)}{z}}`$
is positive.
Hence, the average entropy change is greatest when the final state has a uniform mean $`𝔐_\alpha (z)`$. In the particular case of pure thermal convection this result can be found in Cashwell and Everett Cashwell (1967). However, their result cannot be generalized to more general processes involving deformations because the average change in the internal energy, derived from the first law, is not comparable to the average change in entropy.
The counterpart of maximum entropy in the state of uniform mean $`𝔐_\alpha (z)`$ in an $`L`$-equilbration is minimum $`L`$ in a state of a uniform mean $`𝔐_{\alpha 1}`$ that results in an $`S`$-equilibration. In an $`S`$-equilibrating transition
$`{\displaystyle \frac{1}{𝔐_{\alpha 1}(z)}}`$ $`{\displaystyle \underset{i=1}{\overset{}{}}}`$ $`p_i{\displaystyle _{z_i}^{𝔐_{\alpha 1}(z)}}𝑑L(z)`$
$`<`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}p_i{\displaystyle _{z_i}^{𝔐_{\alpha 1}(z)}}{\displaystyle \frac{dL(z)}{z}}`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{>}{}}}p_i{\displaystyle _{𝔐_{\alpha 1}(z)}^{z_i}}{\displaystyle \frac{dL(z)}{z}}`$
$`<`$ $`{\displaystyle \frac{1}{𝔐_{\alpha 1}(z)}}{\displaystyle \underset{i=1}{\overset{>}{}}}p_i{\displaystyle _{𝔐_{\alpha 1}(z)}^{z_i}}𝑑L(z),`$
implying that $`\overline{\mathrm{\Delta }L}<0`$ when $`\overline{\mathrm{\Delta }S}=0`$. The minimum property follows from the fact that the last term in
$`{\displaystyle \underset{i=1}{\overset{n}{}}}p_i{\displaystyle _{z_i}^{\zeta _i}}𝑑L(z)=`$
$`{\displaystyle \underset{i=1}{\overset{n}{}}}p_i{\displaystyle _{z_i}^{𝔐_{\alpha 1}(z)}}𝑑L(z){\displaystyle \underset{i=1}{\overset{n}{}}}p_i{\displaystyle _{\zeta _i}^{𝔐_{\alpha 1}(z)}}𝑑L(z)`$
is always negative, regardless of the initial $`n`$-tuple $`\stackrel{}{z}`$.
Hence, the mean $`𝔐_{\alpha 1}(z)`$ minimizes $`L`$. Although the $`L`$-potential is nonextensive, it nevertheless shares properties in common to a free energy in manifesting a spontaneous tendency to decrease in the presence of irreversible processes. We shall return to this in the final section of this paper.
## III Bounds on Mean Temperatures and Volumes
Carathéodory’s principle states that there are always neighboring states to any given state that are inaccessible to it by an adiabatic process, whether it be quasi-static or not Carathéodory (1909). This guarantees that there are sets of surfaces $`\psi _i(z)=\text{const}`$., where $`\psi _i(z)`$ is either $`S(z)`$ or $`L(z)`$, that do not intersect with each other. However, it does not say what those states are; this must come from another principle.
This principle states that those states which are adiabatically accessible from any given state must increase the average internal energy Buchdahl (1966), i.e.,
$`\overline{\mathrm{\Delta }E}`$ $`=`$ $`L(z){\displaystyle \underset{i=1}{\overset{n}{}}}p_i{\displaystyle _{V_i}^{V_f}}𝑑V^s`$ (5)
$`=`$ $`L(z)\left\{V_f^s{\displaystyle \underset{i=1}{\overset{n}{}}}p_iV_i^s\right\}0`$
or, equivalently,
$`\overline{\mathrm{\Delta }E}`$ $`=`$ $`S(z){\displaystyle \underset{i=1}{\overset{n}{}}}p_i{\displaystyle _{t_i}^{t_f}}𝑑t^r`$ (6)
$`=`$ $`S(z)\left\{t_f^r{\displaystyle \underset{i=1}{\overset{n}{}}}p_it_i^r\right\}0,`$
since they are linked by the adiabatic constraints $`z_i=\text{const}`$., $`i`$.
According to (5), an adiabatic transition to a state with a final volume, $`V_f`$, is possible so long as there is an increase in the volume. In the limit as $`s\mathrm{}`$, i.e., $`qr`$,
$$V_f\underset{s\mathrm{}}{lim}\left(\underset{i=1}{\overset{n}{}}p_iV_i^s\right)^{1/s}V_{\mathrm{min}},$$
where $`V_{\mathrm{min}}`$ is the volume of the smallest cell. Thus, the larger the adiabatic index, $`s`$, the more weight is given to cells of smaller size. And since the process is adiabatic, inequality (6) says that the final temperature
$$t_f\left(\underset{i=1}{\overset{n}{}}p_it_i^r\right)^{1/r},$$
cannot be inferior to the minimum mean temperature, $`𝔐_r(t)`$.
By contrast, in a process where all the cells have a common empirical temperature, $`t`$, the maximum work that is performed on the system is one where the final volume, $`V_f`$, satisfies
$`\overline{\mathrm{\Delta }L}`$ $`=`$ $`cT^\alpha {\displaystyle \underset{i=1}{\overset{n}{}}}p_i{\displaystyle _{V_i}^{V_f}}𝑑V^{\alpha s}`$
$`=`$ $`cT^\alpha \left\{V_f^{\alpha s}{\displaystyle \underset{i=1}{\overset{n}{}}}p_iV_i^{\alpha s}\right\}0,`$
since in the limit $`s\mathrm{}`$
$$V_f\left(\underset{i=1}{\overset{n}{}}p_iV_i^{\alpha s}\right)^{1/\alpha s}V_{\mathrm{max}}.$$
The upper limit to the final mean volume is the volume of the largest cell, $`V_{\mathrm{max}}`$.
Finally, in a process where no work is performed, and the cells all have the same volume, $`V_i=V`$,
$$\overline{\mathrm{\Delta }L}=cV^{\alpha s}\underset{i=1}{\overset{n}{}}p_i_{t_i}^{t_f}𝑑t^q=cV^{\alpha s}\left\{t_f^q\underset{i=1}{\overset{n}{}}p_it_i^q\right\}0.$$
Since only processes of pure heat conduction are involved, the average energy change manifests the same trend to decrease. The highest attainable temperature is one where there is an $`L`$-conserving equilibration, and
$$t_f=\left(\underset{i=1}{\overset{n}{}}p_it_i^q\right)^{1/q}.$$
The largest mean temperature $`𝔐_q(t)`$ is proportional to the internal energy, and the difference between this mean temperature and the minimum mean temperature $`𝔐_r(t)`$ is related to the system’s capability of performing work, as we shall now discuss.
## IV Metric space of power mean differences
It has long been appreciated that thermodynamic surfaces of convex energy or concave entropy lack the important topological element of a metric. There is nothing in classical thermodynamics that would play the role of a distance function and would satisfy the triangle inequality.
However, the absolute difference of power means have been shown to represent a distance, or metric, on the set of all continuous and monotonic functions in the domain $`[z_c,z_h]`$ Cargo *et. al.* (1969).
Consider the average change in the $`L`$-potential
$$\overline{\mathrm{\Delta }L}=c\left(\underset{i=1}{\overset{n}{}}p_iz_i^{\alpha 1}\right)^{\alpha /(\alpha 1)}c\underset{i=1}{\overset{n}{}}p_iz_i^\alpha ,$$
(7)
where the first term is the result of an $`S`$-equilibration. Setting $`x_i=z_i^\alpha `$, (7) can be written as the difference of two power means
$$\overline{\mathrm{\Delta }L}=c\left(\underset{i=1}{\overset{n}{}}p_ix_i^\gamma \right)^{1/\gamma }c\underset{i=1}{\overset{n}{}}p_ix_i<0,$$
where $`\gamma =(\alpha 1)/\alpha `$, and the inequality follows from (1).
Alternatively, in an $`L`$-equilibration, the average entropy increases by an amount
$$\overline{\mathrm{\Delta }S}=\frac{c}{\gamma }\left\{\left(\underset{i=1}{\overset{n}{}}p_iy_i^{1/\gamma }\right)^\gamma \underset{i=1}{\overset{n}{}}p_iy_i\right\}>0,$$
for the same reason.
The distance $`d(f,g)`$ between the generators $`f`$ and $`g`$ of the power means is defined by
$$d(f,g):=\underset{z}{sup}\left\{\right|𝔐_f(z)𝔐_g(z)|:z_czz_h\}.$$
(8)
For pure thermal conduction, the inequality
$$\left|𝔐_f(z)𝔐_g(z)\right|z_hz_c,$$
shows that the maximum distance between $`f`$ and $`g`$ is bounded by the Carnot efficiency
$$\frac{\left|𝔐_f(z)𝔐_g(z)\right|}{z_h}\eta _t.$$
(9)
All other inequalities that we will derive will be sharper, and, consequently, less efficient.
The distance (8) is obviously symmetric, and satisfies the triangle inequality
$`d(f,g)`$ $`=`$ $`\left|{\displaystyle \underset{i=1}{\overset{n}{}}}p_ix_i\left({\displaystyle \underset{i=1}{\overset{n}{}}}p_ix_i^\gamma \right)^{1/\gamma }\right|`$ (10)
$``$ $`\left|{\displaystyle \underset{i=1}{\overset{n}{}}}p_ix_i{\displaystyle \underset{i=1}{\overset{n}{}}}p_ix_i^\gamma \right|`$
$`+`$ $`\left|{\displaystyle \underset{i=1}{\overset{n}{}}}p_ix_i^\gamma \left({\displaystyle \underset{i=1}{\overset{n}{}}}p_ix_i^\gamma \right)^{1/\gamma }\right|.`$
If $`f`$ denotes the generator of the weighted arithmetic mean and $`g`$ the generator of the mean of order $`\gamma `$, then
$`\left|{\displaystyle \underset{i=1}{\overset{n}{}}}p_ix_i{\displaystyle \underset{i=1}{\overset{n}{}}}p_ix_i^\gamma \right|`$
$`{\displaystyle \underset{i=1}{\overset{n}{}}}p_ifg=fg,`$
where $``$ denotes the norm,
$$fg=\underset{x}{sup}\{|f(x)g(x)|:x_cxx_h\}.$$
Likewise, setting $`w_i=x_i^\gamma `$, the second term in (10) is bounded from above by
$`\left|{\displaystyle \underset{i=1}{\overset{n}{}}}p_iw_i\left({\displaystyle \underset{i=1}{\overset{n}{}}}p_iw_i\right)^{1/\gamma }\right|`$
$`\left|{\displaystyle \underset{i=1}{\overset{n}{}}}p_iw_i{\displaystyle \underset{i=1}{\overset{n}{}}}p_iw_i^\gamma \right|{\displaystyle \underset{i=1}{\overset{n}{}}}p_ifg=fg,`$
since $`𝔐_\gamma (w)𝔐_1(w)`$ for $`\gamma <1`$. The equality sign pertains to the case where all the $`w_i`$ are equal.
This proves that
$$\left|\overline{\mathrm{\Delta }L}\right|2cfg,$$
(11)
and the distance is induced by the norm. The properties of the metric space of equivalent classes on the set of all continuous and montonic functions have been elucidated in Cargo *et. al.* (1969). In particular, the metric space is separable since there is a countable subset everywhere dense in it, like that of the real line.
In some cases it is possible to derive analytic inequalities which are sharper than (9). If a point $`q`$ does not lie in the convex hull of the curve $`\{(f(x),g(x)):axb\}`$, then according to an extension of Carathéodory’s theorem Eggleston (1958), there exist two distinct points $`(f(X_1),g(X_1))`$ and $`(f(X_2),g(X_2))`$, where $`X_1,X_2[a,b]`$, such that the line segment joining them contains the point $`q`$. Thus, there exists positive numbers $`P_1`$ and $`P_2`$, with $`P_1+P_2=1`$, such that
$$\left|\overline{\mathrm{\Delta }L}\right|=c\left|P_1X_1+P_2X_2\left(P_1X_1^\gamma +P_2X_2^\gamma \right)^{1/\gamma }\right|.$$
As an example, consider the case $`\alpha =2`$ or $`\gamma =\frac{1}{2}`$, i.e., $`q=2`$ and $`r=1`$. Then,
$`|\overline{\mathrm{\Delta }L}|`$ $`=`$ $`c\left|P_1X_1+P_2X_2\left(P_1X_1^{1/2}+P_2X_2^{1/2}\right)^2\right|`$
$``$ $`c\left|\frac{1}{2}(a+b)\frac{1}{4}\left(\sqrt{a}+\sqrt{b}\right)^2\right|`$
$`=`$ $`c\frac{1}{4}\left|\left(\sqrt{a}\sqrt{b}\right)^2\right|=\frac{1}{4}(z_hz_c)^2,`$
which, for pure thermal conduction, becomes
$$\left|\overline{\mathrm{\Delta }L}\right|\frac{1}{4}cV^2(t_ht_c)^2=\frac{1}{4}cz_h^2\eta _t^2.$$
For processes involving pure deformations, $`|\overline{\mathrm{\Delta }L}|/L(z_h)\frac{1}{4}\eta _v^2`$, which like thermal conduction places the square of the mechanical efficiency, $`\eta _v`$, as the upper bound on the mean absolute deviation.
## V Tchebychef’s Inequality and Order Statistics
In treating the irreversible transfer of heat between any two cells in the system, say $`i`$ and $`j`$, what is transferred from $`i`$, $`dQ_i`$, is absorbed by $`j`$, $`dQ_j`$, or $`dQ_i=dQ_j`$. For any pair of heat transfers, there results
$$V_idQ_i+V_jdQ_j=(V_iV_j)dQ_i0,$$
(12)
for the $`L`$-potential, while
$$\frac{dQ_i}{T_i}+\frac{dQ_j}{T_j}=\left(\frac{1}{T_i}\frac{1}{T_j}\right)dQ_i0,$$
(13)
for the $`S`$-potential. Consequently, if there is a transfer of heat from $`ij`$, then both $`V_i<V_j`$ and $`T_i<T_j`$, meaning that the temperatures and volumes of the $`n`$-cells are similarly ordered.
Treating temperature and spatial dimension on the same level, we resort to empirical temperatures, $`t_i`$, and linear dimensions, $`R_i`$. The similar ordering of the $`n`$-tuples, $`\stackrel{}{t}`$ and $`\stackrel{}{R}`$, result in the inequality
$`𝔐_r(tR)`$ $`=`$ $`\left({\displaystyle \underset{i=1}{\overset{n}{}}}p_it_i^rR_i^r\right)^{1/r}`$ (14)
$``$ $`\left({\displaystyle \underset{i=1}{\overset{n}{}}}p_it_i^r\right)^{1/r}\left({\displaystyle \underset{i=1}{\overset{n}{}}}p_iR_i^r\right)^{1/r}`$
$`=`$ $`𝔐_r(t)𝔐_r(R),`$
with equality iff all temperatures, $`t_i`$, and linear dimensions, $`R_i`$, are equal.
Since the mean of order $`r`$ is the arithmetic mean of $`z^r`$ raised to the power $`1/r`$, i.e., $`𝔐_r(z)=𝔐_1(z^r)^{1/r}`$, it suffices to consider $`r=1`$. We then have
$`{\displaystyle \underset{i=1}{\overset{n}{}}}p_it_iR_i{\displaystyle \underset{i=1}{\overset{n}{}}}p_it_i{\displaystyle \underset{i=1}{\overset{n}{}}}p_iR_i`$
$`=`$ $`\frac{1}{2}{\displaystyle \underset{i,j=1}{\overset{n}{}}}\left\{p_ip_jt_jR_jp_ip_jt_iR_i+p_ip_jt_iR_ip_ip_jt_jR_i\right\}`$
$`=`$ $`\frac{1}{2}{\displaystyle \underset{i,j=1}{\overset{n}{}}}p_ip_j(t_it_j)(R_iR_j)0,`$
since $`\stackrel{}{t}`$ and $`\stackrel{}{R}`$ are similarly ordered. This is precisely Tchebychef’s inequality Hardy *et. al.* (1952). If $`\stackrel{}{t}`$ and $`\stackrel{}{R}`$ were oppositely ordered then the inequality in (14) would be reversed, thereby violating the second laws, (12) and (13).
In other words, the second laws assert that the averages of the product of the temperatures and linear dimensions of the cells, which are in thermal and mechanical contact, cannot be inferior to the product of their means. This is Tchebychef’s inequality (14). Even though these variables have been assumed to be independent, the second laws, (12) and (13), introduce correlations through heat transfers by similarly ordering them.
## VI From thermodynamics to multifractals and information theory
The ICG limit also allows a connection to be made with multifractals and information theory Lavenda (1998), for isothermal processes occurring in nonextensive systems.
Employing L’Hôpital’s rule we get
$$\overline{\mathrm{\Delta }S}=\mathrm{ln}\left(\underset{i=1}{\overset{n}{}}p_iR_i^{D\gamma }\right)D\gamma \underset{i=1}{\overset{n}{}}p_i\mathrm{log}R_i,$$
(15)
for the average change in entropy in the ICG limit, where we have set the exponent, $`r=D\gamma `$. The exponent $`D`$ is the Hausdorff dimension, defined as
$$\underset{i=1}{\overset{n}{}}R_i^D=1.$$
(16)
Condition (16) plays a role analogous to the Kraft (in)equality for a uniquely decipherable code.
Considering the exponent $`\gamma <1`$, we can apply Hölder’s inequality in the reverse form Hardy *et. al.* (1952)
$`\left[{\displaystyle \underset{i=1}{\overset{n}{}}}\left(p_i^{1/\gamma }R_i^D\right)^\gamma \right]^{1/\gamma }\left[{\displaystyle \underset{i=1}{\overset{n}{}}}p_i^{(1/\gamma )\gamma /(\gamma 1)}\right]^{(\gamma 1)/\gamma }`$
$`{\displaystyle \underset{i=1}{\overset{n}{}}}R_i^D=1,`$
on the strength of (16). With $`\gamma =(\alpha 1)/\alpha <1`$ and $`\alpha >1`$, Hölder’s inequality becomes
$$\underset{i=1}{\overset{n}{}}p_iR_i^{D\gamma }\left(\underset{i=1}{\overset{n}{}}p_i^\alpha \right)^{1/\alpha }.$$
(17)
It is easy to see that we have equality in (17) iff
$$R_i^D=\frac{p_i^\alpha }{_{i=1}^np_i^\alpha },$$
or
$$D\mathrm{ln}R_i=\alpha \mathrm{log}p_i\mathrm{ln}\underset{i=1}{\overset{n}{}}p_i^\alpha ,$$
(18)
which also satisfies the definition of the Hausdorff dimension, (16). Multiplying (18) through by $`p_i`$, summing, and introducing the result into (15) give
$$\overline{\mathrm{\Delta }S}=(\alpha 1)\left(S_1S_\alpha \right)>0.$$
(19)
The entropies, $`S_1`$ and $`S_\alpha `$, are the Shannon-Gibbs,
$$S_1=\underset{i=1}{\overset{n}{}}p_i\mathrm{log}p_i,$$
and Rényi,
$$S_\alpha =\frac{1}{1\alpha }\mathrm{ln}\underset{i=1}{\overset{n}{}}p_i^\alpha ,$$
entropies of order $`1`$ and $`\alpha `$, respectively.
Inequality (19) is result of the fact that for $`\alpha >1`$, the Shannon entropy is greater than the Rényi entropy, while, for $`\alpha <1`$, the converse is true. This can easily be seen by setting $`\beta =\alpha 1`$. Then, on the strength of (1) we have
$$\left(\underset{i=1}{\overset{n}{}}p_ip_i^\beta \right)^{1/\beta }>\underset{i=1}{\overset{n}{}}p_i^{p_i},$$
for $`\beta >0`$, and the reverse inequality for $`\beta <0`$. Hence, inequality (19) is always satisfied. Moreover, in the limit as $`\alpha 1`$, l’Hôpital rule shows that
$$S_1=\underset{\alpha 1}{lim}S_\alpha =\underset{i=1}{\overset{n}{}}p_i\mathrm{log}p_i.$$
In this limit, the average entropy difference, (19), vanishes.
Consequently, (19) shows that in the limit of an isothermal ICG, the average entropy difference is always proportional to the absolute value of the difference between the Shannon and Rényi entropies, when $`D`$ is identified as the Hausdorff dimension.
The generalization of the Hausdorff dimension to multifractals, where the generator of cell sizes of lengths $`R_i`$ with probabilities, $`p_i`$, require *two* exponents Halsey *et. al.* (1986),
$$\underset{i=1}{\overset{n}{}}p_i^\alpha R_i^{D_\alpha (1\alpha )}=1,$$
(20)
where $`D_\alpha `$ is supposed to be some generalization of the Hausdorff dimension, $`D`$.
If $`\{p_i\}`$ is a complete probability distribution, and $`\alpha `$ is restricted to the open interval $`(0,1)`$, in order to ensure that the Rényi entropy be concave, then the usual Hölder inequality, and condition (20), give
$$\underset{i=1}{\overset{n}{}}R_i^{D_\alpha }1.$$
(21)
Since for $`\alpha =1`$, (18) becomes
$$D_1=\frac{S_1}{_{i=1}^np_i\mathrm{log}(1/R_i)},$$
it was thought Halsey *et. al.* (1986) that $`D_\alpha `$ should be related to the Rényi entropy in a similar form, viz.,
$$D_\alpha =\frac{S_\alpha }{_{i=1}^np_i\mathrm{log}(1/R_i)}.$$
This can be derived by averaging
$$D_\alpha (1\alpha )\mathrm{ln}R_i=\mathrm{ln}\underset{i=1}{\overset{n}{}}p_i^\alpha .$$
(22)
Exponentiating both sides, multiplying by $`p_i^\alpha `$, and summing do give (20). But, since the right side of (22) is independent of the index $`i`$, so too must be the left side. This means that all the cell sizes have the same length
$$R^{D_\alpha }=\left(\underset{i=1}{\overset{n}{}}p_i^\alpha \right)^{1/(1\alpha )}=e^{S_\alpha }.$$
In view of condition (21), this would imply
$$S_0S_\alpha .$$
The entropy is largest in either the state of equal probabilities, or in the state of greatest geometrical regularity. This entropy is the Hartley-Boltzmann entropy, $`S_0=\mathrm{ln}n`$, which depends on the number of copies considered, and not on their individual frequencies.
## VII A measure of the tendency to uniformity
In this section we show that the internal energy, or the $`L`$-potential, is a Schur convex function of the empirical temperature and serves as a measure of the tendency of the system to evolve toward a more uniform distribution in temperature.
There is no reason to exclude the possibility that the cells in the final state of thermal equilibrium will have different probabilities to be at a given temperature than those in the initial state at the moment the cells have been placed in thermal contact. Let us therefore introduce a second complete set of probabilities $`q_1,\mathrm{},q_n`$, which are the probabilities that the final temperatures of the cells will be $`\tau _1,\mathrm{},\tau _n`$. If $`\stackrel{~}{t}`$ is any intermediate temperature, we can write the average change in the $`L`$-potential density, $`\mathrm{}=L/V`$ as
$`\overline{\mathrm{\Delta }\mathrm{}}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}p_i{\displaystyle _{t_i}^{\stackrel{~}{t}}}𝑑\mathrm{}(t)+{\displaystyle \underset{i=1}{\overset{n}{}}}q_i{\displaystyle _{\stackrel{~}{t}}^{\tau _i}}𝑑\mathrm{}(t)`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}q_i\mathrm{}(\tau _i){\displaystyle \underset{i=1}{\overset{n}{}}}p_i\mathrm{}(t_i).`$
The process of placing the cells in thermal contact will initiate a process of heat exchange that will ultimately lead to a more uniform temperature distribution. This process of homogenization can be represented mathematically by a doubly stochastic matrix, $`\left(\omega _{ij}\right)`$, whose rows and columns sum to unity Mirsky (1963).
The doubly stochastic matrix relates the set of initial temperatures to the set of final temperatures according to
$$\tau _i=\underset{j=1}{\overset{n}{}}\omega _{ij}t_j,$$
(23)
for $`i=1,\mathrm{},n`$. In other words, the doubly stochastic matrix represents a smoothing operation, and, if the temperatures are ordered in either an increasing or decreasing order, the restriction $`\omega _{ij}[0,\frac{1}{2}]`$ will preserve that order.
The same doubly stochastic matrix, $`(\omega _{ij})`$, relates the new probability distribution $`\{q_i\}`$ to the old one, $`\{p_i\}`$, in an inverse fashion
$$p_j=\underset{i=1}{\overset{n}{}}\omega _{ij}q_i,$$
(24)
for $`j=1,\mathrm{},n`$, to that relating the old to the new temperatures, (23). Multiplying (23) by $`q_i`$, and summing give
$$\underset{i=1}{\overset{n}{}}q_i\tau _i=\underset{i=1}{\overset{n}{}}q_i\underset{j=1}{\overset{n}{}}\omega _{ij}t_j=\underset{j=1}{\overset{n}{}}p_jt_j.$$
This says that the system has the same average temperature before and after the cells have been placed in thermal contact. No heat has been transferred between the system and the environment and no work has been performed so that the same average temperature persists.
Since $`\mathrm{}`$ is Schur convex,
$$\mathrm{}(\tau _i)=\mathrm{}\left(\underset{i=1}{\overset{n}{}}\omega _{ij}t_j\right)\underset{j=1}{\overset{n}{}}\omega _{ij}\mathrm{}(t_j),$$
(25)
where the first equality follows from (23). This is none other than Jensen’s inequality stating that for a convex function: the function of the average can never be greater than the average of the function.
Multiplying both sides of (25) by $`q_i`$ and summing result in
$$\underset{i=1}{\overset{n}{}}q_i\mathrm{}(\tau _i)\underset{i=1}{\overset{n}{}}q_i\underset{j=1}{\overset{n}{}}\omega _{ij}\mathrm{}(t_j)=\underset{j=1}{\overset{n}{}}p_j\mathrm{}(t_j),$$
(26)
on account of (24). This proves that the $`L`$-potential shows a tendency to decrease in the presence of irreversible processes of heat transfer \[vid. (11)\].
Inequality (26) can be most easily established in the case $`n=2`$. Since the two temperature distributions are related by the doubly stochastic matrix, which in the present case is
$$\omega =\left(\begin{array}{cc}\overline{\omega }& \omega \\ \omega & \overline{\omega }\end{array}\right)$$
where $`\overline{\omega }=1\omega `$ for some $`\omega [0,\frac{1}{2}]`$, $`\stackrel{}{t}`$ is said to majorize $`\stackrel{}{\tau }`$, or $`\stackrel{}{t}\stackrel{}{\tau }`$ Marshall *et. al* (1979).
Because $`\mathrm{}`$ is Schur convex
$`q_1\mathrm{}(\tau _1)+q_2\mathrm{}(\tau _2)`$ $`=`$ $`q_1\mathrm{}(\overline{\omega }t_1+\omega (t_2)+q_2\mathrm{}(\omega t_1+\overline{\omega }t_2)`$
$``$ $`q_1[\overline{\omega }\mathrm{}(t_1)+\omega \mathrm{}(t_2)]`$
$`+`$ $`q_2[\omega \mathrm{}(t_1)+\overline{\omega }u(t_2)]`$
$`=`$ $`p_1\mathrm{}(t_1)+p_2\mathrm{}(t_2).`$
When the adiabatic walls between the cells have been replaced by diathermal ones, and the system is left to itself, heat will spontaneously flow from the hotter to the colder cells. In economic jargon these flows can be considered as ‘Robin Hood’ operations Arnold (1986), where there is a transfer of riches from the wealthy to the poorer segments of the population. This transfer is regulated by the doubly stochastic matrix $`\left(\omega _{ij}\right)`$.
The condition that the final, average temperature persist after thermal contact has been made is the same as saying that the riches have remained the same except they have been spread out more equally. If the gas were ICG, the change in the internal energy would vanish since it is linear in the temperature, and the average temperature has not changed. However, due to the property of Schur convexity of the internal energy, or the $`L`$-potential, as a function of temperature for an IGG, they show a net tendency to decrease. This property reflects the tendency of the system to reach a more uniform distribution in temperature. In this respect, an IGG offers a more realistic description of Nature than an ICG. |
warning/0507/astro-ph0507042.html | ar5iv | text | # Detailed Chemical Abundances of Extragalactic Globular Clusters
## 1. Introduction
We are developing a method which will let us measure detailed chemical abundances of Globular Clusters (GCs) from high–resolution, echelle spectra of their integrated light. Our goal is to use the detailed abundance patterns of GC systems in galaxies as distant as 4–5 Mpc to constrain their formation and chemical enrichment histories in the same way as the abundance patterns of old stars have been used to do so in the Milky Way (e.g. Edvardsson et al. 1993; McWilliam et al. 1995; Prochaska et al. 2003; Bensby et al. 2004ab; Fulbright, McWilliam & Rich 2005).
The elements which are the most useful for understanding the processes of galaxy evolution are those which are returned to the ISM primarily through either high–mass, core–collapse SNe (Type II = SNII) or through low–mass, accretion–induced SNe (Type Ia = SNIa). The key is that the former (SNII) evolve on time–scales of order $`10^6`$ years, while the latter (SNIa) generally take of order $`10^9`$ years. Elements produced by SNII — “$`\alpha `$–elements” (e.g. Si, S, Ca) and r–process elements (e.g. Eu) — will therefore build up rapidly over the first few megayears of star formation or in a starburst. Elements produced in SNIa and SNII — Fe and Fe–peak elements (e.g. Sc, V, Cr, Mn, Fe, Co, Ni) — build up over many gigayears as the contribution from SNIa increases. The ratio \[$`\alpha `$/Fe\] relative to \[Fe/H\] is therefore a rough constraint on the star formation rate. Comparison with overall abundance then constrains the timescale of formation.
Unfortunately, the Milky Way is currently the only galaxy for which even a partial history of formation and enrichment can be traced in the stellar fossil record through to the present day. This is because detailed abundance analysis requires high signal–to–noise ratios and high spectral resolutions (S/N$`60`$, $`R>20,000`$). Old stars are too faint to be studied in this way beyond the nearest dwarf spheroidal members of the Local Group ($`100`$ kpc). Even in the nearest gas–rich irregulars of the Local Group, abundance measurements are limited to a few very luminous supergiants (e.g. Venn et al. 2004), and results can only be obtained for a few elements. Moreover, these bright supergiants are very massive (young) and so only identify the recent gas composition. The history of a galaxy can only be learned from long–lived, low–mass stars, which are too faint for high–resolution spectroscopy beyond the nearest dwarf spheroidal galaxies (e.g. Shetrone et al. 2003).
These limitations drive us to target Globular clusters (GCs). Unlike single stars, high–resolution spectra *can* be obtained of unresolved GCs out to $`4`$ Mpc with *current* telescopes. GCs are bright enough ($`10<M_v<6`$ mag) and have low enough velocity dispersions ($`220`$ km/s) that even weak lines ($`15`$ mÅ) could be detected in spectra of their integrated light. Moreover, several lines of evidence from photometry and low–resolution spectroscopy already suggest that GCs trace the star formation and the global formation history of their parent galaxy: the number of GCs per galaxy is roughly constant relative to the total light and mass of a galaxy; young GCs are found in regions of active star formation (e.g. Schweizer & Seitzer 1993, Barth et al. 1995); and the properties of both metal–rich (red) and now also metal–poor (blue) GC systems seem to correlate with the properties of their parent galaxies (e.g. Zepf & Ashman 1993, Geisler, Lee & Kim 1996; Gebhardt & Kissler–Patig 1999; Kundu & Whitmore 2001; Larsen et al. 2001; Strader, Brodie, & Forbes 2004). All of which implies that GCs are bright, observable markers of the chemical enrichment record of normal galaxies.
Of course, significant progress has also been made in understanding the chemical enrichment histories of elliptical galaxies and bulges from analysis of low–resolution, integrated–light (IL) spectra using line indexes, such as the Lick system (Burstein et al. 1984, Faber et al. 1985, Worthey et al. 1994, Worthey & Ottaviani 1997, Trager et al. 1998). These indexes provide estimates of age and composite “metallicity” (“Z”) based on lines from multiple elements such as Fe, Mg, and Ca (see Rose 1984, Worthey et al. 1994, and Worthey & Ottaviani 1997). GCs are ideal targets for these index systems as they are nearly ideal, single–age stellar populations.
A great deal of progress has recently been made with line index systems by using various kinds of principle component analysis to obtain better leverage on Z and also to obtain an “E” parameter, which includes all of the elements (O, Mg, C, etc) typically found to be enhanced in elliptical and bulge (rapidly forming) populations (Proctor & Sansom 2002; Proctor, Forbes, Beasley 2004; Strader & Brodie 2004). However, even for GCs, information from line indexes are limited by two difficulties. First, E and Z do not have the interpretive power of detailed abundances because they include elements which form in multiple sites and include elements which are known to be affected by GC self–enrichment (see discussions in Gratton, Sneden & Carretta 2004). These indexes do not entirely isolate the \[$`\alpha `$/Fe\] ratios which constrain formation timescales. Second, line index results are very sensitive to the calibration of the system, and it can be hard to obtain stellar libraries of the appropriate age and enrichment or to know what calibration to use. This complication has been widely discussed in the literature (e.g. Trippico & Bell 1995; Trager et al. 2000ab; Thomas, Maraston, & Bender 2003; Proctor et al. 2004; Tantalo & Chiosi 2004). Adjustments to the calibration system have been published recently based on stellar models which include a range of enrichments. However, not all of these agree with each other and the models themselves are not well tested against observed populations. This concern is particularly relevant given recent work on M31, which shows that the Milky Way GC system may not be generally representative of GCs in even spiral galaxies; M31 appears to contain a disk GC system (Morrison et al. 2004) which is at least partly composed of young GCs (0.1–0.8 Gyrs; Barmby et al. 2000, Beasley et al. 2004) with different abundance patterns than are seen in the Milky Way (Burstein et al. 2004).
Fortunately, it *is* possible to obtain high–resolution spectra of unresolved GC. A typical GC has velocity dispersions of $`2<\sigma _v<20`$ km/s (Pryor & Meylan 1993, see Figure 1), so that spectra of their integrated light can have line widths in the range $`60,000>\lambda /\mathrm{\Delta }\lambda >6,500`$. For comparison, an elliptical has $`\sigma _v200`$ km/s or greater and limiting line widths of $`\lambda /\mathrm{\Delta }\lambda 1500`$. Absorption features as weak as a $`15`$ mÅ can therefore be seen in high–resolution, integrated light spectra of all but the most massive GCs with S/N$`6090`$.
Detailed abundances have never been obtained for unresolved GCs because current methods of abundance analysis only work for individual stars. We describe here progress that we have made in adapting the basic techniques of single star abundance analysis to spectra of integrated–light.
## 2. A Method for Light–Weighted Abundance Analysis
Abundances are not directly observable quantities; the strength of any given absorption line is a function not only of the abundance of the element, but also of the physical properties (e.g. mass and temperature) of the star. The method we are developing for analyzing IL spectra is based on the standard techniques used to analyze individual stars.
The standard technique for abundance analysis of single stars involves comparing the observed equivalent widths (EWs) of lines for each species with those predicted by model stellar atmospheres and spectral line synthesis. The model atmospheres have the following, observationally constrained parameters: effective temperature ($`T_{\mathrm{eff}}`$, from $`BV`$), specific gravity ($`\mathrm{log}g`$, from luminosity), and microturbulence ($`\xi `$). In principle, the only free parameter is abundance, \[A/H\]. So far, we have used Kurucz models in which Fe alone is varied, so that \[A/H\] is roughly \[Fe/H\]. The Kurucz model atmosphere grids (available from R. Kurucz at http://cfaku5.harvard.edu) give optical depth, temperature, pressure, and electron density in each of 64 layers. The EW of any particular spectral line is then found by radiative transfer through these layers. Given the model atmosphere, the EW of a line is uniquely determined by the wavelength ($`\lambda `$), excitation potential (EP), and gf value of the line, and by the abundance of the element in question, $`[X/H]`$. Of these, only \[X/H\] is adjusted. In standard analysis, the best model atmosphere is found by adjusting \[Fe/H\] in the line synthesis and model atmospheres until a unique, self–consistent solution is found from all Fe I and Fe II lines ($`100`$ lines). A strong constraint on the correct model atmosphere comes from requiring that Fe I lines over the full range of $`\lambda `$, EP, and EW give a stable Fe solution. This requirement is used to adjust $`T_{\mathrm{eff}}`$, $`\mathrm{log}g`$, and $`\xi `$ to find the appropriate stellar atmosphere model. Lines from any other element are then uniquely determined by the input abundance of that element (\[X/H\]) in the line synthesis step. The right abundance is found by adjusting this value until predicted EWs match those observed.
To interpret an IL spectrum, we have developed an original method of detailed chemical abundance analysis which involves producing a light–weighted EW, building on those techniques used for individual stars. For a resolved Galactic GC, the observed CMD tells us the exact fraction of light coming from every stellar type in the cluster. For NGC 6397, for instance, Figure 2 shows the CMD of NGC 6397 with boxes indicating groups of stars which each contain roughly 5% of the cluster’s light. Based on this observed CMD, we can compute a model atmosphere for each box using the Kurucz model atmospheres (ATLAS9) for a range in abundance, \[A/H\]. We have developed code which then synthesizes the individual Fe lines by iteratively calling the spectrum synthesis program MOOG (Sneden 1974) and varying the Fe in the line synthesis only until consistent solutions are obtained. With this Fe abundance, we can then compute a final stellar atmosphere for each box. We then identify the abundances of other elements by adjusting \[X/H\] (in the line synthesis only) to match the observed EWs for a species. In this analysis, a microturbulence scaling law is adopted ($`\xi \mathrm{log}gf`$).
As a first test of this basic strategy, we have analyzed the integrated light spectrum of NGC 6397 using its observed CMD (see Figure 2) to define the contributing stellar atmosphere models and their weights. From this analysis, we successfully obtain abundances for Fe–peak, $`\alpha `$–elements, and light elements which are in good agreement with the results from other groups for single stars in NGC 6397 (see Bernstein & McWilliam 2001 and Bernstein & McWilliam 2005). However, when we target extragalactic GCs, we will not be able to use resolved CMDs to empirically identify the right mix of stellar populations. We must therefore learn to use theoretical CMDs — isochrones based on stellar evolution tracks — to analyze the IL spectra. To do this, and to understand the limitations and impact of stellar evolution models themselves, we have observed a set of GCs in the Milky Way and LMC to serve as a “training set”.
## 3. A Training Set of Galactic and LMC GCs
Our “training set” consists of seven Milky Way and eight LMC GCs. These span the range of velocity dispersions, abundances, ages, and HB morphologies available in the Milky Way and LMC systems (see Figure 1 and Table LABEL:table:tset). For each of the training set GCs, we have obtained integrated light spectra by uniformly scanning the central $`32\times 32`$ arcsec<sup>2</sup> and $`12\times 12`$ arcsec<sup>2</sup> in the Milky Way and LMC GCs, respectively. Crucially, these clusters are all spatially resolved, so detailed color magnitude diagrams (CMDs) can provide age estimates and *a priori* knowledge of the member stellar populations, including HB morphology. We can use this information to refine our methods and determine how such variables affect our results. Moreover, we can also identify individual red giant branch (RGB) stars from these CMDs and use spectra of these individual stars to obtain “fiducial” abundances by standard methods which are consistent with our own atmospheric modeling and synthesis for the integrated light. Below we describe the results we have obtained for one Milky Way GC: NGC 104 (47 Tuc).
To analyze 47 Tuc as we would an unresolved GC, we use isochrones as the template for its stellar population. The analysis here uses isochrones based on the stellar evolution tracks of the Padova group (Girardi et al. 2000) and a Kroupa IMF (Kroupa & Boily 2002). We complete the same analysis for isochrones with a range of age and abundance (1–16 Gyrs, $`2.3<`$\[Fe/H\]$`<0.2`$). The first question one might ask is whether the resulting $`15`$ Gyr, \[Fe/H\]$`0.7`$ isochrone *looks* like the observed CMD of 47 Tuc. Two differences are evident, although not unexpected (e.g. Bergbusch & VandenBerg 2001): the isochrones include lower mass stars than are observed in the clusters and do not reproduce the observed AGB bump. The same problems exist with the BASTI (Pietrinferni et al. 2004) stellar evolution models. To begin, we use the Padova isochrones “as is,” with no adjustments for mass segregation or the AGB bump.
For each isochrone, we split the stellar population into roughly 20 groups each containing roughly 5% of the light, similar to the boxes shown in Figure 2. We then produce a stellar atmosphere for each group and synthesize a light–weighted EW for about 70 Fe I and 10 Fe II lines by combining the synthesized EWs from MOOG for each group, as described in §2. Note that with the isochrones, the parameters of the atmosphere model are completely dictated by the parameters of the isochrone. We adjust only the abundance of Fe in the line synthesis step to obtain a good match between observed and predicted EWs for all Fe I and Fe II line. Figure 3 shows the best–match Fe value for Fe I lines as a function of excitation potential for three different isochrones. The average of these solutions gives us the inferred \[Fe/H\] value (and a statistical uncertainty) for each isochrone. Preferred solutions have small rms scatter and no slope with excitation potential, wavelength, and EW. We then do the same test for Fe II lines. A plot showing the inferred value of \[Fe/H\] from Fe I and Fe II lines for each isochrone is shown in Figure 4. Surprisingly, we find that best–fit \[Fe/H\] solution from Fe I lines is nearly independent of the input isochrone metallicity, while Fe II lines are nearly independent of the input age. Together, the Fe I and Fe II lines provide an excellent constraint on the abundance of the cluster and a way to break the age–metallicity degeneracy that will never be achievable at low resolution.
From the best–fit isochrone, we get \[Fe/H\]=$`0.63\pm 0.03`$ ($`1\sigma `$ statistical uncertainty). We then adjust the luminosity function of one isochrone to remove low–mass stars and add an AGB bump. With this adjusted isochrone, we find \[Fe/H\]$`=0.73`$. Note that these results differ by only 0.1 dex, and that both results are roughly within the range of values obtained recently by different groups from individual stars (Kraft & Ivans 2003, Carretta et al. 2004). Nevertheless, the change in our \[Fe/H\] solution with these adjustments emphasizes the importance of the luminosity function to our results. We will use the rest of the training set to determine the best approach to dealing with these differences between isochrones and real CMDs.
We use the adjusted isochrone to then obtain abundances for a broad range of key elements (see Table LABEL:table:47tuc). Our final abundances are in good agreement with those in the literature for individual 47 Tuc stars. Moreover, we find the expected abundances patterns $`\alpha `$– (enriched relative to solar), Fe–peak (roughly solar), and r–process (enriched) elements that one would expect for an old globular cluster. We also detect enrichment of Na and Al, which is a consistent with self–enrichment through proton–burning in AGB stars and is found for some individual stars (with star–to–star variations) in 47 Tuc (Carretta et al. 2004). This suggests that we will be able to measure such abundance variations in unresolved clusters even if present in only a fraction of the stars in a given cluster! Finally, we note that the age of 47 Tuc is only constrained from our analysis to be greater than $`3`$ Gyrs. We may be able to improve this by calibrating the $`\mathrm{log}gf`$ values of individual Fe lines to the abundance solution for Arcturus; doing so should reduce the scatter in plots like those shown in Figure 3 to $`\pm 0.05`$, as it does in the analysis of individual stars (see Bernstein & McWilliam 2005). However, it is also clear from evolutionary tracks and observed CMDs of Milky Way GCs that stellar populations themselves simply do not vary much with age after a few Gyrs (e.g. Yi et al. 2003).
## 4. Further Spectroscopic Constraints on Integrated Light Abundances
An important issue to explore is the influence of Horizontal Branch (HB) morphology on our abundance analysis because it is not uniquely correlated with age or metallicity and therefore not reliably characterized by the isochrones. Balmer line EWs *and profiles* will be helpful in this regard, because they are very sensitive to the light fraction contributed by hot stars (e.g. Bernstein & McWilliam 2002; Schiavon et al. 2004; Bernstein & McWilliam 2005;). To this end, we have compared observed Balmer profiles for several of our training set clusters with synthesized, light–weighted Balmer profiles based on the isochrones. Our preliminary work suggests that we will need to also synthesize blended lines in the Balmer line wings in order to get accurate Balmer EWs and profiles, particularly in the metal rich clusters. This may be a particular problem for low–resolution spectra. Similar issues have already been noted in the literature to the extent that Balmer line indexes in low–resolution spectra seem to be sensitive to the calibration of the line system (e.g. Proctor, Forbes & Beasley 2004). Additional constraints may also come from CaII lines, which have also been shown to track HB color in low–resolution spectra (e.g. Burstein et al. 2004; Proctor, Forbes & Beasley 2004). Independent models for horizontal branch morphology can also be combined with the isochrone models directly to explore the impact on inferred abundances
We are also planning to explore a broader range of parameters in the isochrones themselves as more models become publicly available. For example, stellar evolution tracks with different abundance ratios (e.g. enriched in $`\alpha `$–elements or with different helium fractions) may yield systematically different abundances and should also be tested against the training set.
## 5. Applications
It is already possible to obtain spectra of extragalactic GCs that can be analyzed with this technique. In the local group, abundances of the more distant LMC GCs could be obtained in a few hours with 4–m class telescopes taking IL spectra, instead of requiring many hours of exposure time on 8–m class telescopes observing individual cluster stars. Beyond the LMC, an handful of galaxies within 4–5 Mpc can be observed using using current 6.5–10 m telescopes. We have already obtained a few spectra ($`R20,000`$, S/N$`=60900`$ at $`H\alpha `$) of confirmed GCs in galaxies NGC 1313 and NGC 5128 (both $`4`$ Mpc away) with the MIKE echelle spectrograph (Bernstein et al. 2003) on the Magellan Telescopes. In the future, hundreds of galaxies in the local universe could be observed with the next generation of ground–based 20–30 m telescopes. In addition, the abundances and integrated light spectra we already have in hand can be used to better understand and calibrate the low–resolution data and line–index systems. This is important because line indexes will always be able to reach galaxies at greater distances than the high–resolution technique described here.
### Acknowledgments.
We gratefully acknowledge the help of Las Campanas technical staff and particularly telescope operators Fernando Peralta and Herman Olivares. R. A. B. acknowledges partial support from NASA through Hubble Fellowship grant HF–01088.01–97A awarded by Space Telescope Science Institute, which is operated for NASA by the Association of Universities for Research in Astronomy, Inc. under contract NAS 5–2655. A. M. acknowledges support from NSF grants AST–96–18623 and AST–00–98612. We also thank the organizers for a very enjoyable and stimulating conference. |
warning/0507/hep-ph0507295.html | ar5iv | text | # Factorization and polarization in two charmed-meson B decays
## I Introduction
The study of $`B`$ meson weak decays is of high interest in heavy flavor physics and CP violation. In particular, much attention has been paid to the two-body charmless hadronic $`B`$ decays, but there are relatively less discussions on the decays with charmful mesons, such as the modes with two charmed-meson in which the final states are both heavy. However, the two charmed-meson decays can provide some valuable and unique information which is different from the light meson productions. For example, CP asymmetries in the decays of $`BD^{()+}D^{()}`$ play important roles in testing the consistency of the standard model (SM) as well as exploring new physics SX . Moreover, these decays are ideal modes to check the factorization hypothesis as the phenomenon of color transparency for the light energetic hadron is not applicable. Since the decay branching ratios (BRs), CP asymmetries (CPAs) and polarizations of $`BD^{()}D^{()}(D^{()}D_s^{()})`$ have been partially observed in experiments PDG04 ; BELLE\_DD ; BABAR\_DD , it is timely to examine these heavy-heavy $`B`$ decays in more detail.
At the quark level, one concludes that the two charmed-meson decays are dominated by tree contributions since the corresponding inclusive modes are $`bqc\overline{c}`$ with $`q=s`$ and $`d`$. It is known that the factorization has been tested to be successful in the usual color-allowed processes. However, the mechanism of factorization in heavy-heavy decays is not the same as the case of the light hadron productions. The color transparency argument Bjorken for light energetic hadrons is no longer valid to the modes with heavy-heavy final states. The reason can be given as follows. Due to the intrinsic soft dynamics in the charmed-meson, non-vanishing soft gluon contributions are involved in the strong interactions between an emission heavy meson with the remained $`BD^{()}`$ system. Since the corresponding divergences may not be absorbed in the definition of the hadronic form factor or hadron wave function, the decoupling of soft divergences is broken. This means that the mechanism of factorization has to be beyond the perturbative frameworks, such as QCD factorization BBNS and soft-collinear effective theory SCET . The large $`N_c`$ limit is another mechanism to justify factorization BGR , corresponding to the effective color number $`N_c=\mathrm{}`$ in the naive factorization approach BSW . The understanding of factorization in heavy-heavy decays requires some quantitative knowledge of non-perturbative physics which is not under control in theory. In this paper, we will assume the factorization hypothesis and apply the generalized factorization approach (GFA) GFA1 ; GFA2 to calculate the hadronic matrix elements.
It is known that annihilation contributions and nonfactorizable effects with final state interactions (FSIs) play important role during the light meson production in $`B`$ meson decays. For instance, to get large strong phases for CP asymmetries (CPAs) in $`BK\pi `$ and $`B\pi \pi `$ decays, these effects are included inevitably KLS\_PRD63 ; CCS\_PRD71 . Moreover, they are also crucial to explain the anomaly of the polarizations in $`B\varphi K^{}`$ decays, measured by BABAR babar\_pol and BELLE belle\_pol recently. By the naive analysis in flavor diagrams, one can easily see that the decay modes of $`BD^{()0}\overline{D}^{()0}`$ and $`D_s^{()0}\overline{D}_s^{()0}`$ are annihilation-dominated processes. Therefore, measurements of these decays will clearly tell us the importance of annihilation contributions in the production of two charmed-meson modes.
For the color-allowed decays, since the short-distant (SD) nonfactorizable parts are associated with the Wilson coefficient (WC) of $`C_1/N_c`$, where $`C_1`$ is induced by the gluon-loop and it is much smaller than $`C_21`$, we can see that the effects arising from the SD non-spectator contributions should be small LM . Nevertheless, long-distant (LD) nonfactorizable contributions governed by rescattering effects or FSIs may not be negligible. Inspired by the anomaly of the large transverse perpendicular polarization, denoted by $`R_{}`$, in $`B\varphi K^{}`$ decays, if there exist significant LD effects, we believe that large values of $`R_{}`$ may appear in $`BD^{}D^{}`$ and $`BD^{}D_s^{}`$ too. As we will discuss, the power-law in the two-vector charmed-meson decay leads to a small $`R_{}`$. The recent measurement of the polarization fraction by the BELLE collaboration gives $`R_{}=0.19\pm 0.08\pm 0.01`$ BELLE\_DD in which the central value is about a factor of three comparing with the model-independent prediction within the factorization approach and heavy quark symmetry. Clearly, to get the implication from the data, we need a detailed analysis in two charmful final states of $`B`$ decays.
To estimate the relevant hadronic effects for two-body decays in the $`B`$ system, we use the GFA, in which the leading effects are factorized parts, while the nonfactorized effects are lumped and characterized by the effective number of colors, denoted by $`N_c^{\mathrm{eff}}`$. Note that the scale and scheme dependences in effective WCs $`C_i^{\mathrm{eff}}`$ are insensitive. In the framework of the GFA, since the formulas for decay amplitudes are associated with the transition form factors, we consider them based on heavy quark effective theory (HQET) NeubertHQET . We will also study their $`\alpha _s`$ NeubertQCD and power corrections which break heavy quark symmetry (HQS) NR . In our analysis, we will try to find out the relationship between the HQS and its breaking effects for $`R_{}`$. In addition, we will reexamine the influence of penguin effects, neglected in the literature LR . We will show that sizable CPAs in $`\overline{B}^0D^+D^{}`$ and $`B^{}D^0D^{}`$ decays may rely on FSIs.
This paper is organized as follows. In Sec. II, we give the effective Hamiltonian for the heavy-heavy B decays. The definition of heavy-to-heavy form factors are also introduced. In Sec. III, we show the general formulas for B to two charmful states in the framework of the GFA. The effects of the heavy quark symmetry breaking on the transverse perpendicular polarization are investigated. In Sec. IV, we provide the numerical predictions on the BRs, direct CPAs and the polarization fractions. Conclusions are given in Sec. V. We collect all factorized amplitudes for $`BPP,PV(VP)`$ and $`VV`$ decays in Appendixes.
## II Effective interactions and parametrization of form factors
The relevant effective Hamiltonian for the $`B`$ meson decaying to two charmful meson states is given by BBL ,
$`H_{\mathrm{eff}}(\mathrm{\Delta }B=1)`$ $`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}\{V_{ub}V_{uq}^{}[C_1(\mu )O_1^u+C_2(\mu )O_2^u]+V_{cb}V_{cq}^{}[C_1(\mu )O_1^c+C_2(\mu )O_2^c]`$ (1)
$`V_{tb}V_{tq}^{}{\displaystyle \underset{k=3}{\overset{10}{}}}C_k(\mu )O_k\}+H.c.,`$
where $`q=s`$ and $`d`$, $`V_{ij}`$ denote the Cabibbo-Kobayashi-Masikawa (CKM) matrix elements, $`C_i(\mu )`$ are the Wilson coefficients (WCs) and $`O_i`$ are the four-fermion operators, given by
$`O_1^u=(\overline{q}_iu_j)_{VA}(\overline{u}_jb_i)_{VA},O_2^u=(\overline{q}_iu_i)_{VA}(\overline{u}_jb_j)_{VA},`$
$`O_1^c=(\overline{q}_ic_j)_{VA}(\overline{c}_jb_i)_{VA},O_2^c=(\overline{q}_ic_i)_{VA}(\overline{c}_jb_j)_{VA},`$
$`O_{3(5)}=(\overline{q}_ib_i)_{VA}{\displaystyle \underset{q^{}}{}}(\overline{q}_j^{}q_j^{})_{VA},O_{4(6)}=(\overline{q}_ib_j)_{VA}{\displaystyle \underset{q^{}}{}}(\overline{q}_j^{}q_i^{})_{VA},`$
$`O_{7(9)}={\displaystyle \frac{3}{2}}(\overline{q}_ib_i)_{VA}{\displaystyle \underset{q^{}}{}}e_q^{}(\overline{q}_j^{}q_j^{})_{V\pm A},O_{8(10)}={\displaystyle \frac{3}{2}}(\overline{q}_ib_j)_{VA}{\displaystyle \underset{q^{}}{}}e_q^{}(\overline{q}_j^{}q_i^{})_{V\pm A}.`$ (2)
with $`i`$ and $`j`$ being the color indices, $`O_{36}`$ ($`O_{710}`$) the QCD (electroweak) penguin operators and $`(\overline{q}_1q_2)_{V\pm A}=\overline{q}_1\gamma _\mu (1\pm \gamma _5)q_2`$. In order to cancel the renormalization scale and scheme dependence in the WCs of $`C_i(\mu )`$, the effective WCs are introduced by
$`C(\mu )O(\mu )C^{\mathrm{eff}}O_{\mathrm{tree}}.`$ (3)
Since the matrix element $`O_{\mathrm{tree}}`$ is given at tree level, the effective WCs are renormalization scale and scheme independent. To be more useful, we can define the effective WCs as
$`a_1^{\mathrm{eff}}`$ $`=`$ $`C_2^{\mathrm{eff}}+{\displaystyle \frac{C_1^{\mathrm{eff}}}{(N_c^{\mathrm{eff}})_1}},a_2^{\mathrm{eff}}=C_1^{\mathrm{eff}}+{\displaystyle \frac{C_2^{\mathrm{eff}}}{(N_c^{\mathrm{eff}})_2}},`$
$`a_{3,4}^{\mathrm{eff}(\mathrm{q})}`$ $`=`$ $`C_{3,4}^{\mathrm{eff}}+{\displaystyle \frac{C_{4,3}^{\mathrm{eff}}}{(N_c^{\mathrm{eff}})_{4,3}}}+{\displaystyle \frac{3}{2}}e_q\left(C_{9,10}^{\mathrm{eff}}+{\displaystyle \frac{C_{10,9}^{\mathrm{eff}}}{(N_c^{\mathrm{eff}})_{10,9}}}\right),`$
$`a_{5,6}^{\mathrm{eff}(\mathrm{q})}`$ $`=`$ $`C_{5,6}^{\mathrm{eff}}+{\displaystyle \frac{C_{6,5}^{\mathrm{eff}}}{(N_c^{\mathrm{eff}})_{6,5}}}+{\displaystyle \frac{3}{2}}e_q\left(C_{7,8}+{\displaystyle \frac{C_{8,7}^{\mathrm{eff}}}{(N_c^{\mathrm{eff}})_{8,7}}}\right),`$ (4)
where
$`{\displaystyle \frac{1}{(N_c^{\mathrm{eff}})_i}}{\displaystyle \frac{1}{N_c}}+\chi _i.`$ (5)
with $`\chi _i`$ being the non-factorizable terms. In the GFA, $`1/(N_c^{\mathrm{eff}})_i`$ are assumed to be universal and real in the absence of FSIs. In the naive factorization, all effective WCs $`C_i^{\mathrm{eff}}`$ are reduced to the corresponding WCs of $`C_i`$ in the effective Hamiltonian and the non-factoziable terms are neglected, i.e., $`\chi _i=0`$.
Under the factorization hypothesis, the tree level hadronic matrix element $`O_{\mathrm{tree}}`$ is factorized into a product of two matrix elements of single currents, which are represented by the decay constant and form factors. The $`BH\left(H=D,D^{}\right)`$ transition form factors are crucial ingredients in the GFA for the heavy-heavy decays. To obtain the transition elements of $`BH`$ with various weak vertices, we first parameterize them in terms of the relevant form factors under the conventional forms as follows:
$`D|V_\mu |\overline{B}`$ $`=`$ $`F_1(q^2)\left\{P_\mu {\displaystyle \frac{Pq}{q^2}}q_\mu \right\}+{\displaystyle \frac{Pq}{q^2}}F_0(q^2)q_\mu ,`$ (6)
$`D^{}(ϵ)|V_\mu |\overline{B}`$ $`=`$ $`{\displaystyle \frac{V(q^2)}{m_B+m_D^{}}}\epsilon _{\mu \alpha \beta \rho }ϵ^\alpha P^\beta q^\rho ,`$
$`D^{}(ϵ)|A_\mu |\overline{B}`$ $`=`$ $`i[2m_D^{}A_0(q^2){\displaystyle \frac{ϵ^{}q}{q^2}}q_\mu +(m_B+m_D^{})A_1(q^2)(ϵ_\mu ^{}{\displaystyle \frac{ϵ^{}q}{q^2}}q_\mu )`$ (7)
$`A_2(q^2){\displaystyle \frac{ϵ^{}q}{m_B+m_D^{}}}(P_\mu {\displaystyle \frac{Pq}{q^2}}q_\mu )],`$
where $`V_\mu =\overline{q}\gamma _\mu b`$, $`A_\mu =\overline{q}\gamma _\mu \gamma _5b`$, $`m_{B,D,D^{}}`$ are the meson masses, $`ϵ_\mu `$ denotes the polarization vector of the $`D^{}`$ meson, $`P=p_B+p_{D^{()}}`$, $`q=p_Bp_{D^{()}}`$ and $`Pq=m_B^2m_{D^{()}}^2`$. According to the HQET, it will be more convenient to define the form factors in terms of the velocity of the heavy quark rather than the momentum. The definition of these form factors can be found in Ref. NR and the relation with the conventional ones are given by
$`F_1(q^2)`$ $`=`$ $`{\displaystyle \frac{m_B+m_D}{2\sqrt{m_Bm_D}}}\left[\xi _+(w){\displaystyle \frac{m_Bm_D}{m_B+m_D}}\xi _{}(w)\right],`$
$`F_0(q^2)`$ $`=`$ $`{\displaystyle \frac{m_B+m_D}{2\sqrt{m_Bm_D}}}\zeta _D(q^2)\left[\xi _+(w){\displaystyle \frac{m_B+m_D}{m_Bm_D}}\left({\displaystyle \frac{w1}{w+1}}\right)\xi _{}(w)\right],`$
$`V(q^2)`$ $`=`$ $`{\displaystyle \frac{m_B+m_D^{}}{2\sqrt{m_Bm_D^{}}}}\xi _V(w),A_1(q^2)={\displaystyle \frac{m_B+m_D^{}}{2\sqrt{m_Bm_D^{}}}}\zeta _D^{}\xi _{A_1}(w),`$
$`A_2(q^2)`$ $`=`$ $`{\displaystyle \frac{m_B+m_D^{}}{2\sqrt{m_Bm_D^{}}}}\left[\xi _{A_1}(w)+{\displaystyle \frac{m_D^{}}{m_B}}\xi _{A_2}(w)\right],`$
$`A_3(q^2)`$ $`=`$ $`{\displaystyle \frac{m_B+m_D^{}}{2\sqrt{m_Bm_D^{}}}}\{{\displaystyle \frac{m_B}{m_B+m_D^{}}}(w+1)\xi _{A_1}(w)`$
$`{\displaystyle \frac{m_Bm_D^{}}{2m_D^{}}}[\xi _{A_3}(w)+{\displaystyle \frac{m_D^{}}{m_B}}\xi _{A_2}(w)]\},`$
$`A_0(q^2)`$ $`=`$ $`A_3(q^2)+{\displaystyle \frac{q^2}{4m_Bm_D^{}}}\sqrt{{\displaystyle \frac{m_B}{m_D^{}}}}\left[\xi _{A_3}(w){\displaystyle \frac{m_D^{}}{m_B}}\xi _{A_2}(w)\right].`$ (8)
with $`\omega =(m_B^2+m_H^2q^2)/(2m_Bm_H)`$ and $`\zeta _H(q^2)=1q^2/(m_B+m_H)^2`$. It is known that under the HQS, $`\xi _+=\xi _V=\xi _{A_1}=\xi _{A_3}=\xi (w)`$ and $`\xi _{}=\xi _{A_2}=0`$. In our numerical estimations, we will base on the results of the HQS and include $`\alpha _s`$ and $`1/m_B`$ power corrections as well.
## III Generalized Factorization formulas and polarization fractions of VV modes
By the effective interactions and the form factors defined in the previous chapter, the decay amplitude could be described by the product of the effective WCs and the hadronic matrix elements in the framework of the GFA. For the hadronic matrix elements in $`BPP`$ decays, we will follow the notation of Ref. GFA2 , given by
$`X_1^{(BP_1,P_2)}`$ $``$ $`P_2|(\overline{q}_2q_3)_{VA}|0P_1|(\overline{q}_1b)_{VA}|\overline{B}=if_{P_2}(m_B^2m_{P_1}^2)F_0^{BP_1}(m_{P_2}^2),`$
$`X_2^{(BP_1,P_2)}`$ $``$ $`P_2|(\overline{q}_2q_3)_{S+P}|0P_1|(\overline{q}_1b)_{SP}|\overline{B}=i{\displaystyle \frac{m_{P_2}^2}{m_2+m_3}}f_{P_2}{\displaystyle \frac{m_B^2m_{P_1}^2}{m_bm_1}}F_0^{BP_1}(m_{P_2}^2),`$ (9)
where $`(\overline{q}_1b)_{SP}=\overline{q}_1(1\gamma _5)b`$, $`(\overline{q}_2q_3)_{S+P}=\overline{q}_2(1+\gamma _5)q_3`$ and $`m_{b,1,2,3}`$ correspond to the masses of quarks. The vertex $`(SP)(S+P)`$ is from the Fierz transformation of $`(VA)(V+A)`$. To get the decay constant and form factors for scalar vertices, we have utilized equation of motion for on-shell quarks. Moreover, we use
$`X_1^{(BP,V)}`$ $``$ $`V|(\overline{q}_2q_3)_{VA}|0P|(\overline{q}_1b)_{VA}|\overline{B}=2f_Vm_VF_1^{BP}(m_V^2)(\epsilon ^{}p__B),`$
$`X_1^{(BV,P)}`$ $``$ $`P|(\overline{q}_2q_3)_{VA}|0V|(\overline{q}_1b)_{VA}|\overline{B}=2f_Pm_VA_0^{BV}(m_P^2)(\epsilon ^{}p__B),`$
$`X_2^{(BV,P)}`$ $``$ $`P|(\overline{q}_2q_3)_{S+P}|0V|(\overline{q}_1b)_{SP}|\overline{B}={\displaystyle \frac{2m_P^2}{m_2+m_3}}f_P{\displaystyle \frac{m_V}{m_b+m_1}}A_0^{BV}(m_P^2)(\epsilon ^{}p__B),`$ (10)
and
$`X^{(BV_1,V_2)}`$ $``$ $`V_2|(\overline{q}_2q_3)_{VA}|0V_1|(\overline{q}_1b)_{VA}|\overline{B}`$ (11)
$`=`$ $`if_{V_2}m_{V_2}[(\epsilon _1^{}\epsilon _2^{})(m_B+m_{V_1})A_1^{BV_1}(m_{V_2}^2)(ϵ_1^{}p__2)(ϵ_2^{}p__1){\displaystyle \frac{2A_2^{BV_1}(m_{V_2}^2)}{(m_B+m_{V_1})}}`$
$`+iϵ_{\mu \nu \alpha \beta }\epsilon _2^\mu \epsilon _1^\nu p__2^\alpha p_1^\beta {\displaystyle \frac{2V^{BV_1}(m_{V_2}^2)}{(m_B+m_{V_1})}}],`$
for $`BPV(VP)`$ and $`BVV`$, respectively. We note that the sign difference of $`X_1^{(BP_1,P_2)}`$ and $`X_2^{(BP_1,P_2)}`$ in Eq. (9) will make the penguin effects become non-negligible. On the other hand, the same sign of $`X_1^{(BV,P)}`$ and $`X_2^{(BV,P)}`$ in Eq. (10) leads to the penguin effects negligible. Since the time-like form factors for annihilation contributions are uncertain, we take $`Y_{1(2)}^{(B,M_1M_2)}M_1M_2|(\overline{q}_2q_3)_{VA}|00|(\overline{q}_1b)_{VA}|\overline{B}`$ and $`Y_3^{(B,M_1M_2)}M_1M_2|(\overline{q}_2q_3)_{S+P}|00|(\overline{q}_1b)_{SP}|\overline{B}`$ to represent them, where $`M`$ can be pseudoscalar or vector bosons. Note that due to the identity of $`\epsilon _i(p_i)p_i=0`$, we have used $`V|(\overline{q}_2q_3)_{S+P}|0=0`$. With these notations and associated effective WCs, one can display the decay amplitude for the specific decay mode. We summary the relevant decay amplitudes in Appendixes. Once we get the decay amplitude, denoted by $`A(BM_1M_2)`$, the corresponding decay rate of the two-body mode could be obtained by
$`\mathrm{\Gamma }(BM_1M_2)={\displaystyle \frac{G_Fp}{16\pi m_B^2}}|A(BM_1M_2)|^2.`$ (12)
with $`p`$ being the spatial momentum of $`M_{1,2}`$. Consequently, the direct CPA is defined by
$`A_{CP}={\displaystyle \frac{\overline{\mathrm{\Gamma }}(\overline{B}\overline{M}_1\overline{M}_2)\mathrm{\Gamma }(BM_1M_2)}{\overline{\mathrm{\Gamma }}(\overline{B}\overline{M}_1\overline{M}_2)+\mathrm{\Gamma }(BM_1M_2)}}.`$ (13)
Besides the BRs and CPAs, we can also study the polarizations of the vector mesons in $`BVV`$ decays. To discuss the polarizations, one can write the general decay amplitudes in the helicity basis to be
$`A^{(\lambda )}`$ $`=`$ $`ϵ_{1\mu }^{}(\lambda )ϵ_{2\nu }^{}(\lambda )\left[ag^{\mu \nu }+{\displaystyle \frac{b}{m_{V_1}m_{V_2}}}p_2^\mu p_1^\nu +i{\displaystyle \frac{c}{m_{V_1}m_{V_2}}}ϵ^{\mu \nu \alpha \beta }p_{1\alpha }p_{2\beta }\right].`$ (14)
In this basis, the amplitudes with various helicities can be given as
$`H_0=axb(x^21),H_\pm =a\pm \sqrt{x^21}c.`$
where $`x=(m_B^2m_{V_1}^2m_{V_2}^2)/(2m_{V_1}m_{V_2})`$. In addition, we can define the polarization amplitudes as
$`A_0=H_0,A_{}={\displaystyle \frac{1}{\sqrt{2}}}(H_++H_{})=\sqrt{2}a,`$
$`A_{}={\displaystyle \frac{1}{\sqrt{2}}}(H_+H_{})=\sqrt{2}\sqrt{x^21}c.`$ (15)
Accordingly, the decay rate expressed by helicity amplitudes for the VV mode can be written as
$`\mathrm{\Gamma }={\displaystyle \frac{G_Fp}{16\pi m_B^2}}\left(|A_0|^2+|A_{}|^2+|A_{}|^2\right),`$
and the polarization fractions can be defined as
$`R_i={\displaystyle \frac{|A_i|^2}{|A_0|^2+|A_{}|^2+|A_{}|^2}}.`$ (16)
where $`i=0`$ and $``$ ($``$), representing the longitudinal and transverse parallel (perpendicular) components, respectively, with the relation of $`R_0+R_{}+R_{}=1`$. Under CP parities, $`R_{0,}`$ are CP-even while $`R_{}`$ is CP-odd.
From the hadronic matrix element in Eq. (11), the amplitudes $`a`$, $`b`$ and $`c`$ in the framework of the GFA are expressed by
$`a`$ $`=`$ $`\stackrel{~}{C}_{eff}(m_B+m_{V_1})m_{V_2}f_{V_2}A_1^{BV_1}(m_{V_2}^2),`$
$`b`$ $`=`$ $`\stackrel{~}{C}_{eff}{\displaystyle \frac{2m_{V_1}m_{V_2}^2}{m_B+m_{V_1}}}f_{V_2}A_2^{BV_1}(m_{V_2}^2),`$
$`c`$ $`=`$ $`\stackrel{~}{C}_{eff}{\displaystyle \frac{2m_{V_1}m_{V_2}^2}{m_B+m_{V_1}}}f_{V_2}V^{BV_1}(m_{V_2}^2).`$ (17)
where $`\stackrel{~}{C}_{eff}`$ represents the involved WCs and CKM matrix elements. With the form factors in Eq. (8) and the heavy quark limit, we get that the ratios $`r_b=b/a`$ and $`r_c=c/a`$ are related. Explicitly, we have
$`r_b=r_c={\displaystyle \frac{m_D^{}}{m_B}}{\displaystyle \frac{1}{(1+w)}}0.16,`$ (18)
which are small. From Eqs. (15) and (16), we find that the polarization fractions behave as
$`R_0R_{},R_{}𝒪\left({\displaystyle \frac{m_D^{}^2}{m_B^2}}\right),`$ (19)
which indicate that the power law in the heavy-heavy decays is different from the light-light ones, which are expected to be $`R_01`$, $`R_{}R_{}𝒪(m_V^2/m_B^2)`$. Moreover, $`R_{}`$ is directly related to $`c`$ and can be written as
$`R_{}`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Gamma }_0}}(x^21)|r_c|^2,`$ (20)
with $`\mathrm{\Gamma }_0=1+(x^21)|r_c|^2+\left|x+(x^21)r_b\right|^2/2`$. By comparing to the result in the HQS, we find an interesting relation
$`{\displaystyle \frac{R_{}}{R_{}^{HQS}}}\left[\zeta _D^{}{\displaystyle \frac{V(q^2)}{A_1(q^2)}}\right]^2.`$ (21)
where $`R_{}^{HQS}=0.055`$ denotes the transverse perpendicular fraction under the HQS. As a good approximation, the form factor $`A_2`$-dependent of $`R_{}`$ is decoupled. By the relationship, we can clearly understand the influence of the HQS breaking effects.
## IV Numerical analysis and discussions
### IV.1 Estimations on the annihilation contributions
Since the annihilation contributions relate to time-like form factors and there are no direct experimental measurements, we shall neglect them in our calculations. However, to make sure that the neglected parts are small, we can connect the processes of $`BD^0\overline{D}^0`$ and $`BD_s^+D_s^{}`$ to the decays $`B^0D_s^{}K^+`$ and $`B^0J/\psi \overline{D}^0`$, which are directly associated with annihilation topologies, with the experimental data of $`Br(B^0D_s^{}K^+)=(3.8\pm 1.3)\times 10^5`$ PDG04 and $`Br(B^0J/\psi \overline{D}^0)<1.3\times 10^5`$ JD , respectively. By the flavor-diagram analysis, shown in Fig. 1,
except there appears a CKM suppressed factor $`V_{cd}\lambda `$ (see Fig.1a) for $`D^0\overline{D}^0`$ and $`D_s^+D_s^{}`$ modes, the four modes have the same decay topologies. Hence, by assuming that they have similar hadronic effects, the BRs of $`B^0D^0\overline{D}^0(D_s^+D_s^{})`$ could be estimated to be less than $`𝒪(10^6)`$.
To give a detailed analysis, we can include the character of each mode, governed by the meson distribution amplitudes. For simplicity, we will concentrate on the leading twist effects and take the meson wave functions to be $`\mathrm{\Phi }_Df_Dx(1x)(1+0.8(12x))`$ PQCD , $`\mathrm{\Phi }_{D_s}f_{D_s}x(1x)(1+0.3(12x))`$ LU , $`J/\mathrm{\Psi }f_{J/\mathrm{\Psi }}x(1x)(x(1x)/(12.8x(1x)))^{0.7}`$ BC and $`\mathrm{\Phi }_Kf_Kx(1x)(1+0.51(12x)+0.3[5(12x)^21])`$ BBKT , where $`x`$ is the momentum fraction of the parton inside the meson and $`f_{D,D_s,J/\mathrm{\Psi },K}`$ are the decay constants of $`D`$, $`D_s`$, $`J/\mathrm{\Psi }`$ and $`K`$ mesons, respectively. From these wave functions, we know that the maximum contributions are from $`x_0(0.35,0.43,0.5,0.5)`$ for $`(D,D_s,J/\mathrm{\Psi },K)`$. With the information, we can estimate the decay amplitudes in order of magnitude for $`\overline{B}_dD_s^+K^{}(J/\mathrm{\Psi }D^0)`$ and $`\overline{B}_dD^0\overline{D}^0(D_s^+D_s^{})`$ as shown in Figs. 1a and 1b with $`q^{}=s(c)`$ and $`q=u(s)`$, respectively. Note that there exists a chiral suppression in the factorized parts in annihilation contributions. However, we just consider the nonfactorized effects in the estimations. Therefore, by Fig. 1 with the gluon exchange, the decay amplitude is related to the propagators of the gluon and the light quark, described by $`1/(k_2+k_3)^2/(k_2+k_3)^2`$, where $`k_{2(3)}`$ denote the momenta carried by the spectators. For simplicity, we have neglected the momentum carried by the light quark of the $`B`$ meson. By the momentum fraction, the decay amplitude could satisfy that $`A1/(x_2x_3)^2`$. Hence, the relative size of the decay amplitudes could be given approximately as
$`A(D_s^+K^{}):A(J/\mathrm{\Psi }D^0):A(D^0\overline{D}^0):A(D_s^+\overline{D}_s^{}){\displaystyle \frac{f_{D_s}f_K}{(x_2x_3)^2}}:{\displaystyle \frac{f_{J/\mathrm{\Psi }}f_D}{(x_2x_3)^2}}:{\displaystyle \frac{\lambda f_D^2}{(x_2x_3)^2}}:{\displaystyle \frac{\lambda f_{D_s}^2}{(x_2x_3)^2}}.`$
With the information of maximum contributions, characterized by $`x_0`$ for each mode, we get
$`A(D_s^+K^{}):A(J/\mathrm{\Psi }D^0):A(D^0\overline{D}^0):A(D_s^+\overline{D}_s^{})`$
$`1:{\displaystyle \frac{f_Df_{J/\mathrm{\Psi }}}{f_{D_s}f_K}}\left({\displaystyle \frac{0.43}{0.65}}\right)^2:{\displaystyle \frac{\lambda f_D^2}{f_{D_s}f_K}}\left({\displaystyle \frac{0.50.43}{0.35^2}}\right)^2:{\displaystyle \frac{\lambda f_{D_s}}{f_K}}\left({\displaystyle \frac{0.5}{0.43}}\right)^2.`$ (22)
With the kinetic effects, the ratios of BRs are roughly to be $`BR(D_s^+K^{}):BR(J/\mathrm{\Psi }D^0):BR(D^0\overline{D}^0):BR(D_s^+\overline{D}_s^{})1:0.25:0.39:0.16`$. That is, the BRs of $`\overline{B}D^0\overline{D}^0(D_s^+D_s^{})`$ could be as large as $`𝒪(10^6)`$, which implies that annihilation effects could be neglected in the discussions on the BR of the production for color-allowed two charmful mesons. We note that our estimations are just based on SD effects and at the level of order of magnitude. Clearly, direct experimental measurements are needed to confirm our results.
### IV.2 $`\alpha _s`$, power corrections and the parametrization of Isgur-Wise function
In the HQS limit, the form factors could be related to a single Isgur-Wise function $`\xi (\omega )`$ by $`\xi _+=\xi _V=\xi _{A_1}=\xi _{A_3}=\xi (w)`$ and $`\xi _{}=\xi _{A_2}=0`$. We now include the perturbative QCD corrections induced by hard gluon vertex corrections of $`bc`$ transitions and power corrections in orders of $`1/m_Q`$ with $`Q=b`$ and $`c`$. Consequently, the form factors can be written as
$`\xi _i(w)=\left\{\alpha _i+\beta _i(w)+\gamma _i(w)\right\}\xi (w).`$ (23)
where $`\xi (w)`$ is the Isgur-Wise function, $`\alpha _+=\alpha _V=\alpha _{A_1}=\alpha _{A_3}=1`$, $`\alpha _{}=\alpha _{A_2}=0`$ and $`\beta _i(\omega )`$ and $`\gamma _i(\omega )`$ stand for effects of $`\alpha _s`$ and power corrections, respectively. Explicitly, for the two-body decays in our study, $`\omega 1.3`$ and the values of the other parameters are summarized as follows NeubertQCD ; NR :
$`\begin{array}{ccc}\beta _+=0.043,\hfill & \beta _{}=0.069,\hfill & \beta _V=0.072,\hfill \\ \beta _{A_1}=0.067,\hfill & \beta _{A_2}=0.114,\hfill & \beta _{A_3}=0.028,\hfill \\ \gamma _+=0.015,\hfill & \gamma _{}=0.122,\hfill & \gamma _V=0.224,\hfill \\ \gamma _{A_1}=0.027,\hfill & \gamma _{A_2}=0.093,\hfill & \gamma _{A_3}=0.014.\hfill \end{array}`$ (28)
Clearly, the range of their effects is from few percent to 20% level. In particular, the power corrections to the form factor $`\xi _V`$ (or say $`V(q^2)`$) are the largest, about 20%. The resultant is also consistent with other QCD approaches, such as the constitute quark model (CQM) MS and the light-front (LF) QCD CCH .
After taking care of the corrections, the remaining unknown is the Isgur-Wise function. To determine it, we adopt a linear parametrization to be $`\xi (w)=1\rho _H^2(w1)`$ for the transition $`BH`$, where $`\rho _H^2`$ is called the slope parameter. We use the BRs of semileptonic $`BD^{()}\mathrm{}^{}\overline{\nu }_{\mathrm{}}`$ decays to determine $`\rho _H^2`$. We note that the values of $`\rho _D^2`$ and $`\rho _D^{}^2`$ are not the same as those in $`D`$ and $`D^{}`$ decays. In our approach, the difference is from higher orders and power corrections. Hence, with $`V_{cb}=0.0416`$ and $`BR(BD^{()}\mathrm{}^{}\overline{\nu }_{\mathrm{}})=2.15\pm 0.22(6.5\pm 0.5)\%`$ PDG04 , we obtain $`\rho _D^2=0.90\pm 0.06`$ and $`\rho _D^{}^2=1.09\pm 0.05`$. Since the errors of $`\rho _H^2`$ are small, we will only use the central values in our numerical results.
### IV.3 Results for BRs and polarization fractions
To get the numerical estimations, the input values for the relevant parameters are taken to be as follows PDG04 ; NS ; CKM :
$`f_D=0.20,f_D^{}=0.23,f_{D_s}=0.24,f_{D_s^{}}=0.275\mathrm{GeV};`$
$`V_{cd}=\lambda ,V_{cs}=1\lambda ^2/2,V_{cb}=A\lambda ^2,`$
$`V_{td}=\lambda |V_{cb}|R_te^{i\varphi _1},V_{ts}=A\lambda ^2,V_{tb}=1,`$
$`A=0.83,\lambda =0.224,\varphi _1=23.4^\mathrm{o},R_t=0.91;`$
$`m_u=0.005,m_d=0.01,m_s=0.15,m_c=1.5,m_b=4.5\mathrm{GeV}.`$ (29)
Note that the numerical results are insensitive to light quark masses. As to the WCs, we adopt the formulas up to one-loop corrections presented in Ref. GFA2 and set $`\mu =2.5`$ GeV. As mentioned early, since the nonfactorized contributions are grouped into $`N_c^{\mathrm{eff}}`$, the color number in Eq. (4) will be regarded as a variable. To display their effects, we take the values of $`N_c^{\mathrm{eff}}=2,\mathrm{\hspace{0.33em}3},\mathrm{\hspace{0.33em}5}`$ and $`\mathrm{}`$.
By following the factorized formulas shown in Appendixes, we present the BRs with various $`N_c^{\mathrm{eff}}`$ in Tables 1, 2 and 3 for $`PP`$, $`PV(VP)`$ and $`VV`$ modes, respectively. In order to accord with the experimental data, our predictions of the BRs are given as the CP-averaged values. For comparisons, we also calculate the results in terms of the form factors given by the CQM and LF, which are displayed in Table 4. Since the CPAs are quite similar in different models, in Table 4 we just show the results in our approach. As to the polarization fractions, we present them in Table 5. Therein, to understand the influence of the HQS breaking effects, we separate the results to be HQS and HQS<sub>I(II)</sub>, representing the HQS results and those with $`\alpha _s`$ ($`\alpha _s`$+power) corrections, respectively.
We now present our discussions on the results as follows:
(1) The non-factorizable contributions are not dominant for color-allowed two charmed-meson decays. According to the classification in Refs. GFA1 ; GFA2 , the decay modes displaced in Tables 1, 2 and 3 belong to class I, which are dominated by the external $`W`$-emission. The leading decay amplitudes are proportional to the effective coefficient $`a_1`$, which is stable against the variation of $`N_c^{\mathrm{eff}}`$. Thus, the predicted branching ratios are insensitive to $`N_c^{\mathrm{eff}}`$. This means that annihilation contributions and FSIs, neglected in the GFA, are sub-leading contributions. On the other hand, by varying $`N_c^{\mathrm{eff}}`$ from $`3`$ to $`2`$, or $`3`$ to $`\mathrm{}`$, the branching ratios change by about 10-20%, which should be the same order as annihilation and FSI effects. From Tables 1, 2 and 3, there are no obvious deviations of the theoretical predictions from the experimental data within the present errors. It is also interesting to note that $`N_c^{\mathrm{eff}}=\mathrm{}`$ is not excluded by experiments if considering the uncertainties of decay constants and from factors. Thus, the large $`N_c`$ limit as a mechanism of factorization is not disfavored yet.
(2) The main uncertainties of theory come from the decay constants and form factors. Because the decay amplitudes are proportional to decay constants, it is clear that the theoretical predictions can be changed with different values of the decay constants. For instance, the branching ratio is $`\mathrm{BR}(\mathrm{our}\mathrm{result})\times \left(\frac{f_{D_s}}{0.24}\right)^2`$ for $`\overline{B}^0D^+D_s^{}`$. The recent experiment $`\mathrm{BR}(\overline{B}^0D^+D_s^{})=18.8\pm 0.9\pm 1.7`$ seems to favor a lower $`f_{D_s^{}}0.24`$ than our choice of $`0.275`$. However, this point has to be checked by other processes. For the form factors, the predictions of BRs in our approach are slightly lower than those in other two approaches (CQM and LF). The present experiment data can not distinguish which model is more preferred. More precise data are necessary. Another place to test different approaches is through the transverse polarization $`R_{}`$. From Table 5, $`R_{}`$ is predicted to be $`0.07,0.08`$ and $`0.09`$ in the CQM, LF and HQET, respectively. The larger prediction in the HQET is due to $`\alpha _s`$ corrections. Except the model-dependent calculation of power corrections in different approaches, one advantage of the HQET is that it permits the calculations of perturbative QCD corrections systematically.
(3) The penguin effects can not be neglected in $`BPP`$ decays. By using the decay amplitudes in Appendixes, the definitions of hadronic effects in Eqs. (9) and (10) and the condition of $`\epsilon _i(p_i)p_i=0`$, we know that the effects of penguin ($`P`$) to tree ($`T`$) level, denoted by $`P/T`$, for $`PP`$, $`VP`$ and $`VV`$ modes are proportional to $`(a_4^{\mathrm{eff}(\mathrm{c})}+2a_6^{\mathrm{eff}(\mathrm{c})})/a_1^{\mathrm{eff}}`$, $`(a_4^{\mathrm{eff}(\mathrm{c})}2a_6^{\mathrm{eff}(\mathrm{c})}^{})/a_1^{\mathrm{eff}}`$ and $`a_4^{\mathrm{eff}(\mathrm{c})}/a_1^{\mathrm{eff}}`$, respectively, where $`=m_D^2/[(m_c+m_d)(m_bm_c)]`$ and $`^{}=m_D^2/[(m_c+m_d)(m_b+m_c)]`$ and the CKM matrix elements have been canceled due to $`|V_{tb}V_{ts}^{}||V_{cb}V_{cs}^{}|`$ and $`|V_{tb}V_{td}^{}||V_{cb}V_{cd}^{}|`$. The situations in the $`PV`$ modes are the same as those in the $`VV`$ modes due to the vector meson being factorized out from the $`BP`$ transition. Since the WCs $`a_4^{\mathrm{eff}(\mathrm{c})}`$ and $`a_6^{\mathrm{eff}(\mathrm{c})}`$ have the same sign, we see clearly that penguin effects in the $`PP`$ modes are larger than those in the $`VV`$ modes; however, due to the cancelation between $`a_4^{\mathrm{eff}(\mathrm{c})}`$ and $`a_6^{\mathrm{eff}(\mathrm{c})}`$, penguin effects could be neglected in $`BVP`$ decays. Hence, the ratios $`|P/T|`$ for $`PP`$, $`VP`$ and $`VV(PV))`$ are around $`15\%`$, $`0\%`$ and $`4\%`$, respectively. For the $`PP,VV(PV)`$ modes, our predictions are consistent with the results in Refs. CY ; KKLM . Note that an 4% penguin contribution was obtained for the $`VP`$ modes in CY . The difference is due to that they used a lower charm quark mass ($`m_c=0.95\mathrm{GeV}`$) than ours. For all the decay modes, the electroweak penguin contributions can be negligible (less than 1%).
(4) Without FSIs, we find that the BRs in the neutral and charged modes have the following relationships:
$`{\displaystyle \frac{1}{\tau _{B^0}}}BR(D^{()+}D_s^{()})`$ $``$ $`{\displaystyle \frac{1}{\tau _{B^+}}}BR(D^{()0}D_s^{()}),`$
$`{\displaystyle \frac{1}{\tau _{B^0}}}BR(D^{()+}D^{()})`$ $``$ $`{\displaystyle \frac{1}{\tau _{B^+}}}BR(D^{()0}D^{()}).`$
In addition, the decays with nonstrangeness charmed mesons are Cabibbo-suppressed compared to the decays with the $`D_s^{()}`$ emission and they satisfy
$`BR(BD^{()}D^{()}){\displaystyle \frac{f_{D^{()}}^2}{f_{D_s^{()}}^2}}\lambda ^2BR(BD^{()}D_s^{()}).`$ (30)
Clearly, if large deviations from the equalities in Eq. (30) are observed in experiments, they should arise from FSIs. Of course, if the BRs of $`\overline{B}^0D^{()0}\overline{D}^{()0}`$ and $`\overline{B}^0D_s^{()+}D_s^{()}`$ with $`𝒪(10^4)`$ are seen, it will be another hint for FSIs EFP .
(5) For the decay amplitude, we write
$`A=T+Pe^{i\theta _W}e^{i\delta },`$ (31)
where $`T`$ and $`P`$ represent tree and penguin amplitudes, and we have chosen the convention such that $`T`$ and $`P`$ are real numbers and $`\theta _W`$ and $`\delta `$ are the CP weak and strong phases, respectively. From Eq. (13), the CPA can be described by
$`A_{CP}={\displaystyle \frac{2(P/T)\mathrm{sin}\delta \mathrm{sin}\theta _W}{1+(P/T)^2+2(P/T)\mathrm{cos}\delta \mathrm{cos}\theta _W}}.`$ (32)
According to the discussions in (1), the maximum CPAs in $`PP`$, $`PV`$ and $`VV(VP)`$ are expected to be around $`26\%`$, $`0\%`$ and $`8\%`$, respectively. However, in $`BD^{()}D^{()}`$ decays, due to $`|\theta _W|=|\varphi _1|`$, if we take $`\delta =90^o`$ and $`\varphi =23.4^o`$, the maximum CPAs for $`PP`$ and $`VV(VP)`$ modes are $`10.3\%`$ and $`1.6\%`$, respectively. Clearly, in the SM, the CPA with $`𝒪(10\%)`$ can be reached in $`BDD`$ decays. Due to the associated CKM matrix element being $`V_{ts}A\lambda ^2`$, there are no CPAs in $`BD^{()}D_s^{()}`$ decays. In the GFA, since the strong phases mainly arise from the one-loop corrections which are usually small, our results on CPAs, shown in Table 4, are all at a few percent level. Therefore, if the CPAs of $`𝒪(10\%)`$ are found in $`\overline{B}^0D^+D^{}`$ and $`\overline{B}^+D^0D^{}`$ decays, we can conclude the large effects of the strong phase are from FSIs.
(6) As discussed before, we know that in two charmful decays the polarization fractions satisfy $`R_{}<<R_0R_{}`$. The current experimental data are: $`R_0=0.52\pm 0.05`$ PDG04 for $`B^0D^+D_s^{}`$ and $`R_0=0.57\pm 0.08\pm 0.02`$, $`R_{}=0.19\pm 0.08\pm 0.01`$ BELLE\_DD and $`R_{}=0.063\pm 0.055\pm 0.009`$ BABAR\_DD for $`B^0D^+D^{}`$. We can see that the experimental measurements support the power-law relation. To estimate how large $`R_{}`$ can be in theory, we use the relationship in Eq. (21) and the form factors in Eq. (7) and we obtain
$`{\displaystyle \frac{R_{}}{R_{}^{HQS}}}\left[{\displaystyle \frac{1+\beta _V+\gamma _V}{1+\beta _{A_1}+\gamma _{A_1}}}\right]^2.`$ (33)
With the values in Eq. (28) and $`R_{}^{HQS}=0.055`$, we get $`R_{}10\%`$. The detailed numerical values can be found in Table 5. Interestingly, for the $`\overline{B}^0D^+D^{}`$ decay, the estimated result is close to the upper limit of $`R_{}=0.063\pm 0.055\pm 0.009`$ observed by BABAR BABAR\_DD but close to the lower limit of $`R_{}=0.19\pm 0.08\pm 0.01`$ observed by BELLE BELLE\_DD . We note that our results are different from the PQCD predictions in which $`R_{}0.06`$ LM . From our results, we can conjecture that if large $`R_{}`$, say around $`20\%`$, is observed, large contributions should arise from FSIs.
## V Conclusions
We have presented a detailed study of B decaying into two charmed-mesons in the generalized factorization approach. The penguin contributions have also been taken into account. If the final states are both pseudoscalar mesons, the ratio of penguin and tree contributions is about 10% in the decay amplitude. The direct CP violating asymmetries have been estimated to be a few percent. For the $`B^0D^+D^{},D^+D^{},D^+D^{}`$ decays, the “penguin pollution” is weaker than that in the $`D^+D^{}`$ mode. Thus, these modes provide cleaner places to cross-check the value of $`sin2\beta `$ measured in the $`B^0J/\psi K`$ decays. The weak annihilation contributions have been found to be small. We have proposed to test the annihilation effects in annihilation-dominated processes of $`B^0D^{()0}\overline{D}^{()0}`$ and $`D_s^{()+}D_s^{()}`$.
We have performed a comprehensive test of the factorization in the heavy-heavy B decays. The predictions of branching ratios in theory are consistent with the experimental data within the present level. The variations of branching ratios with the effective color number $`N_c^{\mathrm{eff}}`$ show that the soft FSIs are not dominant. However, we cannot make the conclusion that they are negligible. Their effects can be of order 10-20% for branching ratios as indicated from the variation of $`N_c^{\mathrm{eff}}`$. Since the soft divergences are not canceled in the non-factorizable corrections, this may indicate that the strong interactions at low energy either become weak or are suppressed by some unknown parameters (such as $`N_c`$ in the large $`N_c`$ theory). If the factorization is still a working concept in the heavy-heavy decays, there must be some non-perturbative mechanisms which prefer the factorization of a large-size charmed-meson from an environment of “soft cloud”. A relevant comment on the necessity of non-perturbative QCD justification can be found in LLW .
The polarization structure in the heavy-heavy decays has shown that the transverse perpendicular polarization fraction $`R_{}`$ is the smallest while the other two are comparable in size. This structure follows from the QCD dynamics in the heavy quark limit. We have found one relation between the transverse perpendicular polarization fraction and the ratios of form factors, in particular $`V(q^2)/A_1(q^2)`$. The corrections to the heavy quark limit give an enhancement of $`R_{}`$ from 0.055 to about 0.09. Since the FSIs are not significant, we do not expect that FSIs can change our prediction of $`R_{}`$ substantially. If future measurements confirm $`R_{}0.2`$ as the recent measurement by BELLE, it will be difficult to explain within the HQET and the factorization hypothesis.
In conclusion, our study has shown that the factorization works well in B meson heavy-heavy decays at present. More precise experimental data are desired to give a better justification. For theory, to explain the mechanism of factorization in the heavy-heavy decays is of high interest. The measurement of the transverse perpendicular polarization provides important information on the size of the heavy quark symmetry breaking or the possibility of large non-factorizable effects.
Acknowledgments
We thank Hai-Yang Cheng and Yu-Kuo Hsiao for many valuable discussions. This work is supported in part by the National Science Council of R.O.C. under Grant #s: NSC-93-2112-M-006-010 and NSC-93-2112-M-007-014.
## $`BPP`$ decays
$`A(\overline{B}^0D^+D_s^{})`$ $`=`$ $`V_{cb}V_{cs}^{}a_1^{\mathrm{eff}}X_1^{(BD,D_s)}V_{tb}V_{ts}^{}[(a_4^{\mathrm{eff}(\mathrm{c})}X_1^{(BD,D_s)}2a_6^{\mathrm{eff}(\mathrm{c})}X_2^{(BD,D_s)})`$ (34)
$`+(a_4^{\mathrm{eff}(\mathrm{d})}Y_1^{(B,DD_s)}2a_6^{\mathrm{eff}(\mathrm{d})}Y_3^{(B,DD_s)})],`$
$`A(\overline{B}^0D_s^+D_s^{})`$ $`=`$ $`V_{cb}V_{cd}^{}a_2^{\mathrm{eff}}Y_1^{(B,D_sD_s)}V_{tb}V_{td}^{}[(a_3^{\mathrm{eff}(\mathrm{s})}+a_3^{\mathrm{eff}(\mathrm{c})})Y_1^{(B,D_sD_s)}`$ (35)
$`+(a_5^{\mathrm{eff}(\mathrm{s})}+a_5^{\mathrm{eff}(\mathrm{c})})Y_2^{(B,D_sD_s)}],`$
$`A(B^{}D^0D_s^{})`$ $`=`$ $`V_{cb}V_{cs}^{}a_1^{\mathrm{eff}}X_1^{(BD,D_s)}+V_{ub}V_{us}^{}a_1^{\mathrm{eff}}Y_1^{(B,DD_s)}V_{tb}V_{ts}^{}[(a_4^{\mathrm{eff}(\mathrm{c})}X_1^{(BD,D_s)}`$ (36)
$`2a_6^{\mathrm{eff}(\mathrm{c})}X_2^{(BD,D_s)})+(a_4^{\mathrm{eff}(\mathrm{u})}Y_1^{(B,DD_s)}2a_6^{\mathrm{eff}(\mathrm{u})}Y_3^{(B,DD_s)})],`$
$`A(\overline{B}^0D^+D^{})`$ $`=`$ $`V_{cb}V_{cd}^{}[a_1^{\mathrm{eff}}X_1^{(BD,D)}+a_2^{\mathrm{eff}}Y_1^{(B,DD)}]V_{tb}V_{td}^{}[(a_4^{\mathrm{eff}(\mathrm{c})}X_1^{(BD,D)}`$ (37)
$`2a_6^{\mathrm{eff}(\mathrm{c})}X_2^{(BD,D)})+(a_4^{\mathrm{eff}(\mathrm{d})}+a_3^{\mathrm{eff}(\mathrm{d})}+a_3^{\mathrm{eff}(\mathrm{c})})Y^{(B,DD)}_1`$
$`+(a_5^{\mathrm{eff}(\mathrm{d})}+a_5^{\mathrm{eff}(\mathrm{c})})Y_2^{(B,DD)}2a_6^{\mathrm{eff}(\mathrm{d})}Y_3^{(B,DD)}],`$
$`A(\overline{B}^0D^0\overline{D}^0)`$ $`=`$ $`(V_{cb}V_{cd}^{}+V_{ub}V_{ud}^{})a_2^{\mathrm{eff}}Y_1^{(B,DD)}V_{tb}V_{td}^{}[(a_3^{\mathrm{eff}(\mathrm{u})}+a_3^{\mathrm{eff}(\mathrm{c})})Y_1^{(B,DD)}`$ (38)
$`+(a_5^{\mathrm{eff}(\mathrm{u})}+a_5^{\mathrm{eff}(\mathrm{c})})Y_2^{(B,DD)}],`$
$`A(B^{}D^0D^{})`$ $`=`$ $`V_{cb}V_{cd}^{}a_1^{\mathrm{eff}}X_1^{(BD,D)}+V_{ub}V_{ud}^{}a_1^{\mathrm{eff}}Y_1^{(B,DD)}V_{tb}V_{td}^{}[(a_4^{\mathrm{eff}(\mathrm{c})}X_1^{(BD,D)}`$ (39)
$`2a_6^{\mathrm{eff}(\mathrm{c})}X_2^{(BD,D)})+a_4^{\mathrm{eff}(\mathrm{d})}Y_1^{(B,DD)}2a_6^{\mathrm{eff}(\mathrm{d})}Y_3^{(B,DD)}].`$
## Appendix A $`BPV(VP)`$ decays
$`A(\overline{B}^0D^+D_s^{})`$ $`=`$ $`V_{cb}V_{cs}^{}a_1^{\mathrm{eff}}X_1^{(BD,D_s^{})}V_{tb}V_{ts}^{}[a_4^{\mathrm{eff}(\mathrm{c})}X_1^{(BD,D_s^{})}`$ (40)
$`+(a_4^{\mathrm{eff}(\mathrm{d})}Y_1^{(B,DD_s^{})}2a_6^{\mathrm{eff}(\mathrm{d})}Y_3^{(B,DD_s^{})})],`$
$`A(\overline{B}^0D^+D_s^{})`$ $`=`$ $`V_{cb}V_{cs}^{}a_1^{\mathrm{eff}}X_1^{(BD^{},D_s)}V_{tb}V_{ts}^{}[(a_4^{\mathrm{eff}(\mathrm{c})}X_1^{(BD^{},D_s)}2a_6^{\mathrm{eff}(\mathrm{c})}X_2^{(BD^{},D_s)})`$ (41)
$`+(a_4^{\mathrm{eff}(\mathrm{d})}Y_1^{(B,D^{}D_s)}2a_6^{\mathrm{eff}(\mathrm{d})}Y_3^{(B,D^{}D_s)})],`$
$`A(\overline{B}^0D_s^+D_s^{})`$ $`=`$ $`V_{cb}V_{cd}^{}a_2^{\mathrm{eff}}Y_1^{(B,D_sD_s^{})}V_{tb}V_{td}^{}[(a_3^{\mathrm{eff}(\mathrm{s})}+a_3^{\mathrm{eff}(\mathrm{c})})Y_1^{(B,D_sD_s^{})}`$ (42)
$`+(a_5^{\mathrm{eff}(\mathrm{s})}+a_5^{\mathrm{eff}(\mathrm{c})})Y_2^{(B,D_sD_s^{})}],`$
$`A(\overline{B}^0D_s^+D_s^{})`$ $`=`$ $`V_{cb}V_{cd}^{}a_2^{\mathrm{eff}}Y_1^{(B,D_s^{}D_s)}V_{tb}V_{td}^{}[(a_3^{\mathrm{eff}(\mathrm{s})}+a_3^{\mathrm{eff}(\mathrm{c})})Y_1^{(B,D_s^{}D_s)}`$ (43)
$`=(a_5^{\mathrm{eff}(\mathrm{s})}+a_5^{\mathrm{eff}(\mathrm{c})})Y_2^{(B,D_s^{}D_s)}],`$
$`A(B^{}D^0D_s^{})`$ $`=`$ $`V_{cb}V_{cs}^{}a_1^{\mathrm{eff}}X_1^{(BD,D_s^{})}+V_{ub}V_{us}^{}a_1^{\mathrm{eff}}X_1^{(B,DD_s^{})}V_{tb}V_{ts}^{}[(a_4^{\mathrm{eff}(\mathrm{c})}X_1^{(BD,D_s^{})}`$ (44)
$`+(a_4^{\mathrm{eff}(\mathrm{u})}Y_1^{(B,DD_s^{})}2a_6^{\mathrm{eff}(\mathrm{u})}Y_3^{(B,DD_s^{})})],`$
$`A(B^{}D^0D_s^{})`$ $`=`$ $`V_{cb}V_{cs}^{}a_1^{\mathrm{eff}}X_1^{(BD^{},D_s)}+V_{ub}V_{us}^{}a_1^{\mathrm{eff}}Y_1^{(B,D^{}D_s)}V_{tb}V_{ts}^{}[(a_4^{\mathrm{eff}(\mathrm{c})}X_1^{(BD^{},D_s)}`$ (45)
$`2a_6^{\mathrm{eff}(\mathrm{c})}X_2^{(BD^{},D_s)})+(a_4^{\mathrm{eff}(\mathrm{u})}Y_1^{(B,D^{}D_s)}2a_6^{\mathrm{eff}(\mathrm{u})}Y_3^{(B,D^{}D_s)})],`$
$`A(\overline{B}^0D^+D^{})`$ $`=`$ $`V_{cb}V_{cd}^{}[a_1^{\mathrm{eff}}X_1^{(BD,D^{})}+a_2^{\mathrm{eff}}Y_1^{(B,DD^{})}]V_{tb}V_{td}^{}[a_4^{\mathrm{eff}(\mathrm{c})}X_1^{(BD,D^{})}`$ (46)
$`+(a_4^{\mathrm{eff}(\mathrm{d})}+a_3^{\mathrm{eff}(\mathrm{d})}+a_3^{\mathrm{eff}(\mathrm{c})})Y_1^{(B,DD^{})}+(a_5^{\mathrm{eff}(\mathrm{d})}+a_5^{\mathrm{eff}(\mathrm{c})})Y_2^{(B,DD^{})}`$
$`2a_6^{\mathrm{eff}(\mathrm{d})}Y_3^{(B,DD^{})}],`$
$`A(\overline{B}^0D^+D^{})`$ $`=`$ $`V_{cb}V_{cd}^{}[a_1^{\mathrm{eff}}X_1^{(BD^{},D)}+a_2^{\mathrm{eff}}Y_1^{(B,D^{}D)}]V_{tb}V_{td}^{}[(a_4^{\mathrm{eff}(\mathrm{c})}X_1^{(BD^{},D)}`$ (47)
$`2a_6^{\mathrm{eff}(\mathrm{c})}X_2^{(BD^{},D)})+(a_4^{\mathrm{eff}(\mathrm{d})}+a_3^{\mathrm{eff}(\mathrm{d})}+a_3^{\mathrm{eff}(\mathrm{c})})Y^{(B,D^{}D)}_1`$
$`+(a_5^{\mathrm{eff}(\mathrm{d})}+a_5^{\mathrm{eff}(\mathrm{c})})Y_2^{(B,D^{}D)}2a_6^{\mathrm{eff}(\mathrm{d})}Y_3^{(B,D^{}D)}],`$
$`A(\overline{B}^0D^0\overline{D}^0)`$ $`=`$ $`(V_{cb}V_{cd}^{}+V_{ub}V_{ud}^{})a_2^{\mathrm{eff}}Y_1^{(B,DD^{})}V_{tb}V_{td}^{}[(a_3^{\mathrm{eff}(\mathrm{u})}+a_3^{\mathrm{eff}(\mathrm{c})})Y_1^{(B,DD^{})}`$ (48)
$`+(a_5^{\mathrm{eff}(\mathrm{u})}+a_5^{\mathrm{eff}(\mathrm{c})})Y_2^{(B,DD^{})}],`$
$`A(\overline{B}^0D^0\overline{D}^0)`$ $`=`$ $`(V_{cb}V_{cd}^{}+V_{ub}V_{ud}^{})a_2^{\mathrm{eff}}Y_1^{(B,D^{}D)}V_{tb}V_{td}^{}[(a_3^{\mathrm{eff}(\mathrm{u})}+a_3^{\mathrm{eff}(\mathrm{c})})Y_1^{(B,D^{}D)}`$ (49)
$`+(a_5^{\mathrm{eff}(\mathrm{u})}+a_5^{\mathrm{eff}(\mathrm{c})})Y_2^{(B,D^{}D)}],`$
$`A(B^{}D^0D^{})`$ $`=`$ $`V_{cb}V_{cd}^{}a_1^{\mathrm{eff}}X_1^{(BD,D^{})}+V_{ub}V_{ud}^{}a_1^{\mathrm{eff}}Y_1^{(B,DD^{})}V_{tb}V_{td}^{}[a_4^{\mathrm{eff}(\mathrm{c})}X_1^{(BD,D^{})}`$ (50)
$`+a_4^{\mathrm{eff}(\mathrm{d})}Y_1^{(B,DD^{})}2a_6^{\mathrm{eff}(\mathrm{d})}Y_3^{(B,DD^{})}],`$
$`A(B^{}D^0D^{})`$ $`=`$ $`V_{cb}V_{cd}^{}a_1^{\mathrm{eff}}X_1^{(BD^{},D)}+V_{ub}V_{ud}^{}a_1^{\mathrm{eff}}Y_1^{(B,D^{}D)}V_{tb}V_{td}^{}[(a_4^{\mathrm{eff}(\mathrm{c})}X_1^{(BD^{},D)}`$ (51)
$`2a_6^{\mathrm{eff}(\mathrm{c})}X_2^{(BD^{},D)})+a_4^{\mathrm{eff}(\mathrm{d})}Y_1^{(B,D^{}D)}2a_6^{\mathrm{eff}(\mathrm{d})}Y_3^{(B,D^{}D)}].`$
## Appendix B $`BVV`$ decays
$`A(\overline{B}^0D^+D_s^{})`$ $`=`$ $`V_{cb}V_{cs}^{}a_1^{\mathrm{eff}}X^{(BD^{},D_s^{})}V_{tb}V_{ts}^{}[a_4^{\mathrm{eff}(\mathrm{c})}X_1^{(BD^{},D_s^{})}`$ (52)
$`+(a_4^{\mathrm{eff}(\mathrm{d})}Y_1^{(B,D^{}D_s^{})}2a_6^{\mathrm{eff}(\mathrm{d})}Y_3^{(B,D^{}D_s^{})})],`$
$`A(\overline{B}^0D_s^+D_s^{})`$ $`=`$ $`V_{cb}V_{cd}^{}a_2^{\mathrm{eff}}Y_1^{(B,D_s^{}D_s^{})}V_{tb}V_{td}^{}[(a_3^{\mathrm{eff}(\mathrm{s})}+a_3^{\mathrm{eff}(\mathrm{c})})Y_1^{(B,D_s^{}D_s^{})}`$ (53)
$`+(a_5^{\mathrm{eff}(\mathrm{s})}+a_5^{\mathrm{eff}(\mathrm{c})})Y_2^{(B,D_s^{}D_s^{})}],`$
$`A(B^{}D^0D_s^{})`$ $`=`$ $`V_{cb}V_{cs}^{}a_1^{\mathrm{eff}}X^{(BD^{},D_s^{})}+V_{ub}V_{us}^{}a_1^{\mathrm{eff}}Y_1^{(B,D^{}D_s^{})}V_{tb}V_{ts}^{}[a_4^{\mathrm{eff}(\mathrm{c})}X_1^{(BD^{},D_s^{})}`$ (54)
$`+(a_4^{\mathrm{eff}(\mathrm{u})}Y_1^{(B,D^{}D_s^{})}2a_6^{\mathrm{eff}(\mathrm{u})}Y_3^{(B,D^{}D_s^{})})],`$
$`A(\overline{B}^0D^+D^{})`$ $`=`$ $`V_{cb}V_{cd}^{}[a^{\mathrm{eff}}X^{(BD^{},D^{})}+a_2^{\mathrm{eff}}Y_1^{(B,D^{}D^{})}]V_{tb}V_{td}^{}[(a_4^{\mathrm{eff}(\mathrm{c})}X^{(BD^{},D^{})}`$ (55)
$`+(a_4^{\mathrm{eff}(\mathrm{d})}+a_3^{\mathrm{eff}(\mathrm{d})}+a_3^{\mathrm{eff}(\mathrm{c})})Y_1^{(B,D^{}D^{})}+(a_5^{\mathrm{eff}(\mathrm{d})}+a_5^{\mathrm{eff}(\mathrm{c})})Y_2^{(B,D^{}D^{})}`$
$`2a_6^{\mathrm{eff}(\mathrm{d})}Y_3^{(B,D^{}D^{})}],`$
$`A(\overline{B}^0D^0\overline{D}^0)`$ $`=`$ $`(V_{cb}V_{cd}^{}+V_{ub}V_{ud}^{})a_2^{\mathrm{eff}}Y_1^{(B,D^{}D^{})}V_{tb}V_{td}^{}[(a_3^{\mathrm{eff}(\mathrm{u})}+a_3^{\mathrm{eff}(\mathrm{c})})Y_1^{(B,D^{}D^{})}`$ (56)
$`+(a_5^{\mathrm{eff}(\mathrm{u})}+a_5^{\mathrm{eff}(\mathrm{c})})Y_2^{(B,D^{}D^{})}],`$
$`A(B^{}D^0D^{})`$ $`=`$ $`V_{cb}V_{cd}^{}a^{\mathrm{eff}}X^{(BD^{},D^{})}+V_{ub}V_{ud}^{}a_1^{\mathrm{eff}}Y_1^{(B,D^{}D^{})}V_{tb}V_{td}^{}[a_4^{\mathrm{eff}(\mathrm{c})}X^{(BD^{},D^{})}`$ (57)
$`+a_4^{\mathrm{eff}(\mathrm{d})}Y_1^{(B,DD)}2a_6^{\mathrm{eff}(\mathrm{d})}Y_3^{(B,DD)}].`$ |
warning/0507/cond-mat0507315.html | ar5iv | text | # Operation of a superconducting nanowire quantum interference device with mesoscopic leads
## I Introduction
The Little-Parks effect concerns the electrical resistance of a thin cylindrically-shaped superconducting film and, specifically, the dependence of this resistance on the magnetic flux threading the cylinder LP ; Tinkham . It is found that the resistance is a periodic function of the magnetic field, with period inversely proportional to the cross-sectional area of the cylinder. Similarly, in a DC SQUID, the critical value of the supercurrent is periodic in magnetic field, with period inversely proportional to the area enclosed by the SQUID ring Tinkham .
In this Paper, we consider a mesoscopic analog of a DC SQUID. The analog consists of a device composed of a thin superconducting film patterned into two mesoscopic leads that are connected by a pair of (topologically) parallel, short, weak, superconducting wires. Thus, we refer to the device as an NQUID (superconducting nanowire quantum interference device). The only restriction that we place on the wires of the device is that they be thin enough for the order parameter to be taken as constant over each cross-section of a wire, varying only along the wire length. In principle, this condition of one-dimensionality is satisfied if the wire is much thinner than the superconducting coherence length $`\xi `$. In practice, it is approximately satisfied provided the wire diameter $`d`$ is smaller than $`4.4\xi `$ LikharevRMP . For thicker wires, vortices can exist inside the wires, and such wires may not be assumed to be one dimensional.
By the term mesoscopic we are characterizing phenomena that occur on length-scales larger than the superconducting coherence length $`\xi `$ but smaller than the electromagnetic penetration depth $`\lambda _{}`$ associated with magnetic fields applied perpendicular to the superconducting film. We shall call a lead mesoscopic if at least one of its two long dimensions is in the mesoscopic regime; the other dimension may be either mesoscopic or macroscopic. Thus, a weak magnetic field applied perpendicular to a mesoscopic lead will penetrate the lead without appreciable attenuation and without driving the lead from the homogeneous superconducting state to the Abrikosov vortex state. This is similar to the regime of operation of superconducting wire networks; see e.g., Ref. AB1 . The nanowires connecting the two leads are taken to be topologically parallel (i.e. parallel in the sense of electrical circuitry): these nanowires and edges of the leads define a closed geometrical contour, which will be referred to as the Aharonov-Bohm (AB) contour. In our approach, the nanowires are considered to be links sufficiently weak that any effects of the nanowires on the superconductivity in the leads can be safely ignored.
The theory presented here has been developed to explain experiments conducted on DNA-templated NQUIDs us\_in\_science . These experiments measure the electrical resistivity of a pair of superconducting nanowires suspended between long superconducting strips (see Fig. 1). In them, a current source is used to pass DC current from a contact on the far end of the left lead to one on the far end of the right lead. The voltage between the contacts is measured (and the resistance is hence determined) as a function of the magnetic field applied perpendicular to the plane of the strips.
In the light of the foregoing remarks, the multiple-connectedness of the device suggests that one should anticipate oscillations with magnetic field, e.g., in the device resistance. Oscillations are indeed observed. But they are distinct from the resistance oscillations observed by Little and Parks and from the critical current oscillations observed in SQUID rings. What distinguishes the resistance oscillations reported in Ref. us\_in\_science from those found, e.g., by Little and Parks? First, the most notable aspect of these oscillations is the value of their period. In the Little-Parks type of experiment, the period is given by $`\mathrm{\Phi }_0/2ab`$, where $`\mathrm{\Phi }_0(hc/2e)`$ is the superconducting flux quantum, $`2a`$ is the bridge separation, and $`b`$ is the bridge length, i.e., the superconducting flux quantum divided by the area of the AB contour (see Fig. 2). In a high-magnetic-field regime, such periodic behavior is indeed observed experimentally, with the length of the period somewhat shorter but of the same order of magnitude as in the AB effect us\_in\_science . However, in a low-magnetic-field regime, the observed period is appreciably smaller (in fact by almost two orders of magnitude for our device geometry). Second, because the resistance is caused by thermal phase fluctuations (i.e. phase slips) in very narrow wires, the oscillations are observable over a wide range of temperatures ($`1\text{K}`$). Third, the Little-Parks resistance is wholly ascribed to a rigid shift of the $`R(T)`$ curve with magnetic field, as $`T_\text{c}`$ oscillates. In contrast, in our system we observe a periodic broadening of the transition (instead of the Little-Parks—type rigid shift) with magnetic field. Our theory explains quantitatively this broadening via the modulation of the barrier heights for phase slips of the superconducting order parameter in the nanowires.
In the experiment, the sample is cooled in zero magnetic field, and the field is then slowly increased while the resistance is measured. At a sample-dependent field ($`5\text{mT}`$) the behavior switches sharply from a low-field to a high-field regime. If the high-field regime is not reached before the magnetic field is swept back, the low-field resistance curve is reproduced. However, once the high-field regime has been reached, the sweeping back of the field reveals phase shifts and hysteresis in the $`R(B)`$ curve. The experiments us\_in\_science mainly address rectangular leads that have one mesoscopic and one macroscopic dimension. Therefore, we shall concentrate on such strip geometries. We shall, however, also discuss how to extend our approach to generic (mesoscopic) lead shapes. We note in passing that efficient numerical methods, such as the boundary element method (BEM) BEM , are available for solving the corresponding Laplace problems.
This paper is arranged as follows. In Section II we construct a basic picture for the period of the magnetoresistance oscillations of the two-wire device, which shows how the mesoscopic size of the leads accounts for the anomalously short magnetoresistance period in the low-field regime. In Section III we concentrate on the properties of mesoscopic leads with regard to their response to an applied magnetic field, and in Section IV we extend the LAMH model to take into account the inter-wire coupling through the leads. Analytical expressions are derived for the short- and long-wire limits, whilst a numerical procedure is described for the general case. The predictions of the model are compared with data from our experiment in Section V, and we give some concluding remarks in Section VI. Certain technical components are relegated to the appendix, as is the analysis of example multiwire devices.
## II Origin of magnetoresistance oscillations
Before presenting a detailed development of the theory, we give an intuitive argument to account for the anomalously-short period of the magnetoresistance in the low-magnetic-field regime, mentioned above.
### II.1 Device geometry
The geometry of the devices studied experimentally is shown in Fig. 2. Five devices were successfully fabricated and measured. The dimensions of these devices are listed in Table 1, along with the short magnetoresistance oscillation period. The perpendicular penetration depth $`\lambda _{}`$ for the films used to make the leads is roughly $`70\mu \text{m}`$, and coherence length $`\xi `$ is roughly $`5\text{nm}`$.
### II.2 Parametric control of the state of the wires by the leads
The essential ingredients in our model are (i) leads, in which the applied magnetic field induces supercurrents and hence gradients in the phase of the order parameter, and (ii) the two wires, whose behavior is controlled parametrically by the leads through the boundary conditions imposed by the leads on the phase of the order parameters in the wires. For now, we assume that the wires have sufficiently small cross-sections that the currents through them do not feed back on the order parameter in the leads. (In Section III.4 we shall discuss when this assumption may be relaxed without altering the oscillation period.) The dissipation results from thermally activated phase slips, which cause the superconducting order parameter to explore a discrete family of local minima of the free energy. (We assume that the barriers separating these minima are sufficiently high to make them well-defined states.) These minima (and the saddle-point configurations connecting them) may be indexed by the net (i.e. forward minus reverse) number of phase slips that have occurred in each wire ($`n_1`$ and $`n_2`$, relative to some reference state). More usefully, they can be indexed by $`n_s=\mathrm{min}(n_1,n_2)`$ (i.e. the net number of phase slips that have occurred in both wires) and $`n_v=n_1n_2`$ (i.e. the number of vortices enclosed by the AB contour, which is formed by the wires and the edges of the leads). We note that two configurations with identical $`n_v`$ but distinct $`n_s`$ and $`n_s^{}`$ have identical order parameters, but differ in energy by
$$IV𝑑t=\frac{\mathrm{}}{2e}I\dot{\mathrm{\Theta }}𝑑t=\frac{h}{2e}I(n_s^{}n_s),$$
(1)
due to the work done by the current source supplying the current $`I`$, in which $`V`$ is the inter-lead voltage, $`\mathrm{\Theta }`$ is the inter-lead phase difference as measured between the two points half-way between the wires, and the Josephson relation $`\dot{\mathrm{\Theta }}=2eV/\mathrm{}`$ has been invoked. In our model, we assume that the leads are completely rigid. Therefore the rate of phase change, and thus the voltage, is identical at all points inside one lead. For sufficiently short wires, $`n_v`$ has a unique value, as there are no stable states with any other number of vortices.
Due to the screening currents in the left lead, induced by the applied magnetic field $`B`$ (and independent of the wires), there is a field-dependent phase $`\delta _{21,L}(B)=_1^2𝑑\stackrel{}{r}\stackrel{}{}\phi (B)`$ (computed below) accumulated in passing from the point at which wire 1 (the top wire) contacts the left (L) lead to the point at which wire 2 (the bottom wire) contacts the left lead (see Fig. 3). Similarly, the field creates a phase accumulation $`\delta _{21,R}(B)`$ between the contact points in the right (R) lead. As the leads are taken to be geometrically identical, the phase accumulations in them differ in sign only: $`\delta _{21,L}(B)=\delta _{21,R}(B)`$. We introduce $`\delta (B)=\delta _{21,L}(B)`$. In determining the local free-energy minima of the wires, we solve the Ginzburg-Landau equation for the wires for each vortex number $`n_v`$, imposing the single-valuedness condition on the order parameter,
$$\theta _{1,LR}\theta _{2,LR}+2\delta (B)=2\pi n_v.$$
(2)
This condition will be referred to as the phase constraint. Here, $`\theta _{1,LR}=_R^L𝑑\stackrel{}{r}\stackrel{}{}\phi (B)`$ is the phase accumulated along wire 1 in passing from the right to the left lead; $`\theta _{2,LR}`$ is similarly defined for wire 2.
Absent any constraints, the lowest energy configuration of the nanowires is the one with no current through the wires. Here, we adopt the gauge in which $`𝑨=By𝒆_x`$ for the electromagnetic vector potential, where the coordinates are as shown in Fig. 2. The Ginzburg-Landau expression for the current density in a superconductor is
$$𝑱\left(\mathbf{}\phi (𝒓)\frac{2e}{\mathrm{}}𝑨(𝒓)\right).$$
(3)
For our choice of gauge, the vector potential is always parallel to the nanowires, and therefore the lowest energy state of the nanowires corresponds to a phase accumulation given by the flux through the AB contour, $`\theta _{1,LR}=\theta _{2,LR}=2\pi Bab/\mathrm{\Phi }_0`$. As we shall show shortly, for our device geometry (i.e. when the wires are sufficiently short, i.e., $`bl`$), this phase accumulation may be safely ignored, compared to the phase accumulation $`\delta (B)`$ associated with screening currents induced in the leads. As the nanowires are assumed to be weak compared to the leads, to satisfy the phase constraint (2), the phase accumulations in the nanowires will typically deviate from their optimal value, generating a circulating current around the AB contour. As a consequence of LAMH theory, this circulating current results in a decrease of the barrier heights for phase slips, and hence an increase in resistance. The period of the observed oscillations is derived from the fact that whenever the magnetic field satisfies the relation
$$2\pi m=2\pi \frac{2abB}{\mathrm{\Phi }_0}+2\delta (B)$$
(4)
\[where $`m`$ is an integer and the factor of $`2`$ accompanying $`\delta (B)`$ reflects the presence of two leads\], there is no circulating current in the lowest in energy state, resulting in minimal resistance. Furthermore, the family of free energy-minima of the two-wire system (all of which, in thermal equilibrium, are statistically populated according to their energies) is identical to the $`B=0`$ case. The mapping between configurations at zero and nonzero $`B`$ fields is established by a shift of the index $`n_vn_vm`$. Therefore, as the sets of physical states of the wires are identical whenever the periodicity condition (4) is satisfied, at such values of $`B`$ the resistance returns to its $`B=0`$ value.
### II.3 Simple estimate of the oscillation period
In this subsection, we will give a “back of the envelope” estimate for the phase gain $`\delta (B)`$ in a lead by considering the current and phase profiles in one such lead. According to the Ginzburg-Landau theory, in a mesoscopic superconductor, subjected to a weak magnetic field, the current density is given by Eq. (3). Now consider an isolated strip-shaped lead used in the device. Far from either of the short edges of this lead, $`𝑨=By𝒆_x`$ is a London gauge London , i.e., along all surfaces of the superconductor $`𝑨`$ is parallel to them; $`𝑨0`$ in the center of the superconductor; and $`\mathbf{}𝑨=0`$. In this special case, the London relation <sup>1</sup><sup>1</sup>1Consider the case in which $`𝑨`$ is a London gauge everywhere (with our choice of gauge, $`𝑨=By𝒆_x`$, this is the case for an infinitely long strip). By using the requirement that $`\mathbf{}𝑨=0`$, together with Eq. (13b), we see that $`\varphi `$ satisfies the Laplace equation. We further insist that no current flows out of the superconductor, i.e., along all surfaces the supercurrent density, Eq. (12), is always parallel to the surface. Together with the requirement that along all surfaces $`𝑨`$ is parallel to them, this implies the boundary condition that $`𝒏\mathbf{}\varphi =0`$. Next, it can be shown that this boundary condition implies that $`\varphi `$ must be a constant function of position in order to satisfy the Laplace equation, and therefore Eq. (12) simplifies to read $`𝑱=(c/8\pi \lambda _{\text{eff}}^2)𝑨`$, which is known as the London relation. states that the supercurrent density is proportional to the vector potential in the London gauge. Using this relation, we find that the supercurrent density is $`𝑱(2e/\mathrm{})𝑨=(2e/\mathrm{})By𝒆_x`$, i.e., there is a supercurrent density of magnitude $`(2e/\mathrm{})Bl`$ flowing to the left at the top (long) edge of the strip and to the right at the bottom (long) edge. At the two short ends of the strip, the two supercurrents must be connected, so there is a supercurrent density of magnitude $`(2e/\mathrm{})Bl`$ flowing down the left (short) edge of the strip and up the right (short) edge (see Fig. 4).
Near the short ends of the strips, our choice of gauge no longer satisfies the criteria for being a London gauge, and therefore $`\mathbf{}\varphi `$ may be nonzero. As, in our choice of gauge, $`𝑨`$ points in the $`𝒆_x`$ direction, the supercurrent on the ends of the strip along $`𝒆_y`$ must come from the $`_y\varphi `$ term. Near the center of the short edge $`_y\varphi =2\pi c_1l/\mathrm{\Phi }_0B`$. The phase difference between the points $`(L,a)`$ and $`(L,a)`$ is therefore given by
$$\delta (B)=_a^a_y\varphi dy=\frac{2\pi c_1}{\mathrm{\Phi }_0}B\mathrm{\hspace{0.17em}2}al,$$
(5)
where we have substituted $`2\pi /\mathrm{\Phi }_0`$ for $`2e/\mathrm{}`$ and $`c_1(a/l)`$ is a function of order unity, which accounts for how the current flows around the corners. As we shall show, $`c_1`$ depends only weakly on $`a/l`$, and is constant in the limit $`al`$.
Finally, we obtain the magnetoresistance period by substituting Eq. (5) into Eq. (4):
$`\mathrm{\Delta }B=\left[\left({\displaystyle \frac{\mathrm{\Phi }_0}{c_1\mathrm{\hspace{0.17em}4}al}}\right)^1+\left({\displaystyle \frac{\mathrm{\Phi }_0}{2ab}}\right)^1\right]^1.`$ (6)
Thus, we see that for certain geometries the period is largely determined not by the flux threading through the geometric area $`2ab`$ but by the response of the leads and the corresponding effective area $`4al`$, provided the nanowires are sufficiently short (i.e. $`bl`$), justifying our assumption of ignoring the phase gradient induced in the nanowires by the magnetic field.
In fact, we can also make a prediction for the periodicity of the magnetoresistance at high magnetic fields, i.e., when vortices have penetrated the leads (see Section III.1). To do this, we should replace $`l`$ in Eq. (6) by the characteristic inter-vortex spacing $`r`$. Note that if $`r`$ is comparable to $`b`$, we can no longer ignore the flux through the AB contour. Furthermore, if $`rb`$ then the flux through the AB contour determines periodicity and one recovers the usual Aharonov-Bohm type of phenomenology.
## III Mesoscale superconducting leads
In this section and the following one we shall develop a detailed model of the leads and nanowires that constitute the mesoscopic device.
### III.1 Vortex-free and vorticial regimes
Two distinct regimes of magnetic field are expected, depending on whether or not there are trapped (i.e. locally stable) vortices inside the leads. As described by Likharev Likharev , a vortex inside a superconducting strip-shaped lead is subject to two forces. First, due to the the currents induced by the magnetic field there is a Magnus force pushing it towards the middle of the strip. Second, there is a force due to image vortices (which are required to enforce the boundary condition that no current flows out of the strip and into the vacuum) pulling the vortex towards the edge. When the two forces balance at the edge of the strip, there is no energy barrier preventing vortex penetration and vortices enter. Likharev has estimated of the corresponding critical magnetic field to be
$$H_\mathrm{s}\frac{\mathrm{\Phi }_0}{\pi d}\frac{1}{\xi a(1)},$$
(7)
where $`d(2l)`$ is the width of the strip and $`a(1)1`$ for strips that are much narrower than the penetration depth (i.e. for $`d\lambda `$).
Likharev has also shown that, once inside a strip, vortices remain stable inside it down to a much lower magnetic field $`H_{c1}`$, given by
$$H_{\text{c1}}=\frac{\mathrm{\Phi }_0}{\pi d}\frac{2}{d}\mathrm{ln}\left(\frac{d}{4\xi }\right).$$
(8)
At fields above $`H_{\text{c1}}`$ the potential energy of a vortex inside the strip is lower than for one outside (i.e. for a virtual vortex ref:topology\_note ). Therefore, for magnetic fields in the range $`H_{\text{c1}}<H<H_\text{s}`$ vortices would remain trapped inside the strip, but only if at some previous time the field were larger than $`H_\text{s}`$. This indicates that hysteresis with respect to magnetic field variations should be observed, once $`H`$ exceeds $`H_\text{s}`$ and vortices become trapped in the leads.
In real samples, in addition to the effects analyzed by Likharev, there are also likely to be locations (e.g. structural defects) that can pin vortices, even for fields smaller than $`H_{c1}`$, so the reproducibility of the resistance vs. field curve is not generally expected once $`H_\text{s}`$ has been surpassed.
As magnetic field at which vortices first enter the leads is sensitive to the properties of their edges, we expect only rough agreement with Likharev’s theory. For sample 219-4, using Likharev’s formula, we estimate $`H_\text{s}=11\text{mT}`$ (with $`\xi =5\text{nm}`$). The change in regime from fast to slow oscillations is found to occur at $`3.1\text{mT}`$ for that sample us\_in\_science . It is possible to determine the critical magnetic fields $`H_\text{s}`$ and $`H_{\text{c1}}`$ by the direct imaging of vortices. Although we do not know of such a direct measurement of $`H_\text{s}`$, $`H_{\text{c1}}`$ was determined by field cooling niobium strips, and found to agree in magnitude to Likharev’s estimate Martinis2004 .
### III.2 Phase variation along the edge of the lead
In the previous section it was shown that the periodicity of the magnetoresistance is due to the phase accumulations associated with the currents along the edges of the leads between the nanowires. Thus, we should make a precise calculation of the dependence of these currents on the magnetic field, and this we now do.
#### III.2.1 Ginzburg-Landau theory
To compute $`\delta (B)`$, we start with the Ginzburg-Landau equation for a thin film as our description of the mesoscopic superconducting leads:
$$\alpha \psi +\beta |\psi |^2\psi +\frac{1}{2m^{}}\left(\frac{\mathrm{}}{i}\mathbf{}\frac{e^{}}{c}𝑨\right)^2\psi =0.$$
(9)
Here, $`\psi `$ is the Ginzburg-Landau order parameter, $`e^{}`$ ($`=2e`$) is the charge of a Cooper pair and $`m^{}`$ is its mass, and $`\alpha `$ and $`\beta `$ may be expressed in terms of the coherence length $`\xi `$ and critical field $`H_\text{c}`$ via $`\alpha =\mathrm{}^2/2m^{}\xi ^2`$ and $`\beta =4\pi \alpha ^2/H_\text{c}^2`$.
The assumptions that the magnetic field is sufficiently weak and that the lead is a narrow strip (compared with the magnetic penetration depth) allow us to take the amplitude of the order parameter in the leads to have the value appropriate to an infinite thin film in the absence of the field. By expressing the order parameter in terms of the (constant) amplitude $`\psi _0`$ and the (position-dependent) phase $`\varphi (𝒓)`$, i.e.,
$$\psi (𝒓)=\psi _0e^{i\varphi (𝒓)},$$
(10)
the Ginzburg-Landau formula for the current density,
$$𝑱=\frac{e^{}\mathrm{}}{2m^{}i}\left(\psi ^{}\mathbf{}\psi \psi \mathbf{}\psi ^{}\right)\frac{e_{}^{}{}_{}{}^{2}}{m^{}c}\psi ^{}\psi 𝑨(𝒓),$$
(11)
becomes
$$𝑱=\frac{e^{}}{m^{}}\psi _0^2\left(\mathrm{}\mathbf{}\varphi (𝒓)\frac{e^{}}{c}𝑨(𝒓)\right),$$
(12)
and \[after dividing by $`e^{i\varphi (𝒓)}`$\] the real and imaginary parts of the Ginzburg-Landau equation become
$`0`$ $`=\left[\alpha \psi _0+\beta \psi _0^3+{\displaystyle \frac{1}{2m^{}}}\psi _0\left|\mathrm{}\mathbf{}\varphi (𝒓){\displaystyle \frac{e^{}}{c}}𝑨(𝒓)\right|^2\right],`$ (13a)
$`0`$ $`={\displaystyle \frac{\mathrm{}^2}{2m^{}i}}\psi _0\left(^2\varphi (𝒓){\displaystyle \frac{e^{}}{\mathrm{}c}}\mathbf{}𝑨(𝒓)\right).`$ (13b)
As long as any spatial inhomogeneity in the gauge-covariant derivative of the phase is weak on the length-scale of the coherence length $`[`$ i.e. $`\xi \left|\mathbf{}\varphi (𝒓)\frac{e^{}}{\mathrm{}c}𝑨(𝒓)\right|1`$ $`]`$, the third term in Eq. (13a) is much smaller than the first two and may be ignored, fixing the amplitude of the order parameter at its field-free infinite thin film value, viz., $`\overline{\psi _0}\sqrt{\alpha /\beta }`$. To compute $`\varphi (𝒓)`$ we need to solve the imaginary part of the Ginzburg-Landau equation.
#### III.2.2 Formulation as a Laplace problem
We continue to work in the approximation that the amplitude of the order parameter is fixed at $`\overline{\psi _0}`$. Starting from Eq. (13b), we see that for our choice of gauge, $`𝑨=By𝒆_x`$, the phase of the order parameter satisfies the Laplace equation, $`^2\varphi =0`$. We also enforce the boundary condition that no current flows out of the superconductor on boundary surface $`\mathrm{\Sigma }`$, whose normal is $`𝐧`$:
$`𝒏𝒋|_\mathrm{\Sigma }=0,`$ (14a)
$`𝒋\left(\mathbf{}\varphi {\displaystyle \frac{2\pi }{\mathrm{\Phi }_0}}𝑨\right).`$ (14b)
#### III.2.3 Solving the Laplace problem for the strip geometry
To solidify the intuition gained via the physical arguments given in Section II, we now determine the phase profile for an isolated superconducting strip in a magnetic field. This will allow us to determine the constant $`c_1`$ in Eq. (6), and hence obtain a precise formula for the magnetoresistance period. To this end, we solve Laplace’s equation for $`\varphi `$ subject to the boundary conditions (14). We specialize to the case of a rectangular strip <sup>2</sup><sup>2</sup>2This specialization is not necessary, but it is convenient and adequately illustrative.
In terms of the coordinates defined in Fig. 2, we expand $`\varphi (x,y)`$ as the superposition
$$\varphi (x,y)=\mathrm{\Theta }_{\text{L/R}}+\underset{k}{}\left(A_ke^{kx}+B_ke^{kx}\right)\mathrm{sin}(ky),$$
(15)
which automatically satisfies Laplace’s equation, although the boundary conditions remain to be satisfied. $`\mathrm{\Theta }_{\text{L(R)}}`$ is the phase at the the point in the left (right) lead located half-way between the wires. In other words $`\mathrm{\Theta }_L=\varphi (Lb,0)`$ and $`\mathrm{\Theta }_R=\varphi (L,0)`$ in the coordinate system indicated in Fig. 2. $`\mathrm{\Theta }_{\text{L/R}}`$ are not determined by the Laplace equation and boundary conditions, but will be determined later by the state of the nanowires.
We continue working in the gauge $`𝑨=By𝒆_x`$. The boundary conditions across the edges at $`y=\pm l`$ (i.e. the long edges) are $`_y\varphi (x,y=\pm l)=0`$. These conditions are satisfied by enforcing $`k_n=\pi (n+\frac{1}{2})/l`$, where $`n=0,1,2,\mathrm{}`$. The boundary conditions across the edges at $`x=\pm L`$ (i.e. the short edges) are $`_x\varphi (x=\pm L,y)=hy`$ (where $`h2\pi B/\mathrm{\Phi }_0`$). This leads to the coefficients in Eq. (15) taking the values
$`B_k`$ $`=A_k={\displaystyle \frac{h}{k_n^3l}}{\displaystyle \frac{(1)^n}{\mathrm{cosh}(k_nL)}}`$ $`(n=0,1,\mathrm{}),`$ (16)
and hence to the solution
$`\varphi (x,y)={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^n\mathrm{\hspace{0.17em}2}h}{k_n^3l\mathrm{cosh}(k_nL)}}\mathrm{sin}(k_ny)\mathrm{sinh}(k_nx).`$ (17)
Figure 5 shows the phase profiles in the leads, in the region close to the trench that separates the leads.
### III.3 Period of magnetoresistance for leads having a rectangular strip geometry
Using the result for the phase that we have just established, we see that the phase profile on the short edge of the strip at $`x=L`$ is given by
$$\varphi (L,y)=\frac{2hl^2}{\pi ^2}\underset{n=0}{\overset{\mathrm{}}{}}\frac{(1)^n}{(n+\frac{1}{2})^3}\mathrm{sin}\frac{\pi \left(n+\frac{1}{2}\right)y}{l},$$
(18)
where we have taken the limit $`L\mathrm{}`$. We would like to evaluate this sum at the points $`(x,y)=(L,\pm a)`$. This can be done numerically. For nanowires that are close to each other (i.e. for $`al`$), an approximate value can be found analytically by expanding in a power series in $`a`$ around $`y=0`$:
$$\begin{array}{c}\hfill \varphi (L,a)=\varphi (L,0)+a\frac{}{y}\varphi (L,y)|_{y=0}\\ \hfill +\frac{a^2}{2}\frac{^2}{y^2}\varphi (L,y)|_{y=0}+O(a^3).\end{array}$$
(19)
The first and third terms are evidently zero, as $`\varphi `$ is an odd function of $`y`$. The second term can be evaluated by changing the order of summation and differentiation. (Higher-order terms are harder to evaluate, as the changing of the order of summation and differentiation does not work for them.) Thus, to leading order in $`a`$ we have
$`\varphi (L,a){\displaystyle \frac{8G}{\pi ^2}}hla,`$ (20)
where $`G_{n=0}^{\mathrm{}}\frac{(1)^n}{(2n+1)^2}0.916`$ is the Catalan number (see Ref. catalan\_n ). This linear approximation is plotted, together with the actual phase profile obtained by the numerical evaluation of Eq. (18), in Fig. 6. Hence, the value of $`c_1`$ in Eq. (5) becomes $`c_1=8G/\pi ^20.74`$, and Eq. (6) becomes
$`B={\displaystyle \frac{\mathrm{\Phi }_0}{2\pi }}h={\displaystyle \frac{\pi ^2}{8G}}{\displaystyle \frac{\mathrm{\Phi }_0}{4al}}.`$ (21)
To obtain this result we used the relation $`\delta (B)/2=\varphi (L,a)`$.
### III.4 Bridge-lead coupling
In order to simplify our analysis we have assumed that the nanowires do not exert any influence on the order parameter in the leads. We examine the justification for this assumption in the setting of the experiment that we are attempting to describe us\_in\_science .
The assumption will be valid if the bending of the phase of the order parameter, in order to accommodate any circulating current around the AB contour, occurs largely in the nanowires. As the phase of the order parameter in the leads satisfies the Laplace equation, which is linear, we can superpose the circulating-current solution with the previously-obtained magnetic-field-induced solution. The boundary conditions on the right lead for the circulating-current solution are $`𝒏\mathbf{}\varphi =0`$ everywhere, except at the two points where the nanowires are attached to the lead \[i.e. at $`(x,y)=(L,\pm a)`$\]. Treating the nanowires as point current sources, the boundary condition on the short edge of the right lead is $`_x\varphi =I(\mathrm{\Phi }_0/H_\text{c}^2s\xi )(\delta (y+a)\delta (ya)`$), where $`I`$ is the current circulating in the loop, and $`H_\text{c}`$, $`s`$, $`\xi `$ are the film critical field, thickness, and coherence length. By using the same expansion as before, Eq. (15), we obtain the coefficients of the Fourier-series in the long strip limit:
$`A_k=I(\mathrm{\Phi }_0/H_\text{c}^2s\xi ){\displaystyle \frac{2\mathrm{sin}(ka)}{kl\mathrm{exp}(kl)}}.`$ (22)
Having the coefficients of the Fourier series, we can find the phase difference in the right lead between the two points at which the nanowires connect to the right lead, induced in this lead by the current circulating in the loop:
$`\delta _{cc}=2I(\mathrm{\Phi }_0/H_\text{c}^2s\xi ){\displaystyle \underset{n=0}{\overset{k=1/w}{}}}{\displaystyle \frac{\mathrm{sin}^2(ka)}{kl}}\mathrm{ln}(2l/\pi w).`$ (23)
Here, we have introduced a large wave-vector $`k`$ cut-off at the inverse of the width $`w`$ of the wire. On the other hand, the current flowing through the wire is
$$\frac{\xi H_\text{c}^2}{\mathrm{\Phi }_0}ws\frac{\mathrm{\Delta }\theta }{b},$$
(24)
where $`\xi `$, $`H_\text{c}`$, and $`s`$ are the wire coherence length, critical field, and height (recall that $`b`$ is the wire length). To support a circulating current that corresponds to a phase accumulation of $`\mathrm{\Delta }\theta `$ along one of the wires, the phase difference between the two nanowires in the lead must be on the order of
$$\delta _{cc}=\mathrm{\Delta }\theta \frac{w}{b}\frac{\left(H_\text{c}^2s\xi \right)_{\text{wire}}}{\left(H_\text{c}^2s\xi \right)_{\text{film}}}\mathrm{ln}\left(\frac{2l}{\pi w}\right).$$
(25)
For our experiments us\_in\_science , we estimate that the ratio of $`\delta _{cc}`$ to $`\mathrm{\Delta }\theta `$ is always less than $`20\%`$, validating the assumption of weak coupling.
### III.5 Strong nanowires
We remark that the assumption of weak nanowires is not obligatory for the computation the magnetoresistance period. Dropping this assumption would leave the period of the magnetoresistance oscillations unchanged.
To see this, consider $`\varphi _{11}`$, i.e., the phase profile in the leads that corresponds to the lowest energy solution of the Ginzburg-Landau equation at field corresponding to the first resistance minimum \[i.e. at B being the first non-zero solution of Eq. (4)\]. For this case, and for short wires, the phase gain along the wires is negligible, whereas the phase gain in the leads is $`2\pi `$, even for wires with large critical current. Excited states, with vortices threading the AB contour, can be constructed by the linear superposition of $`\varphi _{11}`$ with $`\varphi _{0n_v}`$, where $`\varphi _{0n_v}`$ is the phase profile with $`n_v`$ vortices at no applied magnetic field.
This construction requires that the nanowires are narrow, but works independently of whether nanowires are strong or weak, in the limit that $`HH_\text{c}`$. The energy of the lowest energy state always reaches its minimum when the applied magnetic field is such that there is no phase gain (i.e. no current) in the nanowires. By the above construction, it is clear that the resistance of the device at this field is the same as at zero field, and therefore the minimum possible.
Therefore, our calculation of the period is valid, independent of whether the nanowires are weak or strong. However, the assumption of weak nanowires is necessary for the computation of magnetoresistance amplitude, which we present in the following section.
## IV Parallel superconducting nanowires and intrinsic resistance
In this section we consider the intrinsic resistance of the device. We assume that this resistance is due to thermally activated phase slips (TAPS) of the order parameter, and that these occur within the nanowires. Equivalently, these processes may be thought of as thermally activated vortex flow across the nanowires. Specifically, we shall derive analytical results for the asymptotic cases of nanowires that are either short or long, compared to coherence length, i.e. Josephson junctions IZ ; AH or Langer-Ambegaokar-McCumber-Halperin (LAMH) wires LA ; MH ; see also Ref. little . We have not been able to find a closed-form expression for the intrinsic resistance in the intermediate-length regime, so we shall consider that case numerically.
There are two (limiting) kinds of experiments that may be performed: fixed total current and fixed voltage. In the first kind, a specified current is driven through the device and the time-averaged voltage is measured. Here, this voltage is proportional to the net number of phase slips (in the forward direction) per unit time, which depends on the height of the free-energy barriers for phase slips. Why do we expect minima in the resistance at magnetic fields corresponding to $`2\delta =2m\pi `$ and maxima at $`2\delta =(2m+1)\pi `$ for $`m`$ integral, at least at vanishingly small total current through both wires? For $`2\delta =2m\pi `$ the nanowires are unfrustrated, in the sense that there is no current through either wire in the lowest local minimum of the free energy. On the other hand, for $`2\delta =(2m+1)\pi `$ the nanowires are maximally frustrated: there is a nonzero circulating current around the AB contour. Quite generally, the heights of the free-energy barriers protecting locally stable states decrease with increasing current through a wire, and thus the frustrated situation is more susceptible to dissipative fluctuations, and hence shows higher resistance. Note, however, that due to the inter-bridge coupling caused by the phase constraint, the resistance of the full device is more subtle than the mere addition of the resistances of two independent, parallel nanowires, both carrying the requisite circulating current.
In the second kind of experiment, a fixed voltage is applied across the device and the total current is measured. In this situation, the inter-lead voltage is fixed, and therefore the phase drop along each wire is a fixed function of time. Hence, there is no inter-bridge coupling in the fixed voltage regime. Therefore, the resistance of the device would not exhibit magnetic field dependence. If the voltage is fixed far away from the wires, but not in the immediate vicinity of the wires, so that the phase drop along each wire is not rigidly fixed, then some of the magnetic field dependence of the resistance would be restored. In our experiments on two-wire devices, we believe that the situation lies closer to the fixed current limit than to the fixed voltage limit, and therefore we shall restrict our attention to the former limit.
In the fixed-current regime, the relevant independent thermodynamic variable for the device is the total current through the pair of wires, i.e., $`II_1+I_2`$. Therefore, the appropriate free energy to use, in obtaining the barrier heights for phase slips, is the Gibbs free energy $`G(I)`$, as discussed by McCumber M68 , rather than the Helmholtz free energy $`F(\mathrm{\Theta })`$ <sup>3</sup><sup>3</sup>3Recall that the Helmholtz free energy is obtained by minimizing the Ginzburg-Landau free energy functional with respect to the order parameter function $`\psi (𝒓)`$, subject to the phase accumulation constraint $`_\text{L}^\text{R}𝑑𝒓\mathbf{}\varphi =\theta `$.. In the Helmholtz free energy the independent variable can be taken to be $`\mathrm{\Theta }\mathrm{\Theta }_\text{L}\mathrm{\Theta }_\text{R}`$, i.e., the phase difference across the center of the “trench,” defined modulo $`2\pi `$. $`G(I)`$ is obtained from $`F(\mathrm{\Theta })`$ via the appropriate Legendre transformation:
$$G(I)=F(\mathrm{\Theta })\frac{\mathrm{}}{2e}I\mathrm{\Theta },$$
(26)
where the second term represents the work done on the system by the external current source. $`F(\mathrm{\Theta })`$ is the sum of the Helmholtz free energies for the individual nanowires:
$$F(\mathrm{\Theta })=F_1(\theta _1)+F_2(\theta _2),$$
(27)
where $`F_{1(2)}(\theta _{1(2)})`$ is the Ginzburg-Landau free energy for first (second) wire and a simplified notation has been used $`\theta _1\theta _{1,LR}`$ and $`\theta _2\theta _{2,LR}`$. $`\theta _1`$ and $`\theta _2`$ are related to each other and to $`\mathrm{\Theta }`$ through the phase constraint Eq. (2).
### IV.1 Short nanowires: Josephson junction limit
If the nanowires are sufficiently short, they may be treated as Josephson junctions. Unlike the case of long nanowires, described in the following subsection, in this Josephson regime there is no metastability, i.e., the free energy of each junction is a single-valued function of the phase difference, modulo $`2\pi `$, across it. The phase constraint then implies that there is a rigid difference between the phases across the two junctions. As a consequence, $`n_v`$ can be set to zero. The Gibbs free energy in such a configuration is then
$$\begin{array}{c}\hfill G(I)=\frac{\mathrm{}}{2e}\left(I_{\text{c1}}\mathrm{cos}(\theta _1)+I_{\text{c2}}\mathrm{cos}(\theta _2)+I\mathrm{\Theta }\right),\end{array}$$
(28)
where $`I_{\text{c1}}`$ and $`I_{\text{c2}}`$ are the critical currents for the junctions. In thermodynamic equilibrium, the Gibbs free energy must be minimized, so the dependent variable $`\mathrm{\Theta }`$ must be chosen such that $`G(I)/\mathrm{\Theta }=0`$.
Using $`\theta _1=\mathrm{\Theta }+\delta `$ and $`\theta _2=\mathrm{\Theta }\delta `$, $`G(I)`$ may be rewritten in the form
$$\begin{array}{c}\hfill \stackrel{~}{G}(I)=\frac{\mathrm{}}{2e}\left(\sqrt{(I_{\text{c1}}+I_{\text{c2}})^2\mathrm{cos}^2\delta +(I_{\text{c1}}I_{\text{c2}})^2\mathrm{sin}^2\delta }\mathrm{cos}(\vartheta )+I\vartheta _1\right),\end{array}$$
(29)
where we have shifted the free energy by an additive constant $`\left(\frac{\mathrm{}}{2e}\right)I\mathrm{tan}^1\left[\left(\frac{I_{\text{c1}}I_{\text{c2}}}{I_{\text{c2}}+I_{\text{c1}}}\right)\mathrm{tan}\delta \right]`$, and $`\vartheta \mathrm{\Theta }+\mathrm{tan}^1\left[\left(\frac{I_{\text{c1}}I_{\text{c2}}}{I_{\text{c2}}+I_{\text{c1}}}\right)\mathrm{tan}\delta \right]`$. In this model, the option for having $`I_{\text{c1}}I_{\text{c2}}`$ is kept open. Equation (29) shows that, up to an additive constant, the free energy of the two-junction device is identical to that of an effective single-junction device with an effective $`I_\text{c}`$, which is given by
$$I_\text{c}=\sqrt{(I_{\text{c1}}+I_{\text{c2}})^2\mathrm{cos}^2\delta +(I_{\text{c1}}I_{\text{c2}})^2\mathrm{sin}^2\delta }.$$
(30)
Thus, we may determine the resistance of the two-junction device by applying the well-known results for a single junction, established by Ivanchenko and Zil’berman IZ and by Ambegaokar and Halperin AH :
$`R`$ $`=R_\text{n}{\displaystyle \frac{2(1x^2)^{1/2}}{x}}\mathrm{exp}\left(\gamma (\sqrt{1x^2}+x\mathrm{sin}^1x)\right)\mathrm{sinh}(\pi \gamma x/2),`$ (31a)
$`x`$ $`I/I_\text{c},\gamma \mathrm{}I_\text{c}/ek_\text{B}T,`$ (31b)
where $`R_\text{n}`$ is the normal-state resistance of the two-junction device. This formula for $`R`$ holds when the free-energy barrier is much larger than $`k_\text{B}T`$, so that the barriers for phase slips are high. References IZ ; AH provide details on how to calculate the resistance in the general case of an over damped junction, which includes that of shallow barriers. Figure 7 shows the fits to the resistance, computed using Eqs. (3031), as a function of temperature, magnetic field, and total current for sample 219-4. Observe that both the field- and the temperature-dependence are in good agreement with experimental data. In Section V.2.2 we make more precise contact between theory and experiment, and explain how the data have been fitted. We also note that, as it should, our Josephson junction model exactly coincides with our extension of the LAMH model in the limit of very short wires and for temperatures for which the barrier-crossing approximation is valid.
### IV.2 Longer nanowires: LAMH regime
In this section we describe an extension of the LAMH model of resistive fluctuations in a single narrow wire LA ; MH , which we shall use to make a quantitative estimate of the voltage across the two-wire device at a fixed total current. In this regime the nanowires are sufficiently long that they behave as LAMH wires. We shall only dwell on two-wire systems, but we note in passing that the model can straightforwardly be extended to more complicated sets of lead interconnections, including periodic, grating-like arrays (see Appendix A).
As the sample is not simply connected, i.e., there is a hole inside the AB contour, it is possible that there are multiple metastable states that can support the total current. These states differ by the number of times the phase winds along paths around the AB contour. The winding number $`n_v`$ changes whenever a vortex (or an anti-vortex) passes across one of the wires.
In the present theory, we include two kinds of processes that lead to the generation of a voltage difference between the the leads; see Fig. 8. In the first kind of process (Fig. 8a), two phase slips occur simultaneously: a vortex passes across the top wire and, concurrently, an anti-vortex passes across the bottom wire (in the opposite direction), so that the winding number remains unchanged. In the second kind of process (Fig. 8b), the phase slips occur sequentially: a vortex (or anti-vortex) enters the AB contour by passing across the top (or bottom) wire, stays inside the contour for some time-interval, and then leaves the AB contour through the bottom (or top) wire ref:topology\_note .
Our goal is to extend LAMH theory to take into account the influence of the wires on each other. In Appendix C, we review some necessary ingredients associated with the LAMH theory of a single wire. As the wires used in the experiments are relatively short (i.e. 10 to 20 zero-temperature coherence lengths in length), we also take care to correctly treat the wires as being of finite length.
Recall that we are considering experiments performed at a fixed total current, and accordingly, in all configurations of the order parameter this current must be shared between the top and bottom wires. We shall refer to this sharing,
$$I=I_1+I_2,$$
(32)
as the total current constraint. Let us begin by considering a phase-slip event in a device with an isolated wire. While the order parameter in that wire pinches down, the end-to-end phase accumulation must adjust to maintain the prescribed value of the current through the wire. Now consider the two-wire device, and consider a phase slip event in one of the wires. As in the single-wire case, the phase accumulation will adjust, but in so doing it will alter the current flowing through other wire. Thus, in the saddle-point configuration of the two-wire system the current splitting will differ from that in the locally stable initial (and final) state.
Taking into account the two kinds of phase-slip processes, and imposing the appropriate constraints (i.e. the total current constraint and the phase constraint), we construct the possible metastable and saddle-point configurations of the order parameter in the two-wire system. Finally, we compute the relevant rates of thermally activated transitions between these metastable states, construct a Markov chain markov , and determine the steady-state populations of these states. Thus, we are able to evaluate the time-average of the voltage generated between the leads at fixed current due to these various dissipative fluctuations. We mention that we have not allowed for wires of distinct length or constitution (so that the Ginzburg-Landau parameters describing them are taken to be identical). This is done solely to simplify the analysis; extensions to more general cases would be straightforward but tedious.
#### IV.2.1 Parallel pair of nanowires
The total Gibbs free energy for the two-wire system is given by
$$G(I)=F_1(\theta _1)+F_2(\theta _2)4\mathrm{\Theta }(J_1+J_2).$$
(33)
Here, we have followed MH by rewriting the current-phase term in terms of dimensionless currents in wires $`i=1,2`$, i.e., $`J_i`$ defined via $`I_i=8\pi cJ_i/\mathrm{\Phi }_0`$. Moreover, $`\frac{H_\text{c}^2\xi \sigma }{8\pi }`$ is the condensate energy density per unit length of wire, and $`F_i(\theta _i)`$ is the Helmholtz free energy for a single wire along which there is a total phase accumulation of $`\theta _i`$. The precise form of $`F_i(\theta _i)`$ depends on whether the wire is in a metastable or saddle-point state.
We are concerned with making stationary the total Gibbs free energy at specified total current $`I`$, subject to the phase constraint, Eq. (2). This can be accomplished by making stationary the Helmholtz free energy on each wire, subject to both the total current constraint and the phase constraint, but allowing $`\theta _1`$ and $`\theta _2`$ to vary so as to satisfy these constraints—in effect, adopting the total current $`I`$ as the independent variable. The stationary points of the Helmholtz free energy for a single wire are reviewed in Appendix C as implicit functions of $`\theta _i`$, i.e., the end-to-end phase accumulation along the wire. The explicit variable used there is $`J_i`$, which is related to $`\theta _i`$ via Eq. (74).
#### IV.2.2 Analytical treatment in the limit of long nanowires
In the long-wire limit, we can compute the resistance analytically by making use of the single-wire free energy and end-to-end phase accumulation derived by Langer and Ambegaokar LA (and extended by McCumber M68 for the case of the constant-current ensemble). Throughout the present subsection we shall be making an expansion in powers of $`1/b`$, where $`b`$ is the length of the wire measured in units of the coherence length, keeping terms only to first order in $`1/b`$. Thus, one arrives at formulæ for the end-to-end phase accumulations and Helmholtz free energies for single-wire metastable (m) and saddle-point (sp) states M68 :
$`\theta _\text{m}(\kappa )`$ $`=\kappa b,`$ (34a)
$`\theta _{\text{sp}}(\kappa )`$ $`=\kappa b+2\mathrm{tan}^1\left({\displaystyle \frac{13\kappa ^2}{2\kappa ^2}}\right)^{1/2},`$ (34b)
$`F_\text{m}(\kappa )`$ $`=\left(b(1\kappa ^2)^2\right),`$ (34c)
$`F_{\text{sp}}(\kappa )`$ $`=\left(b(1\kappa ^2)^2{\displaystyle \frac{8\sqrt{2}}{3}}\sqrt{13\kappa ^2}\right),`$ (34d)
where $`\kappa `$ is defined via $`J_i=\kappa _i(1\kappa _i^2)`$. In the small-current limit, one can make the further simplification that $`J_i\kappa _i`$; henceforth we shall keep terms only up to first order in $`\kappa `$. To this order, the phase difference along a wire in a saddle-point state becomes
$$\theta _{\text{sp}}=\kappa b+\pi 2\sqrt{2}\kappa .$$
(35)
Next, we make use of these single-wire LAMH results to find the metastable and saddle-point states of the two-wire system, and use them to compute the corresponding barrier heights and, hence, transition rates. At low temperatures, it is reasonable to expect that only the lowest few metastable states will be appreciably occupied. These metastable states, as well as the saddle-point states between them, correspond to pairs, $`\kappa _1`$ and $`\kappa _2`$, one for each wire, that satisfy the total current constraint as well as the phase constraint:
$`\kappa _1+\kappa _2=J,`$ (36)
$`\theta _1(\kappa _1)\theta _2(\kappa _2)=2\pi n_v+2\delta ,`$ (37)
where we need to substitute the appropriate $`\theta _{\text{m}/\text{sp}}(\kappa _i)`$ from Eqs. (34a35) for $`\theta _i(\kappa _i)`$.
In the absence of a magnetic field (i.e. $`\delta =0`$), the lowest energy state is the one with no circulating current, and the current split evenly between the two wires. This corresponds to the solution of Eqs. (3637) with $`n=0`$, together with the substitution (34a) for $`\theta _i(\kappa _i)`$ for both wires (i.e. $`\theta _1=\kappa _1b`$ and $`\theta _2=\kappa _2b`$). Thus we arrive at the solution:
$`\kappa _1`$ $`=J/2,`$ $`\theta _1`$ $`=bJ/2,`$ (38a)
$`\kappa _2`$ $`=J/2,`$ $`\theta _2`$ $`=bJ/2.`$ (38b)
If we ignore the lowest (excited) metastable states then only a parallel phase-slip process is allowed. The saddle point for a parallel phase slip corresponds to a solution of Eqs. (3637) with $`n=0`$ and the substitution (35) for $`\theta _i(\kappa _i)`$ for both wires:
$`\kappa _1`$ $`=J/2,`$ $`\theta _1`$ $`=bJ/2+\pi 2\sqrt{2}J/2,`$ (39a)
$`\kappa _2`$ $`=J/2,`$ $`\theta _2`$ $`=bJ/2+\pi 2\sqrt{2}J/2.`$ (39b)
The change in the phase difference across the center of the trench, $`\mathrm{\Delta }\mathrm{\Theta }[\mathrm{\Theta }_{\text{sp}}\mathrm{\Theta }_\text{m}]`$, is $`\pi 2\sqrt{2}\kappa `$ for a forward phase slip, and $`\pi 2\sqrt{2}\kappa `$ for a reverse phase slip. The Gibbs free-energy barrier for the two kinds of phase slips, computed by subtracting the Gibbs free energy for the ground state from that of the saddle-point state, is
$$\mathrm{\Delta }G=\left(\frac{16\sqrt{2}}{3}\pm 4J\pi \right).$$
(40)
The former free-energy is obtained by substituting Eq. (34c) into Eq. (33) for both wires; the latter one is obtained by substituting Eq. (34d) into Eq. (33) for both wires. We note that the Gibbs free-energy barrier heights for parallel phase slips (in both the forward and reverse directions) are just double those of the LAMH result for a single wire. From the barrier heights, we can work out the generated voltage by appealing to the Josephson relation, $`V=(\mathrm{}/2e)\dot{\mathrm{\Theta }}`$, and to the fact that each phase slip corresponds to the addition (or subtraction) of $`2\pi `$ to the phase. Hence, we arrive at the current-voltage relation associated with parallel phase slips at $`\delta =0`$:
$`V_{\delta =0\text{, par}}={\displaystyle \frac{\mathrm{}}{e}}\mathrm{\Omega }e^{\beta \frac{16\sqrt{2}}{3}}\mathrm{sinh}\left(I/I_0\right),`$ (41)
where the prefactor $`\mathrm{\Omega }`$ may be computed using time-dependent Ginzburg-Landau theory or extracted from experiment, and $`I_0=4e/\beta h`$.
If we take into account the two lowest excited states, which we ignored earlier, then voltage can also be generated via sequential phase slips (in addition to the parallel ones, treated above). To tackle this case, we construct a diagram in which the vertices represent the metastable and saddle-point solutions of Eqs. (3637), and the edges represent the corresponding free energy barriers; see Fig. LABEL:0\_diag. Pairs of metastable-state vertices are connected via two saddle-point-state vertices, corresponding to a phase slip on either the top or the bottom wire. To go from one metastable state to another, the system must follow the edge out of the starting metastable state leading to the desired saddle-point state. We assume that, once the saddle-point state is reached, the top of the barrier has been passed and the order parameter relaxes to the target metastable state. (To make the graph more legible, we have omitted drawing the edge that corresponds to this relaxation process.) To find the Gibbs free-energy difference between a metastable state and a saddle-point state, we need to know the phase difference across the center of the trench. To resolve the ambiguity of $`2\pi `$ in the definition of $`\mathrm{\Theta }`$, the phase difference can be found by following the wire with no phase slip. To further improve the legibility of Fig. LABEL:0\_diag, the free-energy barriers are listed in a separate table to the right. Note, that a phase slip on just one of the wires, being only half of the complete process, can be regarded to a gain in phase of $`\pm \pi `$ for the purposes of calculating voltage, as indicated in both the graph and the table.
Once the table of barrier heights has been computed, we can construct a Markov chain on a directed graph, where the metastable states are the vertices—in effect, an explicit version of our diagram. In general, each pair of neighboring metastable states, $`s_n`$ and $`s_{n+1}`$, are connected by four directed edges:
$`s_n`$ $`\underset{\text{top}}{\overset{}{}}s_{n+1}`$ $`s_n`$ $`\underset{\text{bottom}}{\overset{}{}}s_{n+1}`$ (42a)
$`s_n`$ $`\underset{\text{top}}{\overset{}{}}s_{n+1}`$ $`s_n`$ $`\underset{\text{bottom}}{\overset{}{}}s_{n+1}`$ (42b)
where the probability to pass along a particular edge is given by $`P()=\mathrm{exp}\beta \mathrm{\Delta }G_{()}`$, in which $`\mathrm{\Delta }G_{()}`$ may be read off from the table in Fig. LABEL:0\_diag.
We denote the occupation probability of the $`n^{\text{th}}`$ metastable state by $`o_n`$, where $`n`$ corresponds to the $`n`$ in the phase constraint (2). $`o_n`$ may be computed in the standard way, by diagonalizing the matrix representing the Markov chain markov . Each move in the Markov chain can be associated with a gain in phase across the device of $`\pm \pi `$, as specified in Fig. LABEL:0\_diag. Thus, we may compute the rate of phase-gain, and hence the voltage:
$$V=\frac{\mathrm{\Omega }\mathrm{}}{4e}\underset{nm}{}\frac{o_n}{g_{n,m}}\left(P(s_n\underset{\text{top}}{\overset{}{}}s_m)P(s_n\underset{\text{bot}}{\overset{}{}}s_m)\right),$$
(43)
where the rate prefactor $`\mathrm{\Omega }`$ is to be determined, $`nm`$ indicates that the sum runs over neighboring states only, and $`g_{n,m}`$ keeps track of the sign of the phase-gain for reverse phase-slips:
$$g_{n,m}=\{\begin{array}{cc}\hfill 1,& \text{ if }m>n,\\ \hfill 1,& \text{ if }m<n.\end{array}$$
(44)
For the case $`\delta =0`$, and keeping the bottom three states only, the voltage generated via sequential phase slips turns out to be
$`V_{\delta =0\text{, seq}}={\displaystyle \frac{2\mathrm{}}{e}}\mathrm{\Omega }e^{\beta \left(\frac{8\sqrt{2}}{3}+\frac{\pi ^2}{b}\right)}\mathrm{sinh}(I/2I_0).`$ (45)
Having dealt with the case of $`\delta =0`$ (and hence obtained the value of the resistance at magnetic fields corresponding to resistance minima), we now turn to the case of $`\delta =\pi /2`$, i.e., resistance maxima.
In this half-flux quantum situation, there are two degenerate lowest-energy states, with opposite circulating currents. These states are connected by saddle-point states in which a phase-slip is occurring on either the top or bottom wire. The diagram of the degenerate ground states and the saddle-point states connecting them is shown in Fig. 10. By comparing the diagram with the associated Table, it is easy to see that the free-energy barriers are biased by the current, making clockwise traversals of Fig. 10 more probable than counter-clockwise traversals. As there are only two metastable states being considered, and as they are degenerate, it is unnecessary to go through the Markov chain calculation; clearly, the two states each have a population of $`1/2`$. The voltage being generated by the sequential phase-slip is then given by
$`V_{\delta =\pi /2\text{, seq}}={\displaystyle \frac{\mathrm{}}{2e}}\mathrm{\Omega }e^{\beta \left(\frac{8\sqrt{2}}{3}\frac{\pi ^2}{b}\right)}\mathrm{sinh}(I/2I_0).`$ (46)
$`V_{\delta =\pi /2\text{, seq}}`$ is larger than the sum of $`V_{\delta =0\text{, seq}}`$ and $`V_{\delta =0\text{, par}}`$, so, as expected, the resistance is highest at magnetic fields corresponding to $`\delta =\pi /2`$. For very long wires, the perturbation of one wire when a phase slip occurs in the other is very small, and therefore we expect that the dependence of resistance on magnetic field will decrease with wire length. Indeed, for very long wires, the difference in barrier heights to sequential phase slips between the $`\delta =0`$ and $`\delta =\pi /2`$ cases disappears (i.e. Eq. (45) and Eq. (46) agree when $`b1`$).
#### IV.2.3 Numerical treatment for intermediate-length nanowires
Instead of using the long-wire approximation, Eqs. (34a-34d), we can use the exact functions for the end-to-end phase accumulation along a wire $`\theta (J(\kappa ))`$, and the Helmholtz free energy $`F_{\text{m}/\text{sp}}(J(\kappa ))`$. By dropping the long-wire approximation, as the temperature approaches $`T_\text{c}`$ and the coherence length decreases the picture correctly passes to the Josephson limit. In this approach, the total current and the phase constraints must be solved numerically, as $`\theta (J(\kappa ))`$ is a relatively complicated function. Figure 12 provides an illustration of how, for a single wire, the function $`J(\theta )`$ depends on its length. We shall, however, continue to use the barrier-crossing approximation. Because the barriers get shallower near $`T_\text{c}`$, our results will become unreliable (and, indeed, incorrect) there.
The form of the order parameter that satisfies the Ginzburg-Landau equation inside the wire is expressed in Eqs. (73b59). Therefore, to construct the functions $`\theta (J)`$ and $`F_{\text{m}/\text{sp}}(J)`$ \[i.e. Eqs. (7475)\], we need to find $`u_0(J)`$, i.e., the squared amplitude of the order parameter in the middle of the wire. Hence, we need to ascertain suitable boundary conditions obeyed by the order parameter at the ends of the wire. For thin wires, a reasonable hypothesis is that the amplitude of the order parameter at the ends of the wire matches the amplitude in the leads:
$$f(z=\pm b/2)^2=\frac{H_{\text{c}\text{film}}^2(T)\xi _{\text{film}}^2(T)}{H_{\text{c}\text{wire}}^2(T)\xi _{\text{wire}}^2(T)}.$$
(47)
For wires made out of superconducting material the same as (or weaker than) the leads, this ratio is always larger than unity <sup>4</sup><sup>4</sup>4In finding $`u_0`$ there is a minor numerical difficulty. As the amplitude of the order parameter is expressed via the JacobiSn function, and $`\text{JacobiSn}[z\sqrt{u_2/2},u_1/u_2]`$ is a doubly periodic function in the first variable, it is not obvious whether $`\pm (b/2)\sqrt{u_2/2}`$ lies in the first period, as can be seen from Fig. 11. As the trajectory must be simply periodic, $`z\sqrt{u_2/2}`$ must intersect either a zero or a pole in the first unit quarter cell of the JacobiSn function. Now, we are only interested in trajectories that escape to $`f\mathrm{}`$ \[as $`f(\pm b/2)`$ is assumed to be greater than or equal to unity\], so a pole must be intersected. (However, being outside the first period is unphysical, as it means that somewhere along the wire $`f=\mathrm{}`$.) There are exactly two poles in the first unit quarter cell. They are located at $`2v_1+v_2`$ and $`v_2`$, where $`v_1\mathrm{K}(u_1/u_2)`$ and $`v_2i\mathrm{K}(1u_1/u_2)`$, in which $`\mathrm{K}()`$ is the complete elliptic integral. So, instead of checking whether $`\pm (b/2)\sqrt{u_2/2}`$ is outside the unit quarter cell, we can just determine which pole $`z\sqrt{u_2/2}`$ intersects and then see if $`\pm (b/2)\sqrt{u_2/2}`$ lies beyond that pole or not. .
Once we have computed the functions $`\theta (J)`$ and $`F_{\text{m}/\text{sp}}(J)`$ for both saddle-point and metastable states on a single wire, we can use the phase and total current constraints to build the saddle-point and metastable states for the two-wire device. We proceed as before, by constructing a Markov chain for the state of the device, except that now we include in the graph all metastable states of the device. By diagonalizing the Markov chain, we find the populations of the various metastable states and, hence, the rate of gain of $`\mathrm{\Theta }`$.
We plot the typical magnetoresistance curves for various temperatures, obtained numerically, as well as $`dV/dI`$ vs. $`T`$ for various magnetic fields and total currents. Notice that the resistance at $`\delta =0`$ and large total currents can exceed that at $`\delta =\pi /2`$ with low total-current.
## V Connections with experiment
In this section a connection is made between our calculations and our experiments us\_in\_science . First, the predicted period of the magnetoresistance oscillations is compared to the experimentally obtained one. Then, the experimentally-obtained resistance vs. temperature curves are fitted using our extension of the IZAH Josephson junction model (for shorter wires) and our extension of the LAMH wire model (for longer wires).
### V.1 Device fabrication
Four different devices were successfully fabricated and measured. The devices were fabricated by suspending DNA molecules across a trench and then sputter coating them with the superconducting alloy of MoGe. The leads were formed in the same sputter-coating step, ensuring seamless contact between leads and the wires. Next, the leads were truncated lithographically to the desired width. In the case of device 930-1, after being measured once, its leads were further narrowed using focused ion beam milling, and the device was remeasured. For further details of the experimental procedure see Ref. us\_in\_science .
### V.2 Comparison between theory and experiment
#### V.2.1 Oscillation period
The magnetoresistance periods obtained for four different samples are summarized in Table 1. The corresponding theoretical periods were calculated using Eq. (6), based on the geometry of the samples which was obtained via scanning electron microscopy. To test the theoretical model, the leads of one sample, sample 930-1, were narrowed using a focused ion beam mill, and the magnetoresistance of the sample was remeasured. The theoretically predicted periods all coincide quite well with the measured values, except for sample 219-4, which was found to have a “+” shaped notch in one of the leads (which was not accounted for in calculating the period). The notch effectively makes that lead significantly narrower, thus increasing the magnetoresistance period, and this qualitatively accounts for the discrepancy.
For all samples, when the leads are driven into the vortex state, the magnetoresistance period becomes much longer, approaching the Aharonov-Bohm value for high fields. This is consistent with the theoretical prediction that the period is then given by Eq. (6), but with $`l`$ replaced by the field-dependent inter-vortex spacing $`r`$.
#### V.2.2 Oscillation amplitude
We have made qualitative and quantitative estimates of the resistance of two-bridge devices in several limiting cases. For devices containing extremely short wires \[$`b\xi (T)`$\], such as sample 219-4, the superconducting wires cannot support multiple metastable states, and thus they operate essentially in the Josephson junction limit, but with the junction critical current being a function of temperature given by LAMH theory as $`I_\text{c}(T)=I_\text{c}(0)(1T/T_\text{c})^{3/2}`$. A summary of fits to the data for this sample, using the Josephson junction limit, is shown in Fig. 7. On the other hand, for longer wires it is essential to take into account the multiple metastable states, as is the case for sample 930-1, which has wires of intermediate length. A summary of numerical fits for this sample is shown in Fig. 13. In all cases, only the two low total-current magnetoresistance curves were fitted. By using the extracted fit parameters, the high total-current magnetoresistance curves were calculated, with their fit to the data serving as a self-consistency check. As can be seen from the fits, our model is consistent with the data over a wide range of temperatures and resistances. We remark, however, that the coherence length required to fit the data is somewhat larger than expected for MoGe.
## VI Concluding remarks
The behavior of mesoscale NQUIDs composed of two superconducing leads connected by a pair of superconducting nanowires has been investigated. Magnetoresistance measurements us\_in\_science have revealed strong oscillations in the resistance as a function of magnetic field, and these were found to have anomalously short periods. The period has been shown to originate in the gradients in the phase of the superconducting order parameter associated with screening currents generated by the applied magnetic field. The periods for five distinct devices were calculated, based on their geometry, and were found to fit very well with the experimental results
The amplitude of the magnetoresistance has been estimated via extensions, to the setting of parallel superconducting wires, of the IZAH theory of intrinsic resistive fluctuations in a current-biased Josephson junction for the case of short wires and the LAMH theory of intrinsic resistive fluctuations in superconducting wires for pairs of long wires. In both cases, to make the extensions, it was necessary to take into account the inter-wire coupling mediated through the leads. For sufficiently long wires, it was found that multiple metastable states, corresponding to different winding numbers of the phase of the order parameter around the AB contour, can exist and need to be considered. Accurate fits have been made to the resistance vs. temperature data at various magnetic fields and for several devices by suitably tuning the critical temperatures, zero-temperature coherence lengths, and normal-state resistances of the nanowires.
As these device are sensitive to the spatial variations in the phase of the order parameter in the leads, they may have applications as superconducting phase gradiometers. Such applications may include the sensing of the presence in the leads of vortices or of supercurrents flowing perpendicular to lead edges.
Acknowledgments: This work was supported by the U.S. Department of Energy, Division of Materials Sciences under Award No. DEFG02-91ER45439, through the Frederick Seitz Materials Research Laboratory at the University of Illinois at Urbana-Champaign. AB and DH would like to also acknowledge support from the Center for Microanalysis of Material DOE Grant No. DEFG02-96ER45439, NSF CAREER Grant No. DMR 01-34770, and the A.P. Sloan Foundation.
## Appendix A Multi-wire devices
In this appendix we give an example of how to extend the theory Presented in this Paper to the case of devices comprising more than two wires. In our example, we consider an array of $`n`$ identical short wires (i.e. wires in the Josephson junction limit) spaced at regular intervals. We continue to work at a fixed total current and to ignore charging effects. The end-to-end phase accumulations along the wires are related to each other as
$`\begin{array}{cc}\hfill \theta _2& =\theta _1+2\delta ,\hfill \\ \hfill \theta _3& =\theta _1+4\delta ,\hfill \\ & \mathrm{}\hfill \\ \hfill \theta _n& =\theta _1+2(n1)\delta ,\hfill \end{array}`$ (48)
i.e., $`\theta _n\theta _1=2(n1)\delta `$ (for $`n=2,\mathrm{},N`$), where $`\delta `$ is the phase accumulation in one of the leads between each pair of adjacent wires. The Gibbs free energy of the multi-wire subsystem is given by
$$G(I,\theta _1)=\frac{h}{2e}\left(I_\text{c}\underset{m=1}{\overset{n}{}}\mathrm{cos}\left(\theta _1+2(m1)\delta \right)+I\theta _1\right),$$
(49)
where $`I`$ is the total current and we are assuming that the wires have identical critical currents. As for the two-junction case, this junction array, is equivalent to a single effective junction. Figure 14 shows the critical current of this effective junction as a function of $`\delta `$ for devices comprising 2, 5, and 15 wires. The magnetoresistance of such a device then follows from IZAH theory, i.e., Eq. (31).
## Appendix B Physical Scales
It is convenient to express the results of the long-wire model, Eqs. (41,45,46), in terms of macroscopic physical parameters. Following Tinkham and Lau TL , we express the condensation energy scale per coherence length of wire as
$$=0.22k_\text{B}T_\text{c}(1t)^{3/2}\frac{R_\text{q}}{R_\text{N}}\frac{b}{\xi (T=0)},$$
(50)
where $`tT/T_\text{c}`$, $`R_\text{N}`$ is the normal-state resistance of the device, and $`R_\text{q}h/4e^26.5\text{k}\mathrm{\Omega }`$ is the quantum of resistance. The LAMH prefactor for sequential phase slips then becomes
$$\mathrm{\Omega }=\frac{b\sqrt{1t}}{\xi (T=0)}\left(\frac{8\sqrt{2}}{3k_\text{B}T_\text{c}}\right)^{1/2}\frac{8k_\text{B}(T_\text{c}T)}{\pi \mathrm{}},$$
(51)
and for parallel phase slips becomes
$$\mathrm{\Omega }=\left(\frac{b\sqrt{1t}}{\xi (T=0)}\right)^2\left(\frac{16\sqrt{2}}{3k_\text{B}T_\text{c}}\right)^{1/2}\frac{8k_\text{B}(T_\text{c}T)}{\pi \mathrm{}}.$$
(52)
The remaining parameters in the model are $`R_N`$, $`T_\text{c}`$ and $`\xi (T=0)`$. The normal-state resistance and the critical temperature may be obtained from the $`R`$ vs. $`T`$ curve. The coherence length may be obtained by comparing $``$ obtained from the critical current at low temperature, via
$$I_\text{c}=\frac{2}{3\sqrt{3}}\frac{16\pi }{\mathrm{\Phi }_0},$$
(53)
with $``$ obtained via Eq. (50).
In experiment, it is expected that the two wires are not identical. The long-wire model can be easily extended to this case. The number of parameters to be fitted would then expand to include the normal-state resistance for each wire (only one of which is free, as the pair are constrained by the normal-state resistance of the entire device, which can be extracted from the $`R`$ vs. $`T`$ curve), a zero-temperature coherence length for each wire, and a critical temperatures for each wire.
## Appendix C LAMH theory for a single bridge
In this appendix we reproduce useful formulas from LA LA , and rewrite them in a way that is convenient for further calculations, especially for numerical implementation. As in the case of single-wire LAMH theory, one starts with the Ginzburg-Landau free energy
$$F=_{b/2}^{b/2}\alpha |\psi |^2+\frac{\beta }{2}|\psi |^4+\frac{\mathrm{}^2}{2m}|\psi |^2dz.$$
(54)
The relationships between the parameters of the Ginzburg-Landau free energy ($`\alpha `$ and $`\beta `$), coherence length $`\xi `$, the condensation energy per unit coherence length $``$, the critical field $`H_\text{c}`$ and the cross-sectional area of the wire $`\sigma `$ are given by $`\frac{\alpha ^2}{\beta }=\frac{H_\text{c}^2\sigma }{8\pi }=/\xi `$ and $`\xi ^2=\frac{\mathrm{}^2}{2m|\alpha |}`$. Following McCumber M68 , it is convenient to work in terms of the dimensionless units obtained using the transformations: $`|\psi |^2\frac{\alpha }{\beta }|\psi |^2`$, $`z\sqrt{\frac{2m|\alpha |}{\mathrm{}^2}}z`$, and $`bb/\xi =\sqrt{\frac{2m|\alpha |}{\mathrm{}^2}}b`$. In terms of these units, the free energy becomes
$`F=2{\displaystyle _{b/2}^{b/2}}\left({\displaystyle \frac{1}{2}}(1|\psi |^2)^2+|\psi |^2\right)𝑑z.`$ (55)
The Ginzburg-Landau equation is obtained by varying the free energy:
$$\delta F=0\psi +|\psi |^2\psi ^2\psi =0.$$
(56)
By writing $`\psi =fe^{i\varphi }`$ and taking the real and imaginary parts of the Ginzburg-Landau equation one obtains
$`f+f^3+(\varphi ^{})^2f`$ $`=f^{\prime \prime },`$ (57)
$`2\varphi ^{}f^{}+\varphi ^{\prime \prime }f`$ $`=0.`$ (58)
From Eq. (58), one finds the current conservation law:
$$f^2\varphi ^{}=J,$$
(59)
where $`J`$ is identified with the dimensionless current $`\frac{1}{2i}(\psi ^{}\psi \psi \psi ^{})`$. The physical current (in stat-amps) is given by $`I=JcH_\text{c}^2\sigma \xi /\mathrm{\Phi }_0`$. Expressing $`\varphi ^{}`$ in terms of J, Eq. (57) becomes
$`f^{\prime \prime }`$ $`=f+f^3+{\displaystyle \frac{J^2}{f^3}}={\displaystyle \frac{d}{df}}U(f),`$ (60)
where the effective potential $`U(f)`$ is given by
$`U(f)={\displaystyle \frac{J^2}{2f^2}}+{\displaystyle \frac{f^2}{2}}{\displaystyle \frac{f^4}{4}}.`$ (61)
Following LA, Eq. (61) can usefully be regarded as the equation of motion for a particle with position $`f(z)`$, where $`z`$ plays the role of time, moving in the potential $`U(f)`$ LA . Before proceeding to find the solution of this equation, we pause to consider the type of trajectories that are possible. Later, it will be demonstrated that at the edge of the wire $`f(\pm b)1`$, so the particle starts to the right of the hump; see Fig. 15. If the total energy of the particle is less than the height of the hump, the particle will be reflected by the hump. If, however, the particle starts with more energy than the height of the hump, it will pass over the hump and be reflected by the $`J^2/2f^2`$ dominated part of $`U(f)`$.
The equation of motion can be solved via the first integral (i.e. multiplying both sides by $`f`$ and integrating with respect to $`f`$):
$$E=\frac{(f^{})^2}{2}+U(f)f^{}=\sqrt{2(EU(f))},$$
(62)
where $`E`$ is a constant of integration (i.e. the energy of the particle in the mechanical analogy), which gives
$`z`$ $`={\displaystyle _{f_0}^f}{\displaystyle \frac{df}{\sqrt{2(EU(f))}}}={\displaystyle _{f_0}^f}{\displaystyle \frac{fdf}{\sqrt{2f^2EJ^2f^4+f^6/2}}}.`$ (63)
It is convenient to apply “initial” conditions at the middle of the wire, where $`f(z=0)=f_0`$, and to integrate towards the edges. We require that the particle come back to its starting point after a “time” $`b`$, i.e. at the edges of the wire the amplitude of the order parameter must match the boundary condition. Therefore, the middle of the wire must be the turning-point for the particle, i.e., at $`z=0`$ we have $`E=U(f_0)`$.
What follows next is a series of manipulations via which one can express solution for $`f(z)`$ in terms of special functions.
Step 1: substitution: $`f^2u`$
$`z`$ $`={\displaystyle \frac{1}{2}}{\displaystyle _{u_0}^u}{\displaystyle \frac{du}{\sqrt{2EuJ^2u^2+u^3/2}}}`$ (64)
Step 2: substitution: $`uu_0+ϵ`$
$`2z`$ $`={\displaystyle _0^{uu_0}}{\displaystyle \frac{dϵ}{\left[ϵ\left(\underset{\alpha }{\underset{}{\left(\frac{J^2}{u_0}u_0+u_0^2\right)}}+\underset{\beta }{\underset{}{\left(\frac{3}{2}u_01\right)}}ϵ+ϵ^2/2\right)\right]^{1/2}}}`$ (65)
$`2z`$ $`={\displaystyle _0^{uu_0}}{\displaystyle \frac{dϵ}{(ϵ(ϵ+\beta ϵ+ϵ^2/2))^{1/2}}}`$ (66)
$`={\displaystyle _0^{uu_0}}{\displaystyle \frac{\sqrt{2}dϵ}{(ϵ(ϵ+\underset{u_1}{\underset{}{\beta +\sqrt{\beta ^22\alpha }}})(ϵ+\underset{u_2}{\underset{}{\beta \sqrt{\beta ^22\alpha }}}))^{1/2}}}`$ (67)
$`={\displaystyle _0^{uu_0}}{\displaystyle \frac{\sqrt{2}dϵ}{(ϵ(ϵu_1)(ϵu_2))^{1/2}}}`$ (68)
Step 3: substitution: $`ϵu_1z^2`$
$`2z`$ $`={\displaystyle \frac{2\sqrt{2}}{\sqrt{u_2}}}{\displaystyle _0^{\sqrt{\frac{uu_0}{u_1}}}}{\displaystyle \frac{d\omega }{((\omega ^21)(\frac{u_1}{u_2}\omega ^21))^{1/2}}}`$ (69)
$`={\displaystyle \frac{2\sqrt{2}}{\sqrt{u_2}}}\text{EllipticF}[\text{ArcSin}\left[\sqrt{{\displaystyle \frac{uu_0}{u_1}}}\right],{\displaystyle \frac{u_1}{u_2}}]`$ (70)
The following definitions have been used:
$`\alpha [u_0]`$ $`J^2/u_0u_0+u_0^2,`$ $`\beta [u_0]`$ $`{\displaystyle \frac{3}{2}}u_01,`$ (71)
$`u_1[\alpha ,\beta ]`$ $`\beta \sqrt{\beta ^22\alpha },`$ $`u_2[\alpha ,\beta ]`$ $`\beta +\sqrt{\beta ^22\alpha }.`$ (72)
By inverting relation (70) one obtains an explicit equation for the amplitude of the order parameter as a function of position along the wire (see Fig. 16):
$`f^2(z)`$ $`=u_0+u_1\mathrm{sin}^2\left[\text{JacobiAmplitude}[z\sqrt{{\displaystyle \frac{u_2}{2}}},{\displaystyle \frac{u_1}{u_2}}]\right]`$ (73a)
$`=u_0+u_1\text{JacobiSn}^2[z\sqrt{{\displaystyle \frac{u_2}{2}}},{\displaystyle \frac{u_1}{u_2}}]`$ (73b)
The end-to-end phase difference along the wire may be found by using the current conservation law. Thus one obtains
$$\theta =_{b/2}^{b/2}\frac{J}{f^2(z)}𝑑z=2J_0^{b/2}\frac{dz}{u_0+u_1\text{JacobiSn}^2[z\sqrt{\frac{u_2}{2}},\frac{u_1}{u_2}]}.$$
(74)
The Helmholtz free energy can be found by substituting the expressions for $`f(z)`$ and $`\varphi ^{}(z)`$ into the expression for the free energy. One then obtains
$$F=4_0^{b/2}𝑑z(\frac{1}{2}2f^2+f^4+J^2/u_0+u_0u_0^2/2),$$
(75)
where $`E`$ was expressed in terms of $`u_0`$. Eqs. (7475) provide expressions for $`\theta `$ and $`F`$ which are true regardless of the length of the wire, and therefore may be used as a starting point for computing the Gibbs free-energy of the various metastable states subject to the total-current and the phase constraints. |
warning/0507/math0507606.html | ar5iv | text | # The Noether inequality for smooth minimal 3-folds
## 1. Introduction
In the 1980’s M. Reid, observing the importance of the Noether inequality: $`K^22p_g4`$ for surfaces of general type, asked the following question
###### Question 1.1.
(M. Reid) What is the 3-dimensional version of Noether’s inequality?
Question 1.1 is obviously a very important aspect of threefold geography, just like the well known Miyaoka-Yau inequality. There have been already several works dedicated to the above question:
> $``$ M. Kobayashi (1992, ) studied Gorenstein minimal 3-folds of general type and found an infinite number of examples (Proposition 3.2 in ) satisfying the equality:
>
> (1.1) $`K^3=\frac{4}{3}p_g\frac{10}{3}.`$
>
> $``$ M. Chen (2004, ) gave effective Noether type inequalities for arbitrary minimal 3-folds of general type.
>
> $``$ M. Chen (2004, ) answered Question 1.1 under the assumption that the 3-fold $`X`$ is smooth with an ample canonical line bundle, proving the sharp inequality: $`K^3\frac{4}{3}p_g\frac{10}{3}`$.
>
> \[In the above three items, $`K^3:=K_X^3`$ is the canonical volume and $`p_g:=p_g(X)`$ is the geometric genus of $`X`$.\]
In this paper, we will generalize the main theorem of . The aim is to answer Question 1.1 under a weaker condition:
###### Theorem 1.2.
Let $`X`$ be a smooth projective minimal 3-fold of general type. Then the sharp Noether inequality:
$$K_X^3\frac{2}{3}(2p_g(X)5)$$
holds.
###### Remark 1.3.
The inequality in Theorem 1.2 is sharp because of M. Kobayashi’s interesting examples (cf. Equation (1.1)).
As an application of our results, we present the following corollary which gives a classification of 3-folds of general type with small ”slope” $`K^3/p_g`$:
###### Corollary 1.4.
Let $`X`$ be a projective minimal (i.e., $`K_X`$ is nef) Gorenstein 3-fold of general type with canonical singularities. Assume $`K_X^3<\frac{7}{5}p_g(X)2`$. Then $`X`$ is canonically fibred by curves of genus 2.
The assumption in Corollary 1.4 is not empty again because of M. Kobayashi’s examples.
###### 1.5.
The set up. Let $`X`$ be a projective minimal Gorenstein 3-fold of general type with canonical singularities. According to the work of M. Reid and Y. Kawamata (Lemma 5.1 of) , there is a minimal model $`Y`$ with a birational morphism $`\nu :YX`$ such that $`K_Y=\nu ^{}(K_X)`$ and such that $`Y`$ is factorial with at worst terminal singularities. Thus we may always assume that $`X`$ is factorial with only (necessarily finitely many) terminal singularities. Observing that $`K_X^32`$ (see 2.1 below), the inequality in Theorem 1.2 is automatically true whenever $`p_g(X)4`$. So the essential argumentation takes place when $`p_g(X)`$ is bigger and we are led to study the canonical map $`\mathrm{\Phi }:=\mathrm{\Phi }_{|K_X|}`$ as in the two dimensional case.
Take a birational modification $`\pi :X^{}X`$, which exists by Hironaka’s big theorem, such that:
(1) $`X^{}`$ is smooth;
(2) the movable part of $`|K_X^{}|`$ is base point free;
(3) $`\pi ^{}(K_X)`$ is supported by a normal crossings divisor (so that we are in a position to apply the Kawamata-Viehweg vanishing theorem ).
We will fix some notation below. Denote by $`g`$ the composition $`\mathrm{\Phi }\pi `$. So $`g:X^{}W^{}^N`$ is a morphism. Let $`g:X^{}\stackrel{𝑓}{}B\stackrel{𝑠}{}W^{}`$ be the Stein factorization of $`g`$ (thus $`B`$ is normal and $`f`$ has connected fibers). We can write:
$$K_X^{}=\pi ^{}(K_X)+E_\pi =M+Z^{},$$
where $`M`$ is the movable part of $`|K_X^{}|`$, $`Z^{}`$ the fixed part and $`E_\pi `$ an effective divisor which is a linear combination of distinct exceptional divisors. We may also write:
$$\pi ^{}(K_X)=M+E^{},$$
where $`E^{}=Z^{}E_\pi `$ is an effective divisor. On $`X`$, one may write $`K_XN+Z`$ where $`N`$ is the movable part and $`Z`$ the fixed part. So
$$\pi ^{}(N)=M+\underset{i=1}{\overset{s}{}}d_iE_i$$
with $`d_i>0`$ for all $`i`$. The above sum runs over all those exceptional divisors of $`\pi `$ that lie over the base locus of $`M`$. On the other hand, one may write $`E_\pi =_{j=1}^te_jE_j`$ where the sum runs over all exceptional divisors of $`\pi `$. One has $`e_j>0`$ for all $`1jt`$ because $`X`$ is terminal. Apparently, one has $`ts`$.
Set $`d:=dim(B)`$. We say that $`X`$ is canonically fibred by surfaces if $`d=1`$. Under this situation, we have an induced fibration $`f:X^{}B`$ onto a smooth curve $`B`$. Denote by $`b:=g(B)`$ the geometric genus of $`B`$.
Notations
| $`K^3`$ | the canonical volume of a 3-fold in question |
| --- | --- |
| $`p_g=h^0(𝒪(K))`$ | the geometric genus |
| $`q(V)=h^1(𝒪_V)`$ | the irregularity of $`V`$ |
| $`h^2(𝒪_V)`$ | the second irregularity of a 3-fold $`V`$ |
| $`\chi (𝒪_V)`$ | the Euler Poincare characteristic of $`V`$ |
| $`(K^2,p_g)`$ | invariants of a minimal surface of general type |
| $`g(B)`$ | the genus of a curve $`B`$ |
| $``$ | numerical equivalence |
| $``$ | linear equivalence |
| $`\mathrm{}\mathrm{}`$ | the round up of $``$ ($`\mathrm{}x\mathrm{}:=min\{n|nx\}`$) |
| $`D|_S`$ | the restriction of the divisor $`D`$ to $`S`$ |
| $`DC`$ | the intersection number of a divisor $`D`$ with a curve $`C`$ |
## 2. Reduction to the surface case and the lower bound of $`K^3`$
###### 2.1.
$`K^3`$ is even. Suppose that $`D`$ is any divisor on a smooth 3-fold $`V`$. The Riemann-Roch theorem (cf. appendix in Hartshorne’s book ) gives:
$$\chi (𝒪_V(D))=\frac{D^3}{6}\frac{K_VD^2}{4}+\frac{D(K_V^2+c_2(V))}{12}+\chi (𝒪_V).$$
A direct calculation shows that
$$\chi (𝒪_V(D))+\chi (𝒪_V(D))=\frac{K_VD^2}{2}+2\chi (𝒪_V).$$
Therefore, $`K_VD^2`$ is an even number.
Now let $`X`$ be a projective minimal Gorenstein 3-fold of general type. Denote by $`\nu :VX`$ a smooth birational modification. Let $`D`$ be any divisor on X. Then $`K_XD^2=K_V(\nu ^{}D)^2`$ is even. Especially, $`K_X^3`$ is even and positive.
###### 2.2.
Known results. Let $`X`$ be a projective minimal factorial 3-fold of general type with terminal singularities. The following Noether type inequalities have already been established, where $`d=dim\mathrm{\Phi }_{|K_X|}(X)`$.
> $``$ if $`d=3`$, then $`K_X^32p_g(X)6`$ (cf. M. Kobayashi’s Main Theorem in );
>
> $``$ if $`d=2`$, then $`K_X^3\mathrm{}\frac{2}{3}(g1)\mathrm{}(p_g(X)2)`$ (cf. Chen’s Theorem 4.1(ii) in ), where $`g`$ is the genus of a general fiber of the induced fibration $`f:X^{}B`$; if furthermore $`X`$ is smooth, then $`K_X^3\frac{2}{3}(2p_g(X)5)`$ (cf. Chen’s Theorem 4.3 in );
>
> $``$ if $`d=1`$ and the general fiber $`S`$ of the induced fibration $`f:X^{}B`$ is not a surface of type $`(K^2,p_g)=(1,2)`$, then $`K_X^32p_g(X)4`$ (cf. Chen’s Theorem 4.1(iii) in ).
In order to prove Theorem 1.2, we have to treat the remaining case (in the above third item) where $`X`$ is canonically fibred by surfaces of type $`(1,2)`$. Note that Theorem 1.2 was proved in only under the stronger assumption of $`K_X`$ being ample. Assuming only the nefness of $`K_X`$, we can see that the method in is no longer effective and the situation could be more complicated. It is the aim of this paper to overcome this obstacle and prove our Theorem 1.2.
The rest of this section is devoted to deducing several key inequalities through the $``$-divisor method.
###### 2.3.
Key inequalities. Keep the same notation as in 1.5 and assume that $`K_X`$ is nef and big. Suppose, from now on, $`d=1`$ and $`p_g(X)3`$. We have an induced fibration $`f:X^{}B`$. Denote by $`S`$ a general fiber of $`f`$. Let $`\sigma :SS_0`$ be the contraction onto the minimal model. Suppose $`(K_{S_0}^2,p_g(S_0))=(1,2)`$.
By Lemma 4.5 of , we have two cases exactly:
$$q(X)=b=1\text{and}h^2(𝒪_X)=0,$$
$$q(X)=b=0\text{and}h^2(𝒪_X)1.$$
One may write $`M=_{i=1}^aS_i`$ as a disjoint union of distinct smooth fibers of $`f`$, where $`a=p_g(X)1`$ if $`b=0`$, or $`a=p_g(X)`$ otherwise. Noting that $`\pi ^{}(K_X)_{|S}K_S`$ is a nef and big Cartier divisor and that $`\sigma ^{}(K_{S_0})`$ is the positive part of the Zariski decomposition of $`K_S`$, so $`\pi ^{}(K_X)_{|S}^2=\sigma ^{}(K_{S_0})^2=1`$, and $`\pi ^{}(K_X)_{|S}\sigma ^{}(K_{S_0})`$ by the uniqueness of the Zariski decomposition. According to the construction of $`\pi `$, we know that $`E_{|S}^{}\pi ^{}(K_X)_{|S}`$ is a normal crossing divisor for a general fiber $`S`$.
Now let us assume $`\alpha _3(0,1)`$ be a real number such that
$$h^0(S,K_S+\mathrm{}\alpha E_{|S}^{}\mathrm{})3$$
for all $`\alpha >\alpha _3`$. We may now write $`a`$ as $`a=m_2+m_3+1`$, where $`m_2,m_3`$ are non-negative integers and
$$\frac{am_3}{a}>\alpha _3.$$
Such integers exist: for instance, one may take $`m_3=0`$ and $`m_2=a1`$. What we will show in next sections is that we can find a nontrivial decomposition of $`a`$, i.e., with $`m_3>0`$.
Once we have the above setting, we may deduce an interesting inequality as follows. Write
$$MS_0+\underset{i=1}{\overset{m_2}{}}S_{2,i}+\underset{j=1}{\overset{m_3}{}}S_{3,j}.$$
Since
$$\pi ^{}(K_X)\underset{j=1}{\overset{m_3}{}}S_{3,j}\frac{m_3}{a}E^{}(1\frac{m_3}{a})\pi ^{}(K_X)$$
is nef and big and has normal crossings, the Kawamata-Viehweg vanishing theorem () yields
$$H^1(X^{},K_X^{}+\mathrm{}\pi ^{}(K_X)\underset{j=1}{\overset{m_3}{}}S_{3,j}\frac{m_3}{a}E^{}\mathrm{})=0$$
and hence the exact sequence:
$`0`$ $`H^0(X^{},K_X^{}+\mathrm{}\pi ^{}(K_X){\displaystyle \underset{j=1}{\overset{m_3}{}}}S_{3,j}{\displaystyle \frac{m_3}{a}}E^{}\mathrm{})`$
$`H^0(X^{},K_X^{}+\mathrm{}\pi ^{}(K_X){\displaystyle \frac{m_3}{a}}E^{}\mathrm{})`$
$`_{j=1}^{m_3}H^0(S_{3,j},K_{S_{3,j}}+\mathrm{}(1{\displaystyle \frac{m_3}{a}})E^{}\mathrm{}_{|S_{3,j}})0.`$
In the above sequence, we obviously have
$$\mathrm{}(1\frac{m_3}{a})E^{}\mathrm{}_{|S_{3,j}}\mathrm{}(1\frac{m_3}{a})E_{|S_{3,j}}^{}\mathrm{}$$
and
$$(1\frac{m_3}{a})E_{|S_{3,j}}^{}\frac{am_3}{a}\pi ^{}(K_X)_{|S_{3,j}}.$$
So one has
$$h^0(S_{3,j},K_{S_{3,j}}+\mathrm{}(1\frac{m_3}{a})E_{|S_{3,j}}^{}\mathrm{})3$$
for sufficiently general $`S_{3,j}`$ as a fiber of $`f`$ by our definition of $`\alpha _3`$. The above sequence then gives the inequality
(2.1) $`h^0(X^{},K_X^{}+\mathrm{}\pi ^{}(K_X){\displaystyle \frac{m_3}{a}}E^{}\mathrm{})`$
$`h^0(K_X^{}+\mathrm{}\pi ^{}(K_X){\displaystyle \underset{j=1}{\overset{m_3}{}}}S_{3,j}{\displaystyle \frac{m_3}{a}}E^{}\mathrm{})+3m_3.`$
It is obvious that one has
$`h^0(K_X^{}+\mathrm{}\pi ^{}(K_X){\displaystyle \underset{j=1}{\overset{m_3}{}}}S_{3,j}{\displaystyle \frac{m_3}{a}}E^{}\mathrm{})`$
$`h^0(K_X^{}+\mathrm{}\pi ^{}(K_X){\displaystyle \underset{j=1}{\overset{m_3}{}}}S_{3,j}{\displaystyle \frac{m_2+m_3}{a}}E^{}\mathrm{}).`$
Similarly, because
$$\pi ^{}(K_X)\underset{i=1}{\overset{m_2}{}}S_{2,i}\underset{j=1}{\overset{m_3}{}}S_{3,j}\frac{m_2+m_3}{a}E^{}\frac{1}{a}\pi ^{}(K_X)$$
is nef and big and with normal crossings, the vanishing theorem gives
$$H^1(K_X^{}+\mathrm{}\pi ^{}(K_X)\underset{i=1}{\overset{m_2}{}}S_{2,i}\underset{j=1}{\overset{m_3}{}}S_{3,j}\frac{m_2+m_3}{a}E^{}\mathrm{})=0.$$
So we have the following exact sequence:
$`0`$ $`H^0(X^{},K_X^{}+\mathrm{}\pi ^{}(K_X){\displaystyle \underset{i=1}{\overset{m_2}{}}}S_{2,i}{\displaystyle \underset{j=1}{\overset{m_3}{}}}S_{3,j}{\displaystyle \frac{m_2+m_3}{a}}E^{}\mathrm{})`$
$`H^0(X^{},K_X^{}+\mathrm{}\pi ^{}(K_X){\displaystyle \underset{j=1}{\overset{m_3}{}}}S_{3,j}{\displaystyle \frac{m_2+m_3}{a}}E^{}\mathrm{})`$
$`_{i=1}^{m_2}H^0(S_{2,i},K_{S_{2,i}}+\mathrm{}{\displaystyle \frac{am_2m_3}{a}}E^{}\mathrm{}|_{S_{2,i}})0.`$
The above exact sequence gives
$$h^0(S_{2,i},K_{S_{2,i}}+\mathrm{}\frac{am_2m_3}{a}E^{}\mathrm{}|_{S_{2,i}})p_g(S_{2,i})=2$$
and
(2.2) $`h^0(X^{},K_X^{}+\mathrm{}\pi ^{}(K_X){\displaystyle \underset{j=1}{\overset{m_3}{}}}S_{3,j}{\displaystyle \frac{m_2+m_3}{a}}E^{}\mathrm{})`$
$`h^0(K_X^{}+\mathrm{}\pi ^{}(K_X){\displaystyle \underset{i=1}{\overset{m_2}{}}}S_{2,i}{\displaystyle \underset{j=1}{\overset{m_3}{}}}S_{3,j}{\displaystyle \frac{m_2+m_3}{a}}E^{}\mathrm{})+2m_2.`$
We shall go on studying the group
$$H^0(X^{},K_X^{}+\mathrm{}\pi ^{}(K_X)\underset{i=1}{\overset{m_2}{}}S_{2,i}\underset{j=1}{\overset{m_3}{}}S_{3,j}\frac{m_2+m_3}{a}E^{}\mathrm{}).$$
Apparently, it is slightly bigger than $`H^0(X^{},K_X^{}+S_0)`$.
We set $`\delta :=2h^2(𝒪_X).`$ By looking at the exact sequence:
$$0𝒪_X^{}(K_X^{})𝒪_X^{}(K_X^{}+S_0)𝒪_{S_0}(K_{S_0})0,$$
one has
(2.3) $`h^0(K_X^{}+S_0)p_g(X)+\delta .`$
Combining the above inequalities (2.1)$``$(2.3), we have
(2.4) $`P_2(X)=h^0(K_X^{}+\pi ^{}(K_X))3m_3+2m_2+p_g(X)+\delta .`$
Applying Reid’s plurigenus formula (see the last section of and Lemma 8.3 of ):
$$P_2(X)=\frac{1}{2}K_X^33\chi (𝒪_X)=\frac{1}{2}K_X^33(1b+h^2(𝒪_X)p_g(X)),$$
we get the Noether type inequality:
(2.5) $`K_X^36m_3+4m_24p_g(X)+4h^2(𝒪_X)6b+10.`$
###### 2.4.
A problem on surfaces. As we have seen, the general fiber $`S`$ has the invariants $`(K_{S_0}^2,p_g(S_0))=(1,2)`$. We have a divisor $`E_{|S}^{}\sigma ^{}(K_{S_0})`$ which has normal crossings. So there is a divisor $`D_0|K_{S_0}|`$ with $`E_{|S}^{}=\sigma ^{}(D_0)`$. We expect to find a real number $`\alpha _3(0,1)`$ such that $`h^0(S,K_S+\mathrm{}\alpha E_{|S}^{}\mathrm{})3`$ for all $`\alpha >\alpha _3`$. Furthermore we hope $`\alpha _3`$ to be as small as possible.
## 3. The rounding up problem for (1,2) surfaces
Assume that $`Y`$ is the canonical model of a surface of general type with $`p_g(Y)=2,K_Y^2=1`$, that $`\tau :S_0Y`$ is its minimal model, and finally that $`f:SS_0`$ is a sequence of point blow ups.
We set up the following notation and assumptions:
* $`\mathrm{\Gamma }Y`$ is a canonical divisor
* $`D`$ is the full transform $`\tau ^{}(\mathrm{\Gamma })`$
* we assume that $`f^{}(D)`$ is a normal crossing divisor
* for $`t(0,1)`$ we consider the round up divisor $`\mathrm{\Delta }_t:=\mathrm{}tf^{}(D)\mathrm{}`$
###### Remark 3.1.
Observing that since $`H^1(𝒪_Y)=H^1(𝒪_S)=H^1(K_S)=0`$ (cf. ), one has $`h^0(K_S+\mathrm{\Delta }_t)=p_g(S)+h^0(\omega _{\mathrm{\Delta }_t})`$ where $`\omega _{\mathrm{\Delta }_t}:=𝒪_{\mathrm{\Delta }_t}(K_S+\mathrm{\Delta }_t)`$.
###### Theorem 3.2.
Assume that $`p_g(Y)=2,K_Y^2=1`$, that $`f^{}(D)`$ is a normal crossings divisor, and that $`3/10<t`$. Then $`h^0(K_S+\mathrm{\Delta }_t)=2+h^0(\omega _{\mathrm{\Delta }_t})3.`$
###### Proof.
1) Since $`K_Y^2=1`$, and $`K_Y`$ is ample, $`\mathrm{\Gamma }`$ is irreducible. (Note that $`|K_Y|`$ has one smooth and simple base point and the general member of $`|K_Y|`$ is a smooth curve of genus 2 (cf. page 225 in ). It is well known and easy to show that $`Y`$ is a hypersurface of degree $`10`$ in the weighted projective space $`(1,1,2,5)`$, so $`Y`$ is a finite double cover of $`(1,1,2)`$ and the involution $`\sigma `$ on $`Y`$ induced by the hyperelliptic involution of those genus 2 curves has exactly one isolated fixed point – the base point of $`|K_Y|`$. We shall also denote by the same symbol $`\sigma `$ its lift to a biregular involution on $`S_0`$, observing that again there is exactly one isolated fixed point – the base point of $`|K_{S_0}|`$.
The quotient $`Q_2=Y/\sigma =(1,1,2)`$ is isomorphic to a quadric cone in $`^3`$ and $`\mathrm{\Gamma }`$ is isomorphic to a double cover of $`^1`$ branched in a point $`P_{\mathrm{}}`$ and in a disjoint sub-scheme of length $`5`$ (cf. , page 231, a construction due to Horikawa).
2) Observe that if $`DD^{}`$, and $`\mathrm{\Delta }_t^{}:=\mathrm{}tf^{}(D^{})\mathrm{}`$, then $`h^0(K_S+\mathrm{\Delta }_t)h^0(K_S+\mathrm{\Delta }_t^{})`$.
3) Set $`K:=K_{S_0}`$. Write $`D=\stackrel{~}{\mathrm{\Gamma }}+\stackrel{~}{Z}`$, where $`\stackrel{~}{\mathrm{\Gamma }}`$ is the strict transform of $`\mathrm{\Gamma }`$. Thus $`\stackrel{~}{\mathrm{\Gamma }}K=1,\stackrel{~}{Z}K=0`$. Since $`\mathrm{\Gamma }`$ is a Cartier divisor, it follows that the support of $`\stackrel{~}{Z}`$ is a union of the support of certain fundamental cycles $`Z_i`$ corresponding to the rational double points $`P_iX`$ such that $`P_i\mathrm{\Gamma }`$, and moreover $`\stackrel{~}{Z}=_i\stackrel{~}{Z}_i`$, where $`\stackrel{~}{Z}_iZ_i`$.
4) If we take an effective decomposition $`D=D^{}+W`$, where $`D^{}K=1`$, then $`(D^{})^2=D^{}(KW)=1D^{}W1`$, since a canonical curve is 2-connected (, VII (6.2)).
5) If $`Z^{}K=0`$, and $`Z^{}`$ is (effective and) reduced, then $`(Z^{})^2=2k`$, where $`k`$ is the number of connected components of $`Z^{}`$. In fact, it suffices to prove the formula for $`Z^{}`$ connected, but $`Z^{}`$ is contained in a fundamental cycle, and corresponds therefore to a rational subtree of the Dynkin diagram. Thus, if $`n`$ is the number of edges of the subtree, then $`(Z^{})^2=2(n+1)+2n=2.`$
We pass now to the strategy of proof:
\[S1\] if the arithmetic genus $`p(\stackrel{~}{\mathrm{\Gamma }})1`$ then we pick $`D^{}=\stackrel{~}{\mathrm{\Gamma }}`$ (see point 2)).
Observe now that $`p(\stackrel{~}{\mathrm{\Gamma }})1`$ is equivalent, since $`\stackrel{~}{\mathrm{\Gamma }}K=1`$, to $`\stackrel{~}{\mathrm{\Gamma }}^21`$, or to $`D=\stackrel{~}{\mathrm{\Gamma }}`$, in view of 4). If the first strategy is not allowed, this means that $`\stackrel{~}{\mathrm{\Gamma }}^2=3`$, and $`\stackrel{~}{\mathrm{\Gamma }}^1.`$
If $`\stackrel{~}{\mathrm{\Gamma }}^1`$ we consider the reduced divisor $`\stackrel{~}{\mathrm{\Gamma }}+Z_i^{}`$, where $`Z_i^{}=(Z_i)_{red}`$ is the reduced curve corresponding to one of the divisors $`\stackrel{~}{Z}_i`$ appearing in 3). By 5) and 4) it follows that the odd number $`(\stackrel{~}{\mathrm{\Gamma }}+Z_i^{})^2=5+2(\stackrel{~}{\mathrm{\Gamma }}Z_i^{})`$ equals $`3`$ or $`1`$, accordingly $`(\stackrel{~}{\mathrm{\Gamma }}Z_i^{})=1\mathrm{or}2.`$
\[S2\] If $`(\stackrel{~}{\mathrm{\Gamma }}Z_i^{})=2`$, there are four cases:
\[S2.1\] $`\stackrel{~}{\mathrm{\Gamma }}+Z_i^{}`$ is a normal crossing divisor (of arithmetic genus $`1`$), and we pick $`D^{}=\stackrel{~}{\mathrm{\Gamma }}+Z_i^{}`$.
\[S2.2\] $`\stackrel{~}{\mathrm{\Gamma }}`$ is tangent to a smooth (-2)-curve $`AZ_i^{}`$, and then we take $`D^{}=\stackrel{~}{\mathrm{\Gamma }}+A`$.
\[S2.3\] A fundamental cycle $`Z_1<\stackrel{~}{Z}`$ is of type $`A_4`$ and $`\stackrel{~}{\mathrm{\Gamma }}`$ passes through the central point transversally. Take $`D^{}=D`$ (see the claim below).
\[S2.4\] A fundamental cycle $`Z_1<\stackrel{~}{Z}`$ is of type $`A_2`$ and $`\stackrel{~}{\mathrm{\Gamma }}`$ passes through the central point transversally. Take $`D^{}=\stackrel{~}{\mathrm{\Gamma }}+Z_2`$ (see the claim below).
###### Claim 3.3.
(1) Cases \[S2.1\] – \[S2.4\] are the only possible cases if \[S2\] holds.
(2) In Case \[S2.3\], one has $`K_{S_0}D=\stackrel{~}{\mathrm{\Gamma }}+\stackrel{~}{Z}`$ with $`\stackrel{~}{Z}=A_1+2A+2A^{}+A_4`$, so that $`Z_1=A_1+A+A^{}+A_4`$ is a fundamental cycle of type $`A_4`$.
(3) In Case\[S2.4\], there is another fundamental cycle $`Z_2<\stackrel{~}{Z}`$ of type $`A_m`$ which together with $`\stackrel{~}{\mathrm{\Gamma }}`$ forms a rational loop (of arithmetic genus 1).
###### Proof.
(of the claim) If $`\stackrel{~}{\mathrm{\Gamma }}+Z_i^{}`$ is not a normal crossings divisor, then, the intersection number being $`2`$, either \[S2.2\] holds or $`\stackrel{~}{\mathrm{\Gamma }}`$ meets $`Z_i^{}`$ at a singular point $`P`$ where two components $`A,A^{}`$ meet, and all intersections are transversal. We observed that on $`S_0`$ we have a canonical biregular involution $`\sigma `$, induced from the hyperelliptic involution on the (genus two) canonical curves.
$`P`$ is then a fixed point for the involution $`\sigma `$, which has only the point lying over $`P_{\mathrm{}}`$ as isolated fixed point. Since $`P`$ lies in a fundamental cycle, $`P`$ is a different point than the above isolated fixed point. So there is a $`\sigma `$-fixed curve $`C`$ (on $`S_0`$) through $`P`$. If both $`A,A^{}`$ are $`\sigma `$-stable, then the action $`\sigma _{}`$ on the tangent space at $`P`$ will have three eigenvectors (along $`A,A^{},\stackrel{~}{\mathrm{\Gamma }}`$) and hence it equals $`(1)id`$, contradicting the fact that $`P`$ is not an isolated $`\sigma `$-fixed point.
Thus $`\sigma `$ must interchange $`A`$ and $`A^{}`$.
Let $`\stackrel{~}{Z}_1`$ contain $`A,A^{}`$. Then $`\sigma `$ acts on the graph of $`\stackrel{~}{Z}_1`$ fixing $`P=AA^{}`$. So $`\stackrel{~}{Z}_1`$ is of Dynkin type $`A_{2n}`$ ($`n1`$) and $`P`$ is the central point of $`\stackrel{~}{Z}_1`$. Therefore, $`A,A^{}`$ are the inverse images in the double cover $`S_0Q_2`$ of the last exceptional curve of the blow up of a singular point $`P^{}`$ of the branch curve $`B`$ on $`Q_2`$. Indeed, $`P^{}B`$ is a cusp of type $`(2,2n+1)`$ with fibre $`F`$ the only tangent at $`P^{}B`$. By point 1) follows that $`5(F.B)_P^{}=2n+1`$. Thus $`n=1,2`$. This proves the first assertion.
The second assertion follows from point 1).
Concerning assertion (3), by point 1) and observing that $`\stackrel{~}{\mathrm{\Gamma }}^1`$, our $`F`$ has one further intersection point $`P_2`$ with $`B`$, with $`(F.B)_{P_2}=2`$, and with $`P_2`$ a singular point for $`B`$ of type $`A_n`$. Then assertion (3) follows. ∎
6) If strategies \[S1\] and \[S2\] are both not allowed, this means that $`\stackrel{~}{\mathrm{\Gamma }}^1`$, and that $`(\stackrel{~}{\mathrm{\Gamma }}Z_i^{})=1`$ for each $`i`$.
7) Consider the intersection number $`\stackrel{~}{\mathrm{\Gamma }}\stackrel{~}{Z_i}`$ which equals $`K\stackrel{~}{Z_i}(\stackrel{~}{Z_i})^2=(\stackrel{~}{Z_i})^2`$ as is therefore a strictly positive even number. By 4), since $`(\stackrel{~}{\mathrm{\Gamma }}+\stackrel{~}{Z_i})^2=(\stackrel{~}{\mathrm{\Gamma }})^2(\stackrel{~}{Z_i})^21`$ this number equals $`2`$, or $`\stackrel{~}{Z_i}=\stackrel{~}{Z}`$, in which case we get $`4`$ (indeed, note that $`1=\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{\mathrm{\Gamma }}+\stackrel{~}{Z})`$).
8) Assume still that strategies \[S1\] and \[S2\] are both not allowed, thus $`1=K^2=\stackrel{~}{\mathrm{\Gamma }}^2+_i(2(\stackrel{~}{\mathrm{\Gamma }}\stackrel{~}{Z_i})+\stackrel{~}{Z_i}^2)`$, which, by 7), equals $`3+_i2`$ if there is more than one fundamental cycle. Therefore we conclude that $`\stackrel{~}{\mathrm{\Gamma }}`$ intersects precisely one or two fundamental cycles, and in the former case $`\stackrel{~}{\mathrm{\Gamma }}\stackrel{~}{Z_i}=4.`$
9) Let us consider first the case where there are two fundamental cycles intersecting $`\stackrel{~}{\mathrm{\Gamma }}`$, and observe the following inequalities: $`2=\stackrel{~}{\mathrm{\Gamma }}\stackrel{~}{Z_i}\stackrel{~}{\mathrm{\Gamma }}Z_i\stackrel{~}{\mathrm{\Gamma }}Z_i^{}=1`$, and write $`\stackrel{~}{Z_i}=Z_i+W_i`$. We have $`(\stackrel{~}{\mathrm{\Gamma }}+\stackrel{~}{Z_i})Z_i=0=\stackrel{~}{\mathrm{\Gamma }}Z_i+W_iZ_i+Z_i^2`$. By the well known properties of a fundamental cycle, we have $`Z_i^2=2`$, and $`W_iZ_i0`$, therefore $`\stackrel{~}{\mathrm{\Gamma }}Z_i2`$, and we conclude by the previous inequality that $`\stackrel{~}{\mathrm{\Gamma }}Z_i=2`$.
10) By 6) and 8) it follows that if we have two fundamental cycles which are intersected by $`\stackrel{~}{\mathrm{\Gamma }}`$, both are not reduced. By the standard classification of fundamental cycles, this means that the corresponding rational double points are not of type $`A_n`$, or, equivalently, that on the fibre $`F^1`$ of which $`\mathrm{\Gamma }`$ is the inverse image, we have two triple points. This however contradicts 1), and shows that one of the cases \[S1\] or \[S2\] occurs.
11) Let us consider then the former case in 8), where $`\stackrel{~}{\mathrm{\Gamma }}\stackrel{~}{Z_1}=4`$, and there is only one fundamental cycle which is intersected by $`\stackrel{~}{\mathrm{\Gamma }}`$, so we have $`\stackrel{~}{Z_1}=\stackrel{~}{Z}`$ and we may write accordingly $`Z`$ for the fundamental cycle, and $`Z^{}=Z_{red}.`$ Since $`\stackrel{~}{Z}^2=4`$, $`Z^2=2`$, we can write as in 9) $`\stackrel{~}{Z}=Z+W`$, and $`4=\stackrel{~}{Z}^2=Z^2+W^2+2WZ`$, and we get a sum of non positive terms, where the first two are even and strictly negative. Hence follows that $`2=W^2,WZ=0,\stackrel{~}{\mathrm{\Gamma }}Z=\stackrel{~}{\mathrm{\Gamma }}W=2`$ (note that $`0=ZK_{S_0}=Z(\stackrel{~}{\mathrm{\Gamma }}+Z+W)`$).
Thus again the fundamental cycle corresponds to a triple point of the branch curve, and $`\stackrel{~}{\mathrm{\Gamma }}`$ intersects $`Z^{}`$ in a smooth point, belonging to a (-2)-curve $`A`$ which appears with multiplicity $`2`$ in both $`Z`$ and $`W`$. Write $`W=r_iA_i`$ with $`r_i0`$, then $`A_iZ=0`$ for all $`i`$. So the equation $`A_i(\stackrel{~}{\mathrm{\Gamma }}+Z+W)=0`$ implies $`A_iW=A_i\stackrel{~}{\mathrm{\Gamma }}0`$. Also we have seen that $`W`$ is a sum of only those $`A_i`$’s which are orthogonal to $`Z`$. Moreover, since the point $`P=A\stackrel{~}{\mathrm{\Gamma }}`$ is invariant under the involution $`\sigma `$, we see that $`A`$ is pointwise $`\sigma `$-fixed. Indeed, both $`A`$ and $`\stackrel{~}{\mathrm{\Gamma }}`$ are $`\sigma `$-stable and their tangents are eigenvectors of the action $`\sigma _{}`$ on the tangent space at $`P`$, but $`\stackrel{~}{\mathrm{\Gamma }}`$ is not pointwise fixed and if also $`A`$ were not we would have an isolated fixed point, a contradiction.
Thus after we divide by the involution we obtain a (-4)-curve $`E`$, image of $`A`$, such that $`\stackrel{~}{\mathrm{\Gamma }}`$ is the inverse image of a transversal curve $`\stackrel{~}{F}`$ meeting $`E`$ precisely in the point $`p`$ image of $`P`$.
12) Let us analyse this last case in terms of the double covering $`\mathrm{\Gamma }F`$, where $`F^1`$. Since, on $`S_0`$, $`\stackrel{~}{\mathrm{\Gamma }}`$ is smooth of genus $`0`$, it follows that this covering is branched on the point $`P_{\mathrm{}}`$ and on another point $`PFB`$, where the branch locus $`B`$ of the double covering meets $`F`$ with intersection multiplicity $`(BF)_P=5`$ (observe also that $`B`$ does not contain $`F`$ as a component, else $`K_{S_0}2\stackrel{~}{\mathrm{\Gamma }}`$, absurd.)
Because $`Y`$ has only Rational Double Points as singularities,the branch curve $`B`$ of the double covering has only simple singularities (see ). Since the fundamental cycle is not reduced, then $`B`$ has a triple point at $`P`$. After blowing up $`P`$ we get a (-1)-curve $`E_1`$ and the full transform of $`F`$ is then $`F^{}+E_1`$, and the new branch locus is $`B^{}+E_1`$, where $`B^{}`$ is the proper transform of $`B`$. We know that the curve $`E`$ occurring in the normal crossing resolution of the branch locus has multiplicity $`2`$ in the full transform of $`F`$ (since $`A`$ has multiplicity $`4`$ in $`D=\stackrel{~}{\mathrm{\Gamma }}+\stackrel{~}{Z}`$), so $`E`$ cannot be the proper transform of $`E_1`$, and the new branch locus has a point of multiplicity $`3`$ at the point $`P^{}=F^{}E_1`$, and we must blow up $`P^{}`$, obtaining a (-1)-curve $`E_2`$ which is part of the new branch locus, together with the strict transform $`B^{\prime \prime }`$ of $`B^{}`$ and the strict transform $`E_1^{}`$ of $`E_1`$.
Since we had $`B^{}F^{}=2`$, and $`B^{}`$ has multiplicity $`2`$ at $`P^{}`$, it follows that now $`F^{\prime \prime }B^{\prime \prime }=\mathrm{}`$. Moreover also $`F^{\prime \prime }E_1^{}=\mathrm{}`$, therefore $`E`$ is the strict transform of $`E_2`$. Since $`E`$ has self intersection equal to $`4`$ we need three further blow ups of points (possibly infinitely near) on $`E_2`$.
13) Since the proper transform $`B^{}`$ is singular we can exclude that we have a rational double point of type $`E_6`$ or of type $`E_7`$. The other two cases are separated accordingly as follows (see , II (8.1) and III (7.1) for the one-to-one correspondence between the type of curve singularity of the branch locus $`BQ_2`$ of the double cover $`YQ_2`$ and the type of surface singularity at the corresponding point on the canonical model $`Y`$).
\[S3.1\] $`(B^{}E_1)_P^{}=2`$ implies that we have the $`D_n`$ case, since $`B`$ has then two distinct tangents at $`P`$.
\[S3.2\] $`(B^{}E_1)_P^{}=3`$ implies that we have the $`E_8`$ case, since $`B`$ has then only one tangent at $`P`$.
14) In both cases, we observe that $`\stackrel{~}{Z}`$ is the pull-back of $`E_1^{}+2E_2`$, i.e., the pull back of the maximal ideal of $`P`$ plus the pull back of the maximal ideal of $`P^{}`$. Moreover, in case \[S3.2\], there is only one (-2)-curve $`A`$ which occurs in $`Z`$ with multiplicity two, and such that $`AZ=0`$. In case \[S3.1\] we see instead that $`A`$ is the curve corresponding to the vertex at distance three from the asymmetrical end (observe that our assumptions imply $`n6`$).
15) We proceed by observing that it suffices to verify the statement for one blow up of $`S_0`$ where we have normal crossings for $`C:=f^{}(D)`$.
###### Lemma 3.4.
Let $`CS`$ be a normal crossing divisor, and let $`g:S^{}S`$ the blow up of a point $`P`$. Then, if we set $`C^{}=g^{}(C)`$, and $`\mathrm{\Theta }_t:=\mathrm{}tC\mathrm{}`$, $`\mathrm{\Theta }_t^{}:=\mathrm{}tC^{}\mathrm{}`$, then $`h^0(K_S+\mathrm{\Theta }_t)=h^0(K_S^{}+\mathrm{\Theta }_t^{})`$, , for all $`t(0,1)`$.
More generally, let $`C=_in_iC_i`$ be the decomposition of $`C`$ as a sum of irreducible analytic branches at $`P`$, and let $`m_i:=mult_P(C_i)`$, then the above equality holds if there are two smooth local branches, or just one branch ( i.e., $`n_1=1,n_j=0j2`$) of multiplicity $`m=2`$, provided $`t(0,1)`$.
###### Proof.
Let $`E`$ be the exceptional divisor of $`g`$, let $`C=_in_iC_i`$ be the decomposition of $`C`$ as a sum of irreducible divisors: then
$$C^{}=g^{}(C)=\underset{i}{}n_iC_i^{}+\underset{i}{}n_im_iE,$$
where $`C_i^{}`$ is the proper transform of $`C_i`$, and $`m_i:=mult_P(C_i)`$.
Taking the round up, we obtain
$`\mathrm{\Theta }_t^{}`$ $`=\mathrm{}tC^{}\mathrm{}={\displaystyle \underset{i}{}}\mathrm{}tn_i\mathrm{}C_i^{}+\mathrm{}t{\displaystyle \underset{i}{}}m_in_i\mathrm{}E`$
$`=g^{}(\mathrm{\Theta }_t)+(\mathrm{}t{\displaystyle \underset{i}{}}m_in_i\mathrm{}{\displaystyle \underset{i}{}}\mathrm{}tn_i\mathrm{}m_i)E.`$
Since $`K_S^{}=g^{}(K_S)+E`$, $`K_S^{}+\mathrm{\Theta }_t^{}=g^{}(K_S+\mathrm{\Theta }_t)+[1+\mathrm{}t_im_in_i\mathrm{}_i\mathrm{}tn_i\mathrm{}m_i]E`$ and it suffices that the integer in the square brackets is non negative in order to conclude the desired equality. Notice that the calculation is entirely local, so that we can replace the global decomposition by the local decomposition in analytic branches.
The normal crossings case is a special case of the one where all the multiplicities satisfy $`m_i=1`$ : in this case we want the inequality
(\**) $`1+\mathrm{}{\displaystyle \underset{i}{}}tn_i\mathrm{}{\displaystyle \underset{i}{}}\mathrm{}tn_i\mathrm{}`$
to hold. This is obvious if there are exactly two terms, since for any two real numbers $`a,b`$ holds $`1+\mathrm{}a+b\mathrm{}\mathrm{}a\mathrm{}+\mathrm{}b\mathrm{}`$.
For only one branch, we want $`1+\mathrm{}tm\mathrm{}\mathrm{}t\mathrm{}m`$, and this is true for $`m=2`$, since $`t<1`$.
Case \[S1\]: we take $`C=\stackrel{~}{\mathrm{\Gamma }}`$: it is irreducible of arithmetic genus equal to $`p\{1,2\}`$, therefore, for each $`t(0,1)`$ $`C`$ is equal to the round up of $`tC`$, and $`h^0(\omega _C)=p`$. If $`C`$ has normal crossings, we are done by the previous lemma (choose $`D^{}=C`$ in 2)).
Assume the contrary and assume first $`p=1`$: then $`C`$ has an ordinary cusp, thus the hypothesis of the lemma above applies. After a blow up we get two smooth tangent branches ( $`n_1=1,n_2=2`$), and the lemma still applies. We then get three smooth transversal branches where $`n_1=1,n_2=2,n_3=3`$: the inequality (\**) holds, provided $`1/6<t1/3`$ (since it is equivalent then to $`\mathrm{}6t\mathrm{}2`$) and we are done, since after this blow up we get global normal crossings for the full transform.
Assume now that $`C`$ does not have normal crossings, and that $`p=2`$: we have just verified that an ordinary cusp gives no problem ( as well as a node). We use now the fact that $`C`$ has only double points as singularities, so we have to verify that a tacnode $`y^2=x^4`$ and a higher cusp $`y^2=x^5`$ give no problem (higher singularities are excluded by point 1)).
For a tacnode we get two smooth branches, so the lemma applies, and after the first blow up we get three smooth transversal branches , with $`n_1=1,n_2=1,n_3=2`$, thus (\**) applies again if $`1/4<t1/3`$ (since (\**) is equivalent then to $`\mathrm{}4t\mathrm{}2`$) and after this blow up we get normal crossings.
In the case of the higher cusp, we get one branch of multiplicity $`2`$, so the lemma applies; after the first blow up we get a reduced ordinary cusp transversal to a smooth branch, occurring with multiplicity $`2`$. In this case we have to verify that $`1+\mathrm{}4t\mathrm{}\mathrm{}2t\mathrm{}+2\mathrm{}t\mathrm{}`$ , but this clearly holds for $`1/4<t1/3`$.
After a further blow up, we get two smooth branches, tangent, and with $`n_1=1,n_2=4`$, so the lemma applies. A further blow up, the last before we get normal crossings, yields a point where three smooth branches meet transversally, and $`n_1=1,n_2=4,n_3=5`$: we have to verify whether (\**) holds, i.e., $`1+\mathrm{}10t\mathrm{}\mathrm{}t\mathrm{}+\mathrm{}4t\mathrm{}+\mathrm{}5t\mathrm{}`$. But this holds clearly for $`3/10<t1/3`$ (else , for $`1/5<t3/10`$ there is a loss by $`1`$, which however would not trouble us since we started with $`p=2`$, and we only want the arithmetic genus above to be at least $`1`$). We now treat the remaining cases one by one, using 2).
Case \[S2.1\]: $`D^{}`$ already has normal crossings, and is reduced, thus there is nothing to prove.
Case \[S2.2\]: here $`D^{}`$ consists of two smooth tangent divisors $`^1`$, so its arithmetic genus is $`p=1`$. This is exactly the case of the tacnode, which we already treated, thus this case is also settled.
Case \[S2.3\]: $`D`$ ($`K_{S_0}`$) now has arithmetic genus 2, and does not have normal crossings exactly at the point where $`A,A^{},\stackrel{~}{\mathrm{\Gamma }}^1`$ meet transversally. The local multiplicities are $`2,2,1`$, thus for $`1/5<t1/3`$ we obtain $`1+\mathrm{}5t\mathrm{}\mathrm{}t\mathrm{}+\mathrm{}2t\mathrm{}+\mathrm{}2t\mathrm{}`$. Thus we are done as in the above lemma.
Case \[S2.4\]: $`D^{}`$ already has normal crossings and has arithmetic genus 1. So there is nothing to prove.
Case \[S3.1\]: an explicit calculation, probably well known, (cf. , page 65, lines 2-3) shows that the full transform of the maximal ideal of $`P`$ is the fundamental cycle $`Z`$ of $`D_n`$, while the full transform of the maximal ideal of $`P^{}`$, which is then $`W`$, is the fundamental cycle of the $`D_{n2}`$ configuration obtained by deleting the asymmetric end and its neighbour.
In this case let us choose as $`D^{}=\stackrel{~}{\mathrm{\Gamma }}+2W<D=\stackrel{~}{\mathrm{\Gamma }}+\stackrel{~}{Z}`$: since all the multiplicities of the components of $`D^{}`$ are then either $`1`$, $`2`$ or $`4`$, it follows that for $`1/4<t1/3`$ the round up $`\mathrm{\Delta }_t^{}:=\mathrm{}tD^{}\mathrm{}`$ equals $`\stackrel{~}{\mathrm{\Gamma }}+W`$. Since $`W^2=2`$, $`\stackrel{~}{\mathrm{\Gamma }}W=2`$, the self intersection $`(\mathrm{\Delta }_t^{})^2=1`$, thus $`\mathrm{\Delta }_t^{}`$ has arithmetic genus $`1`$ (topologically it is of elliptic type $`D_{n2}^{}`$ or $`I_{n6}^{}`$ in Kodaira’s notation) and this case is settled by virtue of 2).
Case \[S3.2\]: Also in this case an explicit calculation, probably well known, (cf. , page 65, lines 2-3) shows that the full transform of the maximal ideal of $`P`$ is the fundamental cycle $`Z`$ of $`E_8`$, while the full transform of the maximal ideal of $`P^{}`$, which is then $`W`$, is the fundamental cycle of the $`E_7`$ configuration obtained by deleting the furthest end.
In this case we write the multiplicities for the components of $`\stackrel{~}{Z}`$ starting from left to right (i.e., from middle length end (i.e., $`A`$) to longest end), and then we give the multiplicity for the shortest end: we get the sequence $`4,7,10,8,6,4,2`$ and then $`5`$. We may choose for convenience $`D^{}`$ as $`\stackrel{~}{\mathrm{\Gamma }}+\stackrel{~}{Z}`$ minus twice the longest end and minus its neighbor, i.e., we change the sequence to $`4,7,10,8,6,3,0,5`$. If we now choose $`3/10<t1/3`$, one can easily calculate that the round up $`\mathrm{\Delta }_t^{}:=\mathrm{}tD^{}\mathrm{}`$ equals $`\stackrel{~}{\mathrm{\Gamma }}+W`$ (topologically, it is of elliptic type $`E_7^{}`$ or $`III^{}`$ in Kodaira’s notation), and we are done as in the previous case. ∎
## 4. The Noether inequality
###### Theorem 4.1.
Let $`X`$ be a minimal Gorenstein 3-fold of general type with canonical singularities. Assume either $`p_g(X)2`$ or that $`|K_X|`$ is composed with a pencil of surfaces of type (1,2). Then
$$K_X^3\frac{7}{5}p_g(X)2.$$
###### Proof.
As we have seen in 1.5, we may take $`X`$ to be factorial with only terminal singularities. Because $`K_X^32`$, the inequality is automatically true for $`p_g(X)2`$.
We may suppose, from now on, that $`p_g(X)3`$. Denote by $`f:X^{}B`$ the fibration induced from $`\mathrm{\Phi }_{|K_X|}`$. Let $`S`$ be a general fiber of $`f`$. Lemma 4.5 of says $`0b=g(B)1`$. Theorem 3.2 says that we may take $`\alpha _3=\frac{3}{10}`$ for a general fiber $`S`$ of $`f`$; see the first part of 2.3.
Case 1. $`b=1`$. We may write $`a=p_g(X)=10m+c`$ where $`m0`$ and $`0c9`$ (obviously, for $`m=0`$ we have $`3c9`$).
When $`m>0`$, we take $`m_3:=7m+\mathrm{}_c`$, where $`\mathrm{}_c:=1`$, 0, 1, 2, 2, 3, 4, 4, 5, 6 respectively when $`0c9`$. Then one sees that
$$1\frac{m_3}{a}>\alpha _3=\frac{3}{10}.$$
Take $`m_2=a1m_3=3m+\mathrm{}_c`$, where $`\mathrm{}_c:=c1\mathrm{}_c`$. Then the inequality (2.5) gives
$`K_X^3`$ $`6(7m+\mathrm{}_c)+4(3m+\mathrm{}_c)4p_g(X)+4`$
$`=54m4p_g(X)+6\mathrm{}_c+4\mathrm{}_c+4`$
$`={\displaystyle \frac{7}{5}}p_g(X){\displaystyle \frac{7}{5}}c+2\mathrm{}_c`$
$`{\displaystyle \frac{7}{5}}p_g(X)2.`$
When $`m=0`$, we have $`3a=c9`$. Take $`m_3=`$2, 2, 3, 4, 4, 5, 6 respectively when $`3c9`$. We may easily check that $`1\frac{m_3}{a}>\frac{3}{10}`$. Take $`m_2:=a1m_3=`$0, 1, 1, 1, 2, 2, 2 respectively when $`3c9`$. By inequality (2.5), we may verify case by case that $`K_X^3>\frac{7}{5}p_g(X)2.`$
Case 2. $`b=0`$. We may write $`a=p_g(X)1=10m+c`$ where $`m0`$ and $`0c9`$ (for $`m=0`$ we have $`2c9`$).
Again when $`m>0`$, we take $`m_3:=7m+\mathrm{}_c`$ where $`\mathrm{}_c=1`$, 0, 1, 2, 2, 3, 4, 4, 5, 6 respectively when $`0c9`$. Then the calculation is similar to Case 1. Take $`m_2=a1m_3=3m+\mathrm{}_c`$ where $`\mathrm{}_c=c1\mathrm{}_c`$. Then the inequality (2.5) gives
$`K_X^3`$ $`6(7m+\mathrm{}_c)+4(3m+\mathrm{}_c)4p_g(X)+4h^2(𝒪_X)+10`$
$`54m4p_g(X)+6\mathrm{}_c+4\mathrm{}_c+10`$
$`={\displaystyle \frac{7}{5}}p_g(X){\displaystyle \frac{7}{5}}c+2\mathrm{}_c+{\displaystyle \frac{3}{5}}`$
$`{\displaystyle \frac{7}{5}}p_g(X){\displaystyle \frac{7}{5}}.`$
When $`m=0`$ and $`2a=c9`$, one can in a similar way verify that $`K_X^3>\frac{7}{5}p_g(X)\frac{7}{5}.`$
###### 4.2.
Proof of the main results.
###### Proof.
Now both Theorem 1.2 and Corollary 1.4 follow directly from 2.2 and Theorem 4.1. ∎
###### Remark 4.3.
A quite natural problem left to us is the possibility of generalizing Theorem 1.2 to the case where $`X`$ is Gorenstein minimal. Unfortunately the method of Theorem 4.3 of only works when $`X`$ is smooth. One needs a new method to treat the difficult case where $`X`$ is canonically fibred by curves of genus 2. However we would like to put forward the following:
###### Conjecture 4.4.
The Noether inequality
$$K^3\frac{2}{3}(2p_g5)$$
holds for any projective minimal Gorenstein 3-fold of general type $`X`$.
Acknowledgment. This paper was begun while Zhang was visiting Fudan University at the end of 2004. It was finished when Chen was visiting both Universit$`\ddot{\text{a}}`$t Bayreuth and Universit$`\ddot{\text{a}}`$t Duisburg-Essen in the summer of 2005 financially supported by the joint Chinese-German project on ”Komplexe Geometrie” (DFG & NSFC). The authors would like to thank these 3 universities for their support. Especially Chen would like to thank Ingrid Bauer, Hélène Esnault, Thomas Peternell, Eckart Viehweg for their hospitality and effective discussion. |
warning/0507/astro-ph0507506.html | ar5iv | text | # Sunyaev-Zeldovich Cluster Counts as a Probe of Intra-Supercluster Gas
## 1 Introduction
As much as $`3040\%`$ of all baryons are believed to be in warm gas in large scale filamentary structures connecting clusters and superclusters (SCs) of galaxies. Observational evidence for the so-called *warm-hot intergalactic medium* (WHIM) is still quite rudimentary. First systematic attempts to detect the low X-ray surface brightness enhancement that could possibly be expected from the WHIM were conducted by Persic et al. (1988, 1990) who searched the HEAO-1 A2 database in directions to a sample of SCs. In both these analyses, as well as a similar analysis of Ginga observations of the Coma-A1367 SC (Tawara et al. 1993), only upper limits were obtained on the flux from the putative intra-SC (ISC) gas. More recently, evidence for soft X-ray excess emission in regions around several clusters has been claimed by Kaastra et al. (2003), who carried out a detailed analysis of XMM-Newton line and continuum measurements of extended regions around several clusters.
X-ray emission may be expected, of course, also from the WHIM in the Local SC (LSC); analysis of the HEAO-1 A2 database led Boughn (1999) to conclude that this emission was actually detected. Adopting a simple model for the morphology of the LSC, and assuming uniform temperature and density LSC gas, Boughn deduced the mean electron density to be $`n_e=2.510^6\left(a/20Mpc\right)^{1/2}\left(kT_e/10keV\right)^{1/4}cm^3`$, when using the specified scalings for the values of the electron temperature, $`kT_e`$, and the length of the semi-major axis of the LSC, $`a`$. More generally, Kaastra et al. (2003) have reported the detection of line and continuum WHIM emission in XMM-Newton observations of extended regions around five cluster. While the observational results on the properties of LSC gas are quite uncertain, gaseous filamentary structures seem to be a ubiquitous feature of the large scale mass distribution described by hydrodynamical simulations (e.g. Jenkins et al. 1998). Of considerable interest are the statistics and morphologies of these filaments; these are currently being quantified (e.g. Colberg et al. 2004).
The scant observational information on the WHIM and its properties in the LSC provides strong motivation for assessing the feasibility of probing it by measurements of the Sunyaev-Zeldovich (S-Z) effect (Sunyaev & Zeldovich 1972). The effect - a spectral change of the Planck spectrum due to Compton scattering of the cosmic microwave background (CMB) by electrons in clusters of galaxies - is a major cluster and cosmological probe (as reviewed by Rephaeli 1995, Birkinshaw 1999, and by Carlstrom et al. 2002). In fact, the possibility of detectable S-Z signals in directions to SCs has been considered long ago (Persic, Rephaeli & Boldt 1988, Rephaeli & Persic 1992, Rephaeli 1993), but given the substantial uncertainty in the integrated pressure of ISC, no definite predictions could be made. Measurement capabilities have greatly advanced since then, so the prospects of detecting weak S-Z signals have correspondingly improved.
Interest in the impact of ISC gas on the CMB has also increased. In particular, it has been suggested that gas in the Local Group may be responsible for the power suppression of primary CMB temperature anisotropy on large angular scales. This suppression was deduced from analyses of the COBE/DMR (Hinshaw et al. 1996) and WMAP databases (Bennett et al. 2003). However, the estimated low Comptonization parameter of LG gas makes this suggestion quite unlikely (Rasmussen & Pedersen 2001). More recently, Abramo & Sodré (2003) proposed that warm LSC gas may be responsible for the deduced CMB low multipole power suppression, presumably caused by coincidental alignment of a hot spot in the CMB with the line of sight connecting the Galaxy with the Virgo cluster. At the relevant spectral range of the DMR and WMAP experiments ($`2090GHz`$), the S-Z diminution lowers the power in the low multipoles. Clearly, for this to occur the integrated electron pressure has to be sufficiently high.
WHIM in the LSC may significantly screen our view of the CMB and S-Z sky. In addition to the possible impact on the power spectrum of CMB primary anisotropy, it will also affect the ability to measure the S-Z induced anisotropy and cluster number counts. WHIM in the LSC may appreciably increase S-Z cluster number counts. Moreover, the marked ellipsoidal shape of the LSC, and the far off-center position of the Galaxy, may result in a substantially anisotropic distribution of S-Z counts across the sky. This anisotropy in S-Z counts is obviously in addition to that generated by our motion in the CMB frame (the CMB kinematic dipole) whose effect on cluster counts was recently explored by Chluba et al. (2004).
This paper is arranged as follows: the LSC gas model, the method for calculating the directional distribution of S-Z cluster counts, and the relevant expressions for calculating the S-Z angular power spectrum, are described in section § 2. The main results of the calculations are presented in § 3. Additional aspects of the LSC gas model pertaining to S-Z measurements towards the Virgo cluster are discussed in § 4. A general discussion follows in § 5.
## 2 Model and Method of Calculation
The physical properties of ISC have not yet been well determined; what can be done at present is an attempt to incorporate the established observational results in a reasonable model in order to provide useful insight to guide upcoming observations, mostly of the S-Z effect. Relevant LSC properties that seem to have been roughly determined are the total mass, baryonic mass fraction, and its overall configuration. Knowing the baryon fraction in clusters (e.g. Carlstrom et al. 2002), and adopting the scaling to obtain the global baryonic fraction (based on hydrodynamical simulations; Evrard 1997), $`1014\%`$, the total LSC mass, $`10^{15}M_{}`$, then yields an estimate for the LSC WHIM mass, $`35.610^{13}M_{}`$, assuming (based on current expectations) that it comprises some $`3040\%`$ of the total mass. The distribution of galaxies and clusters in the LSC can be described by a spheroid of semi-axes measuring $`A=20Mpc`$, $`B=6.7Mpc`$, $`C=3.3Mpc`$ (Tully 1982); the Galaxy is on the major axis of the LSC at a distance of $`15Mpc`$ from the center of the spheroid, where the Virgo cluster is located.
Gas in the LSC is likely to be differently distributed than the galaxies, judging by the irregular filamentary structures seen in hydrodynamical simulations. However, it is reasonable to expect that most of the WHIM is concentrated along the major axis of the LSC mass distribution. Our specific model for gas in the LSC is based on the findings of Colberg et al. (2004), who analyzed results of N-body simulations (by Kauffmann et al. 1999) in order to characterize morphologically the filaments observed in the simulations. They found that most filaments have cross sections of $`12Mpc`$, with the gas density decreasing outward from the major axis of the filament as $`r^2`$ beyond some scale radius $`r_s`$. The LSC gas is assumed to have an ellipsoidal configuration with semi-axes (in Mpc) of either $`20\times 1\times 1`$ or $`20\times 2\times 2`$. The highly prolate ellipsoid represents the gas ‘filamentary’ structure. The gas mean density in these structures are $`n_e=2.410^5cm^3`$, or a factor $`4`$ lower, respectively; its temperature is scaled to a value of $`1keV`$).
Clearly, the expected anisotropic morphology of ISC gas in the LSC will be reflected in anisotropic cluster counts when S-Z surveys are conducted in different sky directions. The Comptonization parameter towards the Virgo cluster, a direction which in spherical coordinates corresponds to angles $`(\varphi ,\theta )=(0,\pi /2)`$, is
$$y=\sigma _T\frac{kT_e}{m_ec^2}n_e𝑑\mathrm{}=3.410^6\left(\frac{kT_e}{1keV}\right),$$
(1)
whereas the corresponding value calculated for the opposite direction $`(\pi ,\pi /2)`$ is seven times lower. The non-negligible level of the Comptonization parameter of LSC gas may modify the observed distribution of S-Z clusters across the sky; the effective y-parameter measured along a los to a cluster will have contributions from both intracluster (IC) and ISC gas, giving rise to an intensification of the apparent S-Z signal due to the cluster alone. In spherical coordinates, the path length through the ISC gas configuration as a function of position angle $`\widehat{n}(\varphi ,\theta )`$ is
$$r(\varphi ,\theta )=\frac{\frac{30\mathrm{sin}\theta \mathrm{cos}\varphi }{A^2}+\sqrt{\left(\frac{30\mathrm{sin}\theta \mathrm{cos}\varphi }{A^2}\right)^24\left(\frac{15^2}{A^2}1\right)\left[\frac{\mathrm{sin}^2\theta \mathrm{cos}^2\varphi }{A^2}+\frac{\mathrm{sin}^2\theta \mathrm{sin}^2\varphi }{B^2}+\frac{\mathrm{cos}^2\theta }{C^2}\right]}}{2\left[\frac{\mathrm{sin}^2\theta \mathrm{cos}^2\varphi }{A^2}+\frac{\mathrm{sin}^2\theta \mathrm{cos}^2\varphi }{B^2}+\frac{\mathrm{cos}^2\theta }{C^2}\right]}.$$
(2)
The Comptonization parameter measured along this direction is then simply
$$y(\varphi ,\theta )=\sigma _T\frac{kT_{SC}}{m_ec^2}n_{SC}r(\varphi ,\theta ).$$
(3)
where $`T_{SC}`$ and $`n_{SC}`$ denote the uniform temperature and electron density of the LSC gas.
The flux received from a cluster lying in the direction $`(\varphi ,\theta )`$ is a sum of the cluster flux and the S-Z signal generated in the LSC; assuming a $`\beta `$-King profile for the IC gas density with $`\beta =2/3`$ and an isothermal temperature distribution, the total flux at a frequency $`\nu `$ is
$`\mathrm{\Delta }F(\varphi ,\theta )\mathrm{\Delta }F_C(\varphi ,\theta )+\mathrm{\Delta }F_{SC}(\varphi ,\theta )=`$
$`{\displaystyle \frac{2\left(kT_\gamma \right)^3}{\left(hc\right)^2}}\left({\displaystyle \frac{k\sigma _T}{m_ec^2}}\right)g(x)[2n_0T_0r_c{\displaystyle _0^{p\theta _c}}{\displaystyle \frac{\mathrm{tan}^1\sqrt{\frac{p^2\left(\theta /\theta _c\right)^2}{1+\left(\theta /\theta _c\right)^2}}}{\sqrt{1+\left(\theta /\theta _c\right)^2}}}d\mathrm{\Omega }`$ (4)
$`+n_{SC}T_{SC}{\displaystyle _0^{2\pi }}{\displaystyle _0^{\sigma _b}}r(\varphi ,\theta )\theta d\theta d\varphi ],`$
where $`T_{cmb}=2.726^{}K`$, $`T_0`$, $`n_0`$, $`T_{SC}`$ and $`n_{SC}`$ denote the temperature and density of the cluster and LSC gas, respectively, $`pR_{vir}/r_c`$, $`\sigma _b`$ is the beam width, and $`g(x)`$ is the non-relativistic spectral form of the thermal S-Z effect, valid here owing to the relatively low temperatures of both the ISC and Virgo IC gas. At higher temperatures a relativistically correct expression (e.g. Shimon & Rephaeli 2004) must be taken into account.
The calculations were carried out in a standard $`\mathrm{\Lambda }`$CDM model, specified by the parameters $`\mathrm{\Omega }_m=0.3,\mathrm{\Omega }_\mathrm{\Lambda }=0.7,h=0.7,n=1,\sigma _8=1`$. We use the cluster mass function of Sheth & Tormen (2001) mass, and the temperature-mass relation
$$T=1.3\left(\frac{M}{10^{15}h^1M_{}}\right)^{2/3}\left(\mathrm{\Delta }_cE^2\right)^{1/3}\left(12\frac{\mathrm{\Omega }_\mathrm{\Lambda }(z)}{\mathrm{\Delta }_c}\right)keV,$$
(5)
where $`E^2=\mathrm{\Omega }_m(1+z)^3+\mathrm{\Omega }_\mathrm{\Lambda }`$. A non-evolving gas fraction $`f_0=0.1`$ is assumed. And since our treatment here is most relevant for the Planck all sky survey, we have selected the spectral band to be the HFI $`\nu _0`$=353 GHz, with $`\mathrm{\Delta }\nu =116GHz`$; correspondingly, the beam profile is taken to be Gaussian with FWHM of $`7.1^{}`$.
Number counts are calculated based on the expression
$$N(>\mathrm{\Delta }\overline{F}_\nu )=r^2\frac{dr}{dz}𝑑z_{\mathrm{\Delta }\overline{F}_\nu (M,z)}N(M,z)𝑑M,$$
(6)
where the lower limit is assigned such that the flux received from a cluster with mass $`M`$ situated at redshift $`z`$ exceeds the flux limit of the survey. In order to estimate the number of clusters that may be observed above a flux limit within the framework of the LSC model, the LSC S-Z flux is mapped using the second term in equation (4); the results are illustrated in the upper left-(model 1) and right-hand (model 2) panels of figure 1. The number of clusters with flux exceeding the limit $`\mathrm{\Delta }F_\nu `$ are then calculated using equation (6). Calculated number counts are then multiplied by $`610^4`$ to yield the number of clusters expected within square patches of sky measuring $`5^{}\times 5^{}`$. These are depicted in the lower panel of figure 1. For the detection of a cluster with a given flux to be possible, the combined S-Z flux from the cluster and the intervening LSC gas must exceed that of the measurement sensitivity. For example, if the latter is 30 mJy, then in a sky region where the S-Z flux due to the LSC is 20 mJy, the minimum detection level for a cluster is 10 mJy. Having determined the minimum flux, the expected cluster counts within a 25 square degree patch of sky in the requested direction can be read off the graph in the lower panel of figure (1).
The angular power spectrum of the S-Z effect due to LSC gas can be calculated using the usual spherical harmonic expansion of the temperature anisotropy. Here, however, the flat sky approximation is invalid due to the large angular scales involved. The basic expression is
$$a_\mathrm{}m=_0^{2\pi }_0^\pi \frac{\mathrm{\Delta }T}{T}(\varphi ,\theta )y_\mathrm{}m^{}(\varphi ,\theta )\mathrm{sin}\theta d\theta d\varphi ,$$
(7)
which, after projecting the three dimensional function $`\mathrm{\Delta }T/T(r,\varphi ,\theta )`$ onto the two dimensional celestial sphere, and some algebraic manipulation, assumes the form
$$a_\mathrm{}m=\frac{2\sigma _Tk_BT_en_e}{m_ec^2}_0^{2\pi }_0^\pi r(\varphi ,\theta )y_\mathrm{}m^{}(\varphi ,\theta )\mathrm{sin}\theta d\theta d\varphi ,$$
(8)
where $`r(\varphi ,\theta )`$ is given by equation (2). The angular power spectrum may then be easily calculated as
$$\mathrm{}(\mathrm{}+1)/2\pi C_{\mathrm{}}=\mathrm{}(\mathrm{}+1)/2\pi \frac{_{m=\mathrm{}}^{\mathrm{}}\left|a_\mathrm{}m\right|^2}{2\mathrm{}+1}.$$
(9)
## 3 Results
For flux limits of $`30mJy`$ and $`60mJy`$ the corresponding all-sky cluster counts sum up to $`\mathrm{40\hspace{0.17em}000}`$ and $`\mathrm{4\hspace{0.17em}300}`$, respectively. Assuming that the cluster population is homogeneously distributed across the sky, this amounts to $`24`$ and $`3`$ clusters in each 25 square degrees (rectangular sky region). It is important to emphasize that the number of potentially detected clusters in future S-Z surveys predicted by numerical studies is quite sensitive to the choice of parameters characterizing the background cosmology and cluster properties, and consequently, these counts merely reflect the specific modeling detailed in the last section. Now suppose that observations are conducted in the directions $`(\varphi ,\theta )=(0^{},90^{})`$ (i.e. towards the longest possible los within the LSC halo), $`(5^{},90^{})`$, $`(180^{},90^{})`$ (designating the direction opposite to the Virgo cluster), and, e.g., $`(180^{},20^{})`$. The S-Z fluxes induced by the LSC gas in these directions are (see figure 1) $`21,7,3,0.5mJy`$, respectively, implying that for a flux detection limit of $`30mJy`$ clusters will be detected if their S-Z signals exceed the complementary values ($`9,23,27,29.5mJy`$, respectively) in these directions. The corresponding limiting values for a mean gas density of $`n_e=610^6cm^3`$ are $`25,27,29.3,29.7mJy`$, respectively.
The number of clusters that could be detected in these directions, i.e., those which meet the detectability criterion, can be easily read off the diagram in figure 1. For example, an experiment with a detection limit of $`30mJy`$ yields $`570,73,50,40`$ clusters per 25 square degrees. The corresponding counts for the larger ellipsoid (with gas density lower by a factor of 4) are $`60,50,40,40`$. A limiting flux of $`60mJy`$ yields counts of $`20,8,7,6`$ in the first case, and $`68`$ for all directions in the second case. While these estimates of the cluster number counts are substantially uncertain, and the fact that there is also considerable variation along different los across 25 square degrees (particularly towards both ends of the spheroid major axis), the results nonetheless demonstrate the possible impact of LSC gas, especially in a survey with $`30mJy`$ flux limit, and in direction to the Virgo cluster. The difference between counts along los in the Virgo region and those in other directions markedly exceed the statistical uncertainty (of a purely Poissonian distribution of clusters).
The angular power spectrum of the S-Z effect induced by LSC gas is illustrated in figure 2. The maximum power level is significantly lower than the corresponding values found by Abramo & Sodré. This is due to our lower values of the density and temperature, each lower by a factor of $`2`$ than their values ($`510^5cm^3`$, and 2 keV), and the very different LSC gas configuration assumed by Abramo & Sodré, which is reflected in a different distribution of power among the explored multipole range. In particular, more power is seen at higher multipoles (with regard to Abramo & Sodré’s results), which is expected owing to its highly prolate morphology. Consequently, given the substantial uncertainty in the properties of ISC gas, one cannot convincingly argue that the observed suppression of the primary CMB temperature power spectrum is due to Compton scattering in the LSC.
## 4 Implications for the Virgo Cluster
The central location of Virgo in the LSC implies that measurement of the S-Z effect in this cluster will likely include a substantial contribution from LSC gas. The Comptonization parameter measured along the direction $`(\varphi ,\theta )`$ is
$$y=\sigma _T\frac{kT_{SC}}{m_ec^2}n_{SC}r(\varphi ,\theta )+2\sigma _T\frac{kT_0}{m_ec^2}n_0r_c\mathrm{}(\varphi ,\theta ),$$
(10)
where $`r(\varphi ,\theta )`$ is given in equation (2) and
$$\mathrm{}(\varphi ,\theta )=\left[1+\left(\frac{\theta }{\theta _c}\right)^2\right]^\delta \sqrt{p^2(\theta /\theta _c)^2}H_2F_1[\frac{1}{2},\delta ,\frac{3}{2},\frac{(p^2(\theta /\theta _c)^2)}{1+(\theta /\theta _c)^2}],$$
(11)
is the integrated line of sight through the Virgo cluster. Here $`\theta _c`$ is the core radius, $`\delta 3/2\beta \gamma `$ where $`\gamma `$ is a polytropic index, $`p`$ is the virial to core radius ratio, and $`H_2F_1`$ is the hypergeometric function. The main gas configuration in Virgo is centered on the giant elliptical M87, to which we refer here. Excluding a relatively small region in the center of M87, the gas is taken to be isothermal with a (King) $`\beta `$ density profile, having (Matsumoto et al. 2000) $`n_0=0.019cm^3`$, $`T_0=2.28keV`$, $`r_c=0.014Mpc`$, $`p=125`$, and $`\beta =0.4`$. The magnitude of $`y`$ along the central Virgo los increases by $`20\%`$ due LSC gas, its gradient across the cluster is appreciably shallower than would be expected from the intrinsic density profile, and the cluster seems slightly larger due to the enhancement of the cluster Comptonization parameter.
The combined motion of the earth towards Virgo and the motion of the LSC in the CMB frame accounts for the dipole term; the deduced velocity is $`600kms^1`$ in the direction $`(\mathrm{}=264^{},b=48^{})`$ in galactic coordinates (Fixen et al. 1996). The temperature change generated by the dipole in the direction of Virgo is $`\mathrm{\Delta }T/T1.810^3`$, whereas the combined S-Z temperature change due to the LSC halo and Virgo cluster peaks at $`\mathrm{\Delta }T/T510^5`$. Clearly, the dipole may interfere with S-Z measurements of the Virgo cluster and should therefore be taken into account. This would generally be the case for large angular diameter clusters, across which the dipole gradient can be significant. The temperature change due to the dipole term is simply
$$\left(\frac{\mathrm{\Delta }T}{T}\right)_{dipole}\frac{T_{obs}T_0}{T_0}=\frac{v}{c}\mathrm{cos}\theta ,$$
(12)
whereas the temperature change due to the thermal S-Z effect in the R-J region is $`\mathrm{\Delta }T/T=s(x)y`$, where $`s(x)=x/\mathrm{tanh}(x/2)4`$. In order to be able to compare the contributions of both effects to the temperature change, the galactic coordinates specified above are converted into equatorial coordinates, and transformed such that the Virgo cluster lies in the direction $`(\varphi =0^{},\theta =90^{})`$. In this coordinate system the velocity vector points at $`(\varphi =340^{},\theta =110^{})`$. Figure 3 illustrates the relative temperature change $`(\mathrm{\Delta }T/T)`$ due to the LSC+Virgo S-Z signal, and the latter combined with the CMB dipole $`(\mathrm{\Delta }T/T)_{dipole}+(\mathrm{\Delta }T/T)_{SC+Virgo}`$ within a $`10^{}\times 10^{}`$ patch of sky centered around the Virgo cluster. The relative temperature change due to the dipole does not change considerably across the patch due to its small angular size and the fact that the angle formed between the dipole and the vector pointing at Virgo lies within the range $`21^{}\gamma 35^{}`$, at a significant angular distance from the direction at which the cosine term of the dipole has the steepest slope, $`\gamma =90^{}`$. The angle between the two vectors can be calculated using
$$\mathrm{cos}\gamma =\mathrm{cos}\theta _1\mathrm{cos}\theta _2+\mathrm{sin}\theta _1\mathrm{sin}\theta _2\mathrm{cos}(\varphi _1\varphi _2).$$
(13)
The combined relative temperature change is still dominated by the dipole term, although its profile is now deformed due to the S-Z induced temperature change. In particular a ’warm spot’ may be seen in the direction of Virgo (which translates into an enhancement of $`3\%`$ of the temperature change due to the dipole term in this direction), by virtue of the fact that at a frequency of 353 GHz the cluster behaves as a source for CMB photons. Consequently, it is obvious that the dipole term must be taken into account in measurements of S-Z signals in nearby clusters.
The impact of LSC gas on measurements of the S-Z towards the central Virgo region is appreciable only when observations are made in a survey mode - such as planned with PLANCK - with relatively short scan time over the region, and consequently insufficient sensitivity to remove such an LSC component. Clearly, the additional signal due to LSC gas can largely be subtracted out when pointed observations are made with a suitably selected beam-through pattern.
## 5 Discussion
Although there is yet no unequivocal observational evidence for the WHIM, theory and hydrodynamic cosmological simulations indicate that it is highly prevalent and constitutes a substantial fraction of the baryonic mass fraction. Currently available observations of WHIM in several cluster regions imply appreciable levels of the Comptonization parameter, $`10^5`$, comparable to values in most clusters. Even though we used a simplified model for the WHIM in the LSC in order to explore some of its consequences, we consider the results presented here to be qualitatively valid. In fact, our approximation of the WHIM filamentary structure by an ellipsoidal configuration (with a large volume) quite likely underestimates its S-Z impact, given the basic premise that it constitutes an appreciable fraction of the baryon mass of the LSC. As indicated by morphological analyses of such filaments seen in hydrodynamical simulations, their density grows towards the central axis, such that a clear preferential direction marked by the axis can generate even higher S-Z signals along this direction. Moreover, the adopted symmetric morphology underestimates the density along the main axis in direction to Virgo as there is no cluster in the opposite side of the LG-Virgo direction.
Several potential observational consequences of the presence of such gaseous component have been discussed in this paper: the angular power spectrum of the induced S-Z signal has been calculated and shown to be rather low, particularly with regard to what has been speculated in the literature as a possible explanation for the suppression of primary CMB power on the lowest multipoles. While we cannot rule out a suppression at the claimed level, we do not find this to be very likely. The effect of LSC gas on directional S-Z cluster counts is proportional to the los crossing the volume of gas in the LSC. On the other hand, a definite prediction can be made that the distribution of S-Z cluster counts across the sky will be determined to be anisotropic close to the level estimated here. The impact is relatively significant owing to the highly non-linear form of the mass function; a relatively low S-Z flux induced by intervening LSC gas may give rise to a considerable increase in the number of potentially detected clusters. This is in particular true in light of the decidedly higher abundance of low-flux S-Z clusters predicted by the mass function. Needless to say, this effect is most noticeable along the main axis towards Virgo. S-Z measurements of the Virgo cluster may also be affected by the presence of the WHIM, as well as by the dipole distribution of the primary CMB temperature anisotropy. This is even more relevant to measurements of the kinematic effect, since the induced temperature change does not depend on the temperature of the LSC gas, which tends to be lower than typical IC temperatures. These effects will have to be taken into account in realistic analyses of the S-Z signal from the Virgo cluster.
Finally, although observational evidence for the WHIM is still scarce, it is interesting to use the currently available data to assess its potential of generating a significant S-Z signal towards other SCs. Kaastra et al. (2003) list five clusters (A1795, Sersic 159-03, MKW3s, A2050, A2199) reported to include a WHIM component, modeled to lie within a spherical halo, be isothermal and of homogeneous density. For each cluster they list the WHIM electron density, temperature, and size, such that the Comptonization parameter along the radius of the spherical halo can be easily calculated. They all turn out to lie within the range $`y(0.40.8)10^5`$, i.e. near the level taken here for the LSC. Thus, the S-Z consequences of ISC gas we considered here are likely to be relevant in other nearby SCs, but obviously on smaller angular scales. |
warning/0507/cond-mat0507482.html | ar5iv | text | # Superconducting Proximity Effect through a Magnetic Domain Wall
(July 20, 2005)
## Abstract
We study the superconducting proximity effect in a superconductor-ferromagnet-superconductor (SFS) heterostructure, containing a domain wall in the ferromagnetic region. For the ferromagnet we assume an alloy with an exchange splitting of the conduction bands comparable to the superconducting gaps. We calculate the modification of the density of states in the center of the domain wall as a result of the proximity effect. We show that the density of states is sensitive to domain wall parameters due to triplet-pairing correlations created in vicinity of the domain wall. We present a theoretical tool which in a very effective way enables retaining the full spatially dependent spin-space structure of the problem.
Introduction: Most promising candidates for mesoscopic devices with novel functionality are hybrid structures containing superconducting elements. The key phenomenon that controls the behavior of such systems is the proximity effect. When a superconducting material is placed in contact with a normal metal (N), the superconducting pair correlations leak over to the normal-metal side, changing its conduction properties in the vicinity of the separating interface. Quite similarly, the properties on the superconducting side are also changed (the energy gap $`\mathrm{\Delta }_0`$ is suppressed) due to the contact to a normal metal. An alternative but equivalent way of thinking about the proximity effect is through Andreev-reflection and processes: an incoming electron from the normal side is transmitted together with another one as a Cooper pair into the superconducting side. This phase-coherent electron-hole conversion results in a nonzero pair amplitude in the normal metal.
In the diffusive limit, the correlations relating to an incident electron with an energy $`E`$ (the range of energies being set by the temperature $`T`$) above the chemical potential extend a characteristic distance of $`\xi _N=\sqrt{D/E}`$ into the normal metal;Kogan82 here $`D`$ is the diffusion constant in N. If the extent of the N region is finite, another energy scale, $`E_TD/L^2`$, enters the problem; $`L`$ denotes the width of N. This so-called Thouless energy has associated with it one of the generic features of diffusive superconductor-normal metal heterostructures, the minigap: the density of states in the normal metal develops a gap around the chemical potential in a manner similar to a superconductor (S) but with a smaller magnitude.
If the normal conductor is replaced by a ferromagnet (F), a multitude of new effects arise due to the emergence of yet another energy scale, that of the exchange splitting $`J`$ of the two spin bands. Both on the theoreticalvol ; kad ; hal ; rad ; fom ; bar ; hue ; morten ; Buzdin05 and on the experimental,kon ; rya ; gui ; gir ; gu ; aum ; gee ; frolov ; beckmann side interest has grown recently in the rich physics of such systems. One source for new behavior is that, in the case with a singlet superconductor, the induced pair amplitude in the ferromagnet is oscillatory.buz However, the exchange splitting also gives rise to dephasing which, in turn, results in the decay of induced correlations over a characteristic distance $`\xi _J=\sqrt{D/(E+J)}`$.Bulaevskii85 Unfortunately, since $`J`$ is of the order of the Fermi energy $`E_F`$ in typical ferromagnetic metals, this distance is very short. Still, experimental indication of the oscillatory behavior has been obtained in thin ferromagnetic layers and, relevant to the present paper, in weakly ferromagnetic alloys with $`JE_F`$. kon ; rya
Another question of current interest in SF proximity systems is the role of equal-spin triplet correlations.ber ; mat If created, e.g. near magnetic inhomogeneities, such correlations would not be affected by the exchange splitting but could penetrate considerably longer distances into F.ber Finally, the importance of domain walls has also been stressed for the Andreev conductance.cht ; mel .
In this paper we study an SFS structure, shown schematically in Fig. 1, in equilibrium. The ferromagnetic region consists of two domains with magnetizations oriented in opposite directions. The domains are separated by a domain wall, where the magnetization rotates continuously between the asymptotic values. While varying in direction, the magnitude $`J`$ is assumed constant throughout the F region. We show, that the local density of states (LDOS) in the F region is strongly modified by the presence of the domain wall. In particular, we show that it can be very sensitive to the thickness of the domain wall in a certain parameter region.
Basic equations: Proximity effect is a spatially inhomogeneous phenomenon. An appropriate theoretical tool to treat such a problem is the quasiclassical theory of superconductivity, eil ; lar which in its diffusive version has been formulated by Usadel.usdl In equilibrium, the physical information is contained in the retarded Green functions $`\widehat{G}(z,E)`$. Here, we assume spatial dependence in the coordinate $`z`$ only, and $`E`$ denotes the energy as measured from the chemical potential. The 4$`\times `$4 matrix structure, arising from particle-hole and spin degrees of freedom, is denoted by the hat ( $`\widehat{}`$ ) accent,
$$\widehat{G}=(\begin{array}{cc}𝒢& \\ \stackrel{~}{}& \stackrel{~}{𝒢}\end{array}).$$
(1)
The off-diagonal elements determine the superconducting pair amplitude. Quantities denoted with the “tilde” are related to those without one through $`\stackrel{~}{𝒜}(z,E)=𝒜(z,E^{})^{}`$. All the elements in Eq. (1) are 2$`\times `$2 spin matrices: e.g. $`𝒢=𝒢_{\alpha \beta }`$ with $`\alpha ,\beta =\{,\}`$. The Green functions satisfy the Usadel equation,usdl
$$[E\widehat{\tau }_3\widehat{\mathrm{\Delta }}𝑱\widehat{𝝈},\widehat{G}]_{}+\frac{D}{\pi }\frac{\mathrm{d}}{\mathrm{d}z}\left(\widehat{G}\frac{\mathrm{d}}{\mathrm{d}z}\widehat{G}\right)=\widehat{0},$$
(2)
where the symbol $``$ denotes matrix multiplication, and $`[\widehat{A},\widehat{B}]_{}=\widehat{A}\widehat{B}\widehat{B}\widehat{A}`$. In writing Eq. (2) we have followed the standard way to describe ferromagnetic materials through a spin-dependent energy shift,Bulaevskii85 which has the form $`E\widehat{\tau }_3E\widehat{\tau }_3𝑱\widehat{𝝈}.`$ Here, $`\tau _3`$ denotes the third Pauli-matrix in Nambu space, the vector $`𝑱`$ denotes the effective exchange field of the ferromagnet, and $`\widehat{\mathrm{\Delta }}`$ is the superconducting order parameter (appropriate for weak-coupling spin-singlet pairing). The components of the vector $`\widehat{𝝈}`$ and the order parameter are given by
$$\widehat{\sigma }_i=\left(\begin{array}{cc}\sigma _i& 0\\ 0& \sigma _i^{}\end{array}\right),\widehat{\mathrm{\Delta }}=\left(\begin{array}{cc}0& \mathrm{\Delta }\\ \mathrm{\Delta }^{}& 0\end{array}\right)$$
(3)
where $`\sigma _i`$ are Pauli spin matrices, $`i=x,y,z`$, and $`\mathrm{\Delta }=\mathrm{\Delta }_0i\sigma _y`$. The above procedure is appropriate for describing situations for which $`JE_F`$, which holds e.g. for the ferromagnetic alloys used in Refs. \[kon, ; rya, \]. In writing Eq. (2), we have chosen the normalization according to
$$\widehat{G}\widehat{G}=\pi ^2\widehat{1}.$$
(4)
Riccati parameterization: The spin-dependent nature of SF proximity systems calls for a formulation of the quasiclassical theory that retains the full spin-space structure, especially in studying situations where the exchange-field-orientation varies in space (such as in a domain wall). Within the Eilenberger theory a very convenient formulation already exists, esc employing the so-called Riccati parameterization. scho ; nag The extension of this method to the Usadel theory was achieved only recently, Eschrig04 and has been applied to non-equilibrium situations, Cuevas05 and to FSF-systems with homogeneous magnetizations.Lofwander05 Here we demonstrate its usefulness by applying it to a SFS system with a spatially changing magnetization in a domain wall, a case where the conventional so-called $`\theta `$-parameterizationbelz is not applicable. The spin-dependent Riccati parameterization,esc
$`\widehat{G}`$ $`=`$ $`i\pi \widehat{N}(\begin{array}{cc}(1+\gamma \stackrel{~}{\gamma })& 2\gamma \\ 2\stackrel{~}{\gamma }& (1+\stackrel{~}{\gamma }\gamma )\end{array}),`$ (7)
with
$$\widehat{N}=(\begin{array}{cc}(1\gamma \stackrel{~}{\gamma })^1& 0\\ 0& (1\stackrel{~}{\gamma }\gamma )^1\end{array})$$
(8)
automatically accounts for the normalization (4), which is essential for practical numerical calculations. It is enough to determine one 2$`\times `$2 matrix in spin space, $`\gamma `$. The other, $`\stackrel{~}{\gamma }`$, follows from the above-mentioned (fundamental) symmetry. The transport equation for $`\gamma `$ follows from Eq. (2), and readsEschrig04
$`{\displaystyle \frac{\mathrm{d}^2\gamma }{\mathrm{d}z^2}}`$ $`+`$ $`\left({\displaystyle \frac{\mathrm{d}\gamma }{\mathrm{d}z}}\right){\displaystyle \frac{\stackrel{~}{}}{i\pi }}\left({\displaystyle \frac{\mathrm{d}\gamma }{\mathrm{d}z}}\right)={\displaystyle \frac{i}{D}}[\gamma \mathrm{\Delta }^{}\gamma .`$ (9)
$`.(E𝑱𝝈)\gamma \gamma (E+𝑱𝝈^{})\mathrm{\Delta }],`$
Here, the expression for $`\stackrel{~}{}`$ is obtained by comparing Eq. (1) with Eqs. (7)–(8).
Boundary conditions: Additionally, boundary conditions are required for the different interfaces of the system. Such conditions have been formulated by Nazarov.naz The outer surfaces ($`z=z_o`$) of the superconductors are assumed to border to an insulating region, and the appropriate condition is $`_z\widehat{G}(z_o,E)=0`$, i.e.
$$\frac{d\gamma }{dz}(z_o,E)=0.$$
(10)
On the other hand, the two inner SF interfaces ($`z=z_i^S`$ for the S side, $`z=z_i^F`$ for the F side) are assumed in the following transparent. The corresponding boundary conditions are in this case $`\widehat{G}(z_i^S,E)=\widehat{G}(z_i^F,E)`$, $`\sigma _S_z\widehat{G}(z_i^S,E)=\sigma _F_z\widehat{G}(z_i^F,E)`$, leading to
$`\gamma (z_i^S,E)`$ $`=`$ $`\gamma (z_i^F,E),`$
$`\sigma _S{\displaystyle \frac{d\gamma }{dz}}(z_i^S,E)`$ $`=`$ $`\sigma _F{\displaystyle \frac{d\gamma }{dz}}(z_i^F,E),`$ (11)
where $`\sigma _S`$ and $`\sigma _F`$ refer to the conductivities of S and F, respectively. For simplicity, we have assumed $`\sigma _S=\sigma _F`$, implying the continuity of the derivative at the interface. With the boundary conditions (10) and (Superconducting Proximity Effect through a Magnetic Domain Wall), we have solved Eq. (9) numerically by an iterative procedure (relaxation method) in the entire SFS system.
SFS system with domain wall: We apply the outlined theory to study the SFS structure of Fig. 1 in equilibrium. Lengths are given in units of the superconducting coherence length, $`\xi =\sqrt{D/\mathrm{\Delta }_0}`$. The spin-singlet superconductors are chosen to have the same gap magnitude. The contact areas at the SF interfaces are assumed to be small enough, so that any gap suppression can be neglected. The two superconducting regions and the intermediate ferromagnet are taken to have fixed lengths of $`d_S=5\xi `$ and $`d_F=2\xi `$. We model the domain wall by a varying direction of $`𝐉=(J_x,J_y,J_z)`$ (keeping the magnitude $`J=|𝐉|`$ constant), with $`J_y=0`$ and
$$\left(\begin{array}{c}J_x\\ J_z\end{array}\right)=J\left(\begin{array}{c}\mathrm{cos}\theta (z)\\ \mathrm{sin}\theta (z)\end{array}\right),\theta (z)=\mathrm{arctan}\frac{zz_0}{d_W}.$$
(12)
Here, $`d_W`$ is an effective domain wall width parameter. In the following we study the influence of the width $`d_W`$ of a domain wall centered in F ($`z_0=d_F/2`$) on the density of states in the center of the domain wall ($`z=z_0`$).
Knowing $`\gamma (z,E)`$ from the solution of Eq. (9) with boundary conditions (10) and (Superconducting Proximity Effect through a Magnetic Domain Wall), the quasiclassical Green function and the (total) LDOS
$$N_{tot}(z,E)=\frac{N_0}{2\pi }\mathrm{Im}\mathrm{Tr}𝒢(z,E),$$
(13)
is determined via Eqs. (7)-(8); $`N_0`$ is the normal-state density of states, and Tr denotes the spin trace.
An important characteristic of the value of the LDOS in superconductor-normal metal proximity systems is the minigap: the density of states in the normal-metal region shows a gap of width $`E_g<\mathrm{\Delta }_0`$ induced by proximity to a superconductor. The energy $`E_g`$ can be thought of as that of the lowest-energy Andreev bound state in a finite normal-metal layer. This convenient physical picture can easily be extended to single-domain ferromagnets. The corresponding spin-dependent energy shift of the quasiparticle and the Andreev-reflected quasihole by $`\pm J`$ leads to a reduction of the energy of the lowest-lying bound state, and correspondingly of the minigap, from the expression for a normal metal by $`J`$, vanishing altogether when $`JE_{g,J=0}`$. This picture is confirmed by our numerical calculations.
In the inhomogeneously magnetized case of Fig. 1, the above picture is modified. The effect of the domain wall on the LDOS is summarized in Fig. 2, which shows $`N_{tot}`$ as a function of energy for different domain-wall widths $`d_W`$. Although the value of $`J=0.5\mathrm{\Delta }_0`$ is here larger than the value of the normal state minigap $`E_g0.25\mathrm{\Delta }_0`$ (as seen from the dotted curve in Fig. 2 for $`J=0`$), the minigap is reduced to zero only for larger domain wall widths $`d_W2\xi `$. For the smallest width $`d_W=0.2\xi `$ the minigap is only reduced by about 40%. The additional states which fill the minigap with increasing domain wall width are due to spin triplet correlations, which are sensitive to the direction of $`𝐉`$. Our calculations show that the influence of equal-spin pairing components created by the domain wall increases. This is reflected by the appearance of additional Andreev bound states inside the gap, modifying the LDOS. The relative importance of the triplet correlations depends on $`J`$ and $`d_W`$: as clearly manifested by Fig. 2, the efficiency of the triplet-inducing mechanism grows with increasing $`d_W`$. The inset of Fig. 2 shows the value of the LDOS at $`E=0`$ as a function of $`d_W`$. The interesting observation here is that the LDOS at the chemical potential is very sensitive to the domain wall width when the latter one is comparable to $`\xi `$.
Finally, with a view towards studying the possible effects of the domain walls on the supercurrent flowing in an SFS structure, we have studied the LDOS in the case where there is a phase difference $`\varphi `$ between the two superconductors. This phase difference adds to the one accumulated by the quasiparticles and quasiholes in the ferromagnetic region and, thus, modifies the spectrum of Andreev bound states. Figures 3 and 4 present the LDOS in the middle of the F region for three domain walls with different widths. In the inset of Fig. 3 we also show the zero-energy LDOS for a domain wall of width $`d_W=0.2\xi `$ as a function of the phase difference. As can be seen in Fig. 4, by a possible tuning of the domain wall width $`d_W`$ one can always find a region of strongest sensitivity for a given phase difference $`\varphi `$ and vice versa. This increases the possibilities of controlling the zero energy density of states in the domain wall. The rich structure exhibited by these results could easily result in highly nontrivial behavior of the transport current both in equilibrium and nonequilibrium situations.
Conclusions: We have studied numerically the LDOS in a heterostructure consisting of a ferromagnetic alloy sandwiched between two singlet superconductors. We find strong modifications of the LDOS caused by the presence of a domain wall. As only triplet superconducting correlations are sensitive to the direction of the exchange field, the strong variations in the LDOS result from the presence of triplet correlations induced by the spatially varying magnetization. We also find a strong dependence of the density of states in the domain wall on a possible phase difference between the superconducting order parameters, giving an additional tool to control its value. This motivates future studies of the interplay of a supercurrent and the domain wall (Josephson effect). We hope that the variety of features observed in our calculations motivates further experimental research on proximity systems involving weak ferromagnets.
This work was supported by Deutsche Forschungsgemeinschaft within the Center for Functional Nanostructures. |
warning/0507/hep-ph0507097.html | ar5iv | text | # The clash of symmetries in a Randall-Sundrum-like spacetime.
## I Introduction
The breaking of symmetries is one of the most important issues in particle physics. Nature has taught us the relevance of symmetry principles, especially through Noether’s connection of global symmetries with conservation laws and the successful description of fundamental interactions via local gauge invariance. But to reconcile symmetric field equations with the real world, one needs that some of the solutions to those equations have a reduced symmetry compared to the equations themselves. This is, of course, what we mean by ‘spontaneous symmetry breaking (SSB)’ or ‘hidden symmetry’.
The precise mechanism by which electroweak SSB happens has yet to be experimentally determined. In the standard model (SM), the default option is to use an elementary scalar or Higgs field doublet with a quartic potential that induces a nonzero vacuum expectation value for one of the components. This will soon be tested at the Large Hadron Collider through its physical-Higgs boson search programme.
The purpose of this paper is to further develop a different method of symmetry breaking called the ‘clash of symmetries’ clash ; so10clash ; gaugeclash ; triniclash (see also pv ). This mechanism is relevant for brane-world extensions to the regular $`3+1`$ dimensional standard model rs ; otherbrane . The important development we report here is the rigorous incorporation of gravity into one of the toy Higgs models explored in the previous papers on this topic gaugeclash . Since gravity is fundamental to brane-world models, this is a very welcome step forward.
The toy model we shall analyse consists of two complex scalar fields $`\mathrm{\Phi }_{1,2}`$ coupled to $`4+1`$ dimensional gravity endowed with a bulk cosmological constant. The action is constructed to be invariant under global U(1)$``$U(1) transformations and the discrete $`\mathrm{\Phi }_1\mathrm{\Phi }_2`$ interchange operation. The Higgs potential we construct is sextic, and its specific form is driven by our desire to find an analytic solution to the coupled Einstein-Higgs field equations. The spacetime metric tensor we posit has the Randall-Sundrum non-factorisable form rs , but with the warp factor a completely smooth function of the extra dimension coordinate. The precise Randall-Sundrum metric emerges by taking a certain limit. The Higgs field configurations we seek are kink-like, associated with the spontaneously broken discrete symmetry.
We also study the gauged version of this model, extending the analysis of Ref.gaugeclash . We find additional non-trivial flat-space solutions, but despite an extensive numerical search we have been unable to find any warped metric solutions with nonzero gauge fields.
Section II reviews the clash of symmetries mechanism, while Sec. III defines and analyses the ungauged version of our toy model. This is followed by an account of the gauged cases in Sec. IV. The final section is a conclusion and outlook for future work.
## II The clash of symmetries
The default SSB mechanism employed in the SM and many extensions to it utilises a simple spatially homogeneous, static and stable solution to the field equations: the Higgs fields $`\varphi _i(x^\mu )`$ are set equal to those constant values corresponding to an absolute minimum of the Higgs potential. The locations of the nonzero constants within the Higgs multiplets then determine the stability group of the vacuum state, what we commonly call the unbroken symmetry after SSB. The nonzero Higgs fields become constant background scalar fields throughout the universe.
There are other types of solutions to Higgs field equations that can also serve as stable, static background fields: topological solitons, such as kinks, strings and monopoles. While there is much interest in such configurations, they are not usually considered as viable candidates for playing a major role in spontaneous symmetry breaking. This is simply because they are spatially inhomogeneous. The energy densities required for electroweak and higher symmetry breaking are so high as to be incompatible with the strong evidence for a universe that is spatially homogeneous at large scales. This objection, however, does not apply to brane-world models, because the non-trivial spatial dependence can be restricted to the extra dimension coordinates only.
For definiteness, consider the case of one extra dimension described by Cartesian coordinate $`w`$ and topologically-stable Higgs configurations $`\varphi _i(w)`$, some of which have kink-form with respect to $`w`$. The pattern of spontaneous breaking then effectively becomes a function of $`w`$. The $`3+1`$-brane is located at, say, $`w=0`$. If the degrees of freedom that are confined to the brane local are confined absolutely, then the unbroken symmetry is the stability group of $`\varphi _i(w=0)`$, to be denoted by $`H(w=0)`$ from now on (and let the internal symmetry group of the action be $`G`$). The symmetry breaking will be communicated to the brane-world denizens through appropriate interactions. But in a quantal (and perhaps even in a classical) world, one would not expect strict confinement to the brane: there should be some leakage off it into the bulk. One therefore expects the brane-localised fields to also couple, but with reduced strength, to $`\varphi _i(w)`$ states where $`w`$ is slightly different from zero. If the stability group of $`\varphi _i(0<|w|<ϵ)`$, call it $`H(wϵ)`$, is different from $`H(w=0)`$, then a rather rich effective symmetry breaking outcome could ensue for the brane world. For example, a hierarchical breaking pattern $`GH(w=0)H(wϵ)`$ was speculated upon in Ref.so10clash , where in that case $`G=`$SO(10), $`H(w=0)=`$SU(4)$``$SU(2)$``$U(1) and $`H(wϵ)=`$SU(3)$``$SU(2)$``$U(1)<sup>2</sup>. The potential application to grand unified models and the gauge hierarchy puzzle is clear.
The ‘clash of symmetries’ phenomenon allows the unbroken symmetry at non-asymptotic points in the extra dimension, $`|w|<\mathrm{}`$, to be smaller than the symmetry holding at $`|w|=\mathrm{}`$ clash ; so10clash ; gaugeclash ; triniclash . The boundary conditions at infinity for kink configurations are set to be the Higgs vacuum expectation values (VEVs). They force the spontaneous breaking $`GH(|w|\mathrm{})`$. In usual SSB, the Higgs VEV configuration is used throughout the whole of space, so the symmetry is broken to the group we are calling $`H(|w|\mathrm{})`$ everywhere, homogeneously. In the clash of symmetries mechanism, this symmetry breaking pattern only holds asymptotically with respect to the extra dimension(s).
The clash of symmetries can occur if the isomorphic subgroups left unbroken at $`w=\mathrm{}`$ and $`w=+\mathrm{}`$, $`H(w\mathrm{})H(\mathrm{})`$ and $`H(w+\mathrm{})H(+\mathrm{})`$ respectively, can be differently embedded within the parent group $`G`$. In the cases examined so far in the literature, the breakdown at finite values of $`w`$ is to the intersection of the asymptotic stability groups,
$$H(|w|<\mathrm{})=H(\mathrm{})H(+\mathrm{})H_{\mathrm{clash}}.$$
(1)
Although $`H(\mathrm{})`$ and $`H(+\mathrm{})`$ are isomorphic, their different embeddings within $`G`$ leads their intersection to be a smaller group.<sup>1</sup><sup>1</sup>1There will also typically be kink configurations that interpolate between identically embedded subgroups. Such kinks do not display the clash of symmetries, and one must ensure that they have higher energy than those displaying the clash. Some special values of $`w`$, most often $`w=0`$, may correspond to stability groups larger than $`H_{\mathrm{clash}}`$ because the spatially varying Higgs configurations may instantaneously pass through special patterns as $`w`$ varies. For instance, a component of a Higgs multiplet may be an odd function of $`w`$ such as $`\mathrm{tanh}(w)`$, so it would vanish at $`w=0`$ and thus the unbroken symmetry would be enhanced at that point. This possibility lies behind our remarks two paragraphs above about an application to the gauge hierarchy problem.
Discrete symmetries also play a vital role: their spontaneous breakdown leads the vacuum manifold to have disconnected pieces. Kinks interpolate between vacuum states from disconnected pieces. Within a given topological class of kinks, the one with lowest energy is then guaranteed to be (topologically) stable.
The easiest example of the clash of symmetries uses $`G=`$SU(3) spontaneously breaking to SU(2). An SU(2) subgroup can be embedded in three different ways in SU(3): $`I`$-spin, $`U`$-spin and $`V`$-spin. Clash of symmetries kinks interpolate between VEVs respecting $`(I,U)`$, $`(U,V)`$ or $`(V,I)`$ at $`w=(\mathrm{},+\mathrm{})`$. The differently embedded asymptotic SU(2) symmetries “clash” at non-asymptotic points, leading $`H_{\mathrm{clash}}`$ to be SU(2)$`{}_{I}{}^{}`$SU(2)$`{}_{U}{}^{}=\{1\}`$ and so on.<sup>2</sup><sup>2</sup>2To actually implement this simple SU(3) scheme requires more work: see Ref.clash for details.
The model-building promise of the clash of symmetries lies in the greater symmetry breaking power of spatially inhomogeneous Higgs configurations over the simple homogeneous alternative of standard SSB. We hope that this will lead to models with simpler Higgs sectors and fewer associated parameters. Some of us have already speculated that the ultimate application might be to an $`E_6`$ model, because of the fact that there are three different $`E_6`$ generators that can serve as electric charge, corresponding to three different embeddings of electromagnetism within the full group clash . Standard $`E_6`$ models can only use one at a time. Even more speculatively, a connection between that threefold structure and the three families of quarks and leptons might exist<sup>3</sup><sup>3</sup>3There is some similarity between the clash of symmetries mechanism and orbifold symmetry breaking orbifold . In particular, Ref. buchmuller constructs a model where different subgroups of SO(10) are left unbroken on different branes, with the overall unrboken symmetry being the intersection of those individual subgroups, and with families also distributed amongst the branes. While there are similarities, it is clear that the mechanisms are not identical, because the clash of symmetries has a physical Higgs field while the orbifold method is Higgsless. clash .
But many preparatory studies are needed before such an ambitious vision can be realised. In this paper, we take another important step along that road: the incorporation of gravity. To facilitate the study of gravity, we begin by simplifying the group theoretic structure as much as possible without losing the essence of the clash of symmetries mechanism: we take $`G=`$U(1)$`{}_{1}{}^{}`$U(1)<sub>2</sub>, with $`H(\mathrm{})=`$U(1)<sub>1</sub> and $`H(+\mathrm{})=`$U(1)<sub>2</sub> leading to $`H_{\mathrm{clash}}=\{1\}`$. The discrete symmetry we need for topological stability interchanges the two Abelian sectors.
In the next section we show that the ungauged version of this model can yield a clash-of-symmetries-style kink in a warped metric spacetime. The section after that revisits the gauged extension in both flat and curved spacetime.
## III The ungauged model
### III.1 The field equations
The model has two complex Higgs fields $`\mathrm{\Phi }_1`$ and $`\mathrm{\Phi }_2`$ coupling to $`4+1`$ dimensional gravity, and it respects the $`\mathrm{\Phi }_1\mathrm{\Phi }_2`$ discrete symmetry. The action is
$$S=\left[\frac{\kappa }{2}R+g^{MN}\left(_M\overline{\mathrm{\Phi }}_1_N\mathrm{\Phi }_1+_M\overline{\mathrm{\Phi }}_2_N\mathrm{\Phi }_2\right)+V(\overline{\mathrm{\Phi }}_1\mathrm{\Phi }_1,\overline{\mathrm{\Phi }}_2\mathrm{\Phi }_2)\right]\sqrt{g}d^5x$$
(2)
where $`g_{MN}`$ is the $`4+1`$ dimensional metric tensor, with $`M,N=0,1,2,3,5`$, $`R`$ is the curvature scalar, and overbars denote complex conjugation. The matter action $`S_m`$ is defined as $`S`$ with the Einstein-Hilbert contribution absent. The Higgs potential $`V`$ is as yet unspecified except that it must be invariant under $`\mathrm{\Phi }_1\mathrm{\Phi }_2`$. It will have a constant piece that serves as the bulk cosmological constant. We employ the $`(,+,+,+,+)`$ signature and our sign conventions are as in Weinberg (apart from our definition of $`g\mathrm{Det}(g)`$) weinberg .
The Einstein equation found by requiring a stationary action under metric variations is
$$G^{MN}=\frac{1}{\kappa }T^{MN}$$
(3)
where the Einstein tensor is, as usual, $`G^{MN}=R^{MN}(1/2)g^{MN}R`$ with $`R^{MN}`$ being the Ricci tensor. The stress-energy tensor,
$$T^{MN}\frac{2}{\sqrt{g}}\frac{\delta S_m}{\delta g_{MN}},$$
(4)
evaluates to
$$T^{MN}=2g^{MP}g^{NQ}t_{PQ}g^{MN}\left(g^{PQ}t_{PQ}+V\right),$$
(5)
where
$$t_{MN}_M\overline{\mathrm{\Phi }}_1_N\mathrm{\Phi }_1+_M\overline{\mathrm{\Phi }}_2_N\mathrm{\Phi }_2.$$
(6)
The Higgs Euler-Lagrange equations are
$$_M\left(\sqrt{g}g^{MN}_N\mathrm{\Phi }_i\right)=\sqrt{g}\frac{V}{\overline{\mathrm{\Phi }}_i},$$
(7)
where $`i=1,2`$, and the complex conjugate equations also hold.
We now specify our solution ansatz. We seek configurations that only depend on the extra dimension $`x^5w`$. The imaginary parts of the Higgs fields are set to zero and we write
$$\mathrm{\Phi }_i=\frac{\varphi _i(w)}{\sqrt{2}},$$
(8)
where $`\varphi _i(w)`$ is real. The metric ansatz, of Randall-Sundrum type, is defined through the $`4+1`$ dimensional line element
$$ds_5^2=dw^2+e^{2f(w)}ds_4^2$$
(9)
where $`f(w)`$ is an as yet unknown function and $`ds_4^2=\eta _{\mu \nu }dx^\mu dx^\nu `$ is the $`3+1`$ dimensional Minkowski line element (our brane will be Minkowski), with $`\mu ,\nu =0,1,2,3`$ or $`t,x,y,z`$. The field equations then reduce to
$`3f^{\prime \prime }+6(f^{})^2`$ $`=`$ $`{\displaystyle \frac{1}{\kappa }}\left[{\displaystyle \frac{1}{2}}\left[(\varphi _{1}^{}{}_{}{}^{})^2+(\varphi _{2}^{}{}_{}{}^{})^2\right]+V\right],`$ (10)
$`6(f^{})^2`$ $`=`$ $`{\displaystyle \frac{1}{\kappa }}\left[{\displaystyle \frac{1}{2}}\left[(\varphi _{1}^{}{}_{}{}^{})^2+(\varphi _{2}^{}{}_{}{}^{})^2\right]V\right],`$ (11)
$`\varphi _{i}^{}{}_{}{}^{\prime \prime }+4f^{}\varphi _{i}^{}{}_{}{}^{}`$ $`=`$ $`{\displaystyle \frac{V}{\varphi _i}},`$ (12)
where prime denotes differentiation with respect to $`w`$. Adding and subtracting Eqs. (10) and (11) gives
$`f^{\prime \prime }+4(f^{})^2`$ $`=`$ $`{\displaystyle \frac{2}{3\kappa }}V,`$ (13)
$`f^{\prime \prime }`$ $`=`$ $`{\displaystyle \frac{1}{3\kappa }}\left[(\varphi _{1}^{}{}_{}{}^{})^2+(\varphi _{2}^{}{}_{}{}^{})^2\right].`$ (14)
Note that Eq. (14) is independent of the potential, a feature we will use below.
We will construct a potential $`V`$ that has degenerate global minima at
$$\mathrm{vacuum}1:\varphi _1=u,\varphi _2=0\mathrm{and}\mathrm{vacuum}2:\varphi _1=0,\varphi _2=u,$$
(15)
where $`u`$ would be the VEV in the standard homogeneous SSB mechanism. The vacua are related by the now spontaneously broken discrete interchange symmetry. The Higgs field boundary conditions will be
$$\varphi _i(\mathrm{})=\mathrm{vacuum}2\mathrm{and}\varphi _i(+\mathrm{})=\mathrm{vacuum}1.$$
(16)
### III.2 Constructing the solution
In $`3+1`$ dimensions without gravity, renormalisability limits us to quartic potentials. With extra dimensions, and with gravity included, this is no longer a well-motivated requirement. In fact, the only general constraint on the form of $`V`$ we will impose is that it be bounded from below. A very important topic for the future of the clash of symmetries mechanism will indeed be the rules by which potentials should be specified.
Our more modest goal here is to provide an existence proof that kinks featuring the clash of symmetries form can be consistently coupled to Randall-Sundrum-like gravity. We will search for an analytically tractable model for reasons of both simplicity and elegance. To do this, it makes sense to invert the problem by first specifying the solutions $`\varphi _i(w)`$ and $`f(w)`$ and then constructing a potential that gives those solutions. This, in itself, is non-trivial and interesting.
It turns out that a good choice, satisfying the boundary conditions of Eq. (16) and also the Higgs-potential-independent Eq. (14), is
$`\varphi _1^s(w)`$ $`=`$ $`{\displaystyle \frac{u}{\sqrt{2}}}\sqrt{1+\mathrm{tanh}\beta w},`$ (17)
$`\varphi _2^s(w)`$ $`=`$ $`{\displaystyle \frac{u}{\sqrt{2}}}\sqrt{1\mathrm{tanh}\beta w},`$ (18)
$`f^s(w)`$ $`=`$ $`{\displaystyle \frac{u^2}{12\kappa }}\mathrm{ln}\left(\mathrm{cosh}\beta w\right),`$ (19)
where $`\beta `$ is a free parameter that specifies the brane thickness, and the superscript $`s`$ denotes “solution”. The Higgs configurations are the square roots of the most naïve choices, and $`f^s(w)`$ is essentially a regularisation of $`|w|`$. The latter is, of course, the corresponding function in the Randall-Sundrum scenario. The smoothing out of the Randall-Sundrum cusp goes hand-in-hand with the existence of Higgs field kinks coupled to gravity, and the brane is dynamically generated rather than put in by hand. Note also that
$$\varphi _1^s(w)^2+\varphi _2^s(w)^2=u^2,$$
(20)
a feature that will turn out to be very useful.
The first inkling for what $`V`$ should be comes from Eq. (13). It suggests that
$$V=\frac{\beta ^2u^4}{24\kappa }+\frac{\beta ^2}{2u^2}\left(1+\frac{u^2}{3\kappa }\right)\varphi _1^2\varphi _2^2+U(\varphi _1^2,\varphi _2^2),$$
(21)
where $`U`$ is a function that vanishes at the level of the solution, that is $`U(\varphi _1^s(w)^2,\varphi _2^s(w)^2)=0`$. Notice that a negative bulk cosmological constant,
$$\mathrm{\Lambda }_5=\frac{\beta ^2u^4}{24\kappa },$$
(22)
follows from the ansatz: the bulk has anti-de Sitter qualities, as usual in Randall-Sundrum-like models, though in our case it is of course not precisely anti-de Sitter.
More information about $`U`$ is supplied by the Higgs field equation (12). To solve it, one must specify that
$$U=\frac{\beta ^2}{u^4}\left(\frac{3}{2}+\frac{u^2}{3\kappa }\right)\varphi _1^2\varphi _2^2(\varphi _1^2+\varphi _2^2u^2)+W(\varphi _1^2+\varphi _2^2u^2),$$
(23)
where the new function $`W`$ obeys
$`W(0)`$ $`=`$ $`0,`$ (24)
$`{\displaystyle \frac{W}{\varphi _{1,2}}}|_{\varphi _1^s,\varphi _2^s}`$ $`=`$ $`0,`$ (25)
so as not to affect the satisfaction of the field equations, Eq. (13) and Eq. (14) respectively. As far as they are concerned, $`W`$ can vanish. The resulting sextic potential, however, would not be bounded from below because the coefficient of the first term in $`U`$ is negative.
We therefore construct $`W`$ so that, as well as obeying Eqs. (24) and (25), it leads to a potential that is both bounded from below and has global minima at Eq. (15). The simplest suitable function is
$$W=\zeta \frac{\beta ^2}{4u^2}\left(\frac{3}{2}+\frac{u^2}{3\kappa }\right)(\varphi _1^2+\varphi _2^2u^2)^2\left(\eta +\frac{\varphi _1^2+\varphi _2^2u^2}{u^2}\right)$$
(26)
where $`\eta `$ and $`\zeta `$ are dimensionless parameters.
The full Higgs potential for the model is specified by Eqs. (21), (23) and (26). We will present a global minimisation analysis of it in the next subsection. (To reinstate the complex nature of the Higgs fields, the substitution $`\varphi _i^22\overline{\mathrm{\Phi }}_i\mathrm{\Phi }_i`$ should be made, though for the forthcoming discussion we need only work with the $`\varphi _i`$.)
Before doing so, let us identify the Randall-Sundrum limit of the model. We need to identify a limit in which the warp factor, $`e^{2f^s}`$, becomes $`e^{a|w|}`$ where $`a`$ is some positive constant. Following Ref.davidsonmannheim , we note that $`(\mathrm{cosh}\beta w)^{(1/\beta )}e^{|w|}`$ as $`\beta \mathrm{}`$. Writing
$$e^{2f^s(w)}=\left[(\mathrm{cosh}\beta w)^{\frac{1}{\beta }}\right]^{\frac{u^2\beta }{6\kappa }},$$
(27)
we see that the Randall-Sundrum warp factor, $`e^{\frac{u^2\beta }{6\kappa }|w|}`$, is obtained in the limit
$$\beta \mathrm{},u0\mathrm{such}\mathrm{that}u^2\beta \mathrm{finite}\mathrm{constant}.$$
(28)
We also require $`\eta `$ and $`\zeta `$ to be unchanged when taking the limit. Note that the Higgs field configurations, Eqs. (17) and (18), become step functions as $`\beta \mathrm{}`$, with their magnitude tending to zero because $`u0`$. The Randall-Sundrum fine-tuning condition between the brane tension and the bulk cosmological constant is a feature of our solution. Defining the brane tension via
$$\mathrm{\Lambda }_b\underset{RS}{lim}_{\mathrm{}}^+\mathrm{}\left[\frac{1}{2}\left[(\varphi _{1}^{s}{}_{}{}^{})^2+(\varphi _{2}^{s}{}_{}{}^{})^2\right]+V(\varphi _1^s,\varphi _2^s)\right]\sqrt{g^s}𝑑w,$$
(29)
it is easy to show that
$$\mathrm{\Lambda }_b^2=\frac{3}{2}\kappa \mathrm{\Lambda }_5,$$
(30)
which is the Randall-Sundrum condition.
### III.3 Minimisation analysis
We now identify the parameter space region where our potential has global minima specified by Eq. (15).
To establish the parameter region where the potential is bounded from below, we look at the sixth-order terms in isolation. They are
$$V_6=\lambda (\varphi _1^2+\varphi _2^2)\left[\zeta (\varphi _1^2+\varphi _2^2)^24\varphi _1^2\varphi _2^2\right],$$
(31)
where $`\lambda (3\beta ^2/8u^4)[1+2u^2/(9\kappa )]>0`$. The substitutions $`\varphi _1=\chi \mathrm{cos}\theta `$ and $`\varphi _2=\chi \mathrm{sin}\theta `$ reveal that $`V_6=\lambda \chi ^6(\zeta 4\mathrm{sin}^2\theta \mathrm{cos}^2\theta )`$, from which the condition
$$\zeta >1$$
(32)
is trivially seen to be the only requirement to produce boundedness from below.
The minimisation analysis is a four-parameter problem. For some of the local minima, both their $`(\varphi _1,\varphi _2)`$ locations and the values of the potential at those points are complicated functions of the parameters. But there is a smart way to evade these unenlightening complications while yielding physically meaningful results: approach the Randall-Sundrum limit defined by Eq. (28). Our task then becomes a two parameter problem in $`(\eta ,\zeta )`$ that is readily visualised.
To see this, note that in Eqs. (21), (23) and (26), the combinations
$$1+\frac{u^2}{3\kappa }\mathrm{and}\frac{3}{2}+\frac{u^2}{3\kappa }$$
(33)
frequently appear. In the Randall-Sundrum limit, which includes $`u0`$, these combinations simplify to $`1`$ and $`3/2`$ respectively. In the exact Randall-Sundrum limit, the meaning to be ascribed to the potential is obscure, so we cannot work exactly in that limit. Instead, we will approach the limit without reaching it, and write an approximation to the potential justified by the above remarks. It is
$`VV_{\mathrm{approx}}`$ $`=`$ $`{\displaystyle \frac{\beta ^2u^4}{24\kappa }}+{\displaystyle \frac{\beta ^2}{2u^2}}\varphi _1^2\varphi _2^2{\displaystyle \frac{3\beta ^2}{2u^4}}\varphi _1^2\varphi _2^2(\varphi _1^2+\varphi _2^2u^2)+\zeta {\displaystyle \frac{3\beta ^2}{8u^2}}(\varphi _1^2+\varphi _2^2u^2)^2\left(\eta +{\displaystyle \frac{\varphi _1^2+\varphi _2^2u^2}{u^2}}\right),`$ (34)
$`=`$ $`{\displaystyle \frac{3\beta ^2u^2}{8}}\left[{\displaystyle \frac{u^2}{9\kappa }}+{\displaystyle \frac{4}{3}}p_1^2p_2^24p_1^2p_2^2(p_1^2+p_2^21)+\zeta (p_1^2+p_2^21)^2\left(\eta +p_1^2+p_2^21\right)\right],`$
where
$$p_i\frac{\varphi _i}{u}$$
(35)
is a dimensionless Higgs field. The fact that $`\beta `$ contributes only to an overall scale, and that $`\kappa `$ appears solely in the constant term, reduces the analysis to a two-parameter problem in $`(\eta ,\zeta )`$. While $`\kappa `$ contributes to the potential values, it does so equally for all local minima in this approximation, which means that only $`\eta `$ and $`\zeta `$ need be considered when comparing the local minima to see which are also the global minima.
The potential has three varieties, denoted $`(I,II,III)`$, of possible local minima (or saddle points),
$`(I)`$ $`p_1^2=1,p_2=0\mathrm{degenerate}\mathrm{with}p_1=0,p_2^2=1(\mathrm{discrete}\mathrm{symmetry}\mathrm{breaking}),`$ (36)
$`(II)`$ $`p_1^2=p_2^20(\mathrm{discrete}\mathrm{symmetry}\mathrm{preserving}),`$ (37)
$`(III)`$ $`p_1=p_2=0(\mathrm{no}\mathrm{SSB}).`$ (38)
Which of these are actual local minima depends on $`(\eta ,\zeta )`$. We want $`I`$ to describe the global minima.
Of the three, it is $`II`$ that has the complicated parameter dependence. In fact it is quite easy to show, even for the full potential of Eqs. (21), (23) and (26), that $`V(I)<V(III)`$ requires that
$$\eta >1$$
(39)
which is thus a necessary but not sufficient condition to ensure that $`I`$ gives the global minima. Sufficiency requires an examination of the parameter dependence of $`V(II)`$ also. The results are shown in Fig. 1, where the shaded area is the Higgs parameter region where $`I`$ is the global minimum.
In the red (or medium-shaded if viewing in black-and-white) region, $`I`$ is a local minimum and $`II`$ a saddle point with $`V(I)<V(II)`$, while $`III`$ is a local maximum. As $`\eta `$ decreases below $`3/2`$, $`III`$ develops into a local minimum also; the ordering is $`V(I)<V(II)<V(III)`$ in the green (light-shaded) region and $`V(I)<V(III)<V(II)`$ in the blue (dark-shaded). For $`\eta <1`$, $`III`$ is the global minimum so there is no SSB. The border line between the shaded and unshaded areas is given by
$$\zeta =\frac{2718\eta +61\eta ^2+\sqrt{729+972\eta 2970\eta ^2+1900\eta ^3375\eta ^4}}{24\eta ^3}.$$
(40)
Lowering $`\zeta `$ has the effect of lowering $`V(II)`$. Below this border line, $`V(II)<V(I)<V(III)`$, so the symmetry preserving vacuum is the global minimum.
The border between the blue (dark-shaded) and green (light-shaded) regions is given by
$$\eta =\frac{8\zeta 3\zeta ^2+2\sqrt{3}\sqrt{4\zeta ^27\zeta ^3+3\zeta ^4}}{3\zeta ^2}.$$
(41)
Figures 2-4 show contour plots of the approximate potential for representative points in the red (medium-shaded), green (light-shaded) and blue (dark-shaded) areas of Fig. 1, respectively. The colour- or shade-coding is explained in the captions. The plots illustrate the pattern of extrema described above.
## IV The gauged model
It is natural to extend our toy model by gauging U(1)$``$U(1), and to study the dependence of the gauge fields on $`w`$ by self-consistently solving all of the coupled Euler-Lagrange equations. In the flat space case, Rozowsky, Volkas and Wali gaugeclash found some interesting solutions where the gauge field coupled to $`\mathrm{\Phi }_i`$ tended towards a linear function of $`w`$ on the side of the wall where $`\mathrm{\Phi }_i0`$, and was exponentially suppressed on the other side (Meissner effect). The linearly rising gauge field implies an asymptotically constant magnetic field, so the physical picture was of two infinite, uniform sheets (one for each sector) of supercurrent density flowing along the domain wall. In this section we will first generalise the Rozowsky et al. flat-space analysis and then consider the curved space problem.
### IV.1 Flat space
The Lagrangian is given by
$$=\frac{1}{4}F_1^{MN}F_{1MN}\frac{1}{4}F_2^{MN}F_{2MN}(D^M\mathrm{\Phi }_1)^{}(D_M\mathrm{\Phi }_1)(D^M\mathrm{\Phi }_2)^{}(D_M\mathrm{\Phi }_2)V,$$
(42)
where
$`D_M\mathrm{\Phi }_1`$ $`=`$ $`(_MieA_{1M}i\stackrel{~}{e}A_{2M})\mathrm{\Phi }_1,`$
$`D_M\mathrm{\Phi }_2`$ $`=`$ $`(_Mi\stackrel{~}{e}A_{1M}ieA_{2M})\mathrm{\Phi }_2,`$ (43)
and we take $`V`$ to be quartic
$$V=\lambda _1(\overline{\mathrm{\Phi }}_1\mathrm{\Phi }_1+\overline{\mathrm{\Phi }}_2\mathrm{\Phi }_2u^2)^2+\lambda _2\overline{\mathrm{\Phi }}_1\mathrm{\Phi }_1\overline{\mathrm{\Phi }}_2\mathrm{\Phi }_2.$$
(44)
In the $`\lambda _{1,2}>0`$ parameter space region, the global minima of $`V`$ are
$$\mathrm{Vacuum}1:\overline{\mathrm{\Phi }}_1\mathrm{\Phi }_1=u^2,\mathrm{\Phi }_2=0\mathrm{and}\mathrm{Vacuum}2:\mathrm{\Phi }_1=0,\overline{\mathrm{\Phi }}_2\mathrm{\Phi }_2=u^2.$$
(45)
The Lagrangian is invariant under the discrete symmetry
$$\mathrm{\Phi }_1\mathrm{\Phi }_2,A_1A_2.$$
(46)
This model generalises that considered in Ref.gaugeclash by having $`\stackrel{~}{e}0`$.
The Euler-Lagrange equations are
$`D^MD_M\mathrm{\Phi }_i`$ $`=`$ $`2\lambda _1\mathrm{\Phi }_i(\overline{\mathrm{\Phi }}_i\mathrm{\Phi }_i+\overline{\mathrm{\Phi }}_j\mathrm{\Phi }_ju^2)+\lambda _2\mathrm{\Phi }_i\overline{\mathrm{\Phi }}_j\mathrm{\Phi }_j,`$
$`_MF_i^{MN}`$ $`=`$ $`2\mathrm{Im}(e\overline{\mathrm{\Phi }}_iD^N\mathrm{\Phi }_i+\stackrel{~}{e}\overline{\mathrm{\Phi }}_jD^N\mathrm{\Phi }_j),`$ (47)
where $`(i,j)=(1,2)`$ or $`(2,1)`$. We now specialise to configurations that depend only on $`w`$, impose Lorentz gauge $`_MA_i^M=0`$, and write $`\mathrm{\Phi }_i=R_ie^{i\theta _i}`$. The equations become,
$`R_i^{\prime \prime }`$ $`=`$ $`R_i\left[(eA_i^\mu +\stackrel{~}{e}A_j^\mu )(eA_{i\mu }+\stackrel{~}{e}A_{j\mu })\right]+2\lambda _1R_i(R_i^2+R_j^2u^2)+\lambda _2R_iR_j^2,`$ (48)
$`A_{i}^{\mu }{}_{}{}^{\prime \prime }`$ $`=`$ $`2eR_i^2(eA_i^\mu +\stackrel{~}{e}A_j^\mu )+2\stackrel{~}{e}R_j^2(eA_j^\mu +\stackrel{~}{e}A_i^\mu ),`$ (49)
$`\theta _i^{}`$ $`=`$ $`(eA_i^5+\stackrel{~}{e}A_j^5).`$ (50)
Observe from the last of these equations that the $`A=5`$ gauge field components are pure gauge, so they and the phase fields $`\theta _i`$ decouple from the problem.
We choose to look for solutions that do not break rotational invariance in the $`x`$, $`y`$ and $`z`$ directions. This means that we can always rotate the $`x,y,z`$ coordinate system so that the only nonzero component of $`(A_i^x,A_i^y,A_i^z)`$ is in, say, the $`x`$-direction (this is tantamount to the requirement that the ratios of these components are $`w`$-independent). We also choose that the same component is nonzero for both $`i=1`$ and $`i=2`$. Equations (48) and (49) now simplify to
$`R_i^{\prime \prime }`$ $`=`$ $`R_i\left[(eA_i^t+\stackrel{~}{e}A_j^t)^2+(eA_i^x+\stackrel{~}{e}A_j^x)^2\right]+2\lambda _1R_i(R_i^2+R_j^2u^2)+\lambda _2R_iR_j^2,`$ (51)
$`A_{i}^{t,x}{}_{}{}^{\prime \prime }`$ $`=`$ $`2eR_i^2(eA_i^{t,x}+\stackrel{~}{e}A_j^{t,x})+2\stackrel{~}{e}R_j^2(eA_j^{t,x}+\stackrel{~}{e}A_i^{t,x})`$ (52)
The final choice is whether to choose the gauge fields to be timelike $`(A_i^t0,A_i^x=0)`$, lightlike $`(A_i^t=A_i^x0)`$, or spacelike $`(A_i^t=0,A_i^x0)`$. As is obvious from Eq. (51), the timelike ansatz gives rise to asymptotic oscillatory behaviour for the $`R`$’s and is thus unacceptable. We first discuss the spacelike case.
As a boundary condition, we require that $`R_{1,2}`$ tend to vacua 1 and 2 on opposite sides of the wall:
$$R_1(+\mathrm{})=u,R_2(+\mathrm{})=0\mathrm{and}R_1(\mathrm{})=0,R_2(\mathrm{})=u.$$
(53)
The appropriate boundary conditions for the gauge fields are then determined by examining Eqs. (52) at the asymptotic points. As $`w+\mathrm{}`$, we find that
$$\left(\begin{array}{c}A_{1}^{}{}_{}{}^{\prime \prime }\\ A_{2}^{}{}_{}{}^{\prime \prime }\end{array}\right)=2u^2\left(\begin{array}{cc}e^2& e\stackrel{~}{e}\\ e\stackrel{~}{e}& \stackrel{~}{e}^2\end{array}\right)\left(\begin{array}{c}A_1\\ A_2\end{array}\right),$$
(54)
where $`A_iA_i^x`$. Inputting trial asymptotic solutions of the form
$$A_{1,2}=a_{1,2}e^{\kappa w},$$
(55)
we obtain an eigenvalue equation yielding
$$\kappa =\sqrt{2}u\sqrt{e^2+\stackrel{~}{e}^2}\mathrm{and}\kappa =0.$$
(56)
The zero eigenvalue means that one linear combination of $`A_1`$ and $`A_2`$ is always unsuppressed, and the vanishing second derivative tells us it grows linearly with $`w`$ (provided its amplitude is nonzero). By the symmetry of the problem, a similar situation obtains as $`w\mathrm{}`$. The exponentially suppressed linear combination embodies the Meissner effect. Numerical solutions are depicted in Fig. 5.
The lightlike case is qualitatively similar to the spacelike case. But it is amusing to point out that for the parameter choice $`\lambda _2=4\lambda _1`$ an analytic solution exists. For the lightlike ansatz, the gauge fields cancel out of Eq. 51, and with the stated parameter choice and $`e=1`$, $`\stackrel{~}{e}=0`$, the resulting equations yield the solutions
$$R_{1,2}(w)=\frac{u}{2}\left[1\pm \mathrm{tanh}(\sqrt{\lambda _1}uw)\right].$$
(57)
Inputting these results into Eq. 52, the analytic solutions obeying the correct boundary conditions for the gauge field functions $`A_i^t=A_i^xA_i`$ are
$$A_{1,2}(w)=\mathrm{const}.\times {}_{2}{}^{}F_{1}^{}(\frac{1+\sqrt{2}\sqrt{3}}{2},\frac{1+\sqrt{2}+\sqrt{3}}{2};\mathrm{\hspace{0.17em}1}+\sqrt{2};(1+e^{\pm 2\sqrt{\lambda _1}uw})^1)(1\mathrm{tanh}(\sqrt{\lambda _1}uw))^{\frac{1}{\sqrt{2}}},$$
(58)
where $`{}_{2}{}^{}F_{1}^{}`$ is a hypergeometric function.
### IV.2 Warped spacetime
We now generalise the above model through the inclusion of gravity gaugeggrav . The action is
$$S=\left[\frac{\kappa }{2}R+g^{MN}t_{MN}+\frac{1}{4}g^{MP}g^{NQ}f_{MNPQ}+V\right]\sqrt{g}\mathrm{d}^5x,$$
(59)
with
$`t_{MN}`$ $``$ $`{\displaystyle \underset{i}{}}\left(D_M\mathrm{\Phi }_i\right)^{}D_N\mathrm{\Phi }_i,`$ (60)
$`D_M`$ $``$ $`_M{\displaystyle \underset{i}{}}\mathrm{i}Q_iA_{iM},`$ (61)
$`f_{MNPQ}`$ $``$ $`{\displaystyle \underset{i}{}}F_{iMN}F_{iPQ},`$ (62)
and the charges $`(Q_1,Q_2)`$ of the scalar fields are $`(e,\stackrel{~}{e})`$ for $`\mathrm{\Phi }_1`$ and $`(\stackrel{~}{e},e)`$ for $`\mathrm{\Phi }_2`$.
The Higgs and gauge field equations of motion are
$`D_M\left(\sqrt{g}g^{MN}D_N\mathrm{\Phi }_i\right)\sqrt{g}{\displaystyle \frac{V}{\mathrm{\Phi }_i^{}}}`$ $`=`$ $`0`$ (63)
$`_M\left(\sqrt{g}g^{MP}g^{NQ}F_{iPQ}\right)+\sqrt{g}g^{NP}2\left(e\mathrm{Im}\left(\mathrm{\Phi }_i^{}D_P\mathrm{\Phi }_i\right)+\stackrel{~}{e}\mathrm{Im}\left(\mathrm{\Phi }_j^{}D_P\mathrm{\Phi }_j\right)\right)`$ $`=`$ $`0,`$ (64)
while the Einstein equations, $`\kappa G_{MN}=T_{MN}`$, feature
$`T_{MN}`$ $`=`$ $`2t_{MN}+g^{PQ}f_{MPNQ}+g_{MN}_M,`$ (65)
$`_M`$ $`=`$ $`g^{MN}t_{MN}\frac{1}{4}g^{MP}g^{NQ}f_{MNPQ}V.`$ (66)
The metric ansatz is generalised to
$$\mathrm{d}s^2=\mathrm{e}^{f(w)}\mathrm{d}t^2+\mathrm{e}^{h(w)}\mathrm{d}x^2+\mathrm{e}^{j(w)}\left(\mathrm{d}y^2+\mathrm{d}z^2\right)+\mathrm{d}w^2,$$
(67)
in order to have enough degrees of freedom to be able to look at sufficiently general gauge field configurations. We look for solutions of the form
$`A_{iM}`$ $`=`$ $`(A_i(w),B_i(w),0,0,Z_i(w)),`$ (68)
$`\mathrm{\Phi }_i`$ $`=`$ $`R_i(w)\mathrm{e}^{\mathrm{i}\alpha _i(w)},`$ (69)
To simplify the algebra it is useful to define the following:
$``$ $`=`$ $`\frac{1}{2}\left(f+h+2j\right),`$ (70)
$`\overline{}`$ $`=`$ $`\frac{1}{2}\left(f+h+2j\right),`$ (71)
$`\overline{\overline{}}`$ $`=`$ $`\frac{1}{2}\left(fh+2j\right),`$ (72)
$`𝒜_i`$ $`=`$ $`eA_i+\stackrel{~}{e}A_j,`$ (73)
$`_i`$ $`=`$ $`eB_i+\stackrel{~}{e}B_j,`$ (74)
$`𝒵_i`$ $`=`$ $`eZ_i+\stackrel{~}{e}Z_j\alpha _i^{}.`$ (75)
In terms of these, the equations of motion for the gauge field components are
$`A_i^{\prime \prime }+\overline{}^{}A_i^{}2eR_i^2𝒜_i2\stackrel{~}{e}R_j^2𝒜_j`$ $`=`$ $`0,`$ (76)
$`B_i^{\prime \prime }+\overline{\overline{}}^{}B_i^{}2eR_i^2_i2\stackrel{~}{e}R_j^2_j`$ $`=`$ $`0,`$ (77)
$`\left(eR_i^2+\stackrel{~}{e}R_j^2\right)𝒵_i`$ $`=`$ $`0.`$ (78)
Equation (78) implies $`𝒵_i=0`$ which we use to simplify the rest of the equations of motion. For the scalar fields we have
$$R_i^{\prime \prime }+^{}R_i^{}+\mathrm{e}^f𝒜_i^2R_i\mathrm{e}^h_i^2R_i\mathrm{e}^{\mathrm{i}\alpha _i}\frac{V}{\mathrm{\Phi }_i^{}}=0,$$
(79)
and for the three metric functions we have
$`f^{\prime \prime }+\frac{1}{2}f^2+\frac{1}{3}f^{}h^{}+\frac{2}{3}f^{}j^{}\frac{1}{3}h^{}j^{}\frac{1}{6}j^2\frac{10}{3\kappa }\mathrm{e}^f\mathrm{\Phi }\frac{2}{3\kappa }\mathrm{e}^h\mathrm{\Psi }+\frac{2}{3\kappa }\mathrm{\Omega }`$ $`=`$ $`0,`$ (80)
$`h^{\prime \prime }+\frac{1}{3}f^{}h^{}\frac{1}{3}f^{}j^{}+\frac{1}{2}h^2+\frac{2}{3}h^{}j^{}\frac{1}{6}j^2+\frac{2}{3\kappa }\mathrm{e}^f\mathrm{\Phi }+\frac{10}{3\kappa }\mathrm{e}^h\mathrm{\Psi }+\frac{2}{3\kappa }\mathrm{\Omega }`$ $`=`$ $`0,`$ (81)
$`j^{\prime \prime }\frac{1}{6}f^{}h^{}+\frac{1}{6}f^{}j^{}+\frac{1}{6}h^{}j^{}+\frac{5}{6}j^2+\frac{2}{3\kappa }\mathrm{e}^f\mathrm{\Phi }\frac{2}{3\kappa }\mathrm{e}^h\mathrm{\Psi }+\frac{2}{3\kappa }\mathrm{\Omega }`$ $`=`$ $`0,`$ (82)
where
$`\mathrm{\Phi }`$ $`=`$ $`𝒜_1^2R_1^2+𝒜_2^2R_2^2+\frac{1}{2}A_1^2+\frac{1}{2}A_2^2,`$ (83)
$`\mathrm{\Psi }`$ $`=`$ $`_1^2R_1^2+_2^2R_2^2+\frac{1}{2}B_1^2+\frac{1}{2}B_2^2,`$ (84)
$`\mathrm{\Omega }`$ $`=`$ $`R_1^2+R_2^2+V.`$ (85)
We performed an extensive numerical search for solutions to Eqs. (767779-82), subject to the boundary conditions that the Higgs fields tend to VEVs, and that the metric functions tend to Randall-Sundrum form. For definiteness, we used the same Higgs potential as in the non-gauged model, namely as given by Eqs. (21), (23) and (26). The result of this investigation was the perhaps disappointing one that all solutions found required the gauge fields to vanish. Since the investigation was numerical, we cannot be certain that no non-trivial solutions exist, although we strongly suspect this to be the case.
## V Conclusion
We have constructed a U(1)$``$U(1) model with two complex scalar fields interchanged by a discrete symmetry coupled to $`4+1`$ dimensional gravity. It has a solution featuring clash-of-symmetries-style Higgs kink configurations in a Randall-Sundrum-like spacetime. Gravity is localised to the dynamically generated (smooth) brane, and the Randall-Sundrum limit of the solution can be defined. For the chosen Higgs kink and metric configurations, the Higgs potential had to be of a certain sextic form whose properties we studied in some depth. The symmetry breaking pattern varies as a function of the extra dimension coordinate $`w`$, and displays the clash of symmetries phenomenon: at all points $`|w|<\mathrm{}`$ both U(1)’s are broken, with alternate U(1)’s restored as $`w+\mathrm{}`$ and $`w\mathrm{}`$. The spontaneous breaking of the discrete interchange symmetry guarantees topological stability for the Higgs-gravity-induced brane. We also examined the gauged version of the theory, but we found that to have nonzero gauge fields the spacetime has to be flat.
This work sets the stage for incorporating gravity into more complicated models displaying the clash of symmetries. The eventual aim is to construct a realistic brane-world model that uses the clash of symmetries to understand spontaneous symmetry breaking in a more satisfactory fashion than obtained with the default mechanism using homogeneous Higgs vacuum expectation values.
###### Acknowledgements.
GD was supported in part by the Commonwealth of Australia, DPG by the Puzey bequest to the University of Melbourne, RRV by the Australian Research Council, and KCW in part by the U.S. Department of Energy (DOE) grant DE-FG02-85ER40237. RRV would like to thank D. Vignaud, A. Goldwurm and the APC group based at the Collège de France for their hospitality while this manuscript was completed. |
warning/0507/gr-qc0507122.html | ar5iv | text | # Regularized expression for the gravitational energy-momentum in teleparallel gravity and the principle of equivalence
## 1 Introduction
The notion of gravitational energy-momentum has been adressed recently in the framework of the teleparallel equivalent of general relativity (TEGR) . The TEGR is an alternative geometrical description of Einstein’s general relativity in terms of tetrad fields $`e_\mu ^a`$ and of the torsion tensor $`T_{\mu \nu }^a=_\mu e_\nu ^a_\nu e_\mu ^a`$ ($`a,b,,\mathrm{}`$ and $`\mu ,\nu ,\mathrm{}`$ are SO(3,1) and space-time indices, respectively). The field equations for the tetrad field $`e_\mu ^a`$ are precisely equivalent to Einstein’s equations. Therefore it is not a new theory for the gravitational field. The torsion tensor is related to the antisymmetric part of the Weitzenböck connection $`\mathrm{\Gamma }_{\mu \nu }^\lambda =e^{a\lambda }_\mu e_{a\nu }`$. The curvature tensor constructed out of this connection vanishes identically. Therefore this connection allows the notion of distant parallelism in space-time. Of course one may construct the Christoffel symbols and consider the physical and geometrical properties of the (nonvanishing) Riemann-Christoffel tensor.
The geometrical framework determined by the tetrad field and torsion tensor has proven to be suitable to investigate the problem of defining the gravitational energy-momentum. A consistent expression developed in the realm of the TEGR shares many features with the expected definition. The gravitational energy-momentum $`P^a`$ obtained in the framework of the TEGR has been investigated in the context of several disctinct configurations of the gravitational field. For asymptotically flat space-times $`P^{(0)}`$ yields the ADM energy . In the context of tetrad theories of gravity, asymptotically flat space-times may be characterized by the asymptotic boundary condition,
$$e_{a\mu }\eta _{a\mu }+\frac{1}{2}h_{a\mu }(1/r),$$
(1)
and by the condition $`_\mu e_\nu ^a=O(1/r^2)`$ in the asymptotic limit $`r\mathrm{}`$. In the asymptotic limit above the quantity $`\eta _{a\mu }`$ in Eq. (1) coincides with the metric tensor of the Minkowski space-time $`\eta _{ab}=(+++)`$. An important property of tetrad fields that satisfy the condition above is that in the flat space-time limit we have $`e_\mu ^a(t,x,y,z)=\delta _\mu ^a`$, and therefore $`T_{\mu \nu }^a=0`$. Hence for the flat space-time we normally consider a set of tetrad fields such that $`T_{\mu \nu }^a=0`$ in any coordinate system. This condition establishes the reference space. However, in general an arbitrary set of tetrad fields that yields the metric tensor for asymptotically flat space-times does not satisfy the asymptotic condition given by Eq. (1). Moreover for such tetrad fields we have in general $`T_{\mu \nu }^a0`$ in the flat space-time. It might be argued, therefore, that the expression for the gravitational energy-momentum mentioned above is restricted to a particular class of tetrad fields, namely, to the class of frames such that $`T_{\mu \nu }^a=0`$ if $`e_\mu ^a`$ represents the flat space-time tetrad field.
The definition $`P^a`$ is invariant under global SO(3,1) transformations. We have argued elsewhere that it makes sense to have a dependence of $`P^a`$ on the frame. The energy-momentum in classical theories of particles and fields does depend on the frame, and we assert that such dependence is a natural property of the gravitational energy-momentum. The total energy of a relativistic body, for instance, depends on the frame. We normally assume that a set of tetrad fields is adapted to an ideal observer in the space-time determined by the metric tensor $`g_{\mu \nu }`$. For a given gravitational field configuration (a black hole, for instance), the infinity of possible observers is related to the infinity of tetrad fields (related by a local SO(3,1) transformation) that yields the metric tensor $`g_{\mu \nu }`$. Let $`x^\mu (s)`$ denote the worldline $`C`$ of an observer, and $`u^\mu (s)=dx^\mu /ds`$ its velocity along $`C`$. We may identify the observer’s velocity with the $`a=(0)`$ component of $`e_a^\mu `$, where $`e_a^\mu e_\nu ^a=\delta _\nu ^\mu `$. Thus, $`u^\mu (s)=e_{(0)}^\mu `$ along $`C`$ . The acceleration of the observer is given by
$$a^\mu =\frac{Du^\mu }{ds}=\frac{De_{(0)}^\mu }{ds}=u^\alpha _\alpha e_{(0)}^\mu .$$
(2)
The covariant derivative is constructed out of the Christoffel symbols. We see that $`e_a^\mu `$ determines the velocity and acceleration along the worldline of an observer adapted to the frame. From this perspective we conclude that a given set of tetrad fields, for which $`e_{(0)}^\mu `$ describes a congruence of timelike curves, is adapted to a particular class of observers, namely, to observers determined by the velocity field $`u^\mu =e_{(0)}^\mu `$, endowed with acceleration $`a^\mu `$. If $`e_\mu ^a\delta _\mu ^a`$ in the limit $`r\mathrm{}`$, then $`e_\mu ^a`$ is adapted to stationary observers at spacelike infinity. We may say, therefore, that $`P^{(0)}`$ yields the ADM energy for such observers.
In this article we will extend the definition $`P^a`$ for the gravitational energy-momentum previously considered for arbitrary tetrad fields, namely, for tetrad fields that satisfy $`T_{\mu \nu }^a0`$ in the flat space-time. The redefinition is the only possible consistent extension of $`P^a`$, valid for tetrad fields that do not satisfy boundary conditions like Eq. (1). We will also argue that an existing version of the principle of equivalence, namely, Einstein’s version of the principle, does not pose any obstacle to the concept of localized gravitational energy. We will show that the usual (textbook) version of the principle was never accepted by Einstein.
## 2 The principle of equivalence and the localizability of the gravitational energy
The concept of energy in classical electrodynamics is very simple and well known. We consider an arbitrary volume in the three-dimensional space and verify the existence of field lines of the electric and/or magnetic field in this region. The electromagnetic energy density is given by the standard expression that consists of the sum of the square of the electric and magnetic fields, and therefore is nonvanishing in a space-time region where the field lines are present. Charged particles in this region experience the Lorentz force. Therefore the manifestation of the Lorentz force is an indication of the existence of electromagnetic energy density.
Unfortunately there is not a simple picture in general relativity that relates gravitational “field lines” to the existence of gravitational energy density. Nevertheless, it is legitimate to expect that the manifestation of gravitational forces in a three-dimensional region is an indication of the existence of gravitational energy-momentum density in this region. However in general relativity there is a point of view according to which the gravitational energy density cannot be localized (see §20.4). The argument is the following. In any given small region of the space-time manifold we can find a coordinate system such that the Christoffel symbols disappear. In terms of this appropriate coordinate system the small region in question is “free of gravitational fields”. In summary, this is the argument that has been endorsed by many authors, who claim that the nonlocalizability of the gravitational energy-momentum is due to the principle of equivalence.
We do not endorse the conclusion above for various reasons. First, it is well known that the vanishing of the Christoffel symbols does not imply the vanishing of tidal forces in any infinitesimal region of the space-time. Therefore, the assertion that this region is free of gravitational fields is questionable, because we should agree that the existence of gravitational forces (e.g., on the worldline of a particle) is due to gravitational fields. It is not reasonable to accept the idea of having a force in a space-time region without the associated field in this same region.
Second, the principle of equivalence that supports the conclusion above is related to Pauli’s version of the principle , but is different from Einstein’s version. According to Pauli’s formulation, for every infinitely small world region there always exist a coordinate system in which gravitation has no influence either on the motion of particles or any other physical processes. The distinction between Einstein’s and Pauli’s formulation of the principle of equivalence has been addressed by Norton . From the point of view of Pauli’s formulation, the vanishing of the Christoffel symbols in a space-time region implies that gravitation has no effect in this region. We know, however, that what really vanishes in such region are the first derivatives of the metric tensor. In our opinion the mathematical feature that consists of the vanishing of the first derivatives of any metric tensor - but not of the second and highest derivatives - along any worldline in a Riemannian or pseudo-riemannian manifold, in any dimension, cannot be taken as a physical principle. It is just a feature of differential geometry.
Paulis’s formulation of the principle of equivalence is different from Eintein’s formulation . The latter is unquestionably considered to be the breakthrough that led Einstein to establish the conditions under which a noninertial frame is equivalent to an inertial one, extending in this way the principle of relativity. In view of the practical difficulties related to the description of arbitrary gravitational fields by means of the principle of equivalence, the latter was abandoned in favour of the principle of general covariance. Nevertheless the importance of the principle has always been recognized by Einstein in the years after the formulation of general relativity .
Einstein’s version of the principle of equivalence consists of considering a reference frame $`K`$ (a “Galilean system”), and a reference frame $`K^{}`$, which is uniformly accelerated with respect to $`K`$. Then one asks whether an observer in $`K^{}`$ must understand his condition as accelerated, or whether there remains a point of view acording to which he can interpret his condition as at “rest”. Einstein concludes that by assuming the existence of a homogeneous gravitational field in $`K^{}`$ it is possible to consider the latter as at rest. In his words: The assumption that one may treat $`K^{}`$ as at rest, in all strictness without any laws of nature not being fulfilled with respect to $`K^{}`$, I call the ‘principle of equivalence’ .
Again considering Ref. , we observe that an important remark about Einstein’s formulation of the principle of equivalence is not widely considered in the literature: Einstein’s formulation is established in Minkowski space-time. The passage from $`K`$ to $`K^{}`$ amounts to a frame transformation in a finite region of the space-time, not a coordinate transformation. Moreover Einstein never endorsed Pauli’s formulation. Einstein objected that in the infinitely small every continuous line is a straight line . He believed that the restriction to infinitesimal regions makes it impossible to distinguish the geodesic world lines of free point masses from other world lines and thus it is impossible to judge whether - in the words of Pauli’s formulation - “gravitation has no influence on the motion of particles”. Quoting Norton : It has rarely been acknowledged that Einstein never endorsed the principle of equivalence which results, here called the “infinitesimal principle of equivalence”. Moreover, his early correspondence contains a devastating objection to this principle: in infinitesimal regions of the space-time manifold it is impossible to distinguish geodesics from many other curves and therefore impossible to decide whether a point mass is in free fall. <sup>1</sup><sup>1</sup>1 It is worthwhile to point out a compact statement of the principle formulated by Einstein in 1918 : Principle of Equivalence: inertia and gravity are wesensgleich (identical in essence). From this and from the results of the special theory of relativity it necessarily follows that the symmetrical tensor $`g_{\mu \nu }`$ determines the metrical properties of space, the inertial behaviour of bodies in it, as well as the gravitational action.
The principle of equivalence follows from the equality of inertial and gravitational masses, and really establishes the equivalence between a noninertial reference frame and an inertial one with the addition of a suitable gravitational field. The nonvanishing of tidal forces in infinitesimal regions of the space-time does not allow us to conclude that such regions can be made free of gravitational fields by means of coordinate transformations. Two nearby particles in free fall undergo geodesic deviation irrespective of whether the metric tensor is reduced to the Minkowski form along their worldlines. We recall that by means of coordinate transformations one cannot reduce the tetrad field or the torsion tensor at a space-time point to their flat space-time values. Any space-time region, infinitesimal or not, is flat and consequently “free of gravitational fields” if and only if the Riemann-Christoffel tensor vanishes in this region. Arguments based on the “infinitesimal principle of equivalence” are not conclusive and cannot be taken to rule out the notion of gravitational energy-momentum density.
## 3 A regularized expression for the gravitational energy-momentum
Let us briefly recall the Lagrangian formulation of the TEGR. The Lagrangian density for the gravitational field in the TEGR in empty space-time is given by
$`L(e_{a\mu })`$ $`=`$ $`ke({\displaystyle \frac{1}{4}}T^{abc}T_{abc}+{\displaystyle \frac{1}{2}}T^{abc}T_{bac}T^aT_a)`$ (3)
$``$ $`ke\mathrm{\Sigma }^{abc}T_{abc},`$
where $`k=1/(16\pi )`$ and $`e=det(e_\mu ^a)`$. The tensor $`\mathrm{\Sigma }^{abc}`$ is defined by
$$\mathrm{\Sigma }^{abc}=\frac{1}{4}(T^{abc}+T^{bac}T^{cab})+\frac{1}{2}(\eta ^{ac}T^b\eta ^{ab}T^c),$$
(4)
and $`T^a=T_b^b^a`$. The quadratic combination $`\mathrm{\Sigma }^{abc}T_{abc}`$ is proportional to the scalar curvature $`R(e)`$, except for a total divergence. The field equations for the tetrad field read
$$e_{a\lambda }e_{b\mu }_\nu (e\mathrm{\Sigma }^{b\lambda \nu })e(\mathrm{\Sigma }_a^{b\nu }T_{b\nu \mu }\frac{1}{4}e_{a\mu }T_{bcd}\mathrm{\Sigma }^{bcd})=0,$$
(5)
It is possible to prove by explicit calculations that the left hand side of Eq. (5) is exactly given by $`\frac{1}{2}e[R_{a\mu }(e)\frac{1}{2}e_{a\mu }R(e)]`$. As usual, tetrad fields convert space-time into Lorentz indices and vice-versa.
The definition for the gravitational energy-momentum has first been obtained in the Hamiltonian formulation of the TEGR . However either the Hamiltonian or Lagrangian field equations may be suitably interpreted as equations that define the gravitational energy-momentum. The momentum canonically conjugated to the tetrad components $`e_{aj}`$ is given by $`\mathrm{\Pi }^{aj}=4ke\mathrm{\Sigma }^{a0j}`$. The latter quantity yields the definition of the gravitational energy-momentum $`P^a`$ contained within a volume $`V`$ of the three-dimensional spacelike hypersurface ,
$$P^a=_Vd^3x_k\mathrm{\Pi }^{ak}.$$
(6)
$`P^a`$ transforms as a vector under the global SO(3,1) group. It describes the gravitational energy-momentum with respect to observers adapted to $`e_\mu ^a`$. These observers are characterized by the velocity field $`u^\mu =e_{(0)}^\mu `$, and by the acceleration $`a^\mu `$ given by Eq. (2).
Let us assume that the space-time is asymptotically flat. The total gravitational energy-momentum is given by
$$P^a=_S\mathrm{}𝑑S_k\mathrm{\Pi }^{ak}.$$
(7)
The field quantities are evaluated on a surface $`S`$ in the limit $`r\mathrm{}`$.
In Eqs. (6,7) it is implicitly assumed that the reference space is determined by a set of tetrad fields $`e_\mu ^a`$ for the flat space-time such that the condition $`T_{\mu \nu }^a=0`$ is satisfied. However in general there exist flat space-time tetrad fields for which $`T_{\mu \nu }^a0`$. In this case we may generalize Eq. (6) by adding a suitable reference space subtraction term, exactly like in the Brown-York formalism . The Brown-York quasi-local energy expression is regularized by subtracting the energy of a flat slice of the flat space-time.
Let us denote $`T_{\mu \nu }^a(E)=_\mu E_\nu ^a_\nu E_\mu ^a`$, and $`\mathrm{\Pi }^{aj}(E)`$ as the expression of $`\mathrm{\Pi }^{aj}`$ constructed out of flat tetrads $`E_\mu ^a`$. The regularized form of the gravitational energy-momentum $`P^a`$ is defined by
$$P^a=_Vd^3x_k[\mathrm{\Pi }^{ak}(e)\mathrm{\Pi }^{ak}(E)].$$
(8)
This definition guarantees that the energy-momentum of the flat space-time always vanishes. The reference space-time is determined by the tetrad fields $`E_\mu ^a`$, obtained from $`e_\mu ^a`$ by requiring the vanishing of the physical parameters like mass, angular momentum, etc.
The total gravitational energy-momentum is obtained by integrating over the whole three-dimensional spacelike section. Assuming again that the space-time is asymptotically flat, we have
$$P^a=_S\mathrm{}𝑑S_k[\mathrm{\Pi }^{ak}(e)\mathrm{\Pi }^{ak}(E)],$$
(9)
where the surface $`S`$ is established at spacelike infinity. Like Eq. (6), the definition above transforms as a vector under the global SO(3,1) group.
The definition given by Eq. (8) is valid also in the context of space-times with an arbitrary topology. It is legitimate to take the tetrad fields $`E_\mu ^a`$ to represent the pure de Sitter or anti-de Sitter spaces, for instance, in which case Eq. (8) represents the gravitational energy-momentum defined about the latter space-times.
## 4 Reference frames in the Kerr space-time and the total gravitational energy
In this section we will apply Eq. (9) to a simple set of tetrad fields that describes the Kerr space-time, in order to illustrate the procedure (of course the analysis of the gravitational energy of the Kerr space-time may be carried out by means of several approaches). For this purpose we will evaluate the total gravitational energy. The asymptotic form of the Kerr metric tensor describes the exterior region of a rotating isolated material system. The set of tetrad fields to be considered allows a straightforward evaluation of connections and curvature (in the context of Riemannian geometry), but it has neither a simple geometrical structure when written in cartesian coordinates, nor an appropriate asymptotic behaviour. Before we carry out this analysis, we will recall the construction of tetrad fields as reference frames. We start by considering the flat Minkowski space-time with cartesian coordinates $`x^\mu `$.
Besides $`x^\mu `$, the flat space-time is endowed with cartesian coordinates $`q^a`$. The coordinate system $`q^a`$ establishes a global reference frame. The transformation matrix that relates the two coordinate systems defines a set of tetrad fields for the Minkowski space-time, $`E_\mu ^a=_\mu q^a`$. The coordinate transformation $`dq^a=E_\mu ^adx^\mu `$ can be globally integrated, and therefore it establishes a holonomic transformation between $`q^a`$ and $`x^\mu `$. Rotations and boosts between $`q^a`$ and $`x^\mu `$ are the two basic SO(3,1) transformations. The condition
$$E_{(i)j}(t,x,y,z)=E_{(j)i}(t,x,y,z),$$
(10)
ensures that $`q^a`$ is not rotating with respect to $`x^\mu `$, because a rotation will give rise to antisymmetric components in the sector $`E_{(i)j}`$ (Latin indices from the middle of the alphabet run from 1 to 3). On the other hand, a boost between $`q^a`$ and $`x^\mu `$ implies that $`E_k^{(0)}0`$ . Therefore by imposing $`E_k^{(0)}=0`$, or, equivalently, the time gauge condition,
$$E_{(i)}^0=0,$$
(11)
we ensure that the two coordinate systems have a unique time scale. Both $`q^a`$ and $`x^\mu `$ describe the flat space-time. Thus we may say that if conditions (10) and (11) are imposed on $`E_\mu ^a(t,x,y,z)`$ the reference space-time with coordinates $`q^a`$ is neither rotating nor undergoing a boost with respect to the space-time with coordinates $`x^\mu `$ . If these conditions are imposed we have $`E_\mu ^a(t,x,y,z)`$ $`=\delta _\mu ^a`$, or
$$E_\mu ^a(t,r,\theta ,\varphi )=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& \mathrm{sin}\theta \mathrm{cos}\varphi & r\mathrm{cos}\theta \mathrm{cos}\varphi & r\mathrm{sin}\theta \mathrm{sin}\varphi \\ 0& \mathrm{sin}\theta \mathrm{sin}\varphi & r\mathrm{cos}\theta \mathrm{sin}\varphi & r\mathrm{sin}\theta \mathrm{cos}\varphi \\ 0& \mathrm{cos}\theta & r\mathrm{sin}\theta & 0\end{array}\right).$$
(12)
From a different but equivalent point of view we may say that
$`E_\mu ^a(t,x,y,z)=\delta _\mu ^a`$ is adapted to stationary observers in space-time, namely, observers that are endowed with the velocity field $`u^\mu =E_{(0)}^\mu =\delta _{(0)}^\mu `$ and acceleration $`a^\mu `$ given by Eq. (2). In this case, $`a^\mu =0`$.
A geometrical interpretation of tetrad fields as an observer’s frame can be given as follows. We consider an arbitrary path $`x^\mu (s)`$ of the observer in Minkowski space-time, where $`s`$ is the proper time of the observer. We identify $`dx^\mu /ds=u^\mu =E_{(0)}^\mu `$, where $`E_{(0)}^\mu `$ is the timelike component of the orthonormal frame (the temporal axis of the observer’s local frame). According to the hypothesis of locality , a noninertial observer at each instant along its worldline is equivalent to an otherwise identical momentarily comoving inertial observer. It follows from the hypothesis of locality that each noninertial observer is endowed with an orthonormal tetrad frame $`E_a^\mu `$, whose derivative along the path is given by
$$\frac{dE_a^\mu }{ds}=\varphi _a^bE_b^\mu ,$$
(13)
where $`\varphi _{ab}`$ is the antisymmetric acceleration tensor (not to be confused with $`\varphi ^{aj}`$ given by Eq. (28)). According to Refs. , in analogy with the Faraday tensor we can identify $`\varphi _{ab}(𝐚,𝛀)`$, where $`𝐚`$ is the translational acceleration ($`\varphi _{(0)(i)}=a_{(i)}`$) and $`𝛀`$ is the frequency of rotation of the local spatial frame with respect to a nonrotating (Fermi-Walker transported ) frame. The invariants constructed out of $`\varphi _{ab}`$ establish the acceleration scales and lengths . It follows from Eq. (13) that
$$\varphi _a^b=E_\mu ^b\frac{dE_a^\mu }{ds}=E_\mu ^bu^\lambda _\lambda E_a^\mu .$$
(14)
Therefore given any set of tetrad fields for an arbitrary gravitational field configuration its geometrical interpretation can be obtained by suitably interpreting the velocity field $`u^\mu =e_{(0)}^\mu `$ and the acceleration tensor $`\varphi _{ab}`$, in case we “switch off” the gravitational field by making $`e_\mu ^aE_\mu ^a`$. In several situations it turns out to be easy to impose conditions (10) and (11) on $`e_\mu ^a`$. However, the proper interpretation of $`\varphi _{ab}`$ along a typical trajectory determined by the velocity vector $`u^\mu `$ of a class of observers adapted to a tetrad field seems to be a condition stronger than Eqs. (10) and (11).
Now we consider the Kerr space-time. The line element is given by
$$ds^2=\frac{\psi ^2}{\rho ^2}dt^2\frac{2\chi \mathrm{sin}^2\theta }{\rho ^2}d\varphi dt+\frac{\rho ^2}{\mathrm{\Delta }}dr^2+\rho ^2d\theta ^2+\frac{\mathrm{\Sigma }^2\mathrm{sin}^2\theta }{\rho ^2}d\varphi ^2,$$
(15)
where
$$\rho ^2=r^2+a^2cos^2\theta ,$$
$$\mathrm{\Delta }=r^2+a^22mr,$$
$$\chi =2amr,$$
$$\mathrm{\Sigma }^2=(r^2+a^2)^2\mathrm{\Delta }a^2\mathrm{sin}^2\theta ,$$
$$\psi ^2=\mathrm{\Delta }a^2\mathrm{sin}^2\theta .$$
Imposition of conditions (10) and (11) yields the following expression for $`e_\mu ^a`$,
$$e_{a\mu }=\left(\begin{array}{cccc}\frac{1}{\rho }\sqrt{\psi ^2+\frac{\chi ^2}{\mathrm{\Sigma }^2}\mathrm{sin}^2\theta }& 0& 0& 0\\ \frac{\chi }{\mathrm{\Sigma }\rho }\mathrm{sin}\theta \mathrm{sin}\varphi & \frac{\rho }{\sqrt{\mathrm{\Delta }}}\mathrm{sin}\theta \mathrm{cos}\varphi & \rho \mathrm{cos}\theta \mathrm{cos}\varphi & \frac{\mathrm{\Sigma }}{\rho }\mathrm{sin}\theta \mathrm{sin}\varphi \\ \frac{\chi }{\mathrm{\Sigma }\rho }\mathrm{sin}\theta \mathrm{cos}\varphi & \frac{\rho }{\sqrt{\mathrm{\Delta }}}\mathrm{sin}\theta \mathrm{sin}\varphi & \rho \mathrm{cos}\theta \mathrm{sin}\varphi & \frac{\mathrm{\Sigma }}{\rho }\mathrm{sin}\theta \mathrm{cos}\varphi \\ 0& \frac{\rho }{\sqrt{\mathrm{\Delta }}}\mathrm{cos}\theta & \rho \mathrm{sin}\theta & 0\end{array}\right).$$
(16)
The transformation $`dq^a=e_\mu ^adx^\mu `$ determined by the expression above cannot be globally integrated, because in this case $`e_\mu ^a_\mu q^a`$. Therefore $`dq^a=e_\mu ^adx^\mu `$ is an anholonomic transformation. An important feature of the equation above is that its expression in the asymptotic limit $`r\mathrm{}`$ is given by Eq. (1). Thus we may say that Eq. (16) is adapted to stationary observers at spacelike infinity. We also note that the flat space-time limit of Eq. (16) yields Eq. (12), and therefore $`T_{\mu \nu }^a(E)=0`$.
Equation (16) above has proven to describe satisfactorily the energy-momentum properties of the Kerr space-time (we note that tetrad fields for the Kerr space-time have also been addressed in Refs. ). However, the line element given by Eq. (15) admits a simple form that is useful for computational purposes, and which reads
$$e_{a\mu }=\left(\begin{array}{cccc}A& 0& 0& 0\\ 0& \frac{\rho }{\sqrt{\mathrm{\Delta }}}& 0& 0\\ 0& 0& \rho & 0\\ B& 0& 0& C\end{array}\right),$$
(17)
where
$`A`$ $`=`$ $`\left({\displaystyle \frac{\chi ^2\mathrm{sin}^2\theta +\psi ^2\mathrm{\Sigma }^2}{\rho ^2\mathrm{\Sigma }^2}}\right)^{\frac{1}{2}},`$
$`B`$ $`=`$ $`{\displaystyle \frac{\chi \mathrm{sin}\theta }{\rho \mathrm{\Sigma }}},`$
$`C`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Sigma }\mathrm{sin}\theta }{\rho }}.`$ (18)
The flat space-time limit of Eq. (17) is given by
$$E_\mu ^a=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& r& 0\\ 0& 0& 0& r\mathrm{sin}\theta \end{array}\right).$$
(19)
The expression above yields three nonvanishing torsion components:
$`T_{(2)12}(E)=1`$, $`T_{(3)13}(E)=\mathrm{sin}\theta `$, and $`T_{(3)23}(E)=r\mathrm{cos}\theta `$. Inspite of its simplicity, this tetrad field has a rather intricate structure when written in cartesian coordinates. It reads
$$E_\mu ^a(t,x,y,z)=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& \frac{x}{r}& \frac{y}{r}& \frac{z}{r}\\ 0& \frac{xz}{r\sqrt{x^2+y^2}}& \frac{yz}{r\sqrt{x^2+y^2}}& \frac{\sqrt{x^2+y^2}}{r}\\ 0& \frac{y}{\sqrt{x^2+y^2}}& \frac{x}{\sqrt{x^2+y^2}}& 0\end{array}\right).$$
(20)
In view of the geometrical structure of the equation above, we see that, differently from Eq. (16), Eq. (17) does not display the asymptotic behaviour determined by Eq. (1). Moreover, in general the tetrad field determined by Eq. (20) is adapted to accelerated observers. In order to verify this fact, let us consider a boost in the $`x`$ direction, say, of Eq. (20). We find
$$E_\mu ^a(t,x,y,z)=\left(\begin{array}{cccc}\gamma & \beta \gamma & 0& 0\\ \beta \gamma & \gamma \frac{x}{r}& \frac{y}{r}& \frac{z}{r}\\ 0& \frac{xz}{r\sqrt{x^2+y^2}}& \frac{yz}{r\sqrt{x^2+y^2}}& \frac{\sqrt{x^2+y^2}}{r}\\ 0& \frac{y}{\sqrt{x^2+y^2}}& \frac{x}{\sqrt{x^2+y^2}}& 0\end{array}\right),$$
(21)
where $`\beta `$ and $`\gamma `$ are constants defined by $`\beta =v/c`$ and $`\gamma =\sqrt{1\beta ^2}`$. It is easy to see that along an observer’s trajectory whose velocity is determined by $`u^\mu =(\gamma ,\beta \gamma ,0,0)`$ the quantities $`\varphi _{(j)}^{(k)}=u^i(E_m^{(k)}_iE_{(j)}^m)`$ constructed out of Eq. (21) are nonvanishing. This fact indicates that along the observer’s path the spatial axis $`E_{(i)}^\mu `$ rotate. Nevertheless Eq. (17) yields a satisfactory value for the total gravitational energy-momentum, as we will see.
We will integrate Eq. (9) over a surface of constant radius $`x^1=r`$, and then we require $`r\mathrm{}`$. Therefore we make $`k=1`$ in Eq. (9). Out of Eq. (17) we evaluate all torsion components $`T_{a\mu \nu }`$. We need the quantity
$$\mathrm{\Sigma }^{(0)01}=e_0^{(0)}\mathrm{\Sigma }^{001}=\frac{1}{2}e_0^{(0)}(T^{001}g^{00}T^1).$$
The calculations are lengthy but straightforward, and therefore they will be omitted here. We find
$$\mathrm{\Pi }^{(0)1}(e)=4ke\mathrm{\Sigma }^{(0)01}=\frac{1}{8\pi }\frac{\sqrt{\mathrm{\Delta }}}{\rho }(_r\mathrm{\Sigma })\mathrm{sin}\theta .$$
(22)
The expression of $`\mathrm{\Pi }^{(0)1}(E)`$ is obtained from Eq. (19) or, equivalently, by just making $`m=a=0`$ in the expression above. It is given by
$$\mathrm{\Pi }^{(0)1}(E)=\frac{1}{4\pi }r\mathrm{sin}\theta .$$
(23)
Thus the gravitational energy contained within a surface $`S`$ of constant radius $`r`$ reads
$`P^{(0)}`$ $`=`$ $`{\displaystyle _S}𝑑S_k[\mathrm{\Pi }^{(0)k}(e)\mathrm{\Pi }^{(0)k}(E)]`$ (24)
$`=`$ $`{\displaystyle _S}𝑑\theta 𝑑\varphi {\displaystyle \frac{1}{4\pi }}\mathrm{sin}\theta \left({\displaystyle \frac{1}{2}}{\displaystyle \frac{\sqrt{\mathrm{\Delta }}}{\rho }}(_r\mathrm{\Sigma })+r\right).`$
In the limit $`r\mathrm{}`$ we have
$$4ke\mathrm{\Sigma }^{(0)01}\frac{1}{4\pi }r\mathrm{sin}\theta (1\frac{m}{r}).$$
(25)
Therefore for the total gravitational energy of the Kerr space-time we obtain
$$P^{(0)}_r\mathrm{}𝑑\theta 𝑑\varphi \frac{1}{4\pi }\mathrm{sin}\theta \left(r(1\frac{m}{r})+r\right)=m,$$
(26)
which is the expected result.
We may also integrate Eq. (24) on the surface of constant radius $`r=r_+`$, where $`r_+`$ is the external horizon of the Kerr black hole. On this surface the function $`\mathrm{\Delta }=r^2+a^22mr`$ vanishes. Therefore we find $`P^{(0)}=r_+=m+\sqrt{m^2a^2}`$, a result that is quite different from the irreducible mass of que Kerr black hole. The localization of gravitational energy in the Kerr space-time is correctly described by Eq. (16), according to the discussion in Ref. . As discussed above, the frame determined by Eq. (16) is adapted to stationary observers at spacelike infinity.
Before we close this section let us recall that by means of simple algebraic manipulations an expression for the gravitational energy-momentum flux was developed in Ref. . This expression follows directly from the field equations (5). It reads
$$\frac{d}{dt}\left[_Vd^3x_j\mathrm{\Pi }^{aj}\right]=_S𝑑S_j\varphi ^{aj},$$
(27)
where
$$\varphi ^{aj}=k[ee^{a\mu }(4\mathrm{\Sigma }^{bcj}T_{bc\mu }\delta _\mu ^j\mathrm{\Sigma }^{bcd}T_{bcd})].$$
(28)
The quantity above represents the $`a`$ component of the flux density in the $`j`$ direction. In Ref. this formalism was applied to the evaluation of energy loss in Bondi’s radiative space-time. In Eqs. (27) and (28) it is assumed that for the flat space-time we have $`T_{a\mu \nu }(E)=0`$. We may address Eq. (27) in the context of the present analysis. Let us assume that $`T_{a\mu \nu }(E)0`$. Since $`E_\mu ^a`$ is also a solution of the field equations (5), Eq. (27) is trivially satisfied for $`E_\mu ^a`$. Therefore we may write
$$\frac{d}{dt}\left[_Vd^3x_j[\mathrm{\Pi }^{aj}(e)\mathrm{\Pi }^{aj}(E)]\right]=_S𝑑S_j[\varphi ^{aj}(e)\varphi ^{aj}(E)],$$
(29)
where $`\varphi ^{aj}(E)`$ is constructed out of $`E_\mu ^a`$. We observe that as long as $`E_\mu ^a`$ (and consequently $`\mathrm{\Pi }^{aj}(E)`$) is time independent, the left hand side of Eq. (29) is simplified and therefore the energy-momentum loss can be easily calculated out of any set of tetrad fields. The vanishing of $`\varphi ^{aj}(e)\varphi ^{aj}(E)`$ at spacelike infinity (a feature that is expected to take place for asymptotically flat space-times) ensures the conservation of the total gravitational energy-momentum.
## 5 Discussion
In this article we have extended the definition for the gravitational energy-momentum previously considered in the framework of the TEGR, which requires $`T_{a\mu \nu }(E)=0`$ for the flat space-time, to the case where the flat space-time tetrad fields $`E_\mu ^a`$ yield $`T_{a\mu \nu }(E)0`$. In the context of the regularized gravitational energy-momentum definition it is not strictly necessary to stipulate asymptotic boundary conditions for tetrad fields that describe asymptotically flat space-times.
We have seen that Eqs. (13) and (14) provide a physical interpretation for a set of tetrad fields in Minkowski space-time, in terms of the linear acceleration and rotation of an observer adapted to the frame $`E_\mu ^a`$, endowed with velocity $`u^\mu =E_{(0)}^\mu `$. We note that all frames obtained from $`E_\mu ^a=\delta _\mu ^a`$ by means of a global SO(3,1) transformation (determined by constant transformation matrices $`\mathrm{\Lambda }_b^a`$) yield $`\varphi _a^b=0`$, according to Eq. (14). Thus the requirement $`\varphi _{ab}=\varphi _{ba}=0`$ seems to be equivalent to conditions (10) and (11).
The definition given by Eq. (8) can be applied to an arbitrary volume $`V`$ in any space-time, with an arbitrary topology. We propose that Eq. (8) represents the gravitational energy-momentum relative to the frame determined by the tetrad field $`e_\mu ^a`$, with $`E_\mu ^a`$ representing the tetrad field when the physical parameters of the metric tensor (mass, angular momentum, etc.) vanish.
Acknowledgements
This work was partially supported by the Brazilian Agency CNPQ. M. V. O. Veiga is supported by a CNPQ fellowship. |
warning/0507/physics0507078.html | ar5iv | text | # Design study of an optical cavity for a future photon-collider at ILC
## 1 Introduction
There is a worldwide consensus that the next particle physics project is a linear accelerator for positron-electron ($`e^+e^{}`$) collisions in the energy range of 500 GeV to $``$ 1 TeV. In Summer of 2004 the International Committee for Future Accelerators (ICFA) recommended to the International Linear Collider Steering Committee (ILCSC) an accelerator based on superconducting radio frequency resonators for setting up the International Linear Collider (ILC) . Such a technology has been developed by the TESLA (TeV Energy Superconducting Linear Accelerator) collaboration. Besides collisions between $`e^+e^{}`$, a second interaction point (IP) for collisions between photons as well as collisions of electrons on photons is foreseen. This second IP is commonly referred to as the Gamma-Gamma ($`\gamma \gamma `$) Collider arm of a linear collider. In this paper we will continue to use the parameters of the superconducting TESLA machine , as the parameters of the ILC are expected to be similar.
## 2 Creation and collision of high energy photons
In a photon collider two opposing pulsed electron beams of $`E_0`$ = 250 GeV energy travel towards the interaction point (IP). According to the architecture of the TESLA design each of them collides a few mm before the IP with a tightly focused laser beam. Its diameter is of the order of a few 10 $`\mu `$m, which will be much larger than the elliptical cross section of the electron bunch size. The latter exhibits 88 nm and 4.3 nm in its two perpendicular half axes for TESLA . The laser photons ($`1`$eV) are backscattered as outlined in Fig. 1.
Compton scattering raises their energy close to that of the initial electrons. The produced $`\gamma `$ propagate in the direction of flight path of the electrons. There is a small additional angular spread of the order of $`1/\gamma `$$``$ 2 $`\mu `$rad, where $`\gamma =E_0/(m_ec_0^2)`$ ($`m_e`$: electron mass, $`c_0`$: speed of light in vacuum). Compton backscattering is the most promising process for efficient creation of highly polarized, $`\gamma `$ beams in the range of several hundred GeV with low background . At the IP they collide with a similar opposite $`\gamma `$ or electron beam. The photon spot size at the IP will then be almost equal to that of the electrons at the IP. Therefore the total luminosity of $`\gamma \gamma `$, $`\gamma e^{}`$ collisions will be similar to that of the basic $`e^{}e^{}`$ beams.
¿From kinematics the maximum photon energy $`\mathrm{}\omega _m`$ is given by
$$\mathrm{}\omega _m=\frac{x}{x+1}E_0,x=\frac{4E_0\mathrm{}\omega _0}{m_e^2c_0^4}\mathrm{cos}^2\frac{\alpha }{2}$$
(1)
$`E_0`$, $`\mathrm{}\omega _0`$ and $`\lambda `$ denote the electron beam energy, photon energy and wavelength of the laser. $`\alpha `$ represents the crossing angle between laser and electron beam. It is desirable to keep $`x`$ below 4.8, since for larger $`x`$ the high energy photons are lost by $`e^+e^{}`$ pair production due to their interaction with the laser beam . For the TESLA beam parameters a wavelength near $`\lambda `$ = $`1\mu `$m coinciding with powerful solid state lasers appears to be promising . Then $`x`$$``$$`4.75`$ and $`\mathrm{}\omega _m`$ = 207 GeV arise for a 250 GeV beam. The resulting Compton spectrum is strongly sensitive to the product of the mean electron helicity $`\lambda _e`$ ($`\lambda _e1/2`$) and that of the laser photons $`P_c`$ ($`P_c1`$). A value $`2\lambda _eP_c`$ close to $`1`$, i.e. a circular polarized laser beam with opposite helicity as that of the electrons should be used to maximize luminosity at high photon energies .
The time structure of the laser system must match the bunch structure of the accelerator with 2820 bunches per train, each of 2.4 ps FWHM ($`\sigma `$ = 1 ps) duration for TESLA. The bunches are 337 ns apart and have a 5 Hz repetition rate . Pulses of several Joule are required to scatter the majority of electrons. The precise energy of each pulse depends strongly on the degree of focusing the optical onto the electron beam, as well as on the crossing angle $`\alpha `$ in respect to the electrons. A Nd:YLF-based laser architecture generating this special time structure and providing a flat and stable train of ultraviolet (UV) synchronized ps pulses over 0.8 ms long time periods is already in use for some time at the TESLA Test Facility (TTF). At a wavelength of 1047 nm it delivers $`200\mu `$J per pulse to the nonlinear crystals for conversion into the UV. At present, a maximum single pulse energy in excess of $`300\mu `$J for bursts containing 800 pulses at 1 MHz repetition rate, corresponding to 30 MW peak-power is generated in the infrared. For bursts of 2400 pulses this reduces to $`140\mu `$J per pulse at 3 MHz repetition rate within the pulse train . This system produces an average power between one and two watts. On the other hand, current solid-state laser technology is clearly in the position of generating short pulses at the 5 TW peak-power level. In the foreseeable future however, the required high average power for a photon collider of several ten kilowatt and diffraction limited pulses can not be produced directly with the output from a laser .
According to the TESLA machine parameters each bunch contains $`2\times \mathrm{\hspace{0.17em}10}^{10}`$ electrons, whereas a 5 J Laser pulse at $`\lambda =1064`$ nm consists of $`2.5\times \mathrm{\hspace{0.17em}10}^{19}`$ photons. Provided that all electrons within the bunch undergo scattering, only one in $`10^9`$ photons is lost during a single laser-electron collision. Thus, it has been proposed to reuse the remaining laser pulse by storing it in a passive resonant optical cavity . This results in a significant reduction of the required single pulse energy to be delivered by the laser. The time structure of the electron bunches as planned for TESLA is particularly well suited for application of such a cavity. Its round-trip time has to be adapted to the bunch-spacing of 337 ns resulting in a circumference $`U`$ of about 100 m. As will be shown, this is sufficient for wrapping the optical cavity around the particle detector.
## 3 Existing applications of passive optical cavities
Passive build-up cavities for generating a region of enhanced intensity are routinely used in combination with both continuous wave (cw) and mode-locked sources in a number of laser-related experiments such as frequency doubling (cw-laser), (mode-locked laser), high-resolution spectroscopy , high sensitive detection of absorption or cavity ring-down absorption spectroscopy . For the latter, linear cavities for micro-pulse energy enhancement in the mid-infrared region at 5.3 $`\mu `$m have been investigated at the Stanford free-electron laser .
A recent application is the amplification of ultrashort light pulses through phase coherent superposition of successive pulses from a mode-locked pulse train in a high-finesse optical cavity and subsequent cavity dumping . An amplification factor in the order of 10 has been obtained.
For optical interferometric detection of gravitational waves a 2000-fold power recycling cavity for the GEO600-project, one of several current first-generation large-scale interferometers, is specified to develop a cw light power of approximately 10 kW within the interferometer . Recently, a power enhancement in the range between 1000 and 1200 has already been obtained with approximately 5 W input from the mode cleaner cavities .
Laser increasingly enter also the field of high energy physics for monitoring the transverse size of electron beams (“laser wire”), as well as for precision measurement of electron beam polarization . Both are based on Compton scattering of low energy photons around 1 eV. The ”laser wire” technique relies on an optical cavity for boosting the laser power.
Most of the current cavities for the purpose of frequency doubling have a ring configuration, providing optical isolation from the laser cavity. This advantage has been pointed out in . The ring configuration is also the preferred geometry for an optical cavity for the photon collider.
Some non-resonant storage rings have been demonstrated in which the energy of a laser pulse that was lost through scattering, diffraction at the mirrors and upon impact with the electron beam is partially replaced by an amplifier . In these cases, the round-trip time of the optical pulse is less than the intervals of the laser pulses. These storage rings operate well below the power density required for the $`\gamma \gamma `$-collider.
## 4 Energy storage properties of the passive optical cavity
A short pulse is a wave packet with a wide frequency spectrum. One therefore has to pursue the transfer function approach, whereas the transient response of the circulating field strength $`_{circ}(t)`$ within the cavity to the incoming TESLA burst-mode time structure $`_{in}(t)`$ is mathematically described by a convolution of the incoming laser pulses with the impulse response function $`H(t)`$ of the cavity :
$$_{circ}(t)=_{in}(t)H(t)=_{\mathrm{}}^+\mathrm{}_{in}(\tau )H(t\tau )𝑑\tau \text{.}$$
(2)
Then $`\left|_{circ}(t)\right|^2`$ is proportional to the circulating intra-cavity power $`𝒫_{circ}(t)`$ and the incident power $`𝒫_{in}(t)`$. $`A(t)`$ := $`𝒫_{circ}(t)/𝒫_{in}(t)`$ defines the enhancement factor of the cavity for the incoming power $`𝒫_{in}(t)`$. For sake of simplicity we assume bursts of $`q`$ rectangular shaped laser pulses, each of duration $`\tau _L`$, that are synchronous with the round-trip time $`\tau _{circ}`$ of the cavity. Such a train of pulses is represented by a convolution of the function describing the characteristics of a single pulse with the sampling function:
$`_{in}(t)`$ $`=`$ $`_0\mathrm{}^{i\mathrm{\hspace{0.17em}2}\pi \nu _0t}rect_{\tau _L}(t){\displaystyle \underset{j=0}{\overset{q1}{}}}\delta (tj\tau _{circ})`$
$`=`$ $`_0\mathrm{}^{i\mathrm{\hspace{0.17em}2}\pi \nu _0t}{\displaystyle \underset{j=0}{\overset{q1}{}}}rect_{\tau _L}(tj\tau _{circ}),\text{with}`$
$$rect_{\tau _L}(t)=\{\begin{array}{cc}1\hfill & \text{if}|t|<\tau _L/2\hfill \\ 0\hfill & \text{if}|t|>\tau _L/2\hfill \end{array}\text{.}$$
As a consequence of the shifting property of the $`\delta `$-function, the circulating field of the cavity from Eq. (2) is a sum of properly shifted and scaled replicas of the injected field $`_{in}(t)`$. The convolution of Eq. (2) can however more conveniently be solved by taking the inverse Fourier transform of the intra-cavity spectrum $`\stackrel{~}{_{circ}(\nu )}`$. Then one benefits from the convolution property of the Fourier transform that states
$$\stackrel{~}{_{circ}(\nu )}=F\left[_{circ}(t)\right]=F\left[_{in}(t)H(t)\right]=F\left[_{in}(t)\right]F\left[H(t)\right]\text{.}$$
(4)
The input spectrum $`\stackrel{~}{_{in}}(\nu )`$ of the burst from Eq. (4) according to the Fourier transform
$$\stackrel{~}{_{in}(\nu )}=F\left[_{in}(t)\right]=_{\mathrm{}}^+\mathrm{}_{in}(t)\mathrm{}^{i\mathrm{\hspace{0.17em}2}\pi \nu t}𝑑t$$
(5)
results in
$`\stackrel{~}{_{in}(\nu )}`$ $`=`$ $`_0{\displaystyle _{\mathrm{}}^+\mathrm{}}\mathrm{}^{i\mathrm{\hspace{0.17em}2}\pi (\nu _0\nu )t}{\displaystyle \underset{j=0}{\overset{q1}{}}}rect_{\tau _L}\left(tj\tau _{circ}\right)dt`$
$`=`$ $`_0{\displaystyle \frac{\mathrm{sin}\left[\pi (\nu _0\nu )\tau _L\right]}{\pi (\nu _0\nu )}}\mathrm{exp}\left[i(q1)\pi (\nu _0\nu )\tau _{circ}\right]`$
$`{\displaystyle \frac{\mathrm{sin}\left[q\pi (\nu _0\nu )\tau _{circ}\right]}{\mathrm{sin}\left[\pi (\nu _0\nu )\tau _{circ}\right]}}\text{.}`$
The transfer function $`\stackrel{~}{H(\nu )}`$ for the power build-up ring cavity is almost identical to that of a Fabry-Perot interferometer. Extending the derivation described in to a ring cavity with loss the transfer function can be expressed as
$$\stackrel{~}{H(\nu )}=F\left[H(t)\right]=\frac{\sqrt{1R_c}\mathrm{exp}\left(i\pi \nu \tau _{circ}\right)}{1\sqrt{R_cV}\mathrm{exp}\left(i\mathrm{\hspace{0.17em}2}\pi \nu \tau _{circ}\right)}\text{.}$$
(7)
In the above equation $`\stackrel{~}{H(\nu )}`$ repesents the amplitude of a monochromatic wave within the cavity at optical frequency $`\nu `$ that is due to a monochromatic input field at same frequency with unit amplitude. $`R_c`$ represents the intensity reflectivity of the coupling mirror and $`V`$ the power loss factor for one round-trip. The latter subsumes all loss mechanisms: intra-cavity absorption, diffraction and scattering as well as the finite reflectivity of all $`N`$ cavity mirrors with exception of the coupling mirror. $`R_c`$ and $`V`$ are related to the corresponding quantities for the field amplitudes (denoted with small letters) by
$$\begin{array}{c}r_j=\sqrt{R_j}\text{,}j=1\mathrm{}N\hfill \\ r_{eff}v=\sqrt{R_c}\underset{=:\sqrt{V}}{\underset{}{\sqrt{R_2}\mathrm{}\sqrt{R_N}v}}=\sqrt{R_cV}\hfill \end{array}\}.$$
(8)
As stated in the square modulus of the transfer function Eq. (7) gives the power response to a unit power monochromatic input field. This turns out to be the Airy shape function with Finesse $``$ as a parameter :
$$|\stackrel{~}{H(\nu )}|^2=\frac{1R_c}{\left(1\sqrt{R_cV}\right)^2}\frac{1}{1+\left[\frac{2}{\pi }\mathrm{sin}\left(\pi \nu \tau _{circ}\right)\right]^2},\text{and}$$
(9)
$$=\frac{\pi \sqrt[4]{R_cV}}{1\sqrt{R_cV}}.$$
(10)
The cavity field is maximized at the resonances of the cavity which occur if $`\nu \tau _{circ}`$ equals a positive integer number $`n`$. In terms of the round-trip path length $`U`$ = $`c_0\tau _{circ}`$ and optical wavelength $`\lambda `$ = $`c_0/\nu `$ with $`c_0`$ representing the speed of light in vacuum, the resonance condition is simply
$$U=n\lambda \text{with a positive integer number }n\text{ .}$$
(11)
The transfer function describes the evolution of the electric field at an arbitrarily chosen point within the cavity. in Eq. (7) it is located a propagation distance of $`U/2`$ behind the coupling mirror. This can be seen by the argument $`\nu \tau _{circ}`$ of the complex exponential in the numerator instead of $`2\nu \tau _{circ}`$ for the complete revolution which appears in the denominator.
The above functions Eq. (4, 7) are combined to the intra-cavity spectrum $`\stackrel{~}{_{circ}(\nu )}`$ = $`\stackrel{~}{_{in}(\nu )}\stackrel{~}{H(\nu )}`$. By subsequent application of the inverse Fourier transform
$$_{circ}(t)=F^1\left[\stackrel{~}{_{circ}(\nu )}\right]=_{\mathrm{}}^+\mathrm{}\stackrel{~}{_{circ}(\nu )}\mathrm{}^{i\mathrm{\hspace{0.17em}2}\pi \nu t}𝑑\nu $$
(12)
one finally arrives at the intra-cavity-field
$`_{circ}(t)`$ $`=`$ $`_0\sqrt{1R_c}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}(R_cV)^{n/2}\mathrm{exp}\left[2\pi i\nu _0\left\{t\left(n+{\displaystyle \frac{1}{2}}\right)\tau _{circ}\right\}\right]`$ (13)
$`{\displaystyle \underset{j=0}{\overset{q1}{}}}rect_{\tau _L}\left[t\left(n+j+{\displaystyle \frac{1}{2}}\right)\tau _{circ}\right]\text{.}`$
Since $`\left|\sqrt{R_cV}\mathrm{exp}\left(2\pi i\nu _0t_{circ}\right)\right|<1`$, the identity
$$\frac{1}{1\sqrt{R_cV}\mathrm{exp}\left(2\pi i\nu \tau _{circ}\right)}=\underset{n=0}{\overset{\mathrm{}}{}}(R_cV)^{n/2}\mathrm{exp}\left(2\pi i\nu \tau _{circ}\right)$$
(14)
together with properties of the Fourier transform could be used in derivation of Eq. (13). Since $`\tau _{circ}\tau _L`$, the round-trip time $`\tau _{circ}`$ introduces a staircase behavior for the resulting intra-cavity power. Two regimes evolve from Eq. (13). The power builds up during the first $`p`$ revolutions of the pulse as long as energy is fed from the burst mode laser into the cavity ($`pq`$). When all $`q`$ laser pulses of the burst had been issued the stored power decays during the subsequent circulations ($`p>q`$).
$$\frac{𝒫_{circ}(t)}{𝒫_{in}(t)}=\{\begin{array}{c}(1R_c)\frac{\left(1\sqrt{R_cV}^p\right)^2+4\sqrt{R_cV}^p\mathrm{sin}^2\left(p\pi \nu _0\tau _{circ}\right)}{\text{}\left(1\sqrt{R_cV}\right)^2+4\sqrt{R_cV}\mathrm{sin}^2\left(\pi \nu _0\tau _{circ}\right)}\hfill \\ \text{if}pq\text{and}(2p+1)\tau _{circ}\tau _L<2t<(2p+1)\tau _{circ}+\tau _L\hfill \\ (1R_c)\frac{\left(1\sqrt{R_cV}^q\right)^2+4\sqrt{R_cV}^q\mathrm{sin}^2\left(q\pi \nu _0\tau _{circ}\right)}{\text{}\left(1\sqrt{R_cV}\right)^2+4\sqrt{R_cV}\mathrm{sin}^2\left(\pi \nu _0\tau _{circ}\right)}\left(R_cV\right)^{pq}\hfill \\ \text{if}p>q\text{and}(2p+1)\tau _{circ}\tau _L<2t<(2p+1)\tau _{circ}+\tau _L\hfill \\ 0\text{if }t\text{ otherwise}\hfill \end{array}$$
(15)
At resonance the enhancement of the incoming laser pulse power $`𝒫_{in}(t)`$ simplifies to
$$\frac{𝒫_{circ}(t)}{𝒫_{in}(t)}=\{\begin{array}{cc}(1R_c)\left[\frac{1\sqrt{R_cV}^p}{1\sqrt{R_cV}}\right]^2\hfill & \text{if}pq\text{and}\hfill \\ \multicolumn{2}{c}{(2p+1)\tau _{circ}\tau _L<2t<(2p+1)\tau _{circ}+\tau _L}\\ (1R_c)\left[\frac{1\sqrt{R_cV}^q}{1\sqrt{R_cV}}\right]^2\left(R_cV\right)^{pq}\hfill & \text{if}p>q\text{and}\hfill \\ \multicolumn{2}{c}{(2p+1)\tau _{circ}\tau _L<2t<(2p+1)\tau _{circ}+\tau _L}\\ 0\text{if }t\text{ otherwise}\hfill & \end{array}$$
(16)
Fig. 2 shows the typical transient behavior of the power within the cavity according to Eq. (16).
For an uninterrupted pulse train ($`p\mathrm{}`$) the enhancement factor $`A_p`$ at completion of $`p`$ round-trips converges to:
$$A_{max}=\frac{1R_c}{\left(1\sqrt{R_cV}\right)^2}.$$
(17)
The same expression arises from a monochromatic, i.e. continuous wave. An efficient enhancement will require $`R_c`$ and $`V`$ close to one. Stationary conditions will then be reached to a very good approximation within each burst of the TESLA time structure.
Besides power reflectivities and losses the power enhancement is determined by the length detuning $`\delta `$ of the cavity. Any arbitrary circumference can be expressed as $`U`$ = $`n\lambda +\delta `$ with $`\delta <\lambda `$. From Eq. (15) follows then for sufficient large $`p`$:
$$A(\delta )A_{max}\frac{1}{1+\left[\frac{2}{\pi }\mathrm{sin}\left(\pi \frac{\delta }{\lambda }\right)\right]^2}.$$
(18)
If the reflectivity of the coupling mirror equals a given loss factor $`V`$, the power enhancement takes its maximum possible value $`1/(1V)`$ and all light is absorbed by the resonant cavity. This is known as impedance matching ($`R_c:=V`$) in analogy to properties of electrical cables and waveguides. In Fig. 3 the dependence of the enhancement factor $`A`$ on the input mirror reflectivity at different values of $`V`$ is plotted. Each curve shows a maximum value $`1/(1V)`$ for $`R_c`$ equal to $`V`$.
The resonant ($`\delta =0`$), impedance-matched power enhancement can also be expressed in terms of the finesse $``$ given by
$$\frac{}{\pi }=\sqrt{A_{max}(A_{max}1)}A_{max}(A_{max}1)$$
(19)
Hence, a Finesse in excess of 300 corresponds to an assumed example of $`A_{max}`$ = 100.
The duration for which a laser pulse will be stored within the cavity is given by the cavity photon lifetime $`\tau _{cav}`$. That is defined by the approximately exponential decay of the power following a sudden turn-off of the injected laser pulses at some time $`t_0`$. Then $`\tau _{cav}`$ turns out to be represented by
$$\tau _{cav}=\tau _{circ}\frac{1}{2\mathrm{ln}\left(1{\displaystyle \frac{1}{A_{max}}}\right)}\tau _{circ}\frac{A_{max}}{2}(A_{max}1)$$
(20)
for the impedance matched cavity. A round-trip time of $`\tau _{circ}`$ = 337 ns and $`A_{max}`$ = 100 result in $`\tau _{cav}`$$``$ 49.75$`\times `$337 ns $``$ 16.8 $`\mu `$s. Using this result the effective number $`n_{rt}`$ of round-trips of an optical pulse within the cavity is determined by
$$n_{rt}=\frac{\tau _c}{\tau _{circ}}=\frac{1}{\mathrm{ln}\left(R_cV\right)}\frac{A_{max}}{2}\frac{}{2\pi }.$$
(21)
This number amounts to $`n_{rt}`$$``$ 50 in the above example.
It can be shown that for an uninterrupted pulse train the power that is reflected off an impedance matched cavity ($`R_c=V`$) declines as
$$𝒫_{reflec}^{(p)}=V^{2p1}𝒫_{in}0(p\mathrm{}).$$
(22)
No power will thus be reflected from an impedance matched cavity in steady state. This behavior can be used as an indicator for the alignment of the cavity in an automatic control system.
## 5 Constraints due to the particle detector and design proposal for an optical cavity
The Compton-interaction requires operation of the cavity in the Ultra High Vacuum (UHV) of the accelerator. When operated at the inevitable high power level optical windows within the cavity would imply the risk of distortion of the circulating optical ps-pulse as a result of the non-vanishing B-integral . For maintaining a sufficient high $`\gamma `$ flux density and hence a luminosity of around $`10^{34}`$ cm<sup>-2</sup> s<sup>-1</sup>, the laser pulse must be focused at the Compton conversion point (CP).
The particle detector for the $`\gamma \gamma `$-option is expected to be almost identical to the one for $`e^+e^{}`$-physics located at the leptonic IP. For TESLA, the path length between its end face and the focus of both optical cavities amounts to about 7.40 m . The end faces extend over 7.45 m and are perpendicular to the beam axis. As a consequence of the required tight optical focus, a widespread beam diameter of more than half a meter will emerge from the detector end-faces. Due to shortage of space any optics for the optical cavity should preferably be positioned outside the environment of the particle detector.
These demands suggest to employ two symmetrical telescopic ring cavities, required for Compton conversion at each of the counter-propagating electron beams, and to interlace them without mutual interference. Fig. 4 depicts an aerial view on two possible spatial embeddings of these cavities.
Their optical paths are enclosed within the associated optical beam pipes which are needed for maintaining the vacuum. In the upper sketch each ring resonator stretches across a plane. Both planes enclose an angle of $`90^{}`$ degrees. This tipping allows an un-obstructed passage of both electron beams directions. In addition, a small tilt angle of approx. $`+1.58^{}`$ and $`1.58^{}`$, respec., in the vertical direction prevents the mirrors of the final focusing optics being to close to the leptonic beams. The two laser beams originate from a separate hall placed directly above the detector. They are pointing in opposing directions with the same vertical tilt for coupling into their respective cavities on top of the detector. For the lower design, the optical beam paths outside the detector are kept in two adjacent parallel planes. They cross each other in the interaction region within the detector.
Fig. 5 describes the optical configuration of an individual cavity.
A telescopic convex-concave mirror arrangement generates a focus in the interaction region and a second, identical telescope re-collimates it again. In addition, they introduce a moderate beam magnification that reduces the beam size within the nearly collimated region of the cavity outside the detector. Here, the laser is coupled to the cavity via M<sub>1</sub> or M<sub>8</sub>.
## 6 Telescopic cavity
Disregarding the final size of the mirrors the focal spot size at CP is determined by the Gaussian eigenmode within this cavity. That is calculated by setting up the round-trip matrix $`\left(\left(M_{\text{ccw}}\right)\right)`$, . We assume total correction of the aberrations of the telescope mirrors. When starting at a reference plane containing CP and proceeding in counter clockwise direction, one obtains for the proposed cavity the following dependence on the geometrical data:
$`\left(\left(M_{\text{ccw}}\right)\right)=\left(\begin{array}{cc}G^{}& \frac{\rho ^{}g_1^{}}{2}(1G^{})\\ 4\frac{g_2^{}}{\rho ^{}}& G^{}\end{array}\right),\text{with these shortcuts}`$ (25)
$$\begin{array}{c}\begin{array}{cc}G^{}=2g_1^{}g_2^{}1\hfill & L_1^{}=\frac{L_1}{\text{}M}+\rho _x\hfill \\ g_1^{}=M\left(1\frac{L_1^{}}{\rho ^{}}\right)\hfill & L_2^{}=ML_2+\rho _c\hfill \\ g_2^{}=\frac{1}{\text{}M}\left(1\frac{L_2^{}}{\rho ^{}}\right)\hfill & \rho ^{}=\frac{\rho _c\rho _x}{\text{}2\delta }\hfill \\ & \\ M=\left|\frac{f_c}{\text{}f_x}\right|\hfill & \delta =t(f_c+f_x),f_c=\rho _c/2,f_x=\rho _x/2,\hfill \end{array}\}\hfill \end{array}$$
(26)
This results from multiplication of 2$`\times `$2-matrices describing the sequence of reflections on $`i`$ mirrors with radii of curvature $`\rho _i`$ and the free space propagations in between. $`L_1`$ represents the distance between both concave mirrors $`M_4`$ and $`M_5`$. Between $`M_6`$ and $`M_3`$ the ring is closed by $`L_2`$ via $`M_7`$, $`M_1`$, $`M_8`$ and $`M_2`$ of Fig. 5 on the lower path. $`t`$ designates the length of the telescope (spacing between $`M_3`$ and $`M_4`$ as well as between $`M_5`$, and $`M_6`$) and $`\delta _{tel}`$ quantifies a detuning from the length $`f_c+f_x`$ for afocal alignment of the telescope set up by concave and convex mirrors of effective focal lengths $`f_c`$ and $`f_x`$, respectively. Due to the inclined angle of incidence of the laser beam on the curved mirrors their focal lengths are modified, which is expressed by the effective value. They should have parabolic surfaces for reducing optical aberrations, predominantly astigmatism . Suitable definitions for a design of the cavity are the beam magnification $`\mu `$ as the ratio of beam radii $`w_c`$ and $`w_x`$ on the concave and convex mirror of one telescope, as well as the final image distance $`L_{image}`$ between the focus and the concave mirror of each telescope when a collimated beam enters the telescope on the side of the convex mirror.
$$\begin{array}{c}\mu :=\frac{w_c}{w_x}=M\frac{\delta _{tel}}{\text{}f_x}L_{image}=\frac{\mu f_c^2}{\text{}\delta _{tel}M}=\frac{f_c\left(tf_x\right)}{\text{}tf_cf_x}=\frac{L_1}{\text{}2}.\hfill \end{array}$$
(27)
The location of the telescope within the cavity determines the beam size $`w_{CP,Gaussian}`$ of the Gaussian focus at the reference plane CP :
$$w_{CP,Gaussian}^2=\frac{\lambda }{\pi }\rho ^{}\sqrt{\frac{g_1^{}}{g_2^{}}\left(1g_1^{}g_2^{}\right)}.$$
(28)
For a stable — i.e. on each revolution reproducing Gaussian beam — an additional confinement-condition has to be fulfilled:
$$1G^{}1.$$
(29)
The diameter of the final focusing concave mirrors determine the minimal collision angle $`\alpha `$ between laser and electron beam. They should be kept small for high yield of $`\gamma `$. This however gives rise to an additional contribution $`LF_{diff}`$ to the total loss factor $`V`$ of the cavity due to diffractive power loss, to a diffraction broadening of the focal spot size $`w_{CP,Gaussian}`$ and to deviations from the Gaussian mode. These were studied numerically using the physical optics code GLAD . A Gaussian beam was injected into a cavity formed by perfect reflecting mirrors and was numerically propagated according the Fox-Li approach throughout many revolutions, until a stationary field distribution was established. Then the beam size $`w_{CP}`$ at the focus was calculated by the second moment of the electric field distribution. $`LF_{diff}`$ resulted from the fractional drop of power during one round-trip.
By reducing the size of the concave mirrors, the beam size $`w_{cc}`$ on these mirrors also decreases, minimizing the optical power that is lost. This requires a growth of the beam waist at the CP. The relative broadening $`\gamma _{CP}:=w_{CP}/w_{CP,Gaussian}`$ of the focal spot size compared to the Gaussian focal beam radius is plotted in Fig. 6
as a function of a normalized radius $`a_{cc}/w_{cc,Gaussian}`$ of the concave mirrors. It is expressed in units of the corresponding Gaussian beam radius at this location. This plot is characteristic for the layout of an optical cavity following Fig. 5 For reasons of self-resemblance the diffractive broadening behaves similar when the size of the focus within the cavity of Fig. 5 is varied by shifting both telescopes either an equal amount towards (reduction) or away (enlargement) from the CP.
As shown later the diffraction loss $`LF_{diff}`$ turns out to be negligible for parameters with acceptable $`\gamma _{CP}`$.
## 7 Laser-electron crossing angle
In order to calculate the laser-electron crossing angle $`\alpha `$ and to specify $`w_{CP}`$, diffraction broadening has to be taken into account. In respect to a high yield of $`\gamma `$, crossing angle $`\alpha `$, mirror diameter $`2a_{cc}`$, waist size $`w_{CP}`$, laser pulse energy $`E_{pulse}`$, as well as laser pulse duration $`\tau _{pulse}`$ are all interdependent parameters. Their respective values were determined by a numerical optimization process using the CAIN Monte Carlo code which assumes charged particles interact with a Gaussian optical beam. The center-of-mass energy was set to 500 GeV. The aperture $`2a_{cc}`$ of the final focusing concave mirrors $`M_4`$, $`M_5`$ at distance $`L_{image}`$ from the conversion point CP in Fig. (5) sets an upper limit for the opening angle $`\theta _{cc}`$ of the laser cone that emerges from the beam waist:
$$\theta _{cc}=\frac{a_{cc}}{L_{image}}=\frac{a_{cc}}{w_{cc}}\theta _L.$$
(30)
Replacing $`L_{image}`$ by the far field divergence angle $`\theta _L`$ originating from a Gaussian beam waist $`w_0`$ results in the latter equation. $`w_{cc}`$ represents the beam radius on each of the concave mirrors. For the TESLA parameters, an additional offset angle $`\beta `$ of 17 mrad occurs in view of the disruption of the electron beam during the Compton scattering and the physical size of the electron beam pipes near the interaction region. The crossing angle $`\alpha `$ is thus expressed as
$$\alpha =\frac{a_{cc}}{w_{cc}}\theta _L+\beta ,\beta =17\text{mrad}.$$
(31)
This relation was encoded into CAIN via the optical Rayleigh length $`z_R=\pi w_0^2/\lambda `$ by using the relation $`\theta _L=w_0/z_R`$ of Gaussian beams . Instead of the Gaussian beam waist, the numerical value for the diffraction broadened waist $`w_{CP}`$ = $`\gamma w_{CP,Gaussian}`$ has been used for $`w_0`$. Fig. 7 shows the resulting luminosity
as a function of the crossing angle $`\alpha `$ for different mirror sizes $`a_{cc}/w_{cc}`$ and a laser pulse duration of $`\tau _L=3.5`$ ps FWHM ($`\sigma `$ = 1.5 ps). In the examined range the luminosity rises with decreasing diameter of the mirrors. A value of $`a_{cc}/w_{cc}=0.75`$ is therefore selected. An acceptable crossing angle is then $`\alpha \mathrm{\hspace{0.17em}55}`$ mrad. This corresponds to $`z_R(w_{CP})`$$``$ 0.63 mm, a diffraction broadened beam waist of $`w_{CP}`$$``$ 14.3 $`\mu `$m ($`1/e^2`$) ($`\sigma `$ = 7.15 $`\mu `$m) and a Gaussian waist $`w_{CP,Gaussian}`$$``$ 6.5 $`\mu `$m ($`1/e^2`$) ($`\sigma `$ = 3.3 $`\mu `$m)<sup>1</sup><sup>1</sup>1($`1/e^2`$) designates the radius that is defined by a drop of the intensity to $`1/e^2`$$``$ 13.5 % of its maximum value at the beam center..
A total luminosity<sup>2</sup><sup>2</sup>2Here $`z_{max}`$ is defined as $`z_{max}=x/(x+1+\xi ^2)`$ consistent with the definition in . of $`(z>0.8z_{max})=1.110^{34}\mathrm{cm}^2\mathrm{s}^1`$ can be achieved for these parameters with a pulse energy of 9 J . A non-linearity parameter $`\xi ^2=0.3`$ can be maintained in accordance with reference . In proportion to the laser pulse energy the required average laser power has also gone up to 9 J $`\times `$ 2820 $`\times `$ 5 Hz $``$ 130 kW. All resulting parameters for the Compton interaction zone of a $`\gamma \gamma `$-collider based on 250 GeV electron beams are compiled in Tab. 1.
The diffraction broadened beam waist size of $`w_{CP}`$$``$ 14.3 $`\mu `$m ($`1/e^2`$) corresponds to a Gaussian beam waist $`w_{CP,Gaussian}`$$``$ 6.5 $`\mu `$m ($`1/e^2`$) ($`\sigma `$ = 3.3 $`\mu `$m).
## 8 Specification and enhancement capability of the cavity
In view of the dimension of the particle detector, $`t`$, $`f_c`$, $`f_x`$ and $`\mu `$ were selected such that $`L_{image}15`$ m. The design parameter of the optical cavity according to Fig. 5 are compiled in Tab. 2.
The mode within the optical cavity produces its largest spot-size at the concave mirrors next to the beam waist. In a hypothetical cavity with mirrors of unlimited size a Gaussian beam of $`w_{cc,Gaussian}`$ = 79.1 cm radius ($`1/e^2`$) would develop. Assuming a clipping aperture of 75 % of the Gaussian beam radius, i.e. $`a_{cc}`$ = 0.75 $``$$`w_{cc,Gaussian}`$ = 59.3 cm $``$ 60 cm, the $`1/e^2`$ spot size on the concave mirrors reduces according to the numerical modelling to $``$ 42.7 cm. This leaves sufficient room for the wings of the electric field distribution as will be seen in the discussion of diffraction loss. The concave mirror should therefore have a diameter of about 120 cm. This yields a f-number $`f_\mathrm{\#}`$ = 12.7 for both focussing telescopes.
In this way the concave mirror is designed to represent the dominant aperture at which diffraction will occur. The telescope reduces the lateral extension of the mode by a factor of $`\sqrt{3}`$. However, a reduction of all other mirrors by this factor $`\sqrt{3}`$ is not advisable without introducing further significant apertures. These would influence the beam size even further and create losses which in the end would reduce the efficiency of the Compton conversion. Therefore, all mirrors should have the same diameter of 120 cm.
A plot of the numerically obtained stationary diffraction loss factor $`LF_{diff}`$ as a function of the diffraction broadened beam waist $`w_{CP}`$ is shown in Fig. 8.
It results from slight variations of the distance $`L_1`$ between both telescopes. At the same time $`L_2`$ has been adjusted in the opposite direction, in order to keep the total path length unchanged. The diameter of all mirrors was thereby adapted to end up with the same ratio $`a_{cc}/w_{cc,Gaussian}=0.75`$ as for the optimized cavity. Diffraction loss declines according to Fig. 8 towards smaller foci within the cavity. $`LF_{diff}`$ = 1 denotes no power loss due to diffraction.
A diffraction loss factor of roughly $`LF_{diff}0.9998`$ was obtained from an extrapolation down towards $`14.3\mu `$m diffraction broadened beam waist. Taking the reflectivity of practical highly reflecting mirrors into account reduces the total loss factor to $`V=LF_{diff}R_{HR}^{}{}_{}{}^{7}`$. $`R_{HR}`$ denotes the reflectivity of all remaining mirrors with exception of the coupling mirror. A reflectivity between $`R_{HR}`$ = 99.99 % and 99.95 % would permit a steady-state impedance matched power enhancement between 1100 and 270, for otherwise perfect conditions. The enhancement becomes the more sensitive against any impedance mismatch, the larger the amount of $`A`$ is. In practice, one probably will have a set of spare mirrors with slightly varying reflectivity and check which one will give the best enhancement.
Before onset of each electron bunch train $`n_{\kappa ,prepulse}`$ additional laser pulses are required for the accumulation of a sufficient large optical pulse energy within the cavity, that will ensure an enhancement of at least the fraction $`\kappa `$ ($`\kappa <1`$) of the stationary amount $`A_{max}`$. The number $`n_{\kappa ,prepulse}`$ is derived from Eq. (16) as follows:
$$n_{\kappa ,prepulse}=\frac{\mathrm{ln}\left(1\sqrt{\kappa }\right)}{\mathrm{ln}(\sqrt{R_cV})}.$$
(32)
As an example, the build-up of power obtained from the numerical propagation of an optical pulse under Gaussian seed-conditions through many revolutions is represented by Fig. 9 for $`R_c`$ = 99 %, V = 0.9998, and $`R_{HR}`$ = 100 %. At perfect resonance 1034 pre-pulses are then expected for obtaining an enhancement of at least 99 % of $`A_{max}`$ = 383, i.e. $`A`$ = 378.8. The stationary enhancement of 353 in Fig. 9 following in accordance with Eq. (32) after $``$ 1000 circulations of the pulse reflects a 92 % coupling of the injected Gaussian mode into the cavity mode. The number of laser pre-pulses for reaching an approximate steady-state declines as the impedance matching condition is violated.
When $`A`$ represents the power enhancement factor of the optical cavity, the required average power of the laser is reduced to $`A/(1+n_{prepulse}/n_{train})`$, where $`n_{train}`$ is the number of electron bunches in one train.
## 9 Damage threshold
The TELSA bunch structure consists of bunch trains of 2820 bunches with about 300 ns bunch spacing and a train repetition rate of 5 Hz. This pulse structure represents an intermediate regime between two identified damage mechanisms . For a pulse duration above 100 ps the damage occurs by melting due to heat deposition, whereas for less than 20 ps the damage site is limited to the region where the intensity is sufficient for a laser generated plasma. This occurs before a significant transfer of energy from the electrons to the lattice has taken place.
For $``$ 3.5 ps (FWHM) pulse duration the threshold for 600 shots at 10 Hz repetition rate and 1053 nm wavelength of uncoated very uniform fused silica samples is of the order of 2 J/cm<sup>2</sup> , whereas just $``$ 5 mJ/cm<sup>2</sup> per pulse results from the parameter for the cavity.
For a rough estimation of the worst case risk for material damage the cumulative effect of all pulses within the train is assumed to be represented by a single pulse with total energy and duration equivalent to that of all optical bunches within a train. The parameters of Tab. 1, 2 result then in a laser energy fluence of around 4 J/cm<sup>2</sup> within 10 ns at the mirrors of the cavity for a hypothetical Gaussian beam and 13 J/cm<sup>2</sup> for the cavity mode with truncated mirrors. At the final focussing concave mirrors this fluence is lower by a factor of 3 due to the beam expansion. According the laser damage threshold for 10 ns pulse duration for various materials including fused silica and coatings is higher and lies around 50 J/cm<sup>2</sup> for the substrat, and above approximately 44 J/cm<sup>2</sup> to 120 J/cm<sup>2</sup> for the mirrors depending on the composition of the multilayer coatings. Following the known scaling laws summarized e.g. in , for a wavelength around 1 $`\mu `$m a further factor of $`(1\text{ms}/10\text{ns})^{0.4}=100`$ for the increase of the threshold for the contribution due to thermal induced damage could probably be anticipated for distribution of the optical bunches within a train of $``$ 1 ms duration as in the TESLA time structure.
This means that the expected fluence is below the damage threshold. However, no data for trains of ps-pulses separated on nano- to microsecond time scales which accumulate to the stated fluences are known to us. For a final judgment an experimental study with a representative of the ILC bunch structure would be required.
## 10 Effects of cavity misalignments
¿From Eq. (18) the maximum acceptable length detuning $`\delta _\kappa `$ for the circumference of the optical cavity for a tolerated fractional decline $`\kappa `$ in the power enhancement factor $`A_{max}`$ from the resonance condition is derived to be given by
$$\delta _\kappa \frac{\lambda }{2}\sqrt{\frac{1}{\kappa }1}(1).$$
(33)
Maintaining the power enhancement factor e.g. above $`\kappa =90`$ % of its optimum value $`A_{max}`$ demands a match of the circumference of the cavity of better 0.57 nm for an assumed value $`A_{max}`$ of 100. Technical solutions for such a precise length stabilization are well-known . Even more stringent restrictions are common in interferometrical detection of gravitational waves.
Any misalignment generally results in displacement and broadening of the intra-cavity beam waist. Mechanical vibrations due to instrumentation in the environment of the particle detector, seismic ground motion as well as slight deviations during assembly might slightly modify position and tilt of the mirrors. According to our calculations, the displacement of the beam waist remains smaller then the Rayleigh length, i.e. the depth of the focus. This shift of the beam waist is hence negligible.
Broadening of the waist modifies the cavity-mode and reduces the mode coupling between laser and cavity. The resulting relative decline of the power enhancement factor $`A_{max}`$ was found to be largely independent of the finesse of the cavity. As an example two identical cavities with enhancement factors of 353 and 76, respectively, were considered. Fig. 10 shows the effect of shifting one telescope
over a range of up to 1 cm from its designated location away from the waist. Its internal alignment was assumed as being preserved. Regarding this vast variation range compared to the Rayleigh length in Tab. 2, the decline for sub-mm shifts is rather modest. It remains below 20 %. At a shift of 1 cm reduces the enhancement to 40 % of $`A_{max}`$. A decline of similar amount occurs e.g. for a deviation of the focal length $`f_x`$ of the convex mirror of 12 mm.
An increase of the focal length $`f_x`$ or shift of the telescope away from CP (the location of the beam waist) is accompanied by an increase of the radius of the beam waist. The latter is depicted by Fig. 11. Under the influence of clipping at the concave mirrors the growth of the diffraction broadened beam size is considerably lower. This allows to shift the operating point of the cavity further away from the stability limit. A zero waist size denotes an instable optical cavity.
## 11 Discussion
The $`\gamma \gamma `$-collider will be operated with trains of optical ps-pulses, whereas the optical interferometric detection of gravitational waves relies on continous-wave (cw) laser radiation. An automatic alignment system has already been developed for gravitational wave detection . Therefore, use of an additional uninterrupted train of weak ps-pulses for generation of error signals resulting from misaligned components and constant control of the cavity seems promising, as the cavity for the $`\gamma \gamma `$-collider can benefit from that knowledge.
The tiny absorption of the laser beam within the substrate and mirror coating is expected to induce a nonuniform temperature increase within the optic due to the amount of circulating optical power. This causes a nonuniform distortion of the optical path length by thermal expansion of the optic’s surface, and a variation of the material’s refractive index with temperature (thermal lensing). In our numerical model the focal length has to stay within a few tenth mm of its exact value for missing the desired beam waist by not more than 5 %. A correction of the curvature radius of imaging mirrors has already been demonstrated by exerting an axially symmetrical mechanical strain within the reflecting surface through radiative heating . The remaining non-axial-symmetric wavefront distortion generated by inhomogeneities in the substrate can be compensated by locally heating the mirror in addition with another laser via further computer controlled scanning mirrors. A proof-of-principle experiment has been performed in .
Moreover, the high misalignment sensitivity of the proposed cavity could be overcome by the introduction of adaptive optics for ensuring and controlling the development of a diffraction limited optical mode. This requires an additional control loop acting on the surface shape of at least one deformable mirror. The flat folding mirrors M<sub>1</sub> and M<sub>2</sub> adjacent to the telescopes in Fig. 5 appear to be especially well suited for this task.
The proposed cavity approximates a half-degenerate cavity which is characterized by a round-trip matrix of -1 times the unity matrix. In a half-degenerate cavity, the intensity distribution at any position along the propagation axis is relay imaged at completion of two circulations. Such a cavity can already be described in good approximation by ray-tracing (geometric-optical imaging). Diffraction has only be taken into account for calculation of the focus.
## 12 Conclusion
The actual enhancement factor of the cavity results from the cumulative effect of many small contributions affecting the total loss factor. Residual abberations tend to enlarge the focused beam at the Compton conversion point. Both are difficult to predict in advance. However, the influence of diffraction loss is for the proposed cavity almost negligible.
According to an estimation based on the known properties of laser induced damage, the suggested size of the mirrors provides still some reserve before the energy fluence of the circulating optical pulse reaches the damage threshold. For a final judgement of the upper limit, an experimental study with a representative of the special ILC bunch structure is required.
The use of adaptive optics appears to be essential for operation of such a cavity.
## 13 Acknowledgements
We would like to thank J. Gronberg, I.N. Ross, A. Stahl, and V. Telnov for many useful discussions. We are also grateful to F. Bechtel for his help during numerical calculations of the luminosity optimization as well as to N. Meyners of DESY’s MEA department for preparing the steric schemes used in Fig. 4. |
warning/0507/cond-mat0507331.html | ar5iv | text | # Screening in Ionic Systems: Simulations for the Lebowitz Length
## Abstract
Simulations of the Lebowitz length, $`\xi _\text{L}(T,\rho )`$, are reported for the restricted primitive model hard-core (diameter $`a`$) 1:1 electrolyte for densities $`\rho 4\rho _c`$ and $`T_cT40T_c`$. Finite-size effects are elucidated for the charge fluctuations in various subdomains that serve to evaluate $`\xi _\text{L}`$. On extrapolation to the bulk limit for $`T10T_c`$ the low-density expansions (Bekiranov and Fisher, 1998) are seen to fail badly when $`\rho >\frac{1}{10}\rho _c`$ (with $`\rho _ca^30.08`$). At higher densities $`\xi _\text{L}`$ rises above the Debye length, $`\xi _\text{D}\sqrt{T/\rho }`$, by 10-30$`\%`$ (upto $`\rho 1.3\rho _c`$); the variation is portrayed fairly well by generalized Debye-Hückel theory (Lee and Fisher, 1996). On approaching criticality at fixed $`\rho `$ or fixed $`T`$, $`\xi _\text{L}(T,\rho )`$ remains finite with $`\xi _\text{L}^c0.30a1.3\xi _\text{D}^c`$ but displays a weak entropy-like singularity.
Understanding the thermodynamic and correlation properties of ionic fluids has challenged both theory and experiment wei:sch . Typical electrolytes exhibit phase separation that is analogous to the gas-liquid transition in simple fluids, albeit at rather low temperatures when appropriately normalized. However, the long range of the Coulomb interactions has hampered understanding especially near criticality wei:sch . One crucial aspect is Debye-Hückel screening. For a $`d`$-dimensional classical fluid system with short-range ion-ion potentials beyond the Coulomb coupling $`z_\sigma z_\tau q^2/r^{d2}`$ (where $`z_\sigma `$ is the valence of ions of species $`\sigma `$ and mole fraction $`x_\sigma `$ while $`q`$ is an elementary charge), the charge-charge correlation function, $`G_{ZZ}(𝒓;T,\rho )`$, decays as $`\mathrm{exp}[|𝒓|/\xi _{Z,\mathrm{}}(T,\rho )]`$ (see, e.g., bek:fis ; aqu:fis ): the asymptotic screening length, $`\xi _{Z,\mathrm{}}`$, approaches the Debye length $`\xi _\text{D}=(k_\text{B}T/4\pi \overline{z}_2^2q^2\rho )^{1/2}`$ when the overall ion density $`\rho `$ approaches zero (with $`\overline{z}_2^2=_\sigma z_\sigma ^2x_\sigma `$ bek:fis ; aqu:fis ).
By contrast, at a critical point of fluid phase separation, the density-density (or composition) correlation length, $`\xi _{N,\mathrm{}}(T,\rho )`$, diverges, as do all the moments of $`G_{NN}(𝒓;T,\rho )`$. What then happens to charge screening near criticality? This question was first posed over a decade ago ste and has been addressed recently via the exact solution of $`(d>\mathrm{\hspace{0.17em}2})`$-dimensional ionic spherical models aqu:fis . As anticipated \[4(b)\], the issue of $`\pm `$ ion symmetry proves central. However, spherical models for fluids display several artificial features (e.g., infinite compressibilities on the phase boundary below $`T_c`$; parabolic coexistence curves, $`\beta \frac{1}{2}`$; etc.). Accordingly, understanding screening near criticality for more realistic models remains a significant task.
To that end we report here on a Monte Carlo study of the restricted primitive model (RPM), namely, hard spheres of diameter $`a`$ carrying charges $`q_\pm =\pm q`$ (so that $`z_+=z_{}=1`$, $`x_+=x_{}=\frac{1}{2}`$). Grand canonical simulations have been used and, to accelerate the computations, a finely discretized ($`\zeta =\mathrm{\hspace{0.17em}5}`$ level) lattice version of the RPM has been adopted pan . For this system the critical behavior is well established as of Ising-type with $`T_c^{}k_\text{B}T_ca/q^2\mathrm{\hspace{0.17em}0.05069}`$ and $`\rho _c^{}\rho _ca^30.079`$ lui:fis:pan . Furthermore, it has been demonstrated that for $`\zeta 3`$ the fine-lattice discretization does not qualitatively affect thermodynamic or finite-size properties kim:fis .
Ideally one would like to calculate $`\xi _{N,\mathrm{}}(T,\rho )`$ and $`\xi _{Z,\mathrm{}}(T,\rho )`$ near criticality; but even in nonionic model fluids, obtaining $`\xi _{N,\mathrm{}}`$ via simulations is hardly feasible. Nevertheless, the low-order moments $`M_{N,k}=|𝒓|^kG_{NN}(𝒓)d^dr`$ for $`k=0,1,2,\mathrm{}`$, are accessible and, by scaling, all the $`\xi _{N,k}(M_{N,k}/M_{N,0})^{1/k}`$ for $`k>0`$ diverge like $`\xi _{N,\mathrm{}}`$. However, for charges the Stillinger-Lovett sum rules bek:fis ; aqu:fis dictate $`M_{Z,0}0`$ (so that $`G_{ZZ}(𝒓)`$ is not of uniform sign) while the second moment satisfies $`M_{Z,2}=2\overline{z}_2^2q^2\rho \xi _\text{D}^2`$ which is fully analytic through $`(T_c,\rho _c)`$. On the other hand, the first moment of $`G_{ZZ}(𝒓)`$ is known bei:fel to be intimately related to charge screening via the so-called “area law” of charge fluctuations.
To explain this, consider a regular subdomain $`\mathrm{\Lambda }`$ with surface area $`A_\mathrm{\Lambda }`$ and volume $`|\mathrm{\Lambda }|`$, embedded in a larger domain, specifically say, the cubical $`L^d`$ simulation box. If $`Q_\mathrm{\Lambda }`$ is the total fluctuating charge in $`\mathrm{\Lambda }`$, electroneutrality implies $`Q_\mathrm{\Lambda }=0`$; but the mean square fluctuation, $`Q_\mathrm{\Lambda }^2`$, will grow when $`|\mathrm{\Lambda }|`$ increases. In the absence of screening one expects $`Q_\mathrm{\Lambda }^2|\mathrm{\Lambda }|`$; however, in a fully screened, bulk $`(L\mathrm{})`$ conducting fluid $`Q_\mathrm{\Lambda }^2`$ is asymptotically proportional to the surface area bei:fel . This was first observed by van Beijeren and Felderhof and later proven rigorously by Martin and Yalcin bei:fel . Following Lebowitz bei:fel one may then define a screening distance proportional to $`M_{Z,1}(T,\rho )`$, which we call the Lebowitz length, $`\xi _\text{L}(T,\rho )`$ bek:fis via
$$Q_\mathrm{\Lambda }^2/A_\mathrm{\Lambda }c_d\rho \overline{z}_2^2q^2\xi _\text{L}(T,\rho )\text{as}|\mathrm{\Lambda }|\mathrm{},$$
(1)
where $`c_d`$ is a numerical constant with $`c_3=\frac{1}{2}`$. Note that, since $`G_{ZZ}(𝒓)`$ is not necessarily of uniform sign, $`\xi _\text{L}(T,\rho )M_{Z,1}(T,\rho )`$ might diverge at $`T_c`$ even though the second moment $`M_{Z,2}\xi _\text{D}^2`$ remains finite!
Clearly, by simulating $`Q_\mathrm{\Lambda }^2`$ in various subdomains one may, as we show here, hope to calculate the Lebowitz length. To our knowledge no numerical results have been reported previously for $`d=3`$ although Levesque et al. lev:wei:leb presented a study (above criticality) for $`d=2`$. An exact low density expansion bek:fis proves that $`\xi _\text{L}/\xi _\text{D}1`$ when $`\rho 0`$ and corrections of order $`\rho ^{1/2}`$, $`\rho \mathrm{ln}\rho `$ and $`\rho `$ have been evaluated. This analysis bek:fis also served to validate the generalized Debye-Hückel (GDH) theory for the correlations lee:fis for small $`\rho `$.
The GDH theory, however, did not generate a $`\rho \mathrm{ln}\rho `$ term: nevertheless, as we find here, the exact expansion fails at very low densities — around $`\rho _c/10`$ even for $`T10T_c`$ — while GDH theory provides a reasonable estimate of $`\xi _\text{L}(T,\rho )`$ at higher densities: see Fig. 3 below. Furthermore, our calculations show that $`\xi _\text{L}`$ remains finite at criticality, exceeding $`\xi _\text{D}^c`$ by only $`33\%`$. Nonetheless, the Lebowitz length does exhibit weak singular behavior that, in accord with general theory, matches that of the entropy.
The first serious computational task is to understand the finite-size effects resulting from the $`L\times L\times L`$ simulation box with periodic boundary conditions. Each simulation at a given $`(T^{},\rho ^{})`$ yields a histogram of the total fluctuating charge $`Q_\mathrm{\Lambda }`$ for 24 different subdomains $`\mathrm{\Lambda }`$. We have used: six small cubes of edges $`\lambda L`$ with $`\lambda =0.3,\mathrm{\hspace{0.17em}0.4},\mathrm{},\mathrm{\hspace{0.17em}0.8}`$; seven ‘rods’ of dimensions $`\lambda L\times \lambda L\times L`$ with $`\lambda =0.2,\mathrm{},\mathrm{\hspace{0.17em}0.8}`$, four ‘slabs’ of dimensions $`\lambda L\times L\times L`$ with $`\lambda =0.2,\mathrm{},\mathrm{\hspace{0.17em}0.5}`$; and seven spheres of radius $`R=\lambda L`$ with $`\lambda =0.15`$-$`\mathrm{\hspace{0.17em}0.45}`$ in increments $`\mathrm{\Delta }\lambda =\mathrm{\hspace{0.17em}0.05}`$. To minimize correlations between these various subdomains, they have been located as far apart as feasible.
While the area law for the charge fluctuation, $`Q_\mathrm{\Lambda }^2`$, is rigorously true for $`L\mathrm{}`$ followed by $`\mathrm{\Lambda }\mathrm{}`$, it is by no means clear how it will be distorted for a finite subdomain $`\mathrm{\Lambda }`$ embedded in a finite system. To understand this Fig. 1 presents $`Q_\mathrm{\Lambda }^2`$, normalized by $`q^2`$, for the six cubic subdomains as a function of the reduced area $`A_\mathrm{\Lambda }/L^2`$ at selected temperatures and densities for box sizes $`L^{}L/a=6`$ and $`12`$.
Surprisingly, at high temperature and moderate density ($`T^{}=0.510T_c^{},\rho ^{}=0.08\rho _c^{}`$), the area law is well satisfied for $`\lambda 0.7`$ even for small systems. For $`L^{}=6`$ the data point for $`\lambda =0.8`$ deviates strongly from the linear fit (dashed line) owing to finite-size effects: indeed, electroneutrality dictates that $`Q_\mathrm{\Lambda }^2`$ should vanish when $`\lambda 1`$, corresponding to $`A_\mathrm{\Lambda }/L^2=6`$. At low densities around $`\frac{1}{3}\rho _c`$, the Debye length $`\xi _\text{D}\sqrt{T/\rho }`$ becomes large but nevertheless we see that the area law is still well satisfied. Furthermore, the area law is found to hold even near criticality: see the lowest plot. Note, however, that the linear fits to the data do not pass through the origin. This reflects finite-size effects which are discussed further below.
Combining (1) with the observations illustrated in Fig. 1, we conclude that charge fluctuations in the cubic subdomains are well described by
$$Q_\mathrm{\Lambda }^2(T,\rho ;L)=A_0(T,\rho ;L)+\frac{1}{2}\rho q^2\xi _\text{L}(T,\rho ;L)A_\mathrm{\Lambda },$$
(2)
where the intercept $`A_0(T,\rho ;L)`$ need not vanish. The (fitted) linear slope serves to define the finite-size Lebowitz length, $`\xi _\text{L}(T,\rho ;L)`$, which should approach the bulk value, $`\xi _\text{L}(T,\rho )`$. But by what route?
To answer this question consider Fig. 2 which displays $`\xi _\text{L}(T,\rho ;L)`$ vs. $`1/L^{}`$ for $`T^{}=\mathrm{\hspace{0.17em}0.5}`$ at various densities. It is rather clear that $`\xi _\text{L}(T,\rho ;L)`$ approaches its bulk limit as $`1/L`$. This can be understood by recalling the Lebowitz picture bei:fel in which the uncompensated charge fluctuations in a subdomain arise only from shells of area $`A_\mathrm{\Lambda }`$ and thickness of order $`\xi _\text{L}`$. By invoking the screening of $`G_{ZZ}(r)`$ one can see that $`\mathrm{\Delta }\xi _\text{L}\xi _\text{L}(L)\xi _\text{L}(\mathrm{})`$ for smooth subdomains decays as $`1/L^2`$. Indeed, by this route van Beijeren and Felderhof bei:fel showed explicitly that fluctuations in a sphere of radius $`R`$ (in an infinite system) approach their limiting behavior as $`1/R^2`$. For spheres in finite systems, we observe similarly that $`\xi _\text{L}(L)`$ approaches the bulk value as $`1/L^2`$. However, for cubes—which have edges and corners—and rods with edges, $`\xi _\text{L}(L)`$ gains a lower order, $`1/L`$ term as seen in Fig. 2. (The intercept $`A_0(L)`$ in (2) is, correspondingly, found to vary as $`L`$.) On the other hand, for slabs, lacking edges and corners, we find that $`\xi _\text{L}(L)`$ obtained via (1) approaches the limit exponentially fast.
Having established the finite-size behavior, let us examine $`\xi _\text{L}(T,\rho )`$ on the $`T^{}=\mathrm{\hspace{0.17em}0.5}`$ isotherm, well above $`T_c`$. Figure 3 shows estimates extrapolated from cubes, spheres and slabs. At moderate densities systems up to $`L^{}=\mathrm{\hspace{0.17em}16}`$ suffice but for $`\rho ^{}\mathrm{\hspace{0.17em}0.025}`$ we went up to $`L^{}=\mathrm{\hspace{0.17em}24}`$. The results may be compared with GDH theory lee:fis (dashed curve) and approximants which reproduce the exact low-density expansion known to order $`\rho `$ bek:fis . For the latter we adopt
$`\xi _\text{L}^{\text{[1,0]}}`$ $`=`$ $`\xi _\text{D}(T,\rho )\left[\mathrm{\hspace{0.17em}1}+a_1(T)\rho ^{}+a_2(T)\rho ^{}\mathrm{ln}\rho ^{}\right],`$ (3)
$`\xi _\text{L}^{\text{[0,1]}}`$ $`=`$ $`\xi _\text{D}(T,\rho )/[1a_1(T)\rho ^{}a_2(T)\rho ^{}\mathrm{ln}\rho ^{}],`$ (4)
shown in Fig. 3 as solid and dotted curves, respectively, where $`a_1(T)`$ and $`a_2(T)`$ follow from bek:fis . The simulations agree well with the low-density expansion but only up to $`\rho ^{}\mathrm{\hspace{0.17em}0.005}`$; thereafter $`\xi _\text{L}`$ rises above the Debye length much more slowly. By contrast, GDH theory captures the overall behavior of $`\xi _\text{L}(T,\rho )`$ over a broad density range, representing the numerical estimates to within a few percent at moderate densities, $`0.01\rho ^{}\mathrm{\hspace{0.17em}0.10}`$, where no exact results are available.
In the critical region the first question is the finiteness of $`\xi _\text{L}(T_c,\rho _c)`$. To answer we study $`\xi _\text{L}`$ on the critical isochore $`\rho =\rho _c`$ as $`TT_c`$. Figure 4 kim2 reveals that $`\xi _\text{L}/a`$ falls increasingly rapidly when $`T^{}`$ drops from $`\mathrm{\hspace{0.17em}0.5}`$ but clearly attains a finite nonzero value at $`T_c`$ that exceeds $`\xi _\text{D}^c/a\mathrm{\hspace{0.17em}0.2260}`$ lui:fis:pan . Owing to the relatively strong finite-size dependence of $`\xi _\text{L}`$ and the excessively large computational requirements near $`(T_c,\rho _c)`$, reliable extrapolation to $`L=\mathrm{}`$ is difficult. Nevertheless we may test for the nonanalytic behavior expected in any finite quantity fis:lan .
On general grounds fis:lan weak, entropy-like behavior is predicted. Thus temperature derivatives at $`\rho =\rho _c`$ should diverge like the specific heat, namely as
$$\rho C_V/k_\text{B}A^+/t^\alpha +A^0,$$
(5)
when $`t=(TT_c)/T_c\mathrm{\hspace{0.17em}0}`$, where $`\alpha \mathrm{\hspace{0.17em}0.109}`$ and $`A^+a^3=\mathrm{\hspace{0.17em}0.50}\pm \mathrm{\hspace{0.17em}0.07}`$ kim with, via a rough fit, $`A^0a^30.37`$. A direct comparison for finite $`L`$ of $`(\xi _\text{L}/\xi _\text{D})/T`$ with the specific heat is shown in Fig. 5 kim2 . Bearing in mind the lack of $`\xi _\text{L}`$ data near $`T_c`$ and its imprecision, the resemblance of the two plots is striking: we accept it as confirmation of the anticipated singularity.
Complementary nonanalytic behavior should arise on the critical isotherm as the reduced chemical potential $`\mu ^{}=[\mu \mu _0(T)]/k_\text{B}T`$ pan:fis varies. This is borne out by the plots in Fig. 6 of $`(\xi _\text{L}/\xi _\text{D})/\mu ^{}`$ and $`\mathbf{(}(\rho ^{}U^{})/\mu ^{}\mathbf{)}/\rho ^k`$ with $`k=\frac{1}{2}`$, where $`U^{}(T,\rho )`$ is the configurational energy per particle; the power $`\rho ^k`$ represents a convenient “$`k`$-locus factor” ork:fis:pan . In the bulk limit both functions should, by scaling, diverge as $`1/|\mu \mu _c|^\psi `$ with $`\psi =(1\beta )/(\beta +\gamma )\mathrm{\hspace{0.17em}0.43}`$ lui:fis:pan ; kim:fis .
Returning to the isochore $`\rho =\rho _c`$, theory indicates
$$\xi _\text{L}(T)=\xi _\text{L}^c\left[1+e_\alpha t^{1\alpha }+e_1t+e_\theta t^{1\alpha +\theta }+e_2t^2+\mathrm{}\right],$$
where $`\theta \mathrm{\hspace{0.17em}0.52}`$ is the leading correction exponent kim . By making allowance for the $`L`$-dependence and fitting over various ranges above $`T_c`$ we conclude $`\xi _\text{L}^c\mathrm{\hspace{0.17em}0.30}a`$ and, with less confidence, $`e_\alpha \mathrm{\hspace{0.17em}2.6}\pm 0.2`$ and $`e_12.2\pm 0.3`$.
In summary, the Lebowitz screening length, $`\xi _\text{L}(T,\rho )`$, has been studied for the restricted primitive model electrolyte via grand canonical Monte Carlo simulations of the charge fluctuations in subdomains. The corresponding area law that is asymptotically valid for large subdomains bei:fel holds surprisingly well even in small simulation boxes, $`L\mathrm{\hspace{0.17em}12}a`$. Finite-size effects can be understood so that the bulk, $`L\mathrm{}`$ limit may be extracted by extrapolation vs. $`1/L`$ for cubic subdomains and $`1/L^2`$ for spheres while the effective, finite-size Lebowitz lengths for slabs converge exponentially fast. Evaluation of $`\xi _\text{L}`$ for $`T\mathrm{\hspace{0.17em}10}T_c`$ over densities from $`0.03\rho _c`$ to $`4\rho _c`$ reveals that the exact low-density expansions bek:fis are effective only for $`\rho \frac{1}{10}\rho _c`$ whereas GDH theory lee:fis reproduces well the general trends. Finally, $`\xi _\text{L}`$ remains finite at criticality but exhibits weak, entropy-like singularities on approaching $`(T_c,\rho _c)`$. This is the first time that charge-charge correlations and a strongly state-dependent screening length have been studied by simulations close to criticality.
National Science Foundation support via Grants CHE 99-81772 and 03-01101 (M.E.F.) and DMR 03-46914 (E.L.) is gratefully acknowledged. |
warning/0507/astro-ph0507238.html | ar5iv | text | # Topological signatures in CMB temperature anisotropy maps
## I Introduction
It is becoming widely recognized that universes with a nontrivial spatial topology may be more natural models for our Universe than the traditional simply connected ones. This naturalness can be invoked from the mathematical point of view by arguing that there is an infinity of locally homogeneous and isotropic multiply connected 3–spaces, while there are only three simply connected ones; or with physical arguments coming from incursions into the nobody’s territories of Quantum Gravity and Quantum Cosmology. On the other hand, from a more pragmatic point of view, we can argue that cosmological models with a nontrivial spatial topology offer a very rich field of research, and are particularly well suited to explain certain reported “anomalous” features in CMB temperature maps, such as the alignments of their low $`\mathrm{}`$–modes Aligne1 ; Aligne2 .
Conversely, the full sky CMB temperature maps produced by the space missions COBE and WMAP provide us with an amazingly rich and high quality amount of data with which we can look for the topology of space. This is very compelling for those who wish to unmask our Universe and see its shape, since Cosmic Topology is at present an almost exclusively observational and phenomenological issue, due to the lack of an accepted fundamental physical theory which can predict the global topology of space.
Theory demands topology of space to leave several different kinds of marks in CMB temperature maps. Two of them have been largely studied and exploited to try to unveil the shape of our Universe, the distorsion of the angular power spectrum with respect to that of a simply connected universe PowSpec Weeks04 , and the existence of “circles in the sky” CinSky ; CSSK04 . Two other closely related signatures, a non–null bipolar power spectrum StatAnis ; SH05 , and alignments of the low $`\mathrm{}`$–modes Aligne1 ; Aligne2 ; Weeks04 , have been only marginally used. Our main motivation for deciding to adventure into cosmic topology with the CMB was the desire to get a deeper understanding of the nature and properties of these alignments as a topological signature.
One indispensable tool for a project like this is a software facility to produce simulated CMB temperature maps in multiply connected universes, so that we could systematically study the effects of different sizes and topologies on these alignments. These simulation procedures exist and have been used in several studies in cosmic topology Aurich , LevinB BPS , so we could expect that this issue of the project would not present any problem.
However, almost all known methods for computing CMB temperature anisotropies in multiply connected universes need to solve the Helmholtz equation in the manifold modelling our 3–dimensional space, the only exception to our knowledge being the work of Bond et al. BPS . To solve the Helmholtz equation is a relatively easy problem in Euclidean manifolds LevinA ; RWULL04 , but a very difficult task in spherical SphEigen ; Lachieze and hyperbolic 3–spaces HypEigen . Indeed, the spherical case has been completely solved analitically only recently by Lachièze–Rey Lachieze , while for compact hyperbolic manifolds the only possible approach is numerical HypEigen .
Among other results, in this paper we present a new approach to the computation of the correlation matrix $`a_\mathrm{}ma_\mathrm{}^{}m^{}^{}`$ of the coefficients of the spherical harmonic decomposition of CMB temperature anisotropies in a universe with nontrivial spatial topology. The main feature of this approach is that it avoids the explicit computation of the solutions of the Helmholtz equation in the spatial sections of spacetime. Instead, we express the correlation matrix in terms of the covering group alone. Incidentally, the idea of generating a CMB map exploiting the symmetries of the quotient space was already suggested time ago by Janna Levin and collaborators LevinB ; LSS98 . In particular we wish to quote a citation of a nontrivial claim in LSS98 (p.2695) which we have, in our opinion, succeeded in achieving: “By understanding the symmetries of the fundamental polyhedron and the identification rules, a CBR pattern can be deduced without the need to explicitly obtain the spectrum mode by mode.”
Our main formal result is a generic decomposition of the form
$$X^\mathrm{\Gamma }=X^{s.c.}+X^{t.s.},$$
(1)
where $`X`$ may be any covariance function which can be related to the two–point correlation function of the Newtonian potential (see Sec. III), $`\mathrm{\Gamma }`$ is the covering group of (the multiply connected) space, $`s.c.`$ stands for “simply connected”, and $`t.s.`$ means “topological signature”.
Thus, in Eq. (1), $`X^\mathrm{\Gamma }`$ is the covariance function computed in the manifold $`M=\stackrel{~}{M}/\mathrm{\Gamma }`$, and $`X^{s.c.}`$ the same covariance function but computed in the universal covering space $`\stackrel{~}{M}`$. It means that all the topological information is encoded in the “perturbative” term $`X^{t.s.}`$, and that is why we refer to it, generically, as the topological signature of $`X^\mathrm{\Gamma }`$. This decomposition is always possible whenever one can express $`X^\mathrm{\Gamma }`$ in terms of $`\mathrm{\Gamma }`$, as for example in the PSH method for detecting multiple copies of standard candles CosCris . In the present case we succeeded in writing the correlation matrix of the $`a_\mathrm{}m`$’s in this way by formally manipulating the two–point correlation function of the Newtonian potential derived by Bond, Pogosyan and Souradeep BPS .
Our approach to compute the correlation matrix $`a_\mathrm{}ma_\mathrm{}^{}m^{}^{}`$ has some methodological advantages in the simulation of CMB temperature maps. Indeed, by means of a suitable decomposition of the covering group $`\mathrm{\Gamma }`$ in cyclic subgroups, we are able to write down a formula for the correlation matrix of a complicated topology in terms of the correlation matrices of the cyclic topologies (topologies with a cyclic covering group) that cover it maximally. Since correlation matrices for cyclic manifolds are relatively easy to compute (we obtain a closed quadrature formula for the cylinder), we expect to obtain in the near future more efficient ways to simulate maps for complicated manifolds.
The decomposition of $`\mathrm{\Gamma }`$ in cyclic subgroups describes in a transparent way the symmetries of the manifold, and this fact gives rise to another advantage of our approach, in this case from the observational side. Universes with a cyclic topology present an alignment along the “direction of the generator isometry”. It follows that a CMB map in a universe with a nontrivial topology might present “patterns of alignment” (one for each $`\mathrm{}`$–mode) characteristic of its shape and size. Indeed, and this shall become “obvious”, symmetries of the quotient manifolds translate into symmetries of their patterns of alignment. We propose a method to search for these patterns by constructing maps of the dispersion of the squares $`|a_\mathrm{}m|^2`$ around the power spectrum. This opens the door to the development of methods to look for topology by searching these patterns, instead of limiting ourselves to considerations concerning only the special directions defined by the alignments.
For the sake of brevity, other advantages of our approach to compute the matrix $`a_\mathrm{}ma_\mathrm{}^{}m^{}^{}`$ are discussed in Sec. VI only. We prefer now to make a few remarks on some limitations of our work. We have considered here a few simplifications to develop the formalism, and worked out the details for the very simplest nontrivial topologies. In fact, we (i) have considered the Sachs–Wolfe effect as the only source of temperature anisotropies in CMB maps, (ii) have written the correlation matrix $`a_\mathrm{}ma_\mathrm{}^{}m^{}^{}`$ formally only for Euclidean 3-spaces, (iii) have worked out the details for homogeneous flat manifolds, (iv) have made a detailed analysis and some simulations only for cylinders ($`T^1`$ topology), and (v) these simulations were done using the Einstein–de Sitter model. We wish to close this introduction by justifying each one of these simplifications.
The shape of space is a global property, thus we expect the topological signatures in CMB to show themselves on very large scales only. Although a proof is missing, we believe that the main features of these signals would observationally appear if we restrict the searches to the low $`\mathrm{}`$–modes in the temperature maps, i.e., we do not need high resolution CMB maps in Cosmic Topology! Since the Sachs–Wolfe effect is the main source of temperature anisotropies at these scales RULW04 , we expect that theoretical explorations considering only this effect will put in evidence the main features of the topological signals that we would observe in a real map. The addition of the missing part of the anisotropies will only modify quantitatively the predictions made with our approximation, and thus, will only be important when adjusting theoretical models with data.
We consider this paper technically hard, so much care has been taken to write it in a clear and pedagogical way. The main features of our formalism and of the topological signatures we predict in CMB can be understood by restricting the presentation to Euclidean topologies. The inclusion of nonzero spatial curvature will only introduce additional technical considerations (and nothing qualitatively new), thus we decided to leave the spherical and hyperbolic cosmological models for a future paper.
The same applies to the lack in the paper of detailed explicit correlation matrices for nonhomogeneous flat manifolds. Indeed, the price we pay for simplicity and transparency in the presentation of the results for each specific cyclic topology is the need for very hard calculations in the middle steps, as can be seen in the appendices. Explicit calculations for nonhomogeneous flat manifolds will only add one page to the main body of the paper, and one or two more appendices to the already large list of them (see the end of the introduction). An exhaustive presentation for all the Euclidean manifolds is left for a future work.
The last but not the least, we performed the simulations with the Einstein–de Sitter model for simplicity. Nothing is lost from the theoretical point of view with this simplification since, as discussed in the paper, the structure of the topological signatures in CMB is captured in this oversimplified and old fashioned model of our Universe. However, more realistic $`\mathrm{\Lambda }`$CDM models will be required to confront theory with observations quantitatively.
We close this introduction by giving a detailed description of the structure of the paper. In Sec. II we briefly review the two most common methods to simulate CMB temperature maps in multiply connected universes, as well as present the method we have developed, and implemented for the Euclidean case. In Sec. III we define the topological signature in a correlation function, perform the decomposition of the covering group of a quotient space in its cyclic subgroups, and write the topological signature in terms of this decomposition. We also show here that the symmetries of a quotient space appear transparently in the decomposition of its covering group in cyclic subgroups.
In Sec. IV we apply our formalism to the homogeneous Euclidean manifolds, which are the simplest. We first compute the correlation matrix and the angular power spectrum for the cylinder, and apply the general results of the previous section in order to write down the correlation matrix of the spherical harmonic coefficients and the angular power spectrum for a generic torus. Finally, as an illustration, we write those expressions explicitly for the chimney ($`T^2`$ topology).
In Sec. V we first show, by means of simulations, that in a universe with the topology of a cylinder the low $`\mathrm{}`$–modes are aligned in a similar fashion as they are in the WMAP data. We then use the results in the previous sections to argue that CMB temperature maps in a universe with a nontrivial topology must present characteristic patterns of alignment, and propose the method of mapping on the sphere the dispersion of the squares $`|a_\mathrm{}m|^2`$ to look for them. Finally, in Sec. VI we discuss in detail the results of this paper and suggest further lines of research.
In brief, the main goal of this paper is to show that our approach to the computation of the correlation matrix $`a_\mathrm{}ma_\mathrm{}^{}m^{}^{}`$ (Sec. II), together with the decomposition of the covering group of a manifold in cyclic subgroups (Sec. III), led to the discovery of a new topological signature in CMB temperature maps, i.e., the “patterns of alignment” (Sec. V). To illustrate this we have used the simplest example, i.e., that of flat homogeneous manifolds (Sec. IV).
The paper has four appendices. In Appendix A we collect standard definitions and results related to spherical harmonics in order to set the conventions used in this paper. The technical calculations needed for the computation of the correlation matrix $`a_\mathrm{}ma_\mathrm{}^{}m^{}^{}`$ for the cylinder are presented in Appendices B and C. In the former we develop from scratch the theory of Clausen $`\phi `$–functions, for which we have not found any suitable reference in the literature. In the latter we compute a function that is the key part for computing efficiently the topological signature of the correlation matrix for the cylinder. Finally, in Appendix D we briefly reproduce known results for the correlation matrix of five out of the six compact orientable Euclidean 3–spaces. These formulae have been obtained previously in the literature by considering explicitly the solutions of the Helmholtz equation in these manifolds, and are written in terms of the $`k`$–modes LevinA ; RWULL04 . Our derivation avoids the need for considering these solutions.
## II Simulating CMB temperature maps
In this section we briefly describe two methods currently available for simulating CMB temperature maps in universes with non–trivial spatial topology, and proceed to develop our own formulation. We consider $`\mathrm{\Lambda }`$CDM universes, where the background metric of spacetime is of the Robertson–Walker type, and include scalar and adiabatic perturbations as the seeds for the temperature anisotropies of the CMB. In the Newtonian gauge we have
$$ds^2=a^2(\eta )\left[(1+2\mathrm{\Phi })d\eta ^2(12\mathrm{\Phi })\gamma _{ij}dx^idx^j\right]$$
for the metric, where $`\eta `$ is the conformal time, $`a(\eta )`$ is the scale factor, $`\mathrm{\Phi }`$ is the Newtonian potential, and
$$\gamma _{ij}=\left(1+\frac{K}{4}(x^2+y^2+z^2)\right)^2\delta _{ij}$$
is the metric of the spatial section of the background with sectional curvature $`K=0`$, $`\pm 1`$.
The matter content consists of radiation ($`\mathrm{\Omega }_r`$), baryonic and cold dark matter ($`\mathrm{\Omega }_m=\mathrm{\Omega }_b+\mathrm{\Omega }_{cdm}`$), and dark energy in the form of a cosmological constant ($`\mathrm{\Omega }_\mathrm{\Lambda }`$). Since we are interested on fluctuations on large angular scales, we make the assumption of instantaneous recombination and do not consider finite thickness effects. The main contribution to the temperature anisotropy observed at the direction $`𝐧`$ comes from the complete (ordinary plus integrated) Sachs–Wolfe effect
$`{\displaystyle \frac{\delta T}{T}}(𝐧)`$ $`=`$ $`{\displaystyle \frac{1}{3}}\mathrm{\Phi }(\eta _{LSS},R_{LSS}𝐧)+`$
$`2{\displaystyle _{\eta _{LSS}}^{\eta _0}}𝑑\eta {\displaystyle \frac{\mathrm{\Phi }}{\eta }}|_{(\eta ,R(\eta )𝐧)},`$
where the index $`LSS`$ stands for ‘last scattering surface’, the index $`0`$ for present time, and $`R(\eta )`$ is the comoving distance at instant $`\eta `$ between a photon, scattered at $`\eta _{LSS}`$, and the observer.
The Newtonian potential is written as
$$\mathrm{\Phi }(\eta ,𝐱)=𝑑qF(\eta ,q)\xi (q,𝐱).$$
(3)
The temporal part satisfies the equation
$`F^{\prime \prime }(\eta )+3(1+c_s^2)F^{}(\eta )+[2^{}+`$ (4)
$`(1+3c_s^2)(^2K)+c_s^2q^2]F(\eta )`$ $`=`$ $`0,`$
where $`c_s`$ is the speed of sound in the fluid and $`=a^{}/a`$ is the Hubble parameter in conformal time. On the other hand, the spatial part consists of solutions of the Helmholtz equation
$$(\mathrm{\Delta }+q^2)\xi (q,𝐱)=0,$$
(5)
where the index $`q`$ has been put as a variable in $`\xi `$ for simplicity of notation.
The integral in Eq. (3) has to be understood in a measure theoretic sense. Indeed, for multiply connected spaces the measure $`dq`$ is not the usual one but a combination of a discrete and an absolutely continuous measures, reducing the integral in (3) to a sum and an integral in the usual sense. In particular, if the space is compact, the measure reduces to a discrete one. This comes from the well–known fact that not every eigenmode of the Laplacian operator in the universal covering space $`\stackrel{~}{M}`$ is also an eigenmode in a quotient space $`M=\stackrel{~}{M}/\mathrm{\Gamma }`$. In fact, only eigenmodes in $`\stackrel{~}{M}`$ satisfying the invariance conditions
$$\xi (q,g𝐱)=\xi (q,𝐱)$$
(6)
for any $`g\mathrm{\Gamma }`$ project to eigenmodes in $`M`$.
The most straightforward way of simulating CMB temperature maps is by solving (4) and (5), performing the sum in (3), and then evaluating the SW effect (II). However, one has to consider that a temperature anisotropy map is a realization of a random field on the 2–sphere, and this randomness is inherited from that of the Newtonian potential (3). There are currently two ways to implement this random character in the simulations, and one goal of this paper is to propose a third one.
The first and most direct method is to consider the randomness in the temporal part of the decomposition (3) of the Newtonian potential. The two–point correlation function of the Newtonian potential at fixed time $`\eta `$ can then be written as
$`\mathrm{\Phi }(\eta ,𝐱)\mathrm{\Phi }(\eta ,𝐱^{})`$ $`=`$ $`{\displaystyle }dqdq^{}f(\eta ,q,q^{})\times `$
$`\xi (q,𝐱)\xi ^{}(q^{},𝐱^{}),`$
where $`f(\eta ,q,q^{})=F(\eta ,q)F(\eta ,q^{})`$ is the two–point correlation function for the amplitudes of the scalar perturbation modes and $`\xi (q,𝐱)`$ are normalized solutions of the Helmholtz equation.
Assuming that the Newtonian potential is a homogeneous and isotropic random field, the two point correlation function (II) reduces to a function of time $`\eta `$ and the distance $`d(𝐱,𝐱^{})`$, and thus we get $`f(\eta ,q,q^{})=P_\mathrm{\Phi }(\eta ,q)\delta (qq^{})`$, where $`P_\mathrm{\Phi }(\eta ,q)`$ is the gravitational power spectrum. If, in addition, the Newtonian potential is assumed to be gaussian, its random character is completely encoded in the variance of the temporal part
$$F^2(\eta ,q)=P_\mathrm{\Phi }(\eta ,q).$$
(8)
Specifying this function, one then takes as an initial condition, $`F(\eta _{init},q)`$, a realization of a normal distribution with zero mean and variance given by Eq. (8), and some suitable condition for the initial first derivative. With these initial conditions one solves for (4), so one can now compute the potential (3). Topology is considered by restricting in (3) to normalized solutions of the Helmholtz equation satisfying the invariance conditions (6). This method has been extensively used in Aurich ; LevinB ; LSS98 , although in the former the authors do not consider the randomness of the function $`F(\eta ,q)`$. Instead the random character of the CMB maps is attributed exclusively to the random character of the eigenmodes of the Laplacian in compact hyperbolic spaces.
The second method to produce simulated maps of CMB temperature anisotropies in universes with nontrivial spatial topology was first described in RULW04 , and used in URLW04 ; RWULL04 . It is based in considering the randomness of the Newtonian potential in the eigenmodes of the Laplacian operator. We begin by decomposing the general solution of the Helmholtz equation, in the universal covering space, as a sum of fundamental solutions
$$\xi (q,𝐱)=\underset{\mathrm{},m}{}\widehat{\xi }_\mathrm{}m(q)𝒴_\mathrm{}m(q,𝐱),$$
(9)
where
$$𝒴_\mathrm{}m(q,𝐱)=\rho _{\mathrm{}}(q,x)Y_\mathrm{}m(𝐧)$$
(10)
is the normalized solution of the Helmholtz equation after separating it in radial and angular variables. Here we have put $`x=|𝐱|`$, $`𝐧`$ is the unit vector in the direction of $`𝐱`$, and the $`Y_\mathrm{}m(𝐧)`$ are the spherical harmonic functions (see appendix A). Since the solutions (10) are normalized, the randomness of the eigenmodes’ amplitudes relies on the coefficients $`\widehat{\xi }_\mathrm{}m(q)`$.
Introducing (3) in (II), and using (9) and (10), we arrive at the decomposition of the temperature anisotropy map in spherical harmonics
$$\frac{\delta T}{T}(𝐧)=\underset{\mathrm{},m}{}a_\mathrm{}mY_\mathrm{}m(𝐧),$$
(11)
with multipole coefficients
$$a_\mathrm{}m=𝑑q\widehat{\xi }_\mathrm{}m(q)G_{\mathrm{}}(q),$$
(12)
and the effects of physical cosmology given by
$`G_{\mathrm{}}(q)`$ $`=`$ $`{\displaystyle \frac{1}{3}}F(\eta _{LSS},q)\rho _{\mathrm{}}(q,R_{LSS})+`$
$`2{\displaystyle _{\eta _{LSS}}^{\eta _0}}𝑑\eta {\displaystyle \frac{F}{\eta }}|_{(\eta ,q)}\rho _{\mathrm{}}(q,R(\eta )).`$
At this point it is convenient to recall how does topology enter in the story. Note that, due to the invariance conditions (6), not every solution of the form (9) is a solution in a quotient space. However, one would expect that any solution in a quotient space could be written in this form, the only problem being to find the correct coefficients $`\widehat{\xi }_\mathrm{}m(q)`$. These coefficients are not independent one from the other, since the invariance conditions (6) establish certain relations among them. In what follows we will assume that these relations can always be found, so that we will always represent an eigenmode in a quotient space by Eq. (9), with suitable coefficients.
A crucial point in RULW04 is the decomposition of these coefficients as
$$\widehat{\xi }_\mathrm{}m(q)=\sqrt{P_\mathrm{\Phi }(q)}\widehat{𝐞}_\mathrm{}m(q),$$
(14)
where $`P_\mathrm{\Phi }(q)`$ is the gravitational initial power spectrum, and the $`\widehat{𝐞}_\mathrm{}m(q)`$ form a multivariate gaussian random variable, with a non–diagonal covariance matrix due to the relations among the coefficients $`\widehat{\xi }_\mathrm{}m(q)`$ coming from the invariance conditions.
The simulation procedure can now be described. First we solve Eq. (4) using the initial condition $`F(\eta _{init},q)=1`$, and a suitable condition for the first derivative, and use this in (II) to compute $`G_{\mathrm{}}(q)`$. Then generate a realization of the random variable $`\widehat{𝐞}_\mathrm{}m(q)`$ and use (14) to compute $`\widehat{\xi }_\mathrm{}m(q)`$. The map is now simulated by computing the coefficients $`a_\mathrm{}m`$ using (12), and performing the sum in (11).
An alternative method of simulation, also proposed in RULW04 , is to construct the covariance matrix of the $`a_\mathrm{}m`$’s as
$`a_\mathrm{}ma_\mathrm{}^{}m^{}^{}`$ $`=`$ $`{\displaystyle }dqdq^{}G_{\mathrm{}}(q)G_{\mathrm{}^{}}(q^{})\times `$
$`\widehat{\xi }_\mathrm{}m(q)\widehat{\xi }_\mathrm{}^{}m^{}^{}(q^{}).`$
The substitution of (14) into (II), and the evaluation of the resulting integral give rise to expressions for the covariance matrix in terms of the eigenvalues and eigenmodes of the Laplacian operator. The multipolar coefficients are then obtained directly as a realization of a gaussian distribution with zero mean and covariance given by (II).
The method we propose in this paper lies along these lines, but we are able to manipulate the correlation function for the $`\widehat{\xi }_\mathrm{}m(q)`$ in a way that avoids the need for an explicit determination of the eigenmodes of the Laplacian. Instead, the final expression after the integration of (II) is given in terms of the isometries of the corresponding covering group.
Our starting point is an expression, obtained by Bond et al. in BPS , that relates the two–point correlation function of the Newtonian potential in a simply connected universe, and that in a multiply connected universe, when both potentials have the same initial power spectrum. For a homogeneous and isotropic random field the expression is
$$\mathrm{\Phi }(\eta ,𝐱)\mathrm{\Phi }(\eta ,𝐱^{})^\mathrm{\Gamma }=\underset{g\mathrm{\Gamma }}{}|g|\mathrm{\Phi }(\eta ,𝐱)\mathrm{\Phi }(\eta ,g𝐱^{})^{s.c.},$$
(16)
where $`|g|=1`$ if $`g`$ is orientation preserving, and $`1`$ otherwise.
We now show how to use Eq. (16) in order to express (II) in terms of the covering group, and for simplicity we will restrict the presentation to flat topologies. In Euclidean space, the most general solution of Eq. (5) is written in the form
$$\xi (q,𝐱)=d^3k\delta (qk)\widehat{\xi }(𝐤)e^{i𝐤𝐱}.$$
(17)
If we now expand the plane waves in spherical harmonics as
$$e^{i𝐤𝐱}=4\pi \underset{\mathrm{},m}{}i^{\mathrm{}}j_{\mathrm{}}(kx)Y_\mathrm{}m^{}(𝐧_𝐤)Y_\mathrm{}m(𝐧),$$
(18)
where $`j_{\mathrm{}}(x)`$ is the spherical Bessel function of order $`\mathrm{}`$, and introduce it in (17) we obtain $`\xi (q,𝐱)`$ expressed as in Eq. (9) with
$$\widehat{\xi }_\mathrm{}m(q)=4\pi i^{\mathrm{}}d^3k\delta (qk)\widehat{\xi }(𝐤)Y_\mathrm{}m^{}(𝐧_𝐤),$$
(19)
where $`𝐧_𝐤`$ is the unit vector in the direction of $`𝐤`$, and $`\rho _{\mathrm{}}(q,x)=j_{\mathrm{}}(qx)`$.
Note that we have not decomposed $`\widehat{\xi }_\mathrm{}m(q)`$ as in Eq. (14). Instead, the decomposition (19) allows us to implement the randomness of the Newtonian potential in the modes $`\widehat{\xi }(𝐤)`$. In fact, introducing (19) in (II), the covariance matrix for the $`a_\mathrm{}m`$’s now reads
$`a_\mathrm{}ma_\mathrm{}^{}m^{}^{}`$ $`=`$ $`(4\pi )^2i^{\mathrm{}\mathrm{}^{}}{\displaystyle }d^3kd^3k^{}G_{\mathrm{}}(k)G_{\mathrm{}^{}}(k^{})\times `$
$`\widehat{\xi }(𝐤)\widehat{\xi }^{}(𝐤^{})Y_\mathrm{}m^{}(𝐧_𝐤)Y_\mathrm{}^{}m^{}(𝐧_𝐤^{}).`$
It is the correlation function of the modes $`\widehat{\xi }(𝐤)`$ that carries all the topological information, as we will see in the following.
Introducing (17) in (3) we obtain
$$\mathrm{\Phi }(\eta ,𝐱)=d^3kF(\eta ,k)\widehat{\xi }(𝐤)e^{i𝐤𝐱},$$
thus the two–point correlation function of the Newtonian potential now reads
$`\mathrm{\Phi }(\eta ,𝐱)\mathrm{\Phi }(\eta ,𝐱^{})`$ $`=`$ $`{\displaystyle }d^3kd^3k^{}F(\eta ,k)F(\eta ,k^{})\times `$ (21)
$`\widehat{\xi }(𝐤)\widehat{\xi }^{}(𝐤^{})e^{i(𝐤𝐱𝐤^{}𝐱^{})}.`$
At this point we have to recall that an Euclidean isometry can always be written as $`g=(R,𝐫)`$, where $`R`$ is an orthogonal transformation and $`𝐫`$ is an Euclidean vector, and that this isometry acts on a vector $`𝐱`$ as $`g𝐱=R𝐱+𝐫`$. It is now an easy task to deduce from Eqs. (16) and (21) that
$$\widehat{\xi }(𝐤)\widehat{\xi }^{}(𝐤^{})^\mathrm{\Gamma }=\underset{g\mathrm{\Gamma }}{}\widehat{\xi }(𝐤)\widehat{\xi }^{}(R𝐤^{})^{s.c}e^{iR𝐤^{}𝐫}.$$
(22)
In most inflationary models the initial perturbations of the gravitational field are homogeneous and isotropic Gaussian random fields, thus the correlation matrix of the $`𝐤`$–modes in a simply connected universe takes the form
$$\widehat{\xi }(𝐤)\widehat{\xi }^{}(𝐤^{})^{s.c}=\frac{P_\mathrm{\Phi }(k)}{k^3}\delta (𝐤𝐤^{}).$$
The use of (22) now yields
$$\widehat{\xi }(𝐤)\widehat{\xi }^{}(𝐤^{})^\mathrm{\Gamma }=\frac{P_\mathrm{\Phi }(k)}{k^3}\underset{g\mathrm{\Gamma }}{}\delta (𝐤R𝐤^{})e^{iR𝐤^{}𝐫},$$
for the correlation matrix of the $`𝐤`$–modes in the quotient space $`M=\stackrel{~}{M}/\mathrm{\Gamma }`$, which when substituted in (II) finally gives
$`a_\mathrm{}ma_\mathrm{}^{}m^{}^{}^\mathrm{\Gamma }`$ $`=`$ $`(4\pi )^2i^{\mathrm{}\mathrm{}^{}}{\displaystyle }{\displaystyle \frac{d^3k}{k^3}}\mathrm{\Psi }_{\mathrm{}\mathrm{}^{}}(k)\times `$
$`\mathrm{{\rm Y}}_\mathrm{}^{}m^{}^\mathrm{\Gamma }(𝐤)Y_\mathrm{}m^{}(𝐧_𝐤),`$
where the physical effects are encoded in
$$\mathrm{\Psi }_{\mathrm{}\mathrm{}^{}}(k)=P_\mathrm{\Phi }(k)G_{\mathrm{}}(k)G_{\mathrm{}^{}}(k),$$
(24)
and the topological information in
$$\mathrm{{\rm Y}}_\mathrm{}m^\mathrm{\Gamma }(𝐤)=\underset{g\mathrm{\Gamma }}{}e^{i𝐤𝐫}Y_\mathrm{}m(𝐧_{R^T𝐤}).$$
(25)
The integration in (II) is over the whole $`𝐤`$–space. The topological information is carried in Eq. (25), which automatically selects the eigenvalues of the Laplacian operator in $`M`$. This can be seen in Appendix D, where $`\mathrm{{\rm Y}}_\mathrm{}m^\mathrm{\Gamma }(𝐤)`$ is expressed in terms of Dirac’s delta functions centered in the eigenvalues of the Laplacian operator in the corresponding quotient spaces.
## III Decomposition of $`\mathrm{\Gamma }`$ in cyclic subgroups
In this section we develop some formal results in order to proceed further. Especifically, we define the topological signature of any covariance function that can be decomposed as
$$X^\mathrm{\Gamma }=\underset{g\mathrm{\Gamma }}{}X^g,$$
as for example, the two–point correlation function of the Newtonian potential and the correlation matrix $`a_\mathrm{}ma_\mathrm{}^{}m^{}^{}`$. Then we work out a suitable decomposition of a covering group in their cyclic subgroups, and write down the topological signature in terms of this decomposition. We also show here that the symmetries of a quotient space appear transparently in the decomposition of its covering group in cyclic subgroups. It follows that the main result of this section is the ellucidation of how these symmetries manifest themselves in the topological signature of CMB temperature anisotropy maps.
Let us begin by writing the obvious decomposition
$$X^\mathrm{\Gamma }=X^{s.c.}+X^{\widehat{\mathrm{\Gamma }}},$$
(26)
where $`\widehat{\mathrm{\Gamma }}=\mathrm{\Gamma }\{id\}`$. The second term in the right hand side is the topological signature in the covariance function. The expressions we present in the following are analogous to (26) and are also rather obvious.
It is convenient to introduce a notation, so natural, that has been used in (26) without any previous definition. Let $`S`$ be any subset of isometries of the covering space $`\stackrel{~}{M}`$, then a superscript $`S`$ in the covariance function means
$$X^S=\underset{gS}{}X^g.$$
Then if $`M=\stackrel{~}{M}/\mathrm{\Gamma }`$ is a quotient space and $`\mathrm{\Gamma }_1\mathrm{\Gamma }`$ is any subset of the covering group, we can immediately write $`X^\mathrm{\Gamma }=X^{\mathrm{\Gamma }_1}+X^{\mathrm{\Gamma }\mathrm{\Gamma }_1}`$. This expression is nothing but the simplest generalization of Eq. (26), which corresponds to the trivial case $`\mathrm{\Gamma }_1=\{id\}`$. We get a further generalization as follows, let $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ be any two subsets of the covering group $`\mathrm{\Gamma }`$, such that $`\mathrm{\Gamma }_1\mathrm{\Gamma }_2=\mathrm{\Gamma }_3`$, then
$$X^\mathrm{\Gamma }=X^{\mathrm{\Gamma }_1}+X^{\mathrm{\Gamma }_2}X^{\mathrm{\Gamma }_3}+X^{\mathrm{\Gamma }(\mathrm{\Gamma }_1\mathrm{\Gamma }_2)}.$$
(27)
We can now write the formal result we are interested in. Consider the subsets $`G_1,\mathrm{},G_n\mathrm{\Gamma }`$ such that for any $`ij`$, $`G_iG_j=H`$, then by induction on (27) we get
$$X^\mathrm{\Gamma }=\underset{i=1}{\overset{n}{}}X^{G_i}(n1)X^H+X^{\mathrm{\Gamma }G},$$
(28)
where $`G=G_i`$.
To move forward, let $`G_1=g_1`$ and $`G_2=g_2`$ be two cyclic subgroups of $`\mathrm{\Gamma }`$, and let $`\mathrm{𝟎}\stackrel{~}{M}`$ be a lift to $`\stackrel{~}{M}`$ of the position of the observer in $`M`$. We will say that $`g_1`$ and $`g_2`$ are conjugate by an isometry that “does not move the observer” if there exists an isometry $`\eta `$ fixing $`\mathrm{𝟎}`$ and such that $`g_1=\eta ^1g_2\eta `$. Note that, as a consequence, we have that $`d(\mathrm{𝟎},g_1\mathrm{𝟎})=d(\mathrm{𝟎},g_2\mathrm{𝟎})`$, where $`d(𝐱,𝐲)`$ is the distance between two points $`𝐱`$ and $`𝐲`$ in $`\stackrel{~}{M}`$. By extension, we will also say that the groups $`G_1`$ and $`G_2`$ are conjugate by an isometry that does not move the observer. In addition, we will say that $`g_1`$ is a minimal distance generator of $`G_1`$ (with respect to the observer) if $`d(\mathrm{𝟎},g_1\mathrm{𝟎})d(\mathrm{𝟎},\gamma \mathrm{𝟎})`$ for any other generator $`\gamma G_1`$.
Now consider the isometries $`g_1,\mathrm{},g_n\mathrm{\Gamma }`$ that generate the cyclic groups $`G_i=g_i`$, and such that if $`ij`$ then $`G_iG_j=\{id\}`$. By using Eq. (28) we immediately obtain that the topological signature of the covariance function can be decomposed as
$$X^{\widehat{\mathrm{\Gamma }}}=\underset{i=1}{\overset{n}{}}X^{\widehat{G}_i}+X^{\mathrm{\Gamma }G}.$$
In the following we will be particularly interested in the case where the $`g_i`$’s are minimal distance generators of the $`G_i`$’s, and the latter form a complete set of groups mutually conjugate by isometries that do not move the observer.
Decomposing the topological signature further along these lines, let $`g_1,\mathrm{},g_n,h_1,\mathrm{},h_m\mathrm{\Gamma }`$ be minimal distance generators of the groups $`G_i=g_i`$ and $`H_j=h_j`$, and let $`G=G_i`$ and $`H=H_j`$. Moreover, suppose that the $`G_i`$’s and the $`H_j`$’s form two complete sets of groups mutually conjugate by isometries that do not move the observer, and such that $`GH=\{id\}`$ and $`d(\mathrm{𝟎},g_1\mathrm{𝟎})d(\mathrm{𝟎},h_1\mathrm{𝟎})`$. Then the topological signature can be decomposed as
$$X^{\widehat{\mathrm{\Gamma }}}=\underset{i=1}{\overset{n}{}}X^{\widehat{G}_i}+\underset{i=1}{\overset{m}{}}X^{\widehat{H}_i}+X^{\mathrm{\Gamma }(GH)}.$$
We can proceed along these lines again and again, and obtain the following decomposition, in cyclic subgroups, of the covering group $`\mathrm{\Gamma }`$,
$$\mathrm{\Gamma }=\underset{i=1}{\overset{\mathrm{}}{}}\underset{j=1}{\overset{k_i}{}}\mathrm{\Gamma }_{ij},$$
(29)
where $`g_{ij}\mathrm{\Gamma }`$ is a minimal distance generator of the cyclic group $`\mathrm{\Gamma }_{ij}`$, and such that
1. For each $`i`$, the set $`\{\mathrm{\Gamma }_{i1},\mathrm{},\mathrm{\Gamma }_{ik_i}\}`$ is a complete set of groups mutually conjugate by isometries that do not move the observer.
2. If $`ii^{}`$, the sets $`_{j=1}^{k_i}\mathrm{\Gamma }_{ij}`$ and $`_{j=1}^{k_i^{}}\mathrm{\Gamma }_{i^{}j}`$ have the identity as the only common element.
3. If $`i<i^{}`$, then $`d(\mathrm{𝟎},g_{i1}\mathrm{𝟎})d(\mathrm{𝟎},g_{i^{}1}\mathrm{𝟎})`$.
Then the topological signature of the covariance function can be written as
$$X^{\widehat{\mathrm{\Gamma }}}=\underset{i=1}{\overset{\mathrm{}}{}}\underset{j=1}{\overset{k_i}{}}X^{\widehat{\mathrm{\Gamma }}_{ij}}.$$
(30)
Thus, to compute the topological signature of the covariance function for any multiply connected manifold, it is enough to know how to compute it for manifolds whose covering groups are cyclic groups. It is now obvious that this decomposition will be particularly useful for calculating the correlation matrix of the $`a_\mathrm{}m`$’s for any compact manifold once we know how to calculate it for cyclic flat (twisted cylinders), spherical (lens spaces), and hyperbolic manifolds.
Let us now show that the decomposition (29) describes transparently the symmetries of the quotient manifold $`M=\stackrel{~}{M}/\mathrm{\Gamma }`$. Actually, this decomposition contains the symmetries of the Dirichlet fundamental polyhedron of $`M`$ centered at the observer’s position $`\mathrm{𝟎}\stackrel{~}{M}`$, which is what one expects to reconstruct with cosmological observations. Recall that the Dirichlet fundamental polyhedron centered at $`\mathrm{𝟎}\stackrel{~}{M}`$ is the set $`𝒟_\mathrm{𝟎}\stackrel{~}{M}`$ defined by (see Beardon )
$$𝒟_\mathrm{𝟎}=\{𝐱\stackrel{~}{M}:d(\mathrm{𝟎},𝐱)d(g\mathrm{𝟎},𝐱)\text{ for any }g\mathrm{\Gamma }\}.$$
The first thing to be noted is that, although the whole covering group enters in this definition, it turns out that, in order to effectively construct the Dirichlet polyhedron, we only need the minimal distance generators (and maybe the first few positive powers) of the first few cyclic groups $`\mathrm{\Gamma }_{ij}`$ and their inverses.
In fact, for each $`g\mathrm{\Gamma }`$ consider the semi–space
$$H_g=\{𝐱\stackrel{~}{M}:d(\mathrm{𝟎},𝐱)d(g\mathrm{𝟎},𝐱)\}.$$
Then it is obvious that the Dirichlet polyhedron is the intersection of all of these semi–spaces. However, there is a high redundancy here, since for a sufficiently large positive power $`n`$, we may have
$$\underset{k=1}{\overset{n1}{}}H_{g_{ij}^k}H_{g_{ij}^n},$$
and so this and further powers of $`g_{ij}`$ do not effectively contribute to the polyhedron $`𝒟_\mathrm{𝟎}`$. Additionally, if some $`H_g`$ effectively contributes to the polyhedron, so does $`H_{g^1}`$, thus the same argument holds for the inverses of the minimal distance generators. Moreover, note that due to condition 3. above, for a sufficiently large $`i`$, it may be the case that the semi–spaces $`H_{g_{ij}}`$ do not contribute effectively to the polyhedron.
The faces of the Dirichlet polyhedron are subsets of the boundary planes of the semi–spaces effectively contributing to it. In fact, for each $`H_g`$ effectively contributing, the corresponding face is orthogonal to the geodesic joining $`\mathrm{𝟎}`$ and $`g\mathrm{𝟎}`$, and cuts it at its middle point. It follows that the decomposition (29) describes the symmetries of the Dirichlet fundamental polyhedron of $`M`$ centered at the observer.
## IV Flat homogeneous manifolds
We have seen in the previous section that, to compute the topological signature of CMB temperature anisotropies in a given manifold, we just need to know how to compute it for the cyclic manifolds that cover it maximally. In the flat orientable case the cyclic manifolds are twisted cylinders, i.e., manifolds with covering group generated by a screw motion. We will now focus on the simplest case, the flat homogeneous manifolds, which are generated by translations only.
The flat homogeneous manifolds are 3–tori or $`T^3`$ manifolds (generated by three linearly independent translations), chimneys or $`T^2`$ manifolds (generated by two linearly independent translations), and cylinders or $`T^1`$ manifolds (generated by one translation). Thus, we first compute the correlation matrix $`a_\mathrm{}ma_\mathrm{}^{}m^{}^{}^\mathrm{\Gamma }`$ for cylinders, and then show how the decomposition (30) is used to compute the topological signature in this matrix for two– and three–dimensional tori. We also show that the computation of the angular power spectrum in tori is greatly simplified by this decomposition.
### IV.1 The cylinder
Let us consider a cylinder orthogonal to the $`z`$–direction, that is with covering group generated by the translation $`g=(I,𝐚)`$, with $`𝐚=L\widehat{𝐞}_z`$, where distances are measured in units of the radius of the last scattering surface $`R_{LSS}`$. This choice of the coordinate system is very convenient since here the cylinder appears invariant under (i) arbitrary rotations around the $`z`$–axis, (ii) the parity transformation, and (iii) the reflection on the $`y=0`$ plane. Thus, according to Sec. A.3, we will end with a real correlation matrix with no $`m`$–dependent correlations and the multipoles $`\mathrm{}`$ and $`\mathrm{}^{}`$ correlated only when both are even or odd.
It is convenient to recall here that our cylinder has injectivity radius equal to $`L/2`$, thus cylinders with $`L<2`$ are “small” and might present topological copies of discrete sources and/or circles in the sky. On the other hand, “large” cylinders, i.e., those with $`L>2`$, have undetectable topology with the methods currently available Detect .
The covering group of the cylinder is labeled by the integers as $`g^n=(I,n𝐚)`$, with $`n`$. We then have from (25) that all the topological information is encoded in
$$\mathrm{{\rm Y}}_\mathrm{}m^\mathrm{\Gamma }(𝐤)=\underset{n}{}e^{ink_zL}Y_\mathrm{}m(𝐧_𝐤),$$
and thus the correlation matrix of the $`a_\mathrm{}m`$’s for the cylinder is simply
$$a_\mathrm{}ma_\mathrm{}^{}m^{}^{}^\mathrm{\Gamma }=(4\pi )^2i^{\mathrm{}\mathrm{}^{}}\frac{d^3k}{k^3}\mathrm{\Psi }_{\mathrm{}\mathrm{}^{}}(k)\left(\underset{n}{}e^{ink_zL}\right)Y_\mathrm{}^{}m^{}(𝐧_𝐤)Y_\mathrm{}m^{}(𝐧_𝐤).$$
(31)
To reduce this integral we may use any of the following two identities, either
$$\underset{n}{}e^{ink_zL}=2\pi \underset{p}{}\delta (k_zL2\pi p),$$
(32)
or
$$\underset{n}{}e^{ink_zL}=1+2\underset{n=1}{\overset{\mathrm{}}{}}\mathrm{cos}(nk_zL).$$
(33)
The first identity is obvious since the left hand side is the Fourier expansion of the right hand side. This option yields a formula of the kind obtained in Appendix D, which expresses the correlation matrix in terms of the eigenvalues of the Laplacian operator on the cylinder. The second one still uses a parametrization in terms of the covering group, and thus can be used to isolate the topological signature.
In fact, using (33) to evaluate (31), and integrating in spherical coordinates, we get
$$a_\mathrm{}ma_\mathrm{}^{}m^{}^{}^\mathrm{\Gamma }=a_\mathrm{}ma_\mathrm{}^{}m^{}^{}^{s.c.}+a_\mathrm{}ma_\mathrm{}^{}m^{}^{}^{\widehat{\mathrm{\Gamma }}},$$
(34)
where the simply connected part is as usual
$$a_\mathrm{}ma_\mathrm{}^{}m^{}^{}^{s.c.}=C_{\mathrm{}}^{s.c.}\delta _{\mathrm{}\mathrm{}^{}}\delta _{mm^{}},$$
(35)
with the (simply connected) angular power spectrum given by
$$C_{\mathrm{}}^{s.c.}=(4\pi )^2_0^{\mathrm{}}\frac{dx}{x}\mathrm{\Psi }_{\mathrm{}\mathrm{}}(x),$$
(36)
and the topological signature for the correlation matrix is given by
$`a_\mathrm{}ma_\mathrm{}^{}m^{}^{}^{\widehat{\mathrm{\Gamma }}}`$ $`=`$ $`(4\pi )^2i^{\mathrm{}\mathrm{}^{}}\delta _{\mathrm{}\mathrm{}^{}}^{\text{mod(2)}}\delta _{mm^{}}\times `$
$`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dx}{x}}\mathrm{\Psi }_{\mathrm{}\mathrm{}^{}}\left({\displaystyle \frac{x}{L}}\right)F_{\mathrm{}\mathrm{}^{}}^m(x),`$
with
$$F_{\mathrm{}\mathrm{}^{}}^m(x)=2\underset{n=1}{\overset{\mathrm{}}{}}_1^1𝑑y\mathrm{cos}(nxy)𝒫_{\mathrm{}}^m(y)𝒫_{\mathrm{}^{}}^m(y),$$
(38)
where $`𝒫_{\mathrm{}}^m(x)`$ is the normalized associated Legendre function (see Appendix A). As expected, we have ended up with a real correlation matrix with factors $`\delta _{\mathrm{}\mathrm{}^{}}^{\text{mod(2)}}`$ and $`\delta _{mm^{}}`$.
After evaluating the series in (38), it turns out that $`F_{\mathrm{}\mathrm{}^{}}^m(x)`$ is a piecewise continuous function. In fact, in each interval $`[2\pi q,2\pi (q+1)]`$, it is a polynomial of degree $`(\mathrm{}+\mathrm{}^{}+1)`$ in $`\pi /x`$. Indeed, the final result is
$`F_{\mathrm{}\mathrm{}^{}}^m(x)`$ $`=`$ $`{\displaystyle \underset{q}{}}_{\mathrm{}\mathrm{}^{}}^m(x,q)\mathrm{\Theta }(x2\pi q)\times `$
$`\mathrm{\Theta }(2\pi (q+1)x),`$
where $`\mathrm{\Theta }(x)`$ is the Heaviside step function, and the form of the polynomial $`_{\mathrm{}\mathrm{}^{}}^m(x,q)`$ in the $`q`$–th interval of length $`2\pi `$ is
$`_{\mathrm{}\mathrm{}^{}}^m(x,q)`$ $`=`$ $`4{\displaystyle \underset{k=0}{\overset{\frac{\mathrm{}+\mathrm{}^{}}{2}}{}}}(1)^k𝒫_\mathrm{}\mathrm{}^{}m^{(2k)}(0)\times `$
$`g_{2k+1}(q)\left({\displaystyle \frac{\pi }{x}}\right)^{2k+1}\delta _{\mathrm{}\mathrm{}^{}}.`$
Here $`g_k(q)`$ are polynomials of degree $`k`$ in $`q`$, and $`𝒫_\mathrm{}\mathrm{}^{}m^{(k)}(0)`$ is the $`k`$–th derivative of the polynomial
$$𝒫_\mathrm{}\mathrm{}^{}m(x)=𝒫_{\mathrm{}}^m(x)𝒫_{\mathrm{}^{}}^m(x)$$
(41)
evaluated at the origin. In Appendix B we present recurrence relations for the polynomials $`g_k(q)`$, and all the technical steps that take us from (38) to (IV.1) can be found in Appendix C.
The integrals appearing in the topological signature (IV.1) can be easily evaluated since Eqs. (IV.1) and (IV.1) allow an exact and very fast computation of the function $`F_{\mathrm{}\mathrm{}^{}}^m(x)`$, and the integrands decay very fast, as illustrated in figs.1 and 2. In this figures we have adopted, for simplicity, a scale invariant Einstein–de Sitter model, thus
$$\mathrm{\Psi }_{\mathrm{}\mathrm{}^{}}(x)j_{\mathrm{}}(x)j_{\mathrm{}^{}}(x).$$
(42)
The nice behavior of the integrands in (IV.1) is not a consequence of this particular choice of $`\mathrm{\Psi }_{\mathrm{}\mathrm{}^{}}(x)`$. Actually, the integrand in (IV.1) always decays very fast because $`\mathrm{\Psi }_{\mathrm{}\mathrm{}^{}}(x)`$ and $`F_{\mathrm{}\mathrm{}^{}}^m(x)`$ are both decaying functions, thus the evaluation of the topological signature for the cylinder is always very efficient.
The computation of the topological signature of the power spectrum reduces to a simple integral. In fact we obtain
$$C_{\mathrm{}}^{\widehat{\mathrm{\Gamma }}}=(4\pi )^2_0^{\mathrm{}}\frac{dx}{x}\mathrm{\Psi }_{\mathrm{}\mathrm{}}\left(\frac{x}{L}\right)f_{\mathrm{}}(x),$$
with
$$f_{\mathrm{}}(x)=\frac{1}{2\mathrm{}+1}\underset{m=l}{\overset{\mathrm{}}{}}F_{\mathrm{}\mathrm{}}^m(x).$$
Using (38) to perform this sum, the Addition Theorem for Spherical Harmonics (see Appendix A) yields immediately
$$C_{\mathrm{}}^{\widehat{\mathrm{\Gamma }}}=2(4\pi )^2_0^{\mathrm{}}\frac{dx}{x^2}\mathrm{\Psi }_{\mathrm{}\mathrm{}}\left(\frac{x}{L}\right)\phi _1(x),$$
(43)
where $`\phi _1(x)`$ is the first Clausen $`\phi `$–function given in Appendix B.
In fig.3a we show the low $`\mathrm{}`$–modes of the topological signature of the angular power spectrum of a cylinder, normalized w.r.t. $`C_{\mathrm{}}^{s.c.}`$, as a function of its size $`L`$. We can see that the topological signature is typically much smaller than the cosmic variance, even for small cylinders which have already been discarded observationally as candidates for the shape of our Universe because of the lack of antipodal matched circles in WMAP data Aligne1 ; CSSK04 . Thus it is apparent that the angular power spectrum is not a good indicator to look for topology in this case.
The correlation matrix given by (34)–(IV.1) corresponds to a cylinder for which the direction of compactification is parallel to the $`z`$–axis. The correlation matrix corresponding to a cylinder with a different orientation can be easily obtained from the previous one by simply rotating the celestial sphere. Thus, parametrizing the rotations with Euler angles, if $`R(\alpha ,\beta ,\gamma )SO(3)`$ takes the $`z`$–axis to the direction of compactification of the cylinder, the topological signature of the corresponding correlation matrix can be computed using the expressions (53)–(55) of Appendix A yielding
$`a_\mathrm{}ma_\mathrm{}^{}m^{}^{}_R^{\widehat{\mathrm{\Gamma }}}`$ $`=`$ $`e^{i(m^{}m)\alpha }\times `$
$`{\displaystyle \underset{m_1}{}}d_{mm_1}^{\mathrm{}}(\beta )d_{m^{}m_1}^{\mathrm{}^{}}(\beta )\times `$
$`a_{\mathrm{}m_1}a_{\mathrm{}^{}m_1}^{}^{\widehat{\mathrm{\Gamma }}},`$
since $`a_\mathrm{}ma_\mathrm{}^{}m^{}^{}^{s.c.}`$ is rotationally invariant. Moreover, the $`\gamma `$ angle does not appear in this expression since $`R_z(\gamma )`$ in (54) does not move the $`z`$–axis, and $`a_\mathrm{}ma_\mathrm{}^{}m^{}^{}^\mathrm{\Gamma }`$ is invariant under such rotations.
It should be noted here that, no matter its orientation, the cylinder is always invariant under parity, thus its correlation matrix will always conserve the factor $`\delta _{\mathrm{}\mathrm{}^{}}^{\text{mod(2)}}`$. On the other hand, the correlation matrix will remain real as far as we perform rotations with $`\alpha =0`$, since in this case we do not rotate the cylinder around the $`z`$–axis, and thus it remains invariant under reflections on the plane $`y=0`$. However, any rotation of the cylinder (other than one with $`\beta =\pi `$) makes it non–invariant under azimuthal rotations, thus the correlation matrix of an arbitrarily oriented cylinder has $`m`$–dependent correlations. All these features can be seen explicitly in (IV.1).
### IV.2 tori
In order to calculate the correlation matrix of the $`a_\mathrm{}m`$’s for a two– or a three–torus we use the decomposition (29) of its covering group in cyclic subgroups. Let $`\mathrm{\Gamma }_{ij}=g_{ij}`$ be the covering group of the cylinder generated by the element $`g_{ij}\mathrm{\Gamma }`$, and let us write $`L_i=d(\mathrm{𝟎},g_{ij}\mathrm{𝟎})`$, $`g_i=(I,L_i\widehat{𝐞}_z)`$, and $`\mathrm{\Gamma }_i=g_i`$. In the Euclidean case, the orientation preserving isometries that do not move the observer are rotations, thus let $`R_{ij}SO(3)`$ be the rotation taking $`\widehat{𝐞}_z`$ to the unit vector along $`g_{ij}\mathrm{𝟎}`$.
Using the decomposition (30), we easily write the topological signature for the torus as a superposition of topological signatures of rotated cylinders. In fact,
$$a_\mathrm{}ma_\mathrm{}^{}m^{}^{}^{\widehat{\mathrm{\Gamma }}}=\underset{i=1}{\overset{\mathrm{}}{}}\underset{j=1}{\overset{k_i}{}}a_\mathrm{}ma_\mathrm{}^{}m^{}^{}_{R_{ij}}^{\widehat{\mathrm{\Gamma }}_i},$$
(45)
where the correlation matrices of the rotated cylinders are written in terms of the Wigner $`D`$–functions and Euler angles, according to (IV.1), as
$`a_\mathrm{}ma_\mathrm{}^{}m^{}^{}_{R_{ij}}^{\widehat{\mathrm{\Gamma }}_i}`$ $`=`$ $`e^{i(m^{}m)\alpha _{ij}}\times `$
$`{\displaystyle \underset{m_1}{}}d_{mm_1}^{\mathrm{}}(\beta _{ij})d_{m^{}m_1}^{\mathrm{}^{}}(\beta _{ij})\times `$
$`a_{\mathrm{}m_1}a_{\mathrm{}^{}m_1}^{}^{\widehat{\mathrm{\Gamma }}_i},`$
and $`(\beta _{ij},\alpha _{ij})`$ are the angular spherical coordinates of the vector $`g_{ij}\mathrm{𝟎}`$, and $`k_i`$ is the number of cylinders of size $`L_i`$. Since any group of translations is invariant under parity, from Sec. A.3 we know that the correlation matrix for a homogeneous flat manifold has always the factor $`\delta _{\mathrm{}\mathrm{}^{}}^{\text{mod(2)}}`$, and this is evident from (45), since it is just a sum of correlation matrices of cylinders.
The power spectrum is rotationally invariant, thus from (45) one can easily write down the expression for the topological signature of the power spectrum of the torus as a superposition of topological signatures of power spectra of cylinders,
$$C_{\mathrm{}}^{\widehat{\mathrm{\Gamma }}}=\underset{i=1}{\overset{\mathrm{}}{}}k_iC_{\mathrm{}}^{\widehat{\mathrm{\Gamma }}_i}.$$
(47)
Let us consider a chimney with square base for the sake of illustration. It is convenient to orient the chimney so that its covering group consists of translations in the horizontal plane. We take as generators of the covering group the translations $`g_1=(I,𝐚)`$ and $`g_2=(I,𝐛)`$, with $`𝐚=L\widehat{𝐞}_x`$ and $`𝐛=L\widehat{𝐞}_y`$.
It is more convenient to reparametrize the cyclic decomposition as follows. Parametrize each cyclic subgroup by a pair of integer numbers $`(p,q)`$ as $`G_{pq}=g_2^qg_1^p`$. Clearly, if the greatest common divisor of $`(p,q)`$ is $`r`$, then
$$G_{pq}<G_{\frac{p}{r}\frac{q}{r}},$$
where ‘$`<`$’ means ‘subgroup of’. Thus we must restrict the labels to pairs $`(p,q)`$ of coprime numbers.
The only exceptions are when (i) $`p=\pm 1`$ and $`q=0`$ and viceversa, and (ii) when $`p=\pm 1`$ and $`q=\pm 1`$. Thus the first two complete sets of cyclic subgroups conjugate by a rotation are $`\{G_{1,0},G_{0,1}\}`$ and $`\{G_{1,1},G_{1,1}\}`$. In both cases the conjugation is performed by a rotation of $`\pi /2`$ around the $`z`$–axis. The compactification lengths of the corresponding cylinders are $`L_{1,0}=L_{0,1}=L`$ and $`L_{1,1}=L_{1,1}=\sqrt{2}L`$ respectively. The Euler angles $`(\beta ,\alpha )`$ to rotate the corresponding cylinders from the $`z`$–axis to their orientation in the chimney, according to (IV.2), are $`\beta =\pi /2`$ in all cases, and $`\alpha _{1,0}=0`$, $`\alpha _{0,1}=\pi /2`$, $`\alpha _{1,1}=\pi /4`$ and $`\alpha _{1,1}=3\pi /4`$, respectively.
To write the remaining complete sets of cyclic subgroups conjugate by a rotation let us define, for a pair of coprime natural numbers $`(p,q)`$, with $`p>q1`$, the groups
$`G_{pq}^{(1)}=G_{pq}=g_2^qg_1^p`$ , $`G_{pq}^{(3)}=G_{q,p}=g_2^pg_1^q,`$
$`G_{pq}^{(2)}=G_{qp}=g_2^pg_1^q`$ , $`G_{pq}^{(4)}=G_{p,q}=g_2^qg_1^p.`$
The compactification lengths are all equal to $`L_{pq}=\sqrt{p^2+q^2}L`$, and the Euler angles $`(\beta ,\alpha )`$ to rotate the corresponding cylinders from the $`z`$–axis to their orientation in the chimney, according to (IV.2), are $`\beta =\pi /2`$ in all cases, and
$`\alpha _{pq}^{(1)}=\mathrm{arctan}{\displaystyle \frac{q}{p}}`$ , $`\alpha _{pq}^{(3)}={\displaystyle \frac{\pi }{2}}+\alpha _{pq}^{(1)},`$
$`\alpha _{pq}^{(2)}={\displaystyle \frac{\pi }{2}}\alpha _{pq}^{(1)}`$ , $`\alpha _{pq}^{(4)}=\pi \alpha _{pq}^{(1)},`$
respectively.
Let us denote by $`\mathrm{\Gamma }_{pq}`$ the covering group of the cylinder with compactification scale $`L_{pq}`$ and oriented along the $`z`$–axis. Then, putting all this together, using (45) and (IV.2), and taking into account the invariance properties derived in Sec. A.3, the topological signature of the chimney with square base is
$`a_\mathrm{}ma_\mathrm{}^{}m^{}^{}^{\widehat{\mathrm{\Gamma }}}`$ $`=`$ $`\delta _{mm^{}}^{\text{mod(4)}}{\displaystyle \underset{m_1}{}}d_{mm_1}^{\mathrm{}}(\pi /2)\times `$
$`d_{m^{}m_1}^{\mathrm{}^{}}(\pi /2)𝒲_{\mathrm{}\mathrm{}^{}m_1}^{m^{}m},`$
with
$$𝒲_{\mathrm{}\mathrm{}^{}m_1}^m=2\left(a_{\mathrm{}m_1}a_{\mathrm{}^{}m_1}^{}^{\widehat{\mathrm{\Gamma }}_{1,0}}+(1)^{m/4}a_{\mathrm{}m_1}a_{\mathrm{}^{}m_1}^{}^{\widehat{\mathrm{\Gamma }}_{1,1}}\right)+4\underset{(p,q)}{}\mathrm{cos}m\alpha _{pq}^{(1)}a_{\mathrm{}m_1}a_{\mathrm{}^{}m_1}^{}^{\widehat{\mathrm{\Gamma }}_{pq}},$$
(49)
where the sum in $`(p,q)`$ is evaluated only for pairs of coprime natural numbers $`(p,q)`$ such that $`p>q1`$.
The topological signature of the power spectrum of the chimney with square base is simply
$$C_{\mathrm{}}^{\widehat{\mathrm{\Gamma }}}=2\left(C_{\mathrm{}}^{\widehat{\mathrm{\Gamma }}_{1,0}}+C_{\mathrm{}}^{\widehat{\mathrm{\Gamma }}_{1,1}}\right)+4\underset{(p,q)}{}C_{\mathrm{}}^{\widehat{\mathrm{\Gamma }}_{pq}}.$$
(50)
Since the topological signature of the power spectrum of a cylinder converges quickly to zero as a function of the compactification scale (see fig.3a), it follows that the sum in (50) also converges quickly.
Moreover, the topological signature of the power spectrum of the chimney is larger than that of the cylinder. This is so because the $`\mathrm{}`$-th mode of the topological signature of the angular power spectrum of the cylinder oscillates very slowly. Thus from (50) this signature is slightly higher in the chimney, as can be seen in fig.3b. Actually, this is a general result that holds for manifolds whose covering groups are not cyclic.
## V Patterns of alignment
The nondiagonal character of the topological signature of the correlation matrix of the $`a_\mathrm{}m`$’s in multiply connected universes and their $`m`$–dependence are manifestations of their globally anisotropic nature. They manifest themselves in statistically anisotropic temperature maps, i.e., realizations of random temperature fluctuations for which mean values of functions of the temperature over ensembles of universes depend on the orientation StatAnis .
In this section we analyze an expected consequence of the topology of space on the temperature anisotropies of the CMB that has not received the deserved attention up to the present, namely the existence of preferred directions in space. We show that the decomposition of the topological signature of the correlation matrix of the $`a_\mathrm{}m`$’s in a universe with a complex topology, in signatures corresponding to cyclic topologies, demands the existence of “patterns of alignments” along these directions. For the sake of simplicity we consider the Einstein–de Sitter model, thus from now on we will take (42) to perform all our calculations.
We want to call attention to the existence of alignments of the low $`\mathrm{}`$–modes of the CMB temperature maps in multiply connected universes. Indeed, in fig.4 we show a low resolution temperature map simulation for a cylinder with $`L=2`$ (in units of $`R_{LSS}`$), together with the maps corresponding to the first four $`\mathrm{}`$–modes. One can see that these $`\mathrm{}`$–maps present alignments along the $`z`$–direction, which in this case is the unique direction of compactification of space.
Similar alignments as those present in our simulations have been reported as being observed in WMAP data, and have been attributed to a possible nontrivial topology of space with the shape of a cylinder Aligne1 . These models have been quickly abandoned due to the lack of circles in the sky which should be present if the Universe were small Aligne1 ; CSSK04 . However our simulations show that even in universes slightly larger, and so not presenting such circles, these alignments should still be observable. Thus whether these observed alignments are a consequence of a nontrivial shape of our Universe is still an open question LocShape .
We will show here that if our Universe had a nontrivial topology, its CMB temperature map will present characteristic patterns of alignment, even if its size is somewhat larger than the observable universe. Moreover, from the observed patterns of alignment, we might be able to reconstruct the shape of space.
In fig.5 we show the topological signature of the low $`\mathrm{}`$–modes of the diagonal part of the correlation matrix of the $`a_\mathrm{}m`$’s, normalized w.r.t. $`C_{\mathrm{}}^\mathrm{\Gamma }`$, for a cylinder oriented along the polar axis, as a function of the size $`L`$ of compactification. It is apparent that, for a given $`\mathrm{}`$–mode, there are multipole coefficients for which their expected values are above the mean (the angular power spectrum), and others for which these expected values are below it. This is the reason why the low $`\mathrm{}`$–modes in a cylinder are aligned. Actually, the expectation values $`|a_\mathrm{}m|^2`$ are all equal to $`C_{\mathrm{}}`$ only in the simply connected case, thus the dispersion around the mean
$$\sigma _{\mathrm{}}=\sqrt{\frac{1}{2\mathrm{}+1}\underset{m}{}\left(|a_\mathrm{}m|^2C_{\mathrm{}}\right)^2}$$
(51)
is null. However, in a multiply connected universe, this dispersion is non zero, and in a particular map, it adds to the cosmic variance. Thus it seems natural to propose the dispersion of the squares $`|a_\mathrm{}m|^2`$ around their mean value as a measure of these alignments in a map.
In fig.6 we show a plot of the dispersion (51), normalized w.r.t. $`C_{\mathrm{}}^\mathrm{\Gamma }`$, for a cylinder oriented along the polar axis, as a function of $`L`$, for low multipoles. Note that even for a large cylinder ($`L2`$) the dispersion is larger than $`15\%`$ of the power for multipoles up to $`\mathrm{}=5`$. Indeed, on these scales the dispersion is of the order of the cosmic variance, and thus might be detectable.
In order to show that this is a good measure of the alignment of multipoles, and that it provides an efficient method to determine the directions of possible alignments in real or simulated maps, let us compute the dispersion of the squares $`|a_\mathrm{}m|^2`$, Eq. (51), for a cylinder which is oriented along a direction making an angle $`\beta `$ with the $`z`$–axis. Each one of these squares can be computed with
$$|a_\mathrm{}m|^2_R^\mathrm{\Gamma }=\underset{m_1}{}\left[d_{mm_1}^{\mathrm{}}(\beta )\right]^2|a_{\mathrm{}m_1}|^2^\mathrm{\Gamma },$$
which is nothing but (IV.1) restricted to the diagonal part.
In fig.7 it is shown this dispersion as a function of $`\beta `$ for different multipole coefficients and for different values of $`L`$. One can see that the dispersion has a maximum when the cylinder is oriented along the $`z`$–axis. Thus in order to look for the alignments in a hypothetical universe with the shape of a cylinder, one should just rotate the celestial sphere around different directions until find those two opposite ones along which the dispersion of the squares $`|a_\mathrm{}m|^2`$ is maximum.
However, in order to collect definitive evidence that the universe is indeed a cylinder, one should map the dispersion of the squares $`|a_\mathrm{}m|^2`$ on the sphere for each $`\mathrm{}`$–mode, i.e. one should determine the dispersion (51) as a function of the orientation of the celestial sphere. If the universe had the topology of a cylinder, these dispersion maps should be axially symmetric around a special direction, where the dispersion is maximum. Moreover, this direction should be identified with the direction of compactification of the cylinder.
If the universe has the topology of a flat homogeneous manifold note, from (45) and (IV.2), that the topological signature is a superposition of rotated cylinders of different sizes. Thus a CMB map for a universe with this kind of topology might present alignments along the directions corresponding to these cylinders. In fact, rotating the celestial sphere and computing the dispersion of the squares $`|a_\mathrm{}m|^2`$, an easy computation shows that whenever we perform the rotation $`R(0,\theta ,\phi )`$, with $`\theta =\beta _{ij}`$ and $`\phi =\alpha _{ij}`$, one has the cylinder labeled by $`(i,j)`$ oriented along the polar axis, and thus *dispersion maps might present local maxima along these directions*.
Whether these local maxima are observable in a given dispersion map will depend on (i) the scale of compactification of the corresponding cylinder $`L_i`$, (ii) the background due to the simply connected part, and (iii) the other cylinders’ topological signatures. For large values of $`L_i`$ the corresponding local maxima will not be observable, however one can expect those maxima corresponding to the smaller cylinders to be detectable. The existence and distribution of these maxima in each dispersion map, together with their relative intensities is what we call a *pattern of alignment*.
It might seem that the problem of constructing dispersion maps for manifolds that are not flat homogeneous is more involved, since general cyclic manifolds do not have axial symmetry as the cylinder has. Eq. (IV.1) depends on two angles only because the cylinder is axially symmetric, but in the general case the expression for the correlation matrix in a rotated frame depends on the three Euler angles. Thus it seems at first sight that in these cases, a dispersion map should be a function on the 3–sphere. Fortunately, the diagonal elements of the rotated correlation matrix depend only on the last two Euler angles as
$`|a_\mathrm{}m|^2_R^\mathrm{\Gamma }`$ $`=`$ $`{\displaystyle \underset{m_1,m_2}{}}e^{i(m_2m_1)\gamma }\times `$
$`d_{mm_1}^{\mathrm{}}(\beta )d_{mm_2}^{\mathrm{}}(\beta )\times `$
$`a_{\mathrm{}m_1}a_{\mathrm{}m_2}^{}^\mathrm{\Gamma },`$
thus the same conclusion holds in the general case. *Dispersion maps on the 2–sphere for low $`\mathrm{}`$–modes should display patterns of alignment showing the symmetries of our Universe if it has a (not too large) nontrivial topology*.
## VI Discussion
In order to study systematically the effects of a nontrivial spatial topology in the temperature fluctuations of the CMB, we need to have the ability to simulate efficiently temperature maps in multiply connected $`\mathrm{\Lambda }`$CDM cosmologies. Almost all the usual methods to perform these simulations use explicitly the solutions of the Helmholtz equation in 3–manifolds with nontrivial topology. The computation of the eigenfunctions and eigenvalues of the Laplacian operator is simple only in Euclidean manifolds, while in spherical and hyperbolic spaces it is a nontrivial problem. In fact, it is only recently that an analytical computation has been achieved for all the spherical manifolds. The hyperbolic cases still have to be done numerically.
In this paper we have developed a simulation procedure that avoids the explicit use of the solutions of the Helmholtz equation. Instead, our results are expressed in terms of the covering group $`\mathrm{\Gamma }`$ of the corresponding manifold. In this section we summarize the details of the method, its efficiency, the simple applications performed here, and discuss future related work.
### VI.1 The formalism
The cornerstone of our method is formula (16), which is the two–point correlation function of the scalar perturbations in a multiply connected universe expressed in terms of the covering group of the manifold BPS . By means of simple formal manipulations we obtain an expression for the correlation matrix of the spherical harmonic expansion coefficients of the temperature maps, Eqs. (II)–(25), which contain all the topological information expressed as a sum over the elements of the covering group.
Former applications of (16) required a regularization procedure in order to account for divergences of the series, as well as some resummation techniques for accelerating the convergence. We do not have these problems here because the divergent series, which are actually distributions, appear only inside integrals. Indeed, on the one hand, we show in Appendix D that our formalism easily reproduces results previously reported in the literature, as well as some simple generalizations, without the need of any regularization procedure. On the other hand, elementary decompositions of the two–point correlation function (16), shown in Sec. III, guarantee that our final expressions are highly convergent, as discussed below.
Two decompositions of a generic covariance function which can be written as a sum over the covering group are crucial for the efficiency of our formalism. The first one, a trivial decomposition given by (26), defines the topological signature of the covariance function. When written for the correlation matrix of the harmonic expansion coefficients, it yields the topological signature in the temperature anisotropy maps, as illustrated for the cylinder by (34). This expression shows that the topological signature is nothing but a “perturbation” of the correlation matrix corresponding to the simply connected case. Since, as discussed in Sec. IV.1, these “perturbations” are small, the efficiency of the calculation follows.
The second decomposition given by (30) allows us to write the topological signature of any manifold in terms of the topological signatures of its maximal covering manifolds with cyclic covering groups. The example of the tori illustrates the power of this approach, since we can write down explicit formulae for a general torus whether its generating translations are orthogonal and/or equal. Trying to do this with the explicit use of eigenfunctions of the Laplacian (or with the method used in Appendix D) turns out to be tedious if not difficult.
Moreover, by construction, this decomposition is invariant under the symmetries of the manifold, thus it carries information on how shall these symmetries manifest in individual CMB temperature anisotropy maps, as will be discussed in Sec. VI.2.
Another advantage of this second decomposition is the simplicity for writing down the power spectrum for complicated manifolds. Expressions like (47) and (50) are computationally very efficient once we have saved the power spectrum for cyclic manifolds as a function of its scale of compactification, since we have just to perform a weighted sum of power spectra for cyclic manifolds at different scales considering the multiplicity of the decomposition.
### VI.2 Topological signatures
A further advantage of splitting the correlation matrix of the multipole coefficients into its simply connected part and its topological signature is that we can identify very easily the geometric features of the signature. Although we have made explicit calculations for flat homogeneous manifolds only, qualitatively these results are general.
Universes with cylindrical topology of size $`L2`$ present clear alignments of their low $`\mathrm{}`$–modes along the direction of compactification. A dispersion map of the squares $`|a_\mathrm{}m|^2`$, for a given low $`\mathrm{}`$, exhibits an axial symmetry around this direction, thus it reduces to a function of the polar angle. These dispersion maps are shown in fig.7.
By decomposing the covering group $`\mathrm{\Gamma }`$ in cyclic subgroups one can see that, whatever the shape of our Universe, and if it is not too large, dispersion maps (one for each individual low multipole) might show patterns of alignment. In the general case such maps are functions on the two–dimensional projective space or, by a lifting, on the 2–sphere. Although we have shown the existence of patterns of alignment explicitly only for homogeneous flat manifolds, it follows from the exposition of the general formalism that the same conclusions hold for any manifold of constant curvature. Thus, we propose the construction of these dispersion maps in the WMAP data, and so the search for patterns of alignment, as a new method for detecting a possible nontrivial topology of our Universe.
It is interesting to comment on some features relating Levin and collaborators’ proposal of pattern formation in CMB temperature maps and the results we present in this paper. The patterns proposed by Levin et al. LevinB ; LSS98 are due to individual eigenmodes ($`k`$–modes), the patterns we have identified here are due to multipole modes ($`\mathrm{}`$–modes). In either case the modes compete to form their patterns in a CMB temperature map, however the observable modes in a map on the sphere are the latter, since spherical harmonics form a base on the space of functions on the sphere. On the other side, the association between real space perturbations and angular temperature fluctuations requires some averaging over the $`k`$–modes Inoue03 , thus these patterns appear mixed in a map and their observation might demand more elaborated techniques.
### VI.3 Further remarks and future research
The formalism we have developed in this paper reveals new insights on the problem of characterizing the marks that topology leaves in CMB maps, and opens up new possibilities for developing further methods for unvealing the shape of our Universe. It makes explicit that the multipole alignments observed in COBE and WMAP full sky CMB temperature maps may be a manifestation of its global shape, provides details of the nature and features of these alignments, and gives at least one methodology to test this hypothesis. As a consequence, further work is much needed.
One line of further research is the implementation of our formalism in the spherical and hyperbolic cases. One way to do this requires first to identify the radial part of the fundamental solution (10) of the Helmholtz equation in the universal covering, as well as the analog of the “plane wave expansion” solution (17) in these geometries, and to write the expansion of the corresponding “plane waves” in spherical harmonics as in (18). The difficult part seems to be expressing the “plane wave expansion” in a suitable form to reproduce the formal steps used in the Euclidean case.
Moreover we have to compute the topological signature of all other cyclic manifolds, in order to extend the computations to any quotient space that could be a candidate for the shape of space. We also need to include acoustic oscillations, and Doppler and finite width effects in $`\mathrm{\Lambda }`$CDM models so that we could determine the relevant angular scales in Cosmic Topology, i.e. the angular scales at which the topological signatures appear. This is a crucial step in order to confront quantitatively the theory with real CMB maps in an efficient and rigurous way. An ultimate goal may be to implement all this methodology in a software package for public use.
The identification of the “topological signature” of the correlation matrix $`a_\mathrm{}ma_\mathrm{}^{}m^{}^{}^\mathrm{\Gamma }`$ also opens up a path for solving a problem raised by Riazuelo et al. in RULW04 . The correlation matrix for a multiply connected universe is non diagonal and, tipically, $`m`$–dependent. In fact, this is the source of the statistical anisotropy of the CMB in these universes. However, for very large manifolds this correlation matrix becomes effectively diagonal, and equal to that corresponding to the universal covering counterpart. A natural question raises up: at what typical scales does the correlation matrix “becomes diagonal”? In terms of our formalism this problem can be stated as finding the scales at where the topological signature becomes observationally negligible compared to the simply connected part. A closed analysis of the topological signature might give some answers to this and related questions. For example, establishing bounds on the integral in (IV.1) might solve the problem for the cylinder.
## Acknowledgments
We wish to thank Cristiane Camilo Hernandez for her unvaluable help with fig.4, Wanderson Wanzeller for his continuous help in computational issues, and Carlos Alexandre Wensche for showing us the papers Aligne1 which triggered our interest in this topic. We would also like to thank the Brazilian federal institutions CBPF and INPE for warm hospitality in several ocassions, and to the participants of the Seminar of Cosmic Topology held monthly at IFT and the Workshop New Physics from Space held every year in Campos do Jordão, São Paulo, where we had lots of opportunities to discuss this work at the several stages of its development during the last two years. W.S. Hipólito–Ricaldi acknowledges CAPES and G.I. Gomero aknowledges FAPESP (contract 02/12328-6) for financial support.
## Appendix A Spherical harmonics
In order to be self contained and to set the notation used in the paper, in this appendix we present basic definitions, some useful formulae of spherical harmonic functions and Wigner rotation matrices, as well as some invariance properties of the correlation matrix $`a_\mathrm{}ma_\mathrm{}^{}m^{}^{}`$ under coordinate transformations. For a complete treatment of spherical harmonic functions the reader can consult VMK .
### A.1 Basic definitions
Let us denote by $`𝐧=(\theta ,\phi )`$ a point in a 2–sphere parametrized in the usual spherical coordinates, then the spherical harmonic functions are defined as
$$Y_\mathrm{}m(𝐧)=\sqrt{\frac{2\mathrm{}+1}{4\pi }\frac{(\mathrm{}m)!}{(\mathrm{}+m)!}}P_{\mathrm{}}^m(\mathrm{cos}\theta )e^{im\phi },$$
where
$$P_{\mathrm{}}^m(x)=(1)^m\left(1x^2\right)^{m/2}\frac{d^m}{dx^m}P_{\mathrm{}}(x)$$
are the associated Legendre functions with non–negative index $`0m\mathrm{}`$, and with
$$P_{\mathrm{}}(x)=\frac{1}{2^{\mathrm{}}\mathrm{}!}\frac{d^{\mathrm{}}}{dx^{\mathrm{}}}\left(x^21\right)^{\mathrm{}}$$
being the Legendre polynomials. The associated Legendre functions with negative index $`m`$ are defined by
$$P_{\mathrm{}}^m(x)=(1)^m\frac{(\mathrm{}m)!}{(\mathrm{}+m)!}P_{\mathrm{}}^m(x).$$
Moreover, it is often convenient to introduce the normalized associated Legendre functions
$$𝒫_{\mathrm{}}^m(x)=\sqrt{\frac{2\mathrm{}+1}{2}\frac{(\mathrm{}m)!}{(\mathrm{}+m)!}}P_{\mathrm{}}^m(x).$$
It can easily be seen that the Legendre polynomial $`P_{\mathrm{}}(x)`$ is an $`\mathrm{}`$–th degree polynomial of parity $`\mathrm{}`$, and thus the associated Legendre function $`P_{\mathrm{}}^m(x)`$ is a function of parity $`\mathrm{}m`$. It follows that the function $`𝒫_\mathrm{}\mathrm{}^{}m(x)`$ defined in (41) is an $`(\mathrm{}+\mathrm{}^{})`$–th degree polynomial of parity $`\mathrm{}+\mathrm{}^{}`$, and thus the expression for $`F_{\mathrm{}\mathrm{}^{}}^m(x)`$ in (38), which is evaluated in Appendix C, contains only even polynomials.
Spherical harmonics form a complete orthonormal set of functions on the sphere, thus their most common application is in the expansion of functions, like a CMB temperature anisotropy map, in multipoles as in (11), where the coefficients $`a_\mathrm{}m`$, called the multipole coefficients, are given by
$$a_\mathrm{}m=_{𝕊^2}𝑑\mathrm{\Omega }\frac{\delta T}{T}(𝐧)Y_\mathrm{}m^{}(𝐧).$$
Since the temperature map is a real function on the sphere, the multipole coefficients obey the constraint $`a_\mathrm{}m^{}=(1)^ma_{\mathrm{},m}`$.
A very useful formula is given by the Addition Theorem for Spherical Harmonics
$$P_{\mathrm{}}(𝐧𝐧^{})=\frac{4\pi }{2\mathrm{}+1}\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}Y_\mathrm{}m(𝐧)Y_\mathrm{}m^{}(𝐧^{}),$$
(52)
which for the particular case $`𝐧=𝐧^{}`$ yields the identity
$$\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}\left[𝒫_{\mathrm{}}^m(x)\right]^2=\frac{2\mathrm{}+1}{2}.$$
### A.2 Wigner $`D`$–functions
In several ocassions it is convenient to rotate the sphere and compute the multipole coefficients in this new coordinate system. This can be achieved by means of the Wigner $`D`$–functions which can be defined operationally as the functions $`D_{mm_1}^{\mathrm{}}(R)`$ such that, for any rotation $`RSO(3)`$, then
$$Y_\mathrm{}m(R𝐧)=\underset{m_1}{}D_{mm_1}^{\mathrm{}}(R)Y_{\mathrm{}m_1}(𝐧).$$
In this case, it can be shown that the multipole coefficients of the temperature anisotropy map in the rotated reference frame are
$$\stackrel{~}{a}_\mathrm{}m=\underset{m_1}{}D_{mm_1}^{\mathrm{}}(R)a_{\mathrm{}m_1}.$$
This expression can be used to compute the correlation matrix of the $`a_\mathrm{}m`$’s in a rotated frame simply as
$$a_\mathrm{}ma_\mathrm{}^{}m^{}^{}_R=\underset{m_1,m_1^{}}{}D_{mm_1}^{\mathrm{}}(R)D_{m^{}m_1^{}}^{\mathrm{}^{}}(R)a_{\mathrm{}m_1}a_{\mathrm{}^{}m_1^{}}^{}.$$
(53)
The Wigner $`D`$–functions take a very simple form when we express the rotation matrix $`R`$ in terms of its Euler angles as
$$R(\alpha ,\beta ,\gamma )=R_z(\alpha )R_y(\beta )R_z(\gamma ).$$
(54)
Indeed, for this decomposition we have
$$D_{mm^{}}^{\mathrm{}}(R(\alpha ,\beta ,\gamma ))=e^{i(m\alpha +m^{}\gamma )}d_{mm^{}}^{\mathrm{}}(\beta ),$$
(55)
where $`d_{mm^{}}^{\mathrm{}}(\beta )=D_{mm^{}}^{\mathrm{}}(R_y(\beta ))`$ is a real matrix with the following symmetries
$`d_{mm^{}}^{\mathrm{}}(\beta )`$ $`=`$ $`(1)^{mm^{}}d_{m^{}m}^{\mathrm{}}(\beta ),`$
$`d_{mm^{}}^{\mathrm{}}(\beta )`$ $`=`$ $`d_{m^{},m}^{\mathrm{}}(\beta ),`$
$`d_{mm^{}}^{\mathrm{}}(\pi \beta )`$ $`=`$ $`(1)^\mathrm{}m^{}d_{m,m^{}}^{\mathrm{}}(\beta ),`$
$`d_{mm^{}}^{\mathrm{}}(\beta )`$ $`=`$ $`(1)^{m^{}m}d_{m,m^{}}^{\mathrm{}}(\beta ).`$
There exist several explicit and recursive formulae to compute these matrices (see VMK ). A very efficient recursive procedure can be found in BFB97 . The following formula will be enough to reproduce the results presented in this paper.
$$d_{mm^{}}^{\mathrm{}}(\beta )=\sqrt{(\mathrm{}+m)!(\mathrm{}m)!(\mathrm{}+m^{})!(\mathrm{}m^{})!}\underset{k}{}(1)^k\frac{\left(\mathrm{cos}\frac{\beta }{2}\right)^{2\mathrm{}2k+mm^{}}\left(\mathrm{sin}\frac{\beta }{2}\right)^{2km+m^{}}}{k!(\mathrm{}+mk)!(\mathrm{}m^{}k)!(m^{}m+k)!},$$
where the sum in $`k`$ is evaluated whenever the arguments inside the factorials are non–negative.
### A.3 Symmetry considerations
Some consequences of the symmetries of the quotient manifold on the invariance structure of the correlation matrix of the $`a_\mathrm{}m`$’s can be deduced directly from the transformation rules of the spherical harmonic functions under coordinate transformations. The results obtained in this way are formal, generic, and are very useful in practical computations. We end this Appendix by deducing the invariance properties the correlation matrix must have, given some symmetries of the corresponding quotient manifold. These invariance properties have been used in RULW04 to simplify the correlation matrix for the 3–torus, however we want to remark here that they are general and do not depend on the geometry of the universal covering space.
Let us begin with the invariance properties of $`a_\mathrm{}ma_\mathrm{}^{}m^{}^{}^\mathrm{\Gamma }`$ under rotations around the $`z`$–axis. Under a rotation $`R_z(\alpha ):\phi \phi +\alpha `$, the function $`Y_\mathrm{}m(𝐧)`$ transforms as
$$Y_\mathrm{}m(R_z(\alpha )𝐧)=e^{im\alpha }Y_\mathrm{}m(𝐧).$$
As a consequence the transformation rules for the multipole coefficients of a CMB temperature map are of the form $`\stackrel{~}{a}_\mathrm{}m=e^{im\alpha }a_\mathrm{}m`$, and so the correlation matrix transforms under this rotation as
$$a_\mathrm{}ma_\mathrm{}^{}m^{}^{}_{R_z(\alpha )}^\mathrm{\Gamma }=e^{i(m^{}m)\alpha }a_\mathrm{}ma_\mathrm{}^{}m^{}^{}^\mathrm{\Gamma }.$$
(56)
We extract two consequences out of (56). First, if the quotient space is invariant under a rotation of $`\alpha =2\pi /s`$ around the $`z`$–axis, then the correlation matrix must be zero unless $`m=m^{}`$ mod $`s`$. Second, if the quotient space is invariant under “any” rotation around the $`z`$–axis, the correlation matrix must be zero unless $`m=m^{}`$. In practice, if we take our coordinate system such that the fundamental polyhedron of the quotient manifold is oriented so that it is invariant under a $`2\pi /s`$ rotation around the polar axis, the correlation matrix will present a factor $`\delta _{mm^{}}^{\text{mod(s)}}`$, and correspondingly, if the orientation is such that the polyhedron is invariant under arbitrary rotations around the $`z`$–axis, the correlation matrix will present a factor $`\delta _{mm^{}}`$.
Let us now take a look at invariance under the inversion transformation $`P:𝐧𝐧`$. Under this transformation the spherical harmonic functions change as
$$Y_\mathrm{}m(P𝐧)=(1)^{\mathrm{}}Y_\mathrm{}m(𝐧),$$
thus the multipole coefficients $`a_\mathrm{}m`$ change as $`\stackrel{~}{a}_\mathrm{}m=(1)^{\mathrm{}}a_\mathrm{}m`$, and as a consequence the transformation rule for the correlation matrix is
$$a_\mathrm{}ma_\mathrm{}^{}m^{}^{}_P^\mathrm{\Gamma }=(1)^\mathrm{}+\mathrm{}^{}a_\mathrm{}ma_\mathrm{}^{}m^{}^{}^\mathrm{\Gamma }.$$
(57)
Thus, if the fundamental polyhedron is oriented such that it appears invariant under the parity transformation, the correlation matrix must be zero unless $`\mathrm{}=\mathrm{}^{}`$ mod 2, i.e., the correlation matrix will present a factor $`\delta _{\mathrm{}\mathrm{}^{}}^{\text{mod(2)}}`$.
To end this section, let us consider the reflection on the $`y=0`$ plane. This operation changes only the azimuthal angle as $`P_y:\phi \phi `$, thus the transformation rule for the spherical harmonics are
$$Y_\mathrm{}m(P_y𝐧)=Y_\mathrm{}m^{}(𝐧),$$
the multipole coefficients $`a_\mathrm{}m`$ change as $`\stackrel{~}{a}_\mathrm{}m=a_\mathrm{}m^{}`$, and thus, the transformation rule for the correlation matrix is
$$a_\mathrm{}ma_\mathrm{}^{}m^{}^{}_{P_y}^\mathrm{\Gamma }=a_\mathrm{}ma_\mathrm{}^{}m^{}^{}^\mathrm{\Gamma }.$$
(58)
It immediately follows that if the fundamental polyhedron is oriented such that it appears invariant under the reflection on the $`y=0`$ plane, the correlation matrix must be real.
## Appendix B Clausen functions
In this appendix we briefly present some computational aspects of the theory of Clausen functions, as far as we need them for our purposes. Clausen functions are periodic functions of period $`2\pi `$. There are two kinds of Clausen functions, the $`\phi `$–class and the $`\psi `$–class. Clausen $`\phi `$–functions can be expressed in terms of polynomials, while Clausen $`\psi `$–functions involve higher transcendental functions, the so–called Clausen integrals. Fortunately, we are interested exclusively in the Clausen $`\phi `$–functions, thus we will develop the details of the theory only for them. The Clausen $`\phi `$–functions are defined as
$`\phi _{2s1}(x)`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{sin}nx}{n^{2s1}}}`$
$`\phi _{2s}(x)`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{cos}nx}{n^{2s}}}`$
for $`s=1,2,\mathrm{}`$, and can be calculated recursively with the formulae,
$`\phi _{2s}(x)`$ $`=`$ $`\zeta (2s){\displaystyle _0^x}\phi _{2s1}(y)𝑑y`$
$`\phi _{2s+1}(x)`$ $`=`$ $`{\displaystyle _0^x}\phi _{2s}(y)𝑑y,`$
where
$$\zeta (s)=\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n^s}$$
is the Riemann Zeta function.
These recurrence relations are complemented by the initial condition
$`\phi _1(x)`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{sin}nx}{n}}`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{q}{}}[(2q+1)\pi x]\times `$
$`\mathrm{\Theta }(x2\pi q)\mathrm{\Theta }(2\pi (q+1)x),`$
where $`\mathrm{\Theta }(x)`$ is the Heaviside step function. Formula (B) can be verified by computing the Fourier series of the second right hand side.
Since the Clausen functions are periodic of period $`2\pi `$, we can write
$`\phi _s(x)`$ $`=`$ $`{\displaystyle \underset{q}{}}f_s(x2\pi q)\times `$
$`\mathrm{\Theta }(x2\pi q)\mathrm{\Theta }(2\pi (q+1)x),`$
with $`f_1(x)=\frac{1}{2}(\pi x)`$. The recurrence formulae (B) yield the following expressions for the Clausen functions in the period $`[0,2\pi ]`$,
$`f_{2s+1}(x)`$ $`=`$ $`{\displaystyle \underset{r=0}{\overset{s1}{}}}{\displaystyle \frac{(1)^r}{(2r+1)!}}\zeta (2(sr))x^{2r+1}+`$
$`{\displaystyle \frac{(1)^s}{2}}\left({\displaystyle \frac{\pi x^{2s}}{(2s)!}}{\displaystyle \frac{x^{2s+1}}{(2s+1)!}}\right)`$
for $`s=0,1,2,\mathrm{}`$, and
$`f_{2s}(x)`$ $`=`$ $`{\displaystyle \underset{r=0}{\overset{s1}{}}}{\displaystyle \frac{(1)^r}{(2r)!}}\zeta (2(sr))x^{2r}+`$
$`{\displaystyle \frac{(1)^s}{2}}\left({\displaystyle \frac{\pi x^{2s1}}{(2s1)!}}{\displaystyle \frac{x^{2s}}{(2s)!}}\right)`$
for $`s=1,2,3,\mathrm{}`$.
From the definitions (B) we get $`f_{2s+1}(\pi )=0`$, which can be used to obtain a recurrence formula for the Riemann Zeta function of even argument,
$`\zeta (2s)`$ $`=`$ $`{\displaystyle \underset{r=1}{\overset{s1}{}}}{\displaystyle \frac{(1)^{r+1}}{(2r+1)!}}\zeta (2(sr))\pi ^{2r}`$
$`{\displaystyle \frac{(1)^ss}{(2s+1)!}}\pi ^{2s}.`$
Writing $`\zeta (2s)=g_{2s}(0)\pi ^{2s}`$, and substituting this into (B) we have
$$g_{2s}(0)=\underset{r=1}{\overset{s1}{}}\frac{(1)^{r+1}}{(2r+1)!}g_{2(sr)}(0)\frac{(1)^ss}{(2s+1)!}.$$
The convenience for introducing this notation will be apparent in what follows.
We will now seek for generalizations of the formulae (B) and (B), i.e., we look for explicit expressions for the Clausen functions in the $`q`$–th interval $`[2\pi q,2\pi (q+1)]`$. Since the Clausen functions satisfy the periodicity conditions $`\phi _{2s1}(2\pi q)=0`$ and $`\phi _{2s}(2\pi q)=\zeta (2s)`$, the recurrence relations (B) can be rewritten in the form
$`\phi _{2s}(x)`$ $`=`$ $`\zeta (2s){\displaystyle _{2\pi q}^x}\phi _{2s1}(y)𝑑y`$
$`\phi _{2s+1}(x)`$ $`=`$ $`{\displaystyle _{2\pi q}^x}\phi _{2s}(y)𝑑y,`$
Defining the polynomials $`f_s^q(x)=f_s(x2\pi q)`$, we notice that $`\phi _s(x)`$ coincides with $`f_s^q(x)`$ in the interval $`[2\pi q,2\pi (q+1)]`$. This fact, and the expressions (B), allow us to write recurrence formulae analog to (B) for the polynomials $`f_s^q(x)`$ as follows
$`f_{2s}^q(x)`$ $`=`$ $`g_{2s}(q)\pi ^{2s}{\displaystyle _0^x}f_{2s1}^q(y)𝑑y,`$
$`f_{2s+1}^q(x)`$ $`=`$ $`g_{2s+1}(q)\pi ^{2s+1}+{\displaystyle _0^x}f_{2s}^q(y)𝑑y,`$
where
$`g_{2s}(q)`$ $`=`$ $`g_{2s}(0)+{\displaystyle \frac{1}{\pi ^{2s}}}{\displaystyle _0^{2\pi q}}f_{2s1}^q(y)𝑑y,`$
$`g_{2s+1}(q)`$ $`=`$ $`{\displaystyle \frac{1}{\pi ^{2s+1}}}{\displaystyle _0^{2\pi q}}f_{2s}^q(y)𝑑y,`$
with initial conditions, given by the first Clausen function, $`f_1^q(x)=g_1(q)\pi \frac{x}{2}`$ and $`g_1(q)=q+\frac{1}{2}`$.
The expressions (LABEL:q-PerClausRecurse) can be written in a unified way as
$$f_s^q(x)=g_s(q)\pi ^s(1)^s_0^xf_{s1}^q(y)𝑑y.$$
Using this expression we readily obtain the explicit formula, which is the generalization of (B) and (B) we were looking for,
$`f_s^q(x)`$ $`=`$ $`{\displaystyle \underset{r=0}{\overset{s1}{}}}{\displaystyle \frac{(1)^{\mu (r,s)}}{r!}}g_{sr}(q)\pi ^{sr}x^r`$
$`{\displaystyle \frac{(1)^{\mu (s,1)}}{2}}{\displaystyle \frac{x^s}{s!}},`$
where
$$\mu (r,s)=\frac{r}{2}+\frac{1+(1)^s}{4},$$
and $`x`$ is the floor function of $`x`$, i.e., the largest integer smaller than $`x`$.
The expressions (LABEL:q-gRecurse) can also be written in a unified way as
$$g_s(q)=g_s(0)+\frac{(1)^s}{\pi ^s}_0^{2\pi q}f_{s1}^q(y)𝑑y,$$
where
$$g_s(0)=\{\begin{array}{cc}\frac{\zeta (s)}{\pi ^s}\hfill & \text{if }s\text{ is even}\hfill \\ 0\hfill & \text{if }s>1\text{ is odd}.\hfill \end{array}$$
From this we get the expression analogous to (B)
$$g_s(q)=g_s(0)+(1)^s\left[\underset{r=1}{\overset{s1}{}}(1)^{\mu (r1,s1)}\frac{2^r}{r!}g_{sr}(q)q^r(1)^{\mu (s1,1)}\frac{2^{s1}}{s!}q^s\right].$$
(69)
The polynomials $`g_s(q)`$ can also be written in the canonical form
$$g_s(q)=\underset{k=0}{\overset{s}{}}A_k^sq^k,$$
(70)
where the coefficients are given by $`A_0^s=g_s(0)`$,
$$A_n^s=(1)^s\underset{r=1}{\overset{n}{}}(1)^{\mu (r1,s1)}\frac{2^r}{r!}A_{nr}^{sr}$$
for $`0<n<s`$, and
$`A_s^s`$ $`=`$ $`(1)^s[{\displaystyle \underset{r=1}{\overset{s1}{}}}(1)^{\mu (r1,s1)}{\displaystyle \frac{2^r}{r!}}A_{sr}^{sr}`$
$`(1)^{\mu (s1,1)}{\displaystyle \frac{2^{s1}}{s!}}],`$
with initial conditions $`A_0^1=\frac{1}{2}`$ and $`A_1^1=1`$. These coefficients are obtained by just introducing (70) into (69) and collecting terms.
## Appendix C The function $`F_{\mathrm{}\mathrm{}^{}}^m(x)`$
In this Appendix we evaluate the function $`F_{\mathrm{}\mathrm{}^{}}^m(x)`$ given by (38). We first observe (see Appendix A) that the function $`𝒫_\mathrm{}\mathrm{}^{}m(x)`$, given by (41), is an even polynomial of $`(\mathrm{}+\mathrm{}^{})`$–degree. Thus, we begin by considering the integral
$$I(\alpha )=_1^1P(y)\mathrm{cos}\alpha ydy,$$
where $`P(y)`$ is an even analytical function. Integrating succesively by parts we get
$`I(\alpha )`$ $`=`$ $`2[{\displaystyle \frac{\mathrm{sin}\alpha }{\alpha }}{\displaystyle \underset{s=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^s}{\alpha ^{2s}}}P^{(2s)}(1)+`$
$`{\displaystyle \frac{\mathrm{cos}\alpha }{\alpha ^2}}{\displaystyle \underset{s=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^s}{\alpha ^{2s}}}P^{(2s+1)}(1)],`$
where $`P^{(k)}(x)`$ is the $`k`$–th derivative of $`P(x)`$.
Making $`\alpha =nx`$ and $`P(x)=𝒫_\mathrm{}\mathrm{}^{}m(x)`$ in (C), substituting (C) in (38), and performing the sum in $`n`$ we get
$`F_{\mathrm{}\mathrm{}^{}}^m(x)`$ $`=`$ $`4{\displaystyle \underset{s=0}{\overset{\frac{\mathrm{}+\mathrm{}^{}}{2}}{}}}(1)^s[{\displaystyle \frac{𝒫_\mathrm{}\mathrm{}^{}m^{(2s)}(1)}{x^{2s+1}}}\phi _{2s+1}(x)+`$
$`{\displaystyle \frac{𝒫_\mathrm{}\mathrm{}^{}m^{(2s+1)}(1)}{x^{2s+2}}}\phi _{2s+2}(x)],`$
where $`\phi _k(x)`$ is the $`k`$–th Clausen $`\phi `$–function defined in Appendix B.
Since the Clausen functions are periodic functions of period $`2\pi `$, analytic in each period, it follows that $`F_{\mathrm{}\mathrm{}^{}}^m(x)`$ is a piecewise continuous function, analytic in each period as well. Thus we will now show how the explicit expression for $`F_{\mathrm{}\mathrm{}^{}}^m(x)`$, in the $`q`$–th interval $`[2\pi q,2\pi (q+1)]`$, given in (IV.1), comes out.
Introducing the explicit form for the Clausen $`\phi `$–functions (B), in the sum of (C) yields a huge expresion, but a close inspection reveals that it is a polynomial in $`\pi /x`$. The independent term is simply
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{s=0}{\overset{\frac{\mathrm{}+\mathrm{}^{}}{2}}{}}}{\displaystyle \frac{(1)^s}{(s+1)!}}𝒫_\mathrm{}\mathrm{}^{}m^{(s)}(1)`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle _1^1}𝒫_\mathrm{}\mathrm{}^{}m(x)𝑑x`$
$`=`$ $`{\displaystyle \frac{1}{4}}\delta _{\mathrm{}\mathrm{}^{}},`$
where the first equality can be deduced by writing the Taylor expansion of the integrand of the right hand side, and integrating. On the other hand, summing up all the coefficients of the $`(r+1)`$–th odd term, and proceeding as before, we have this term equal to
$$(1)^r𝒫_\mathrm{}\mathrm{}^{}m^{(2r)}(0)g_{2r+1}(q)\left(\frac{\pi }{x}\right)^{2r+1},$$
while the $`(r+1)`$–th even term is equal to
$$(1)^r𝒫_\mathrm{}\mathrm{}^{}m^{(2r+1)}(0)g_{2r+2}(q)\left(\frac{\pi }{x}\right)^{2r+2},$$
which by the parity of $`𝒫_\mathrm{}\mathrm{}^{}m(x)`$ is zero. Summing up all the terms we finally get (IV.1) and (IV.1).
## Appendix D Known results for closed flat 3–manifolds
In this section we briefly show how we can obtain the formulae for the correlation matrix of the $`a_\mathrm{}m`$’s and the angular power spectrum, currently available in the literature, for some closed flat manifolds, as well as a simple generalization, i.e., considering the observer out of the axis of rotations of the screw motions of the covering group. We present explicit derivations and formulae for the correlation matrix $`a_\mathrm{}ma_\mathrm{}^{}m^{}^{}`$, as well as for the power spectrum, in order to allow the interested reader to perform their own simulations confidently.
We first give a brief description of flat orientable closed 3–manifolds and their covering groups. The versions of the diffeomorphic and isometric classifications of flat 3–manifolds we present here were given by Wolf in Wolf , and previous descriptions in the context of cosmic topology were given in Gomero (see LevinA ; RWULL04 for alternative descriptions).
There are six diffeomorphic classes of compact orientable Euclidean 3–manifolds. The generators for the covering groups of the first five classes, $`𝒢_1𝒢_5`$, are $`\gamma _1=(I,𝐚)`$, $`\gamma _2=(I,𝐛)`$ and $`\gamma _3=(A_i,𝐜)`$, where $`A_1=I`$ is the identity and
$`A_2=\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right)`$ , $`A_4=\left(\begin{array}{ccc}0& 1& 0\\ 1& 0& 0\\ 0& 0& 1\end{array}\right),`$
$`A_3=\left(\begin{array}{ccc}0& 1& 0\\ 1& 1& 0\\ 0& 0& 1\end{array}\right)`$ , $`A_5=\left(\begin{array}{ccc}0& 1& 0\\ 1& 1& 0\\ 0& 0& 1\end{array}\right),`$
for the classes $`𝒢_1𝒢_5`$ respectively. It is important to remark that these matrices for the rotations are written in the basis formed by the set $`\{𝐚,𝐛,𝐜\}`$ of linearly independent vectors. Thus, the torus $`𝒢_1`$ is generated by three independent translations, while for the other manifolds the generators are two independent translations and a screw motion along a linearly independent direction. The manifold $`𝒢_6`$ is the most involved since their generators are all screw motions. In the following we present some general considerations concerning the classes $`𝒢_1𝒢_5`$ only.
For space forms of the classes $`𝒢_2𝒢_5`$, the following facts are easily derivable (see Gomero for details):
1. The vector $`𝐜`$ is orthogonal to both $`𝐚`$ and $`𝐛`$.
2. The angle between $`𝐚`$ and $`𝐛`$ is a free parameter for the class $`𝒢_2`$, while its value is fixed to be $`2\pi /3`$, $`\pi /2`$ and $`\pi /3`$ for the classes $`𝒢_3`$, $`𝒢_4`$ and $`𝒢_5`$ respectively.
3. Denoting by $`|𝐜|`$ the length of the vector $`𝐜`$, and similarly for any other vector, one has that $`|𝐚|=|𝐛|`$ for the classes $`𝒢_3𝒢_5`$, while both lengths are independent free parameters in the class $`𝒢_2`$. Moreover, in all classes $`𝒢_2𝒢_5`$, $`|𝐜|`$ is an independent free parameter.
4. Denoting the canonical unitary basis vectors in Euclidean space by $`\{\widehat{𝐞}_x,\widehat{𝐞}_y,\widehat{𝐞}_z\}`$, one can always write $`𝐚=|𝐚|\widehat{𝐞}_x`$, $`𝐛=|𝐛|\mathrm{cos}\phi \widehat{𝐞}_x+|𝐛|\mathrm{sin}\phi \widehat{𝐞}_y`$, and $`𝐜=|𝐜|\widehat{𝐞}_z`$, for the basis $`\{𝐚,𝐛,𝐜\}`$, where $`\phi `$ is the angle between $`𝐚`$ and $`𝐛`$, as established in the item 2.
Thus in dealing with manifolds of classes $`𝒢_2𝒢_5`$, the axis of rotation of the generator screw motion can be taken to be the $`z`$–axis, and the orthogonal part of this generator, in the basis $`\{\widehat{𝐞}_x,\widehat{𝐞}_y,\widehat{𝐞}_z\}`$, is
$$A=\left(\begin{array}{ccc}\mathrm{cos}\alpha & \mathrm{sin}\alpha & 0\\ \mathrm{sin}\alpha & \mathrm{cos}\alpha & 0\\ 0& 0& 1\end{array}\right),$$
(75)
with $`\alpha =\pi `$, $`2\pi /3`$, $`\pi /2`$ and $`\pi /3`$ respectively. Since the axis of rotation passes through the origin, the translational part of the generator $`\gamma _3`$ is $`𝐜=(0,0,L_z)`$, where we have put $`|𝐜|=L_z`$ as is usual in cosmic topology.
However, in cosmological applications we need to consider the arbitrariness of the position of the observer inside space. Thus if the axis of rotation is at a distance $`\rho `$ from the origin (the observer), and its intersection with the horizontal plane makes an angle $`\varphi `$ with the positive $`x`$–axis, the translational part of the screw motion $`\gamma _3=(A,𝐜)`$ is
$`𝐜`$ $`=`$ $`\rho [\mathrm{cos}\varphi \mathrm{cos}(\varphi +\alpha )]\widehat{𝐞}_x+`$ (76)
$`\rho [\mathrm{sin}\varphi \mathrm{sin}(\varphi +\alpha )]\widehat{𝐞}_y+L_z\widehat{𝐞}_z.`$
In order to perform calculations of the topological signature of CMB temperature maps, we need to write the covering group for the manifold under study in a compact form. For a torus $`𝒢_1`$ the problem is trivial, since the covering group is generated by three independent translations, and thus any two isometries commute (see Sec.D.1 below), while the covering groups for the other closed flat manifolds are noncommutative since they contain screw motions.
The generators of the covering groups for the classes $`𝒢_2𝒢_5`$ satisfy certain relations of the form
$$\gamma _3\gamma _1^{n_1}\gamma _2^{n_2}=\gamma _1^{m_1}\gamma _2^{m_2}\gamma _3,$$
where $`n_1,n_2,m_1,m_2`$, and they hold whether the axis of rotation passes through the origin or not. It follows that a generic isometry can always be put in the form
$$\gamma =\gamma _1^{n_1}\gamma _2^{n_2}\gamma _3^{n_3},$$
(77)
with $`\gamma _1^{n_1}=(I,n_1𝐚)`$, $`\gamma _2^{n_2}=(I,n_2𝐛)`$, and
$$\gamma _3^{n_3}=(A^h,n_3𝐜_{}+𝒪_h𝐜_{}),$$
where $`A`$ is given by (75), $`\alpha =2\pi /s`$, $`n_3=sq+h`$, with $`q`$ and $`h`$ integers such that $`0<hs`$, the parameter $`s`$ being $`2,3,4`$ and 6 corresponding to $`𝒢_2`$, $`𝒢_3`$, $`𝒢_4`$ and $`𝒢_5`$ respectively,
$$𝒪_h=\underset{j=0}{\overset{h1}{}}A^j,$$
$`𝐜_{}=L_z\widehat{𝐞}_z`$, and $`𝐜_{}=\rho [\mathrm{cos}\varphi \mathrm{cos}(\varphi +\alpha )]\widehat{𝐞}_x+\rho [\mathrm{sin}\varphi \mathrm{sin}(\varphi +\alpha )]\widehat{𝐞}_y`$.
It is now straightforward to compute both the correlation matrix $`a_\mathrm{}ma_\mathrm{}^{}m^{}^{}`$ and the angular power spectrum $`C_{\mathrm{}}`$. They all have a simple structure. We first describe the general procedure for obtaining these expressions and present the results in a unified form. We finally specify each case separately. Note that, due to (76), in all of our calculations we are considering that the observer may be off an axis of rotation of the screw motions of $`\mathrm{\Gamma }`$.
Upon introducing (77) into (25), we transform the series of exponentials in a series of Dirac’s delta functions by using (32). The integration of (II) is then immediate in Cartesian coordinates, the general result being
$`a_\mathrm{}ma_\mathrm{}^{}m^{}^{}`$ $`=`$ $`{\displaystyle \frac{(4\pi )^2}{V}}i^{\mathrm{}\mathrm{}^{}}{\displaystyle \underset{𝐩\widehat{}^3}{}}{\displaystyle \frac{1}{\beta ^3}}\mathrm{\Psi }_{\mathrm{}\mathrm{}^{}}(2\pi \beta )\times `$
$`Y_\mathrm{}^{}m^{}(𝐧_\stackrel{}{\beta })Y_\mathrm{}m^{}(𝐧_\stackrel{}{\beta })f_m^{}^\mathrm{\Gamma }(2\pi \stackrel{}{\beta }),`$
where $`V`$ is the volume of the manifold, and $`\widehat{}^3=^3(0,0,0)`$, since the term corresponding to $`𝐩=0`$ represents a constant perturbation, and thus is neglected.
The function
$$f_m^\mathrm{\Gamma }(𝐤)=\frac{1}{s}\left[1+\underset{h=1}{\overset{s1}{}}\omega _s^{hm}e^{i𝐤𝒪_h𝐜}\right],$$
where $`\omega _s`$ is the first complex $`s`$th root of unity, is a complex modulating term characteristic of the geometry and topology of the spatial section of the universe model, and depends only on the screw motion generators. The vector $`\stackrel{}{\beta }(𝐩)`$ comes from the discretization of the wavevector $`𝐤`$ due to the Dirac’s deltas (each $`2\pi \beta `$ is an eigenvalue of the Laplacian operator), and $`𝐧_\stackrel{}{\beta }`$ is the unit vector in the direction of $`\stackrel{}{\beta }`$.
Using the property $`a_\mathrm{}ma_\mathrm{}^{}m^{}^{}=a_\mathrm{}^{}m^{}a_\mathrm{}m^{}^{}`$ one can easily show, by resumming the series, that the variances of the multipole moments can be put in the general form
$`|a_\mathrm{}m|^2`$ $`=`$ $`{\displaystyle \frac{(4\pi )^2}{V}}{\displaystyle \underset{𝐩\widehat{}^3}{}}{\displaystyle \frac{1}{\beta ^3}}\mathrm{\Psi }_{\mathrm{}\mathrm{}}(2\pi \beta )\times `$
$`\left|Y_\mathrm{}m(𝐧_\stackrel{}{\beta })\right|^2\mathrm{}\left(f_m^\mathrm{\Gamma }(\stackrel{}{2\pi \beta })\right),`$
where $`\mathrm{}`$ stands for the real part of a complex number. The angular power spectrum is then
$$C_{\mathrm{}}=\frac{4\pi }{V}\underset{𝐩\widehat{}^3}{}\frac{1}{\beta ^3}\mathrm{\Psi }_{\mathrm{}\mathrm{}}(2\pi \beta )\mathrm{\Xi }_{\mathrm{}}(2\pi \stackrel{}{\beta }),$$
where
$$\mathrm{\Xi }_{\mathrm{}}(𝐤)=\frac{4\pi }{2\mathrm{}+1}\underset{m=l}{\overset{\mathrm{}}{}}\left|Y_\mathrm{}m(𝐧_𝐤)\right|^2\mathrm{}\left(f_m^\mathrm{\Gamma }(𝐤)\right)$$
can be evaluated using the Addition Theorem for Spherical Harmonics yielding
$$\mathrm{\Xi }_{\mathrm{}}(𝐤)=\frac{1}{s}\left[1+\underset{h=1}{\overset{s1}{}}P_{\mathrm{}}(\mathrm{cos}\theta _{𝐤,h})\mathrm{cos}(𝐤𝒪_h𝐜)\right],$$
where
$$\mathrm{cos}\theta _{𝐤,h}=\mathrm{cos}^2\theta _𝐤+\mathrm{sin}^2\theta _𝐤\mathrm{cos}\frac{2\pi h}{s}.$$
### D.1 Rectangular torus $`𝒢_1`$
The generators for the rectangular torus are the translations $`𝐚=L_x\widehat{𝐞}_x`$, $`𝐛=L_y\widehat{𝐞}_y`$ and $`𝐜=L_z\widehat{𝐞}_z`$, thus a generic isometry of its covering group can be written as $`\gamma =(I,𝐫)`$, with $`𝐫=n_xL_x\widehat{𝐞}_x+n_yL_y\widehat{𝐞}_y+n_zL_z\widehat{𝐞}_z`$, and $`n_x,n_y,n_z`$, i.e., the covering group of $`𝒢_1`$ is parametrized by $`^3`$.
It follows immediately that, for a rectangular torus, the expression (25) takes the form
$`\mathrm{{\rm Y}}_\mathrm{}m^\mathrm{\Gamma }(𝐤)`$ $`=`$ $`{\displaystyle \underset{𝐧^3}{}}e^{i(n_xk_xL_x+n_yk_yL_y+n_zk_zL_z)}Y_\mathrm{}m(𝐧_𝐤)`$
$`=`$ $`(2\pi )^3{\displaystyle \underset{𝐩^3}{}}\delta (k_xL_x2\pi p_x)\times `$
$`\delta (k_yL_y2\pi p_y)\delta (k_zL_z2\pi p_z)Y_\mathrm{}m(𝐧_𝐤).`$
Following the procedure described above one gets $`f_m^\mathrm{\Gamma }(𝐤)=1`$, and $`\beta _x=\frac{p_x}{L_x}`$, $`\beta _y=\frac{p_y}{L_y}`$, and $`\beta _z=\frac{p_z}{L_z}`$. In particular, we have the well known result
$$C_{\mathrm{}}=\frac{4\pi }{V}\underset{𝐩\widehat{}^3}{}\frac{1}{\beta ^3}\mathrm{\Psi }_{\mathrm{}\mathrm{}}(2\pi \beta ).$$
### D.2 Rectangular $`𝒢_2`$
The generators for the rectangular $`𝒢_2`$ are $`\gamma _1=(I,𝐚)`$, $`\gamma _2=(I,𝐛)`$ and $`\gamma _3=(A,𝐜)`$, with $`𝐚=L_x\widehat{𝐞}_x`$, $`𝐛=L_y\widehat{𝐞}_y`$, $`𝐜=2\rho \mathrm{cos}\varphi \widehat{𝐞}_x+2\rho \mathrm{sin}\varphi \widehat{𝐞}_y+L_z\widehat{𝐞}_z`$, and $`A`$ given in (75) with $`\alpha =\pi `$. They satisfy the relations $`\gamma _1\gamma _3\gamma _1=\gamma _3`$ and $`\gamma _2\gamma _3\gamma _2=\gamma _3`$, which allow to write any isometry of the covering group by (77) with
$$\gamma _3^{n_3}=\{\begin{array}{ccc}(I,n_3𝐜_{})\hfill & & \text{if }n_3\text{ is even}\\ (A,n_3𝐜_{}+𝐜_{})\hfill & & \text{if }n_3\text{ is odd}\end{array},$$
(78)
where $`𝐜_{}=2\rho \mathrm{cos}\varphi \widehat{𝐞}_x+2\rho \mathrm{sin}\varphi \widehat{𝐞}_y`$.
It follows from (77) and (78) that the expression (25) takes the form
$`\mathrm{{\rm Y}}_\mathrm{}m^\mathrm{\Gamma }(𝐤)`$ $`=`$ $`{\displaystyle \underset{𝐧^3}{}}e^{i(n_xk_xL_x+n_yk_yL_y+2n_zk_zL_z)}\left[1+(1)^me^{i𝐤𝐜}\right]Y_\mathrm{}m(𝐧_𝐤)`$
$`=`$ $`(2\pi )^3{\displaystyle \underset{𝐩^3}{}}\delta (k_xL_x2\pi p_x)\delta (k_yL_y2\pi p_y)\delta (k_zL_z\pi p_z)Y_\mathrm{}m(𝐧_𝐤)f_m^\mathrm{\Gamma }(𝐤),`$
where we have put $`n_1=n_x`$, $`n_2=n_y`$, and $`n_3=2n_z`$ or $`2n_z+1`$, depending on whether $`n_3`$ is even or odd. The components of $`\stackrel{}{\beta }`$ are $`\beta _x=\frac{p_x}{L_x}`$, $`\beta _y=\frac{p_y}{L_y}`$, and $`\beta _z=\frac{p_z}{2L_z}`$.
### D.3 $`𝒢_3`$
The generators for a manifold of class $`𝒢_3`$ are $`\gamma _1=(I,𝐚)`$, $`\gamma _2=(I,𝐛)`$, and $`\gamma _3=(A,𝐜)`$, with $`𝐚=L\widehat{𝐞}_x`$, $`𝐛=\frac{L}{2}(\widehat{𝐞}_x\sqrt{3}\widehat{𝐞}_y)`$, $`𝐜=\frac{\rho }{2}(3\mathrm{cos}\varphi +\sqrt{3}\mathrm{sin}\varphi )\widehat{𝐞}_x+\frac{\rho }{2}(3\mathrm{sin}\varphi \sqrt{3}\mathrm{cos}\varphi )\widehat{𝐞}_y+L_z\widehat{𝐞}_z`$, and $`A`$ given in (75) with $`\alpha =2\pi /3`$. They satisfy the relations $`\gamma _2^1\gamma _3\gamma _1=\gamma _3`$ and $`\gamma _1\gamma _2\gamma _3\gamma _2=\gamma _3`$, which allow us to write any isometry of the covering group by (77) with
$$\gamma _3^{n_3}=\{\begin{array}{ccc}(I,n_3𝐜_{})\hfill & & \text{if }n_3=0\text{ mod 3}\\ (A,n_3𝐜_{}+𝐜_{})\hfill & & \text{if }n_3=1\text{ mod 3}\\ (A^2,n_3𝐜_{}+𝒪_2𝐜_{})\hfill & & \text{if }n_3=2\text{ mod 3}\end{array},$$
(79)
where $`𝐜_{}=\frac{\rho }{2}(3\mathrm{cos}\varphi +\sqrt{3}\mathrm{sin}\varphi )\widehat{𝐞}_x+\frac{\rho }{2}(3\mathrm{sin}\varphi \sqrt{3}\mathrm{cos}\varphi )\widehat{𝐞}_y`$.
It follows from (77) and (79) that the expression (25) takes the form
$`\mathrm{{\rm Y}}_\mathrm{}m^\mathrm{\Gamma }(𝐤)`$ $`=`$ $`{\displaystyle \underset{𝐧^3}{}}e^{i\left[n_xk_xL+n_y\left(\frac{\sqrt{3}}{2}k_y\frac{1}{2}k_x\right)L+3n_zk_zL_z\right]}\left[1+\omega _3^me^{i𝐤𝐜}+\omega _3^{2m}e^{i𝐤𝒪_2𝐜}\right]Y_\mathrm{}m(𝐧_𝐤)`$
$`=`$ $`(2\pi )^3{\displaystyle \underset{𝐩^3}{}}\delta (k_xL2\pi p_x)\delta \left(\left[\frac{\sqrt{3}}{2}k_y\frac{1}{2}k_x\right]L2\pi p_y\right)\delta \left(k_zL_z\frac{2\pi }{3}p_z\right)Y_\mathrm{}m(𝐧_𝐤)f_m^\mathrm{\Gamma }(𝐤),`$
where we have put $`n_1=n_x`$, $`n_2=n_y`$, and $`n_3=3n_z`$, $`3n_z+1`$ or $`3n_z+2`$ according to (79). We also get $`\beta _x=\frac{p_x}{L}`$, $`\beta _y=\frac{\sqrt{3}}{3L}(2p_y+p_x)`$, and $`\beta _z=\frac{p_z}{3L_z}`$.
### D.4 $`𝒢_4`$
The generators for a manifold of class $`𝒢_4`$ are $`\gamma _1=(I,𝐚)`$, $`\gamma _2=(I,𝐛)`$, and $`\gamma _3=(A,𝐜)`$, with $`𝐚=L\widehat{𝐞}_x`$, $`𝐛=L\widehat{𝐞}_y`$, $`𝐜=\rho (\mathrm{cos}\varphi +\mathrm{sin}\varphi )\widehat{𝐞}_x+\rho (\mathrm{sin}\varphi \mathrm{cos}\varphi )\widehat{𝐞}_y+L_z\widehat{𝐞}_z`$, and $`A`$ given in (75) with $`\alpha =\pi /2`$. They satisfy the relations $`\gamma _2^1\gamma _3\gamma _1=\gamma _3`$ and $`\gamma _1\gamma _3\gamma _2=\gamma _3`$, which allow us to write any isometry of the covering group by (77) with
$$\gamma _3^{n_3}=\{\begin{array}{ccc}(I,n_3𝐜_{})\hfill & & \text{if }n_3=0\text{ mod 4}\\ (A,n_3𝐜_{}+𝐜_{})\hfill & & \text{if }n_3=1\text{ mod 4}\\ (A^2,n_3𝐜_{}+𝒪_2𝐜_{})\hfill & & \text{if }n_3=2\text{ mod 4}\\ (A^3,n_3𝐜_{}+𝒪_3𝐜_{})\hfill & & \text{if }n_3=3\text{ mod 4}\end{array},$$
(80)
where $`𝐜_{}=\rho (\mathrm{cos}\varphi +\mathrm{sin}\varphi )\widehat{𝐞}_x+\rho (\mathrm{sin}\varphi \mathrm{cos}\varphi )\widehat{𝐞}_y`$.
Similarly, using (77) and (80), the expression (25) takes the form
$`\mathrm{{\rm Y}}_\mathrm{}m^\mathrm{\Gamma }(𝐤)`$ $`=`$ $`{\displaystyle \underset{𝐧^3}{}}e^{i(n_xk_xL+n_yk_yL+4n_zk_zL_z)}\left[1+{\displaystyle \underset{h=1}{\overset{3}{}}}\omega _4^{hm}e^{i𝐤𝒪_h𝐜}\right]Y_\mathrm{}m(𝐧_𝐤)`$
$`=`$ $`(2\pi )^3{\displaystyle \underset{𝐩^3}{}}\delta (k_xL2\pi p_x)\delta (k_yL2\pi p_y)\delta \left(k_zL_z\frac{\pi }{2}p_z\right)Y_\mathrm{}m(𝐧_𝐤)f_m^\mathrm{\Gamma }(𝐤),`$
where we have put $`n_1=n_x`$, $`n_2=n_y`$, and $`n_3=4n_z`$, $`4n_z+1`$, $`4n_z+2`$ or $`4n_z+3`$ according to (80), and $`\omega _4`$ is the first complex 4th-rooth of unity. We also get $`\beta _x=\frac{p_x}{L}`$, $`\beta _y=\frac{p_y}{L}`$, and $`\beta _z=\frac{p_z}{4L_z}`$.
### D.5 $`𝒢_5`$
The generators for the rectangular $`𝒢_5`$ are $`\gamma _1=(I,𝐚)`$, $`\gamma _2=(I,𝐛)`$, and $`\gamma _3=(A,𝐜)`$, with $`𝐚=L\widehat{𝐞}_x`$, $`𝐛=\frac{L}{2}(\widehat{𝐞}_x+\sqrt{3}\widehat{𝐞}_y)`$, $`𝐜=\frac{\rho }{2}(\mathrm{cos}\varphi +\sqrt{3}\mathrm{sin}\varphi )\widehat{𝐞}_x+\frac{\rho }{2}(\mathrm{sin}\varphi \sqrt{3}\mathrm{cos}\varphi )\widehat{𝐞}_y+L_z\widehat{𝐞}_z`$, and $`A`$ given in (75) with $`\alpha =\pi /3`$. They satisfy the relations $`\gamma _2^1\gamma _3\gamma _1=\gamma _3`$ and $`\gamma _1\gamma _2^1\gamma _3\gamma _2=\gamma _3`$, which allow us to write any isometry of the covering group by (77) with
$$\gamma _3^{n_3}=\{\begin{array}{ccc}(I,n_3𝐜_{})\hfill & & \text{if }n_3=0\text{ mod 6}\\ (A,n_3𝐜_{}+𝐜_{})\hfill & & \text{if }n_3=1\text{ mod 6}\\ (A^2,n_3𝐜_{}+𝒪_2𝐜_{})\hfill & & \text{if }n_3=2\text{ mod 6}\\ (A^3,n_3𝐜_{}+𝒪_3𝐜_{})\hfill & & \text{if }n_3=3\text{ mod 6}\\ (A^4,n_3𝐜_{}+𝒪_4𝐜_{})\hfill & & \text{if }n_3=4\text{ mod 6}\\ (A^5,n_3𝐜_{}+𝒪_5𝐜_{})\hfill & & \text{if }n_3=5\text{ mod 6}\end{array},$$
(81)
where $`𝐜_{}=\frac{\rho }{2}(\mathrm{cos}\varphi +\sqrt{3}\mathrm{sin}\varphi )\widehat{𝐞}_x+\frac{\rho }{2}(\mathrm{sin}\varphi \sqrt{3}\mathrm{cos}\varphi )\widehat{𝐞}_y`$.
Using (77) and (81), the expression (25) takes the form
$`\mathrm{{\rm Y}}_\mathrm{}m^\mathrm{\Gamma }(𝐤)`$ $`=`$ $`{\displaystyle \underset{𝐧^3}{}}e^{i\left[n_xk_xL+n_y\left(\frac{\sqrt{3}}{2}k_y+\frac{1}{2}k_x\right)L+6n_zk_zL_z\right]}\left[1+{\displaystyle \underset{h=1}{\overset{5}{}}}\omega _6^{hm}e^{i𝐤𝒪_h𝐜}\right]Y_\mathrm{}m(𝐧_𝐤)`$
$`=`$ $`(2\pi )^3{\displaystyle \underset{𝐩^3}{}}\delta (k_xL2\pi p_x)\delta \left(\left[\frac{\sqrt{3}}{2}k_y+\frac{1}{2}k_x\right]L2\pi p_y\right)\delta \left(k_zL_z\frac{\pi }{3}p_z\right)Y_\mathrm{}m(𝐧_𝐤)f_m^\mathrm{\Gamma }(𝐤),`$
where we have put $`n_1=n_x`$, $`n_2=n_y`$, and $`n_3=6n_z`$, $`6n_z+1`$, $`6n_z+2`$, $`\mathrm{}`$, $`6n_z+5`$ according to (81), and $`\omega _6`$ is the first complex 6th-rooth of unity. We also get $`\beta _x=\frac{p_x}{L}`$, $`\beta _y=\frac{\sqrt{3}}{3L}(2p_yp_x)`$, and $`\beta _z=\frac{p_z}{6L_z}`$. |
warning/0507/hep-ph0507147.html | ar5iv | text | # Pion-delta sigma-term
## I introduction
The delta ($`\mathrm{\Delta }`$) plays a very important role in low-energy pion-nucleon ($`\pi N`$) scattering and correlated processes, such as the nucleon-nucleon interaction. Its contribution as an intermediate state, in many instances, supersedes that of the the nucleon. This happens for two main reasons. The first one is that the $`\pi N\mathrm{\Delta }`$ coupling constant is rather large, whereas the other is related to chiral symmetry.
At low energies, pion-hadron interactions are well described by effective theories, in which an approximate $`SU(2)\times SU(2)`$ symmetry, broken by the pion mass $`(\mu )`$, accounts for the smallness of the $`u`$ and $`d`$ quark masses. In this framework, elastic pion-baryon scattering is dominated by diagrams involving both contact terms and propagating states. In order to comply with threshold chiral theorems, the latter are typically given by polynomials in small quantities, such as the pion mass or three-momenta, divided by energy denominators. When the delta is present, the scale of some denominators is given by the quantity $`\omega _\mathrm{\Delta }=(M^2m^2\mu ^2)/2m`$, where $`M`$ and $`m`$ are respectively the delta and nucleon masses. As the difference $`\mathrm{\Delta }Mm`$ is small, one has $`\omega _\mathrm{\Delta }\mathrm{\Delta }`$. Delta contributions are given by ratios of small quantities and may turn out to be large. In such cases, numerical values adopted for $`\mathrm{\Delta }`$ do influence predictions produced by effective theories, especially those that rely on the small scale expansionHHK or the heavy baryon approximationHB .
In chiral perturbation theory, there is a clear conceptual distinction between the bare baryon masses, present in the lagrangian, and their observed values, which include loop corrections. The former should, in principle, be preferred as inputs in the evaluation of theoretical amplitudes. Nevertheless, as there is little knowledge available concerning the bare delta mass, one tends to use physical values in calculations. In most cases, it is reasonable to expect that this would have little numerical importance. On the other hand, in the case of the parameter $`\mathrm{\Delta }`$, which is a small quantity, the influence of loops may become relatively large.
Recently Bernard, Hemmert and MeissnerBHM , have stressed that the value of $`\mathrm{\Delta }_0`$, the delta-nucleon mass splitting in the chiral limit, is an important constraint to lattice data extrapolation. The purpose of the present work is to estimate the delta $`\sigma `$-term, which controls the change induced in $`\mathrm{\Delta }`$ when one goes from bare to physical masses. This $`\sigma `$-term was studied in the framework of a quark model by Lyubovitskij, Gutsche, Faessler and DrukarevLGFD and the reader is referred to their paper for a clear formulation of the problem and earlier works.
According to the Feynman-Hellmann theoremFH the mass $`m_B`$ of a baryon $`B`$ is related to its sigma-term $`\sigma _B`$ by $`\sigma _B=\mu ^2dm_B/d\mu ^2`$. Therefore the sigma-term provides a measure of the shift in the baryon mass due to chiral symmetry breaking. Whenever it is possible to evaluate $`\sigma _B`$ as a function of $`\mu `$, the bare mass $`m_{B_0}`$ can be extracted from the relation
$$m_B=m_{B_0}+_0^{\mu ^2}𝑑\lambda \sigma _B(\lambda )/\lambda .$$
(1)
As the leading term in $`\sigma _B`$ is proportional to $`\mu ^2`$, it enters directly the mass shift and the difference $`m_B\sigma _B`$ already provides a crude estimate for the bare value. In the case of the nucleon, one has $`\sigma _N`$=45 MeVGLS , which amounts to 5% of its physical mass. In chiral perturbation theory, the leading contribution to $`\sigma _N`$ cannot be predicted theoretically. Formally, it is associated with the constant $`c_1`$ of the second order lagrangianGSS ; BL , which can be extracted from empirical subthreshold information. The situation of the delta is much worse, for $`\pi \mathrm{\Delta }`$ scattering data are not available. One is then forced to resort to models in order to calculate the delta $`\sigma `$-term, which is associated with the parameter $`a_1`$ defined in ref.HHK .
In this work we estimate $`\sigma _\mathrm{\Delta }`$ using a model which proved to be successful in the case of the nucleon. Our paper is organized as follows. In section II we review our calculational procedure in the case of the nucleon and present results for the delta in section III, leaving technical details to the appendices. The main expressions for both the nucleon and delta $`\sigma `$-terms in configuration space are given in appendix B, written in terms of the loop integrals defined in appendix A. The consistency of our results with standard chiral counting rules is discussed in appendix C whereas their behaviour in the chiral limit is given in appendix D. A summary is provided in section IV.
## II model for the sigma term
In order to evaluate $`\sigma _\mathrm{\Delta }`$, we follow a procedure used previously in the study of $`\sigma _N(t)`$, the nucleon scalar form factorR , which is briefly reviewed here. The leading contributions to this function is $`𝒪(q^2)`$ whereas the triangle diagram, involving only known masses and coupling constants, gives rise to corrections which begin at $`𝒪(q^3)`$ and are completely determined. At $`𝒪(q^4)`$, on the other hand, interactions incorporate the low energy constants $`c_1,c_2`$ and $`c_3`$. Data on $`\pi N`$ subthreshold coefficients indicate that $`c_2`$ and $`c_3`$ are larger than $`c_1`$ and that their values are approximately saturated by $`\mathrm{\Delta }`$ intermendiate statesBL . Thus, up to $`𝒪(q^4)`$, the function $`\sigma _N(t)`$ can be well represented by the leading tree contribution associated with $`c_1`$, supplemented by the two triangle diagrams shown in fig.1, involving $`N`$ and $`\mathrm{\Delta }`$ intermediate states. In the sequence we will make use of the fact that, in configuration space, contact and loop contributions split apart, since the Fourier transform acts as a filterL . As a result, the theoretically undetermined leading tree term yields a zero-range $`\delta `$-function, whereas the triangle diagrams give rise to spatially distributed structures, fully determined by known parameters.
The nucleon scalar form factor in momentum space is defined by
$$N(p^{})|_{sb}|N(p)=\sigma _N(t)\overline{u}(p^{})u(p),$$
(2)
where $`_{sb}`$ is the symmetry breaking term in the lagrangian and $`t=(p^{}p)^2`$. In terms of the quark degrees of freedom, one has $`_{sb}=\widehat{m}(\overline{u}u+\overline{d}d)`$, with $`\widehat{m}=(m_u+m_d)/2`$. The configuration space scalar form factor is denoted by $`\stackrel{~}{\sigma }_N`$ and given by
$$\stackrel{~}{\sigma }_N(𝒓)=\frac{d^3q}{(2\pi )^3}e^{i𝒒𝒓}\sigma _N(t)$$
(3)
with $`𝒒=(𝒑^{}𝒑)`$, in the Breit frame. The nucleon $`\sigma `$-term, defined as $`\sigma _N\sigma _N(t=0)`$, is given by
$$\sigma _N=4\pi _0^{\mathrm{}}𝑑rr^2\stackrel{~}{\sigma }_N(r).$$
(4)
The contributions from the diagrams of fig.1 to $`\stackrel{~}{\sigma }_N(𝒓)`$ read
$$\stackrel{~}{\sigma }_N(𝒓)=4c_1\mu ^2\delta ^3(𝒓)+\stackrel{~}{\sigma }_{N_N}(r)+\stackrel{~}{\sigma }_{N_\mathrm{\Delta }}(r),$$
(5)
where $`\stackrel{~}{\sigma }_{N_N}(r)`$ and $`\stackrel{~}{\sigma }_{N_\mathrm{\Delta }}(r)`$ are given by eqs.(43) and (44) of appendix B and displayed in fig.2. These functions are based on unregularized loop integrals and diverge for small values of $`r`$. In momentum space, regularization is achieved by means of counterterms, which give rise to polynomials in $`t`$, designed to cancel the divergences of the loop integrals. In configuration space, this regularization procedure amounts to adding $`\delta `$-functions and their derivatives to $`\stackrel{~}{\sigma }_N(r)`$. This gives rise to a regularized form factor which is very large both at $`r=0`$ and in a sizeable vicinity of that point. We argue, in the sequence, that this picture is not consistent with the definition of the form factor given by eq.(2).
Pions are Goldstone bosons, collective states derived from the $`q\overline{q}`$ condensate. The corresponding degrees of freedom are appropriately accomodated into non-linear lagrangians and described by the field $`U=\mathrm{exp}(i𝝉\widehat{𝝅}\theta )`$, where $`\widehat{𝝅}`$ is the isospin direction and $`\theta `$ is the chiral angle. This function can be expressed as $`U=\mathrm{cos}\theta +i𝝉\widehat{𝝅}\mathrm{sin}\theta `$ and the dimensional pion field is given by
$$\mathit{\varphi }=f_\pi \mathrm{sin}\theta \widehat{𝝅}.$$
(6)
Long ago, SkyrmeSky , in a series of papers, considered the possibility of pion fields being either weak or strong. It is worth noting that the words weak and strong, as used here, are akin to the notion of weak and strong electromagnetic fields developed by Schwinger, and not at all related to the nature of the fundamental interactions. In the former case, changes in the $`q\overline{q}`$ condensate are small, one relies on the approximation $`\mathit{\varphi }f_\pi \theta \widehat{𝝅}`$ and can employ perturbative techniques, as in chiral perturbation theory (ChPT). In the latter, disturbances of the QCD vacuum become important and the non-linear nature of pionic interactions manifests itself through the condition $`|\mathit{\varphi }|f_\pi `$. The physical picture behind eq.(6) is that pions, as Goldstone bosons, destroy the $`q\overline{q}`$ condensate in order to exist.
When strong fields are present, constraints also apply to the scalar form factor. The symmetry breaking lagrangian is written in terms of the dimensional pion field as
$$_{sb}=\frac{1}{4}f_\pi ^2\mu ^2Tr\left[U+U^{}\right]=f_\pi ^2\mu ^2\mathrm{cos}\theta .$$
(7)
This structure shows that $`_{sb}`$ is a bound function and definition (2) means that the same necessarily happens with the scalar form factor. The function $`\stackrel{~}{\sigma }_N(x)`$ corresponds to a mass density induced in the vacuum by the presence of the nucleon, which manifests itself in the form of a pion cloud. Far away from the nucleon, eq.(7) yields the density of the condensate, which is negative and equal to $`f_\pi ^2\mu ^2`$. In the description of a nucleon, it is convenient to use a convention for the energy in which the density tends to zero at long distances and $`_{sb}`$ is rewritten as
$$_{sb}=f_\pi ^2\mu ^2(\mathrm{cos}\theta 1).$$
(8)
In this new convention, the density vanishes when $`r\mathrm{}`$ and increases monotonically as one approaches the center of the nucleon as in fig.2. At a critical radius $`R`$ one has $`\mathrm{cos}\theta =1`$, the density becomes that of empty space and the condition
$$\stackrel{~}{\sigma }_N(R)=f_\pi ^2\mu ^2$$
(9)
holds. Beyond this point, a further increase in $`\stackrel{~}{\sigma }_N`$ would correspond to $`\mathrm{cos}\theta >1`$. In order to prevent this behaviour, we assume that the condensate no longer exists in the region $`r<R`$, and that the energy density saturates at $`r=R`$. For this reason, in our previous evaluation of $`\sigma _N`$R , we used the expression
$$\sigma _N=\frac{4}{3}\pi R^3f_\pi ^2\mu ^2+4\pi _R^{\mathrm{}}𝑑rr^2\stackrel{~}{\sigma }_N(r)$$
(10)
instead of eq.(4). This procedure is the basis of our model.
In the numerical determination of $`\sigma _N`$, we use the results of appendix B and consider two possibilities for the $`\pi N\mathrm{\Delta }`$ coupling constant in the lagrangian (41), corresponding to either the $`SU(4)`$ prediction $`g_{\pi N\mathrm{\Delta }}=3g_A/2\sqrt{2}=1.33`$ or $`g_{\pi N\mathrm{\Delta }}=1.47`$, which yields $`\mathrm{\Gamma }`$=120MeV for the $`\mathrm{\Delta }`$ decay width. The corresponding results, given in table 1, are quite close to the value extracted from experiment by Gasser, Leutwyler and SainioGLS , namely $`\sigma _N=45`$MeV.
Consistency with chiral symmetry is an important issue in this problem. Therefore, we note that, although the chiral powers of the pion mass expected from triangle diagrams are not explicit in the expressions of appendix B, the use of covariant relations among integralsHR allows results for partial contributions to the $`\sigma `$-term to be recast in such a way that these powers become apparent, as shown in appendix C. In appendix D we show that the formal chiral expansion of eq.(10) gives rise to the expected non-analytic terms ($`\mathrm{log}\mu `$ and $`\mu ^3`$) and agrees fully with that produced by standard chiral perturbation theoryBL , provided the renormalization scale is identified with $`1/R`$.
## III delta $`\sigma `$-term
The delta scalar form factor is defined as
$$<\mathrm{\Delta }(p^{},s^{})|_{sb}|\mathrm{\Delta }(p,s)>\overline{u}_\mu ^s^{}(𝒑^{})\left[g^{\mu \nu }\sigma _\mathrm{\Delta }(t)+p^\nu p^\mu F_T(t)\right]u_\nu ^s(𝒑),$$
(11)
where $`u_\nu ^s`$ is the $`\mathrm{\Delta }`$ spinorD and $`\sigma _\mathrm{\Delta }`$ and $`F_T`$ are respectively the scalar and tensor form factors. The minus sign on the r.h.s. is associated with the conventions used in the free $`\mathrm{\Delta }`$ lagrangian as in ref.HHK . We assume that the scalar form factor is determined by a short range contact interaction and the two long range two-pion processes shown in fig.1.
In figs. 3 and 4 (zoom in) we display the profile functions for the partial contributions to $`\sigma _\mathrm{\Delta }`$ given by eqs.(48,49) and it is interesting to note that the nucleon contribution oscillates in the outer region, in sharp contrast with fig.2. This behaviour is due to the fact that the delta is unstable and makes $`\sigma _\mathrm{\Delta }`$ to be smaller than $`\sigma _N`$. The structure of partial contributions for $`SU(4)`$ coupling constants is given in table 2, where core and cloud refer respectively to regions inside and outside the cutting radius $`R`$.
Processes containing nucleon intermediate states give rise to an imaginary component $`\sigma _\mathrm{\Delta }^I`$ for the delta $`\sigma `$-term, which can be related to the decay width by means of the Feynman-Hellmann theoremFH :
$$\sigma _\mathrm{\Delta }i\sigma _\mathrm{\Delta }^I=\mu ^2\frac{d(Mi\mathrm{\Gamma }/2)}{d\mu ^2}.$$
(12)
UsingBL
$`\mathrm{\Gamma }={\displaystyle \frac{g_{\pi N\mathrm{\Delta }}^2}{24\pi M^2f_\pi ^2}}q_\mathrm{\Delta }^3\left[(M+m)^2\mu ^2\right],`$
$`q_\mathrm{\Delta }={\displaystyle \frac{1}{2M}}\sqrt{M^4+m^4+\mu ^42m^2M^22\mu ^2M^22\mu ^2m^2},`$ (13)
one finds
$$\sigma _\mathrm{\Delta }^I=\frac{g_{\pi N\mathrm{\Delta }}^2\mu ^2}{48\pi M^2f_\pi ^2}\left\{q_\mathrm{\Delta }^3+3q_\mathrm{\Delta }\frac{M^2+m^2\mu ^2}{4M^2}\left[(M+m)^2\mu ^2\right]\right\}.$$
(14)
The values of the distance $`R`$ for which $`\stackrel{~}{\sigma }_\mathrm{\Delta }(R)/f_\pi ^2\mu ^2=1`$ and of the delta $`\sigma `$-term, calculated by means of eqs.(48, 49), are given in table 3, for different choices of the coupling constants $`g_{\pi N\mathrm{\Delta }}`$ and $`g_{\pi \mathrm{\Delta }\mathrm{\Delta }}`$. The $`SU(4)`$ predictions for these constants are 1.33 and 0.75, whereas the value 1.47 for the former yields the empirical decay width. The value 0.67 for the latter was used in ref.BHM . Results for the real component of $`\sigma _\mathrm{\Delta }`$ are sensitive to the coupling constant $`g_{\pi \mathrm{\Delta }\mathrm{\Delta }}`$ and fully consistent with that given in ref.LGFD , namely $`\sigma _\mathrm{\Delta }=(32\pm 3)`$ MeV. On the other hand, our prediction is larger than that quoted in ref.BHM . The values for the imaginary component $`\sigma _\mathrm{\Delta }^I`$, obtained by means of eq.(48), are identical with those given by eq.(14), as they should.
## IV summary
We have discussed a model aimed at determining $`\sigma `$terms, which consists in cutting off configuration space expressions at the point where the cosine of the chiral angle becomes larger than 1. The model has been used to calculate $`\sigma _N`$ and $`\sigma _\mathrm{\Delta }`$ with success. In the former case, a value very close to that extracted from experiment by Gasser, Leutwyler and SainioGLS , was obtained. In the case of the delta, the prediction 28 MeV $`\sigma _\mathrm{\Delta }`$ 32 MeV, depending on the coupling constants employed, is also very close to the result produced by another groupLGFD . The fact that the delta can decay gives rise to a pion cloud which includes an oscillation and is responsible for both the relation $`\sigma _\mathrm{\Delta }<\sigma _N`$ and the consistency of the imaginary part of $`\sigma _\mathrm{\Delta }`$ with the decay width. Analytic expressions also comply with chiral counting rules and give rise to expected non-analytic terms in the chiral limit. These features suggest that our calculational procedure is sound and can be reliably applied to other systems.
## Appendix A loop integrals
In the triangle diagrams, $`p`$ and $`p^{}`$ are the initial and final baryon momenta, whereas $`k`$ and $`k^{}`$ are the momenta of the exchanged pions. We also employ the variables
$$q=(pp^{}),P=(p+p^{})/2,Q=(k+k^{})/2.$$
(15)
In all diagrams, the external baryon, with mass $`m_e`$, is assumed to be on shell and one has
$`p^2=p^2=m_e^2,`$ $`Pq=0,`$ (16)
$`\overline{u}\overline{)}qu=0,`$ $`\overline{u}\overline{)}Pu=m\overline{u}u,`$ (17)
$`\overline{u}^\mu \overline{)}qu^\nu =0,`$ $`\overline{u}^\mu \overline{)}Pu^\nu =M\overline{u}^\mu u^\nu .`$ (18)
The basic loop integrals needed in this work involve either two or three denominators. We use the definition
$$[\mathrm{}]=\frac{d^4Q}{(2\pi )^4}\frac{1}{[(Q+q/2)^2\mu ^2][(Qq/2)^2\mu ^2]}$$
(19)
and the dimensionless expressions
$`I_{\pi \pi }={\displaystyle [\mathrm{}]}={\displaystyle \frac{i}{(4\pi )^2}}\mathrm{\Pi }_{\pi \pi }^{(00)},`$ (20)
$`I_{\pi \pi }^{\mu \nu }={\displaystyle [\mathrm{}]\frac{Q^\mu Q^\nu }{\mu ^2}}={\displaystyle \frac{i}{(4\pi )^2}}\left[g^{\mu \nu }\overline{\mathrm{\Pi }}_{\pi \pi }^{(00)}+\mathrm{}\right],`$ (21)
$`I_{x\pi \pi }={\displaystyle [\mathrm{}]\frac{2\mu m_e}{[(Q+P)^2m_x^2]}}={\displaystyle \frac{i}{(4\pi )^2}}\mathrm{\Pi }_{x\pi \pi }^{(000)},`$ (22)
$`I_{x\pi \pi }^\mu ={\displaystyle [\mathrm{}]\frac{(Q^\mu /\mu )\mathrm{\hspace{0.33em}\hspace{0.33em}2}\mu m_e}{[(Q+P)^2m_x^2]}}={\displaystyle \frac{i}{(4\pi )^2}}\left[{\displaystyle \frac{P^\mu }{m_e}}\mathrm{\Pi }_{x\pi \pi }^{(100)}+\mathrm{}\right],`$ (23)
$`I_{x\pi \pi }^{\mu \nu }={\displaystyle [\mathrm{}]\frac{(Q^\mu Q^\nu /\mu ^2)\mathrm{\hspace{0.33em}\hspace{0.33em}2}\mu m_e}{[(Q+P)^2m_x^2]}}={\displaystyle \frac{i}{(4\pi )^2}}\left[g^{\mu \nu }\overline{\mathrm{\Pi }}_{x\pi \pi }^{(000)}+\mathrm{}\right],`$ (24)
$`I_{x\pi \pi }^{\mu \nu \rho }={\displaystyle [\mathrm{}]\frac{(Q^\mu Q^\nu Q^\rho /\mu ^3)\mathrm{\hspace{0.33em}\hspace{0.33em}2}\mu m_e}{[(Q+P)^2m_x^2]}}={\displaystyle \frac{i}{(4\pi )^2}}\left[g^{\mu \nu }{\displaystyle \frac{P^\rho }{m_e}}\overline{\mathrm{\Pi }}_{x\pi \pi }^{(100)}+\mathrm{}\right].`$ (25)
where the ellipses indicate terms that do not contribute to the scalar form factors. The usual Feynman techniques for loop integration allow one to write the regular parts of these integrals as
$`\mathrm{\Pi }_{\pi \pi }^{(00)}={\displaystyle _0^1}𝑑a\mathrm{ln}\left({\displaystyle \frac{D_{\pi \pi }}{\mu ^2}}\right),`$ (26)
$`\overline{\mathrm{\Pi }}_{\pi \pi }^{(00)}={\displaystyle _0^1}𝑑a{\displaystyle \frac{D_{\pi \pi }}{2\mu ^2}}\mathrm{ln}\left({\displaystyle \frac{D_{\pi \pi }}{\mu ^2}}\right),`$ (27)
$`\mathrm{\Pi }_{x\pi \pi }^{(k00)}={\displaystyle _0^1}𝑑aa{\displaystyle _0^1}𝑑b[m_e(1a)/\mu ]^k\left({\displaystyle \frac{2\mu m_e}{D_{x\pi \pi }}}\right),`$ (28)
$`\overline{\mathrm{\Pi }}_{x\pi \pi }^{(k00)}={\displaystyle \frac{m_e}{\mu }}{\displaystyle _0^1}𝑑aa{\displaystyle _0^1}𝑑b[m_e(1a)/\mu ]^k\mathrm{ln}\left({\displaystyle \frac{D_{x\pi \pi }}{2\mu m_e}}\right),`$ (29)
with
$`D_{\pi \pi }=\mu ^2a(1a)q^2,`$ (30)
$`D_{x\pi \pi }=a\mu ^2+(1a)m_x^2a(1a)m_e^2a^2b(1b)q^2.`$ (31)
The dimensionless configuration space functions $`S`$ are defined as
$$S=\frac{d𝒌}{(2\pi )^3}e^{i𝒌𝒙}\mathrm{\Pi }.$$
(32)
with $`x=\mu r`$ and $`𝒌=𝒒/\mu `$. Performing the Fourier transforms, we find
$`S_{\pi \pi }^{(00)}`$ $`=`$ $`{\displaystyle \frac{1}{\pi x^2}}K_1(2x),`$ (33)
$`\overline{S}_{\pi \pi }^{(00)}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi x^4}}\left[xK_0(2x)+K_1(2x)\right],`$ (34)
$`\varphi ^2>0S_{x\pi \pi }^{(k00)}`$ $`=`$ $`{\displaystyle \frac{2m_e}{\mu }}{\displaystyle \frac{1}{\pi x}}{\displaystyle _0^1}𝑑a{\displaystyle \frac{[m_e(1a)/\mu ]^k}{a}}K_0(2\varphi x),`$ (35)
$`\varphi ^2<0S_{x\pi \pi }^{(k00)}`$ $`=`$ $`{\displaystyle \frac{m_e}{\mu }}{\displaystyle \frac{1}{x}}{\displaystyle _0^1}𝑑a{\displaystyle \frac{[m_e(1a)/\mu ]^k}{a}}\left[Y_0(2|\varphi |x)+iJ_0(2|\varphi |x)\right],`$ (36)
$`\varphi ^2>0\overline{S}_{x\pi \pi }^{(k00)}`$ $`=`$ $`{\displaystyle \frac{1}{\pi x^2}}{\displaystyle \frac{m_e}{\mu }}{\displaystyle _0^1}𝑑aa[m_e(1a)/\mu ]^k\varphi K_1(2\varphi x),`$ (37)
$`\varphi ^2<0\overline{S}_{x\pi \pi }^{(k00)}`$ $`=`$ $`{\displaystyle \frac{1}{2x^2}}{\displaystyle \frac{m_e}{\mu }}{\displaystyle _0^1}𝑑aa[m_e(1a)/\mu ]^k|\varphi |\left[Y_1(2|\varphi |x)+iJ_1(2|\varphi |x)\right],`$ (38)
with
$$\varphi ^2=[a\mu ^2+(1a)m_x^2a(1a)m_e^2]/(\mu ^2a^2).$$
(39)
## Appendix B scalar form factors
We give here the expressions for $`\stackrel{~}{\sigma }_{B_I}`$, due to triangle diagrams containing external states $`B`$ and an intermediate states $`I`$. The following interaction lagrangiansD ; H ; HHK ; BL are used
$`_{\pi NN}={\displaystyle \frac{g_A}{2f_\pi }}\left\{\overline{N}\gamma _\mu \gamma _5\tau _aN\right\}^\mu \varphi _a,`$ (40)
$`_{\pi N\mathrm{\Delta }}={\displaystyle \frac{g_{\pi N\mathrm{\Delta }}}{f_\pi }}\left\{\overline{\mathrm{\Delta }}^\mu \left[g_{\mu \nu }(Z1/2)\gamma _\mu \gamma _\nu \right]M_aN\right\}^\nu \varphi _a+h.c.,`$ (41)
$`_{\pi \mathrm{\Delta }\mathrm{\Delta }}={\displaystyle \frac{g_{\pi \mathrm{\Delta }\mathrm{\Delta }}}{f_\pi }}\left\{\overline{\mathrm{\Delta }}^\mu \left(g_{\mu \nu }\gamma _\lambda g_{\mu \lambda }\gamma _\nu g_{\lambda \nu }\gamma _\mu \right)\gamma _5T_a\mathrm{\Delta }^\nu \right\}^\lambda \varphi _a,`$ (42)
where $`\varphi `$, $`N`$ and $`\mathrm{\Delta }`$ denote pion, nucleon and delta fields, $`f_\pi `$ is the pion decay constant, $`g_A`$, $`g_{\pi N\mathrm{\Delta }}`$ and $`g_{\pi \mathrm{\Delta }\mathrm{\Delta }}`$ are coupling constants, and $`\tau `$, $`M`$ and $`T`$ are matrices that couple nucleons and deltas into isospin 1 states, with $`\tau _a\tau _a=3`$, $`M_a^{}M_a=2`$, $`M_aM_a^{}=1`$, and $`T_a^{}T_a=15/4`$. For the coupling constants we use $`g_A=1.25`$ and the $`SU(4)`$ results $`g_{\pi N\mathrm{\Delta }}=3g_A/2\sqrt{2}`$ and $`g_{\pi \mathrm{\Delta }\mathrm{\Delta }}=3g_A/5`$. We also use $`f_\pi =93`$ MeV, $`\mu =139.57`$ MeV, $`m=938.27`$ MeV and $`M=1232`$ MeV.
nucleon:
Using the loop integrals $`S`$ defined in appendix A, we obtain the following contributions to the nucleon scalar form factor
$`\stackrel{~}{\sigma }_{N_N}(r)={\displaystyle \frac{3}{4}}\left[{\displaystyle \frac{\mu g_A}{4\pi f_\pi }}\right]^22m(\mu ^3)\left\{S_{\pi \pi }^{(00)}S_{N\pi \pi }^{(100)}\right\},`$ (43)
$`\stackrel{~}{\sigma }_{N_\mathrm{\Delta }}(r)=2\left[{\displaystyle \frac{\mu g_{\pi N\mathrm{\Delta }}}{4\pi f_\pi }}\right]^2{\displaystyle \frac{(m+M)}{6M^2}}(\mu ^3)\{[(m+M)(2Mm)+2\mu ^2+{\displaystyle \frac{m\mu ^2(1\mathbf{}^2/2)}{(m+M)}}]S_{\pi \pi }^{(00)}`$
$`{\displaystyle \frac{2m\mu ^2}{(m+M)}}\overline{S}_{\pi \pi }^{(00)}`$
$`+{\displaystyle \frac{1}{2m\mu }}\left[(m^2M^2)(m+M)(2Mm)+2\mu ^2(m^2M^2)+6M^2\mu ^2(1\mathbf{}^2/2)\right]S_{\mathrm{\Delta }\pi \pi }^{(000)}`$
$`+{\displaystyle \frac{1}{2m}}[(m+M)(4mMM^2m^2)2\mu ^2(2Mm)+{\displaystyle \frac{6M^2\mu ^2}{(m+M)}}(1\mathbf{}^2/2)]S_{\mathrm{\Delta }\pi \pi }^{(100)}\}.`$ (44)
delta:
In the evaluation of the triangle diagram, the external deltas are on shell and one has the constraints $`pu^s(𝒑)=\gamma u^s(𝒑)=\overline{u}^s^{}(𝒑^{})p^{}=\overline{u}^s^{}(𝒑^{})\gamma =0`$. The $`T`$ matrix can be cast in the form
$$iT=\overline{u}_\mu ^s^{}(𝒑^{})\left\{\frac{d^4Q}{(2\pi )^4}\frac{\mathrm{\Theta }^{\mu \nu }}{[(Qq/2)^2\mu ^2][(Q+q/2)^2\mu ^2]}\right\}u_\nu ^s(𝒑),$$
(45)
with
$`\mathrm{\Theta }_N^{\mu \nu }=2\left[{\displaystyle \frac{\mu g_{\pi N\mathrm{\Delta }}}{f_\pi }}\right]^2(Q+q/2)^\mu (Qq/2)^\nu {\displaystyle \frac{m+M+\overline{)}Q}{\overline{p}^2m^2}},`$ (46)
$`\mathrm{\Theta }_\mathrm{\Delta }^{\mu \nu }={\displaystyle \frac{15}{4}}\left[{\displaystyle \frac{\mu g_{\pi \mathrm{\Delta }\mathrm{\Delta }}}{f_\pi }}\right]^2\left[g^{\mu \nu }\left(2M{\displaystyle \frac{4M^2}{\overline{p}^2M^2}}\right)\overline{)}Q{\displaystyle \frac{8}{3}}(Q+q/2)^\mu (Qq/2)^\nu {\displaystyle \frac{M\overline{)}Q}{\overline{p}^2m^2}}\right],`$ (47)
with $`\overline{p}=P+Q`$. Performing the integrals, comparing the results with eq.(2), and going to configuration space, we obtain the contributions
$`\stackrel{~}{\sigma }_{\mathrm{\Delta }_N}(r)=\left[{\displaystyle \frac{\mu g_{\pi N\mathrm{\Delta }}}{4\pi f_\pi }}\right]^2\mu (\mu ^3)\left[{\displaystyle \frac{(m+M)}{2M}}\overline{S}_{N\pi \pi }^{(000)}+{\displaystyle \frac{\mu }{2M}}\overline{S}_{N\pi \pi }^{(100)}\right],`$ (48)
$`\stackrel{~}{\sigma }_{\mathrm{\Delta }_\mathrm{\Delta }}(r)={\displaystyle \frac{15}{4}}\left[{\displaystyle \frac{\mu g_{\pi \mathrm{\Delta }\mathrm{\Delta }}}{4\pi f_\pi }}\right]^22M(\mu ^3)\left[S_{\pi \pi }^{(00)}S_{\mathrm{\Delta }\pi \pi }^{(100)}{\displaystyle \frac{2\mu }{3M}}\overline{S}_{\mathrm{\Delta }\pi \pi }^{(000)}+{\displaystyle \frac{2\mu ^2}{3M^2}}\overline{S}_{\mathrm{\Delta }\pi \pi }^{(100)}\right].`$ (49)
## Appendix C chiral symmetry
In this appendix we show that results (43,44) and (48,49) are fully compatible with standard chiral power counting by means of a covariant chiral expansionHR . It is important to note that these expressions contain a factor $`(\mu ^3)`$, which comes from the definition of the configuration space function $`S`$ and must not be included in the counting. With this previous in mind, we use the following relations among integrals,
$`S_{\pi \pi }^{(00)}=\left[1{\displaystyle \frac{\mu ^2\mathbf{}^2}{4m_e^2}}\right]S_{x\pi \pi }^{(100)}+\left[{\displaystyle \frac{\mu }{2m_e}}(1\mathbf{}^2/2)+{\displaystyle \frac{(m_e^2m_x^2)}{2\mu m_e}}\right]S_{x\pi \pi }^{(000)},`$ (50)
$`2\left[1{\displaystyle \frac{\mu ^2\mathbf{}^2}{4m_e^2}}\right]\overline{S}_{x\pi \pi }^{(000)}=\left[{\displaystyle \frac{\mu }{2m_e}}(1\mathbf{}^2/2)+{\displaystyle \frac{(m_e^2m_x^2)}{2\mu m_e}}\right]S_{\pi \pi }^{(00)}`$
$`\left[{\displaystyle \frac{(m_e^2m_x^2)^2}{4\mu ^2m_e^2}}{\displaystyle \frac{m_e^2+m_x^2}{2m_e^2}}+{\displaystyle \frac{\mu ^2}{4m_e^2}}+{\displaystyle \frac{m_x^2}{4m_e^2}}\mathbf{}^2\right]S_{x\pi \pi }^{(000)},`$ (51)
$`2\left[1{\displaystyle \frac{\mu ^2\mathbf{}^2}{4m_e^2}}\right]\overline{S}_{x\pi \pi }^{(100)}={\displaystyle \frac{1}{3}}(1\mathbf{}^2/4)S_{\pi \pi }^{(00)}`$
$`\left[{\displaystyle \frac{(m_e^2m_x^2)^2}{4\mu ^2m_e^2}}{\displaystyle \frac{m_e^2+m_x^2}{2m_e^2}}+{\displaystyle \frac{\mu ^2}{4m_e^2}}+{\displaystyle \frac{m_x^2}{4m_e^2}}\mathbf{}^2\right]S_{x\pi \pi }^{(100)},`$ (52)
which are obtained by multiplying eqs.(22-24) by $`P_\mu `$, neglecting short range terms with a single pion propagator, and going to configuration space. In the case $`m_e=M`$ and $`m_x=m`$, the expansion of eqs.(22-23) yield the leading order relations
$`\mathrm{\Pi }_{\mathrm{\Delta }\pi \pi }^{(000)}{\displaystyle \frac{(2m_e\mu )}{(m^2M^2)}}\mathrm{\Pi }_{\pi \pi }^{(00)},`$ (53)
$`\mathrm{\Pi }_{\mathrm{\Delta }\pi \pi }^{(100)}{\displaystyle \frac{(2m_e\mu )^2}{(m^2M^2)^2}}\overline{\mathrm{\Pi }}_{\pi \pi }^{(00)}.`$ (54)
Truncating the expansions at $`𝒪(q^4)`$, we find
$`\stackrel{~}{\sigma }_{N_N}(r)={\displaystyle \frac{3}{4}}\left[{\displaystyle \frac{\mu g_A}{4\pi f_\pi }}\right]^2\mu (\mu ^3)\left[(1\mathbf{}^2/2)S_{N\pi \pi }^{(000)}{\displaystyle \frac{\mu }{2m}}\mathbf{}^2S_{\pi \pi }^{(00)}\right],`$ (55)
$`\stackrel{~}{\sigma }_{N_\mathrm{\Delta }}(r)={\displaystyle \frac{4}{3}}\left[{\displaystyle \frac{\mu g_{\pi N\mathrm{\Delta }}}{4\pi f_\pi }}\right]^2{\displaystyle \frac{\mu ^2}{(Mm)}}(\mu ^3)\left[(1\mathbf{}^2/2){\displaystyle \frac{m^2}{3M^2}}(1\mathbf{}^2/4)\right]S_{\pi \pi }^{(00)},`$ (56)
$`\stackrel{~}{\sigma }_{\mathrm{\Delta }_N}(r)={\displaystyle \frac{1}{4}}\left[{\displaystyle \frac{\mu g_{\pi N\mathrm{\Delta }}}{4\pi f_\pi }}\right]^2{\displaystyle \frac{(M+m)}{M}}\mu (\mu ^3)[{\displaystyle \frac{(M+m)}{4\mu M^2}}(M^2m^2)S_{N\pi \pi }^{(100)}`$
$`+(1{\displaystyle \frac{Mm}{6M}})(1{\displaystyle \frac{\mathbf{}^2}{4}})S_{N\pi \pi }^{(000)}],`$ (57)
$`\stackrel{~}{\sigma }_{\mathrm{\Delta }_\mathrm{\Delta }}(r)={\displaystyle \frac{15}{4}}\left[{\displaystyle \frac{\mu g_{\pi \mathrm{\Delta }\mathrm{\Delta }}}{4\pi f_\pi }}\right]^2\mu (\mu ^3)[(1\mathbf{}^2/2)S_{\mathrm{\Delta }\pi \pi }^{(000)}{\displaystyle \frac{\mu \mathbf{}^2}{2M}}S_{\mathrm{\Delta }\pi \pi }^{(100)}`$
$`{\displaystyle \frac{4}{3}}\overline{S}_{\mathrm{\Delta }\pi \pi }^{(000)}+{\displaystyle \frac{4\mu }{3M}}\overline{S}_{\mathrm{\Delta }\pi \pi }^{(100)}].`$ (58)
These results, except for eq. 57, which contains imaginary terms, are compatible with chiral counting rules. Contributions begin at $`𝒪(q^3)`$ for diagrams in which internal and external baryons are identical and at $`𝒪(q^4)`$ when this does not happen.
## Appendix D chiral limit
In this section we show that our model for the nucleon $`\sigma `$-term is consistent with the standard ChPT expansion. In the paper by Becher and LeutwylerBL , one finds, using our notation
$`\sigma _N=4c_1\mu ^2{\displaystyle \frac{9g_A^2\mu ^3}{64\pi f_\pi ^2}}{\displaystyle \frac{3\mu ^4}{16\pi ^2f_\pi ^2m}}\left(g_A^28c_1m+c_2m+4c_3m\right)\mathrm{ln}{\displaystyle \frac{\mu }{m}}`$
$`{\displaystyle \frac{3\mu ^4}{64\pi ^2f_\pi ^2m}}\left(3g_A^28c_1m+4c_3m\right)+2\overline{e}_1,`$ (59)
where $`c_i`$ and $`e_1`$ are, respectively, low enegy constants (LECs) from the $`_N^{(2)}`$ and $`_N^{(4)}`$ lagrangians. The bar over $`e_1`$ indicates that it has been renormalized.
In order to expand our $`\sigma _N`$, we use in eq.(53) the result
$$S_{N\pi \pi }^{(000)}\frac{e^{2x}}{2x^2}+\frac{\mu }{m\pi x^2}\left[xK_0(2x)+K_1(2x)\right],$$
(60)
which holdsHRR for $`\mu /m<<1`$. This allows integrations in eq.(10) to be performed analytically and one finds
$`\sigma _N={\displaystyle \frac{4}{3}}\pi R^3f_\pi ^2\mu ^2+{\displaystyle \frac{3g_A^2\mu ^3}{16\pi f_\pi ^2}}\{({\displaystyle \frac{1}{4}}+{\displaystyle \frac{1}{2\mu R}})e^{2\mu R}{\displaystyle \frac{\mu }{2m\pi }}[4K_0(2\mu R)`$ (61)
$`+(2\mu R+{\displaystyle \frac{6}{2\mu R}})K_1(2\mu R)]\}+{\displaystyle \frac{g_{\pi N\mathrm{\Delta }}^2\mu ^4}{6\pi ^2f_\pi ^2(Mm)}}[K_0(2\mu R)+(3{\displaystyle \frac{m^2}{2M^2}}){\displaystyle \frac{K_1(2\mu R)}{2\mu R}}].`$
An expansion for small values of $`\mu `$ yields
$`\sigma _N={\displaystyle \frac{4}{3}}\pi R^3f_\pi ^2\mu ^2`$ (62)
$`+{\displaystyle \frac{3g_A^2}{16\pi f_\pi ^2}}\left[{\displaystyle \frac{\mu ^2}{2R^2}}\left(R{\displaystyle \frac{1}{m\pi }}\right){\displaystyle \frac{3\mu ^3}{4}}{\displaystyle \frac{\mu ^4}{m\pi }}\mathrm{ln}\mu R+{\displaystyle \frac{\mu ^4(54\gamma )}{4m\pi }}+{\displaystyle \frac{\mu ^4R}{2}}\right]`$
$`+{\displaystyle \frac{g_{\pi N\mathrm{\Delta }}^2}{12\pi ^2f_\pi ^2(Mm)}}\left[{\displaystyle \frac{\mu ^2}{R^2}}\left(3{\displaystyle \frac{m^2}{2M^2}}\right)+\mu ^4\left(4{\displaystyle \frac{m^2}{M^2}}\right)\mathrm{ln}\mu R+\mu ^4\left(3+4\gamma +{\displaystyle \frac{m^2}{2M^2}}(12\gamma )\right)\right],`$
where $`\gamma `$ is the Euler constant. This result reproduces the first three terms of eq.(59), provided one absorbs the factors proportional to $`\mu ^2`$ into the definition of $`c_1`$, uses the delta contributions to the $`c_i`$, which are given by
$$c_1^\mathrm{\Delta }=0,c_2^\mathrm{\Delta }=\frac{4g_{\pi N\mathrm{\Delta }}^2m^2}{9M^2(Mm)},c_3^\mathrm{\Delta }=\frac{4g_{\pi N\mathrm{\Delta }}^2}{9(Mm)},$$
(63)
and chooses the value $`R=1/m`$ for the cutting radius. As the renormalized constant $`\overline{e}_1`$ contains factors proportional to $`\mu ^4`$, terms of this kind need not to coincide.
###### Acknowledgements.
The works by I.P.C. and D.O.S. were supported by CNPq and the work by G.R.S.Z., by FAPESP (Brazilian agencies). |
warning/0507/math0507487.html | ar5iv | text | # Admissible Dirichlet Series
## 1. Introduction
In 1956 Hayman defined admissible functions—they are analytic in a neighborhood of 0 and one can use the saddle point method to estimate the coefficients of the power series expansion of such functions. They include the functions $`e^z`$ and $`\mathrm{exp}{\displaystyle \frac{1}{1z}}`$, are closed under product (of series with the same radius of convergence) and under exponentiation.
In this paper a notion of admissibility for functions that have Dirichlet series expansions is proposed. We believe that this is a viable analog of Hayman’s definition because (1) this notion of admissible generalizes the conditions of Tenenbaum in , (2) there is a fundamental theorem (Theorem 7) that is the analog of Hayman’s fundamental theorem, and (3) a product of admissible Dirichlet series (with the same abscissa of convergence) is again admissible.
## 2. Definition of Admissible
###### Theorem 1.
Suppose the function $`𝐅(s)`$
* has a Dirchlet series expansion $`𝐅(s)=_{n1}f(n)n^s`$, where the coefficients $`f(n)`$ are nonnegative real, with $`f(1)>0`$,
* has abscissa of (absolute) convergence $`\alpha [0,\mathrm{})`$, and
* $`𝐅(s)`$ has no zeros in its halfplane of convergence.
Then there exists a Dirichlet series $`𝐇(s)`$ with real coefficients such that $`𝐅(s)=e^{𝐇(s)}`$ for $`\sigma >\alpha `$ where $`s=\sigma +it`$.
###### Proof.
This is a slight specialization of Theorem 11.14 in Apostol . ∎
Let $``$ be the set of real numbers. Assuming that $`𝐅(s)`$ satisfies (A1)–(A3) and $`𝐅(s)=e^{𝐇(s)}`$, we will need the basic facts about a Taylor series expansion with remainder of $`𝐇(\sigma +it)`$ about $`t=0`$. With
(1)
$$𝐚(s):=𝐇^{}(s)𝐛(s):=𝐇^{\prime \prime }(s)𝐜(s):=𝐇^{\prime \prime \prime }(s)$$
we have, for $`\sigma >\alpha `$ and $`t`$,
(2) $`𝐇(\sigma +it)`$ $`=`$ $`𝐇(\sigma )+i𝐚(\sigma )t{\displaystyle \frac{𝐛(\sigma )}{2}}t^2+𝖱(\sigma +it)`$
with the remainder term given by
(3) $`𝖱(\sigma +it)`$ $`=`$ $`{\displaystyle \frac{i}{2}}{\displaystyle _0^t}𝐜(\sigma +iv)(tv)^2𝑑v.`$
Thus we can write
(4) $`{\displaystyle \frac{𝐅\left(\sigma +it\right)}{𝐅(\sigma )}}`$ $`=`$ $`\mathrm{exp}\left(i𝐚(\sigma )t{\displaystyle \frac{𝐛(\sigma )}{2}}t^2+𝖱\left(\sigma +it\right)\right)`$
(5) $`{\displaystyle \frac{\left|𝐅\left(\sigma +it\right)\right|}{𝐅(\sigma )}}`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \frac{𝐛(\sigma )}{2}}t^2\right)\left|\mathrm{exp}\left(𝖱\left(\sigma +it\right)\right)\right|.`$
###### Definition 2.
Suppose $`𝐅(s)`$ satisfies (A1)–(A3), with $`𝐅(s)=e^{𝐇(s)}`$. Let $`𝐚(s)`$ and $`𝐛(s)`$ be the first two derivatives of $`𝐇(s)`$ as in (1), let $`𝖱(\sigma +it)`$ be the remainder term of the Taylor expansion for $`𝐇(\sigma +it)`$ as in (2), and suppose there is a function $`\delta :(\alpha ,\beta )(0,1)`$, for some $`\beta >\alpha `$, such that as $`\sigma \alpha +`$
* $`\delta (\sigma )0`$
* $`\sigma ^2𝐛(\sigma )\mathrm{}`$
* $`𝐛(\sigma )\mathrm{exp}\left(𝐛(\sigma )\delta (\sigma )^2\right)0`$
* $`𝖱(\sigma +it)0\text{uniformly for }|t|\delta (\sigma )`$
* $`{\displaystyle \frac{\sigma \sqrt{𝐛(\sigma )}}{𝐅(\sigma )}}{\displaystyle _{|t|\delta (\sigma )}}\left|𝐅(\sigma +it)\right|{\displaystyle \frac{dt}{\sigma ^2+t^2}}0`$.
Then we say $`𝐅(s)`$ is admissible, as witnessed by $`\delta (\sigma )`$.
###### Remark 3.
Except for $`\mathrm{\S }`$5 we can replace (A6) and (A8) by the following, giving a more general notion of admissible:
* $`𝐛(\sigma )\delta (\sigma )^2\mathrm{}`$
* $`{\displaystyle \frac{\sigma \sqrt{𝐛(\sigma )}}{𝐅(\sigma )}}{\displaystyle _{|t|\delta (\sigma )}}𝐅(\sigma +it)x^{it}{\displaystyle \frac{dt}{(\sigma +it)(\sigma +1+it)}}0\text{uniformly for }x>0`$
The full strength of (A6) and (A8) are used to prove the product theorem in $`\mathrm{\S }`$5.
## 3. Asymptotic Estimates and Regular Variation
In this section we assume that $`𝐅(s)`$ is admissible, witnessed by $`\delta (\sigma )`$. (A5) implies
(6)
$$𝐛(\sigma )\mathrm{}\text{as }\sigma \alpha +,$$
so there is a $`\beta >\alpha `$ such that
(7)
$$𝐛(\sigma )>0\text{for }\sigma (\alpha ,\beta ).$$
We will consistently use $`\beta `$ as a number in $`(\alpha ,\mathrm{})`$ such that $`𝐛(\sigma )>0`$ on $`(\alpha ,\beta )`$, keeping the original requirement that $`\delta (\sigma )`$ be defined on $`(\alpha ,\beta )`$. From (A6) and (6) we have
(8)
$$𝐛(\sigma )\delta (\sigma )^2\mathrm{}\text{as }\sigma \alpha +.$$
###### Definition 4.
The partial sums of the coefficients of $`𝐅(s)`$ and its integral are denoted as follows:
$`F(x)`$ $`:={\displaystyle \underset{nx}{}}f(n)`$
$`\widehat{F}(x)`$ $`:={\displaystyle _1^x}F(u)𝑑u.`$
Hayman makes a direct application of Cauchy’s integral formula to express the coefficients of a power series. Tenenbaum makes a direct application of Perron’s integral formula to express $`F(x)`$. The next lemma, where the Perron formula is used to express $`\widehat{F}(x)`$, is used to derive a formula that leads to the verification of regular variation at infinity for $`F(x)`$. <sup>1</sup><sup>1</sup>1Oppenheim (, ) appears to state this lemma for a particular choice of $`𝐅(s)`$, the zeta function connected with ‘Factorisatio Numerorum’, but it is a general result. On pages 210–211 of , one of the sources cited for Lemma 5, the several occurrences of $`\mathrm{exp}(𝐒(z))`$ need to be replaced by $`𝐒(z).`$
###### Lemma 5.
For $`x>0`$ and $`c>\alpha `$,
$$\widehat{F}(x)=\frac{1}{2\pi i}_{ci\mathrm{}}^{c+i\mathrm{}}𝐅(s)\frac{x^{s+1}}{s(s+1)}𝑑s.$$
###### Proof.
See or Lemma 11.22 in . ∎
An elementary estimate will also be needed.
###### Lemma 6.
For $`h,\lambda >0`$ and $`\kappa `$,
$$_h^he^{i\kappa \lambda u^2}𝑑u=\sqrt{\frac{\pi }{\lambda }}e^{\kappa ^2/4\lambda }\left(1+\epsilon (h,\kappa ,\lambda )\right),$$
where
$$\left|\epsilon (h,\kappa ,\lambda )\right|<\frac{2}{h\sqrt{\lambda }}.$$
The following gives the fundamental formula for $`\widehat{F}(x)`$. It is this form, rather than the asymptotics that can be obtained by specializing $`\sigma `$ to be the saddle point $`\sigma _x`$, that leads to a verification of regular variation at infinity.
###### Theorem 7.
For $`x>0`$ and $`\sigma >\alpha `$
$$\widehat{F}(x)=\frac{x^{\sigma +1}𝐅(\sigma )}{\sigma (\sigma +1)\sqrt{2\pi 𝐛(\sigma )}}\left(\mathrm{exp}\left(\frac{\left(𝐚(\sigma )+\mathrm{log}x\right)^2}{2𝐛(\sigma )}\right)+R(x,\sigma )\right)$$
where
$$R(x,\sigma )0\text{as }\sigma \alpha +,\text{uniformly for }x>0.$$
###### Proof.
For $`x>0`$ and $`\sigma >\alpha `$
$`\widehat{F}(x)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle _{\sigma i\mathrm{}}^{\sigma +i\mathrm{}}}𝐅(s){\displaystyle \frac{x^{s+1}}{s(s+1)}}𝑑s\text{by Lemma }\text{5}`$
$`=`$ $`{\displaystyle \frac{x^{\sigma +1}}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{𝐅(\sigma +it)x^{it}}{(\sigma +it)(\sigma +1+it)}}𝑑t`$
$`=`$ $`{\displaystyle \frac{x^{\sigma +1}}{2\pi }}\left(J_1(\sigma ,x)+J_2(\sigma ,x)\right)`$
where
$`J_1(\sigma ,x)`$ $`={\displaystyle _{\delta (\sigma )}^{\delta (\sigma )}}{\displaystyle \frac{𝐅(\sigma +it)x^{it}}{(\sigma +it)(\sigma +1+it)}}𝑑t`$
$`J_2(\sigma ,x)`$ $`={\displaystyle _{|t|\delta (\sigma )}}{\displaystyle \frac{𝐅(\sigma +it)x^{it}}{(\sigma +it)(\sigma +1+it)}}𝑑t.`$
Since
$`\left|J_2(\sigma ,x)\right|`$ $`{\displaystyle _{|t|\delta (\sigma )}}\left|𝐅(\sigma +it)\right|{\displaystyle \frac{dt}{\sigma ^2+t^2}}`$
by (A8) we immediately have
$`J_2(\sigma ,x)`$ $`={\displaystyle \frac{\sqrt{2\pi }𝐅(\sigma )}{\sigma (\sigma +1)\sqrt{𝐛(\sigma )}}}\mathrm{o}(1)`$
as $`\sigma \alpha +`$, uniformly for $`x>0`$.
Let us collect some simple facts before estimating $`J_1(\sigma ,x)`$. We easily have
(9) $`{\displaystyle \frac{\sigma +1}{\sigma +1+it}}`$ $`=1+\mathrm{o}(1)`$
as $`\sigma \alpha +`$, uniformly for $`|t|\delta (\sigma )`$, since
$`{\displaystyle \frac{\sigma +1}{\sigma +1+it}}`$ $`=1{\displaystyle \frac{it}{\sigma +1+it}}`$
and
$`\left|{\displaystyle \frac{it}{\sigma +1+it}}\right|`$ $`{\displaystyle \frac{\delta (\sigma )}{\sigma +1}}=\mathrm{o}(1)\text{by }\text{(A4)}.`$
Also
(10) $`\left|{\displaystyle \frac{\sigma }{\sigma +it}}1\right|`$ $`=\left|{\displaystyle \frac{it}{\sigma +it}}\right|{\displaystyle \frac{|t|}{\sigma }}.`$
Let
$$a(\sigma ,x)=𝐚(\sigma )+\mathrm{log}x.$$
Then by (4), (A7), (9) and (10), for $`\sigma (\alpha ,\beta )`$ and $`x>0`$
$`J_1(\sigma ,x)`$ $`=`$ $`{\displaystyle _{|t|\delta (\sigma )}}{\displaystyle \frac{𝐅(\sigma +it)x^{it}}{(\sigma +it)(\sigma +it+1)}}𝑑t`$
$`=`$ $`{\displaystyle \frac{𝐅(\sigma )}{\sigma (\sigma +1)}}{\displaystyle _{|t|\delta (\sigma )}}\mathrm{exp}\left(ia(\sigma ,x)t{\displaystyle \frac{𝐛(\sigma )}{2}}t^2+𝖱(\sigma +it)\right)`$
$`{\displaystyle \frac{\sigma (\sigma +1)}{(\sigma +it)(\sigma +1+it)}}dt`$
$`=`$ $`{\displaystyle \frac{𝐅(\sigma )}{\sigma (\sigma +1)}}{\displaystyle _{|t|\delta (\sigma )}}\mathrm{exp}\left(ia(\sigma ,x)t{\displaystyle \frac{𝐛(\sigma )}{2}}t^2\right)\left(1+\mathrm{o}(1)+\mathrm{O}\left({\displaystyle \frac{|t|}{\sigma }}\right)\right)𝑑t`$
$`=`$ $`\underset{J_{11}(\sigma ,x)}{\underset{}{{\displaystyle \frac{𝐅(\sigma )}{\sigma (\sigma +1)}}{\displaystyle _{|t|\delta (\sigma )}}\mathrm{exp}\left(ia(\sigma ,x)t{\displaystyle \frac{𝐛(\sigma )}{2}}t^2\right)\left(1+\mathrm{o}(1)\right)𝑑t}}`$
$`+\underset{J_{12}(\sigma ,x)}{\underset{}{{\displaystyle \frac{𝐅(\sigma )}{\sigma (\sigma +1)}}{\displaystyle _{|t|\delta (\sigma )}}\mathrm{exp}\left(ia(\sigma ,x)t{\displaystyle \frac{𝐛(\sigma )}{2}}t^2\right)\mathrm{O}\left({\displaystyle \frac{|t|}{\sigma }}\right)𝑑t}}.`$
For $`J_{12}(\sigma ,x)`$ we have, for $`\sigma (\alpha ,\beta )`$ and $`x>0`$,
$`\left|J_{12}(\sigma ,x)\right|`$ $`=\mathrm{O}\left({\displaystyle \frac{𝐅(\sigma )}{\sigma (\sigma +1)}}{\displaystyle _{|t|\delta (\sigma )}}\mathrm{exp}\left({\displaystyle \frac{𝐛(\sigma )}{2}}t^2\right){\displaystyle \frac{|t|}{\sigma }}𝑑t\right)`$
$`=\mathrm{O}\left({\displaystyle \frac{𝐅(\sigma )}{\sigma ^2(\sigma +1)}}{\displaystyle _0^{\mathrm{}}}\mathrm{exp}\left({\displaystyle \frac{𝐛(\sigma )}{2}}t^2\right)t𝑑t\right)`$
$`=\mathrm{O}\left({\displaystyle \frac{𝐅(\sigma )}{\sigma (\sigma +1)\sqrt{𝐛(\sigma )}}}\left({\displaystyle \frac{1}{\sigma \sqrt{𝐛(\sigma )}}}\right)\right).`$
Thus by (A5)
$`J_{12}(\sigma ,x)`$ $`={\displaystyle \frac{\sqrt{2\pi }𝐅(\sigma )}{\sigma (\sigma +1)\sqrt{𝐛(\sigma )}}}\mathrm{o}(1)`$
as $`\sigma \alpha +`$, uniformly for $`x>0`$.
From Lemma 6 we have, for $`\sigma (\alpha ,\beta )`$ and $`x>0`$,
$`J_{11}(\sigma ,x)`$ $`=`$ $`{\displaystyle \frac{𝐅(\sigma )}{\sigma (\sigma +1)}}{\displaystyle _{|t|\delta (\sigma )}}\mathrm{exp}\left(ia(\sigma ,x)t{\displaystyle \frac{𝐛(\sigma )}{2}}t^2\right)\left(1+\mathrm{o}(1)\right)𝑑t`$
$`=`$ $`{\displaystyle \frac{𝐅(\sigma )}{\sigma (\sigma +1)}}\sqrt{{\displaystyle \frac{\pi }{𝐛(\sigma )/2}}}(\mathrm{exp}\left({\displaystyle \frac{a(\sigma ,x)^2}{2𝐛(\sigma )}}\right)`$
$`+\epsilon (\delta (\sigma ),a(\sigma ,x),𝐛(\sigma )/2)+\mathrm{o}(1))`$
$`=`$ $`{\displaystyle \frac{\sqrt{2\pi }𝐅(\sigma )}{\sigma (\sigma +1)\sqrt{𝐛(\sigma )}}}\left(\mathrm{exp}\left({\displaystyle \frac{\left(𝐚(\sigma )+\mathrm{log}(x)\right)^2}{2𝐛(\sigma )}}\right)+\mathrm{o}(1)\right)`$
as $`\sigma \alpha +`$, uniformly for $`x>0`$ since by Lemma 6 and (8)
$$\left|\epsilon (\delta (\sigma ),a(\sigma ,x),𝐛(\sigma )/2)\right|<\frac{2}{\delta (\sigma )\sqrt{𝐛(\sigma )/2}}=\mathrm{o}(1).$$
Combining these results we have
$`J(\sigma ,x)`$ $`:=`$ $`J_1(\sigma ,x)+J_2(\sigma ,x)`$
$`=`$ $`{\displaystyle \frac{\sqrt{2\pi }𝐅(\sigma )}{\sigma (\sigma +1)\sqrt{𝐛(\sigma )}}}\left(\mathrm{exp}\left({\displaystyle \frac{\left(𝐚(\sigma )+\mathrm{log}(x)\right)^2}{2𝐛(\sigma )}}\right)+\mathrm{o}(1)\right)`$
as $`\sigma \alpha +`$, uniformly for $`x>0`$, and the proof of the theorem is completed by observing that
$$\widehat{F}(x)=\frac{x^{\sigma +1}}{2\pi }J(\sigma ,x).$$
###### Corollary 8.
$`𝐚(\sigma )`$ is strictly increasing on $`(\alpha ,\beta )`$ and as $`\sigma \alpha +`$
* $`𝐚(\sigma )\mathrm{}`$
* $`{\displaystyle \frac{𝐚(\sigma )^2}{𝐛(\sigma )}}\mathrm{}`$
* $`\left(\sigma \alpha \right)𝐚(\sigma )\mathrm{}`$.
###### Proof.
We know that $`𝐚^{}(\sigma )=𝐛(\sigma )>0`$ on $`(\alpha ,\beta )`$, so $`𝐚(\sigma )`$ is strictly increasing on $`(\alpha ,\beta )`$. Now $`\widehat{F}(1)=0`$, so by Theorem 7 with $`x=1`$ we have for $`\sigma (\alpha ,\beta )`$
$$0=\frac{𝐅(\sigma )}{\sigma (\sigma +1)\sqrt{2\pi 𝐛(\sigma )}}\left(\mathrm{exp}\left(\frac{𝐚(\sigma )^2}{2𝐛(\sigma )}\right)+R(1,\sigma )\right)$$
so
$$0=\mathrm{exp}\left(\frac{𝐚(\sigma )^2}{2𝐛(\sigma )}\right)+R(1,\sigma ).$$
Therefore
$$\mathrm{exp}\left(\frac{𝐚(\sigma )^2}{2𝐛(\sigma )}\right)0\text{ as }\sigma \alpha +,$$
and thus (b) holds. Since $`𝐚(\sigma )`$ decreases on $`(\alpha ,\beta )`$ as $`\sigma \alpha +`$ and $`𝐛(\sigma )\mathrm{}`$ by (6) we see that (a) follows from (b).
To prove (c) we first use (a) to choose a $`\gamma (\alpha ,\beta )`$ so that $`𝐚(\sigma )<0`$ for $`\sigma (\alpha ,\gamma )`$. Then for $`\alpha <\sigma _1<\sigma <\gamma `$ we have, by the mean value theorem,
$`{\displaystyle \frac{1}{𝐚(\sigma )}}{\displaystyle \frac{1}{𝐚(\sigma _1)}}`$ $`=`$ $`{\displaystyle \frac{𝐚^{}(\xi )}{𝐚(\xi )^2}}(\sigma \sigma _1)\text{for some }\xi (\sigma _1,\sigma )`$
$`=`$ $`{\displaystyle \frac{𝐛(\xi )}{𝐚(\xi )^2}}(\sigma \sigma _1)`$
$`=`$ $`\mathrm{o}(\sigma \sigma _1)\text{as }\sigma \alpha +`$
$`=`$ $`\mathrm{o}(\sigma \alpha )\text{as }\sigma \alpha +.`$
Letting $`\sigma _1\alpha +`$
$$\frac{1}{𝐚(\sigma )}=\mathrm{o}(\sigma \alpha )\text{ as }\sigma \alpha +$$
so $`(\sigma \alpha )𝐚(\sigma )\mathrm{}`$ as $`\sigma \alpha +`$.
In view of Corollary 8(a), from now on we will assume that $`\beta `$ was chosen small enough that $`𝐚(\sigma )<0`$ for $`\sigma (\alpha ,\beta )`$.
Notice that $`𝐚(\sigma )\mathrm{}`$ as $`\sigma \alpha +`$ implies that for $`x`$ sufficiently large the equation $`𝐚(\sigma )+\mathrm{log}x=0`$ has, by the continuity of $`𝐚(\sigma )`$, a solution. In particular since $`𝐚^{}(\sigma )`$ is positive on $`(\alpha ,\beta )`$, for
$$xx_0:=\mathrm{exp}\left(𝐚(\beta )\right)+\mathrm{\hspace{0.17em}1}$$
one has a unique solution in $`(\alpha ,\beta )`$.
###### Definition 9.
For $`xx_0`$ (as just described) let $`\sigma _x`$ be the unique solution for $`\sigma (\alpha ,\beta )`$ to the equation
(11)
$$𝐚(\sigma )+\mathrm{log}x=0.$$
The function $`\sigma _x`$ is strictly decreasing on $`[x_0,\mathrm{})`$ and
(12)
$$\sigma _x\alpha \text{+}\text{as}x\mathrm{}$$
since $`\underset{x\mathrm{}}{lim}𝐚(\sigma _x)=\underset{x\mathrm{}}{lim}\mathrm{log}x=\mathrm{}`$.
Also note that if one puts $`\sigma =\sigma _x`$ (where $`xx_0`$) in the expression for $`\widehat{F}(x)`$ in Theorem 7 then it simplifies to
$$\widehat{F}(x)=\frac{x^{\sigma _x+1}𝐅(\sigma _x)}{\sigma _x(\sigma _x+1)\sqrt{2\pi 𝐛(\sigma _x)}}\left(1+R(x,\sigma _x)\right),$$
where $`R(x,\sigma _x)0`$ as $`x\mathrm{}`$. So we have the following.
###### Corollary 10.
$$\widehat{F}(x)\frac{x^{\sigma _x+1}𝐅(\sigma _x)}{\sigma _x(\sigma _x+1)\sqrt{2\pi 𝐛(\sigma _x)}}\text{as }x\mathrm{}\text{.}$$
The choice of $`\sigma =\sigma _x`$ is what is commonly meant by ‘finding the saddlepoint’, and the resulting formula for $`\widehat{F}(x)`$ is the result of ‘applying the saddlepoint method’. In reality the value $`s=\sigma _x`$ is usually only near a saddle point of the integrand of the integral in Lemma 5, that is, a point $`s`$ where the derivative of the integrand vanishes. By choosing the line of integration of this integral to pass through (a point near) the saddle point one hopes to concentrate the value of the integral in a small neighborhood of the real axis. Indeed, that is what happens for admissible functions. In the proof of Theorem 7, the value of $`\widehat{F}(x)`$ is concentrated in the integral $`J_1(x)`$ when $`\sigma =\sigma _x`$ as $`x\mathrm{}`$, leading to Corollary 10 above.
###### Corollary 11.
As $`\sigma \alpha +`$
* $`{\displaystyle \frac{𝐅(\sigma )}{\sigma \sqrt{𝐛(\sigma )}}}\mathrm{}`$ and
* $`𝐅(\sigma )\mathrm{}`$.
###### Proof.
Note that $`\widehat{F}(2)>0`$ (since $`f(1)>0`$). Then for $`\sigma (\alpha ,\beta )`$, by Theorem 7
(13)
$$\widehat{F}(2)=\frac{2^{\sigma +1}𝐅(\sigma )}{\sigma (\sigma +1)\sqrt{2\pi 𝐛(\sigma )}}\left(\mathrm{exp}\left(\frac{\left(𝐚(\sigma )+\mathrm{log}2\right)^2}{2𝐛(\sigma )}\right)+R(2,\sigma )\right).$$
By Corollary 8(a) there is a $`\gamma (\alpha ,\beta )`$ such that $`𝐚(\sigma )`$ is negative on $`(\alpha ,\gamma )`$, and thus nonzero. For $`\sigma (\alpha ,\gamma )`$ we then have
$`{\displaystyle \frac{\left(𝐚(\sigma )+\mathrm{log}2\right)^2}{𝐛(\sigma )}}`$ $`=`$ $`{\displaystyle \frac{𝐚(\sigma )^2}{𝐛(\sigma )}}\left(1+{\displaystyle \frac{2\mathrm{log}2}{𝐚(\sigma )}}+{\displaystyle \frac{(\mathrm{log}2)^2}{𝐚(\sigma )^2}}\right).`$
By Corollary 8 the right hand side of this equation goes to $`\mathrm{}`$ as $`\sigma \alpha +`$, so
$$\mathrm{exp}\left(\frac{\left(𝐚(\sigma )+\mathrm{log}2\right)^2}{2𝐛(\sigma )}\right)0\text{ as }\sigma \alpha +.$$
From Theorem 7 we know $`R(2,\sigma )0`$ as $`\sigma \alpha +`$; and clearly
$$\frac{2^{\sigma +1}}{(\sigma +1)\sqrt{2\pi }}\frac{2^{\alpha +1}}{(\alpha +1)\sqrt{2\pi }}<\mathrm{}\text{ as }\sigma \alpha +.$$
The left side of (13) is a positive constant, so it follows that part (a) of this Corollary must hold: $`{\displaystyle \frac{𝐅(\sigma )}{\sigma \sqrt{𝐛(\sigma )}}}\mathrm{}`$ as $`\sigma \alpha +.`$ Then part (a) and (A5) give $`𝐅(\sigma )\mathrm{}`$ as $`\sigma \alpha +`$, which is part (b). ∎
The next corollary shows that as $`\sigma \alpha +`$ we have $`𝐅(\sigma )`$ growing much faster than any power of $`𝐚(\sigma )`$ or $`𝐛(\sigma )`$. This leads in turn to the fact that $`𝐅(\sigma )`$ grows much faster than any power of $`\sigma \alpha `$. Consequently $`𝐅(s)`$ cannot have a pole at $`\alpha `$.
###### Corollary 12.
* For all $`\epsilon >0`$,
$$𝐚(\sigma )=\mathrm{o}\left(𝐅(\sigma )^\epsilon \right)\text{and}𝐛(\sigma )=\mathrm{o}\left(𝐅(\sigma )^\epsilon \right)\text{ as }\sigma \alpha +.$$
* For all $`r`$,
$$(\sigma \alpha )^r𝐅(\sigma )\mathrm{}\text{ as }\sigma \alpha +.$$
###### Proof.
We break the proof of (a) into two claims.
*Claim 1*: For all $`\epsilon >0`$ and all $`\gamma (\alpha ,\beta )`$ there is a $`\sigma (\alpha ,\gamma )`$ such that
$$\frac{|𝐚(\sigma )|}{𝐅(\sigma )^\epsilon }<\mathrm{\hspace{0.17em}1}.$$
Assume not. Then we can choose $`\epsilon >0`$ and $`\gamma (\alpha ,\beta )`$ such that for all $`\sigma (\alpha ,\gamma )`$
(14)
$$\frac{|𝐚(\sigma )|}{𝐅(\sigma )^\epsilon }\mathrm{\hspace{0.17em}1}.$$
Then for $`\alpha <\sigma _1<\sigma <\gamma `$ by the mean value theorem
$`{\displaystyle \frac{1}{𝐅(\sigma )^\epsilon }}{\displaystyle \frac{1}{𝐅(\sigma _1)^\epsilon }}`$ $`=`$ $`{\displaystyle \frac{\epsilon 𝐅^{}(\xi )}{𝐅(\xi )^{1+\epsilon }}}(\sigma \sigma _1)\text{for some }\xi (\sigma _1,\sigma )`$
$`=`$ $`\epsilon {\displaystyle \frac{𝐚(\xi )}{𝐅(\xi )^\epsilon }}(\sigma \sigma _1)`$
$`=`$ $`\epsilon {\displaystyle \frac{|𝐚(\xi )|}{𝐅(\xi )^\epsilon }}(\sigma \sigma _1)`$
$``$ $`\epsilon (\sigma \sigma _1)\text{by (}\text{14}\text{).}`$
Letting $`\sigma _1\alpha +`$ gives $`1/𝐅(\sigma )^\epsilon \epsilon (\sigma \alpha )`$; so
(15)
$$𝐅(\sigma )\left(\frac{1}{\epsilon (\sigma \alpha )}\right)^{1/\epsilon }\text{for }\sigma (\alpha ,\gamma ).$$
We can also assume that $`\gamma (\alpha ,\beta )`$ is such that $`(\sigma \alpha )\left|𝐚(\sigma )\right|>2/\epsilon `$ for $`\sigma (\alpha ,\gamma )`$ by Corollary 8(c). So for $`\alpha <\sigma <\sigma _2<\gamma `$
$$\frac{𝐅^{}(u)}{𝐅(u)}=𝐚(u)>\frac{2}{\epsilon (u\alpha )}\text{for }u[\sigma ,\sigma _2]$$
which implies that
$$_\sigma ^{\sigma _2}\frac{𝐅^{}(u)}{𝐅(u)}𝑑u>_\sigma ^{\sigma _2}\frac{2}{\epsilon (u\alpha )}𝑑u,$$
that is,
$$\left(\mathrm{log}𝐅(\sigma _2)\mathrm{log}𝐅(\sigma )\right)>\frac{2}{\epsilon }\mathrm{log}\left(\frac{\sigma _2\alpha }{\sigma \alpha }\right).$$
From this inequality and (15)
$$\mathrm{log}𝐅(\sigma _2)+\frac{2}{\epsilon }\mathrm{log}\left(\frac{\sigma _2\alpha }{\sigma \alpha }\right)<\mathrm{log}𝐅(\sigma )\frac{1}{\epsilon }\left(\mathrm{log}\frac{1}{\epsilon }+\mathrm{log}\frac{1}{\sigma \alpha }\right).$$
Thus
$`{\displaystyle \frac{1}{\epsilon }}\mathrm{log}{\displaystyle \frac{1}{\sigma \alpha }}`$ $`>`$ $`\mathrm{log}𝐅(\sigma _2)+{\displaystyle \frac{2}{\epsilon }}\mathrm{log}(\sigma _2\alpha ){\displaystyle \frac{1}{\epsilon }}\mathrm{log}{\displaystyle \frac{1}{\epsilon }}+{\displaystyle \frac{2}{\epsilon }}\mathrm{log}{\displaystyle \frac{1}{\sigma \alpha }}`$
$`=`$ $`C+{\displaystyle \frac{2}{\epsilon }}\mathrm{log}{\displaystyle \frac{1}{\sigma \alpha }}`$
where $`C`$ is independent of $`\sigma `$. Hence
$$1>\frac{C\epsilon }{\mathrm{log}\left(1/(\sigma \alpha )\right)}+2\mathrm{\hspace{0.33em}2}\text{as }\sigma \alpha +,$$
which is a contradiction, proving Claim 1.
*Claim 2*: For all $`\epsilon >0`$ there is a $`\gamma (\alpha ,\beta )`$ such that
(16)
$$\frac{|𝐚(\sigma )|}{𝐅(\sigma )^\epsilon }<1\text{for }\sigma (\alpha ,\gamma ]\text{.}$$
Let $`\epsilon >0`$ be given. From Claim 1 and Corollary 8(b) we know that there exists a $`\gamma (\alpha ,\beta )`$ such that
(17)
$$\frac{|𝐚(\gamma )|}{𝐅(\gamma )^\epsilon }<1\text{and}\frac{𝐛(\sigma )}{𝐚(\sigma )^2}<\epsilon \text{for }\sigma (\alpha ,\gamma ].$$
We will show this $`\gamma `$ is such that (16) holds. Otherwise there is a $`\sigma (\alpha ,\gamma )`$ such that $`|𝐚(\sigma )|/𝐅(\sigma )^\epsilon 1`$. By the intermediate value theorem there must be a $`\sigma (\alpha ,\gamma )`$ such that $`|𝐚(\sigma )|/𝐅(\sigma )^\epsilon =1`$. Letting $`\sigma _1`$ be the largest such $`\sigma `$ in $`(\alpha ,\gamma )`$ we have
$$\frac{|𝐚(\sigma _1)|}{𝐅(\sigma _1)^\epsilon }=1\text{and}\frac{|𝐚(\sigma )|}{𝐅(\sigma )^\epsilon }<1\text{for }\sigma (\sigma _1,\gamma ].$$
Equivalently
(18)
$$|𝐚(\sigma _1)|𝐅(\sigma _1)^\epsilon =0\text{and}|𝐚(\sigma )|𝐅(\sigma )^\epsilon <0\text{for }\sigma (\sigma _1,\gamma ].$$
As $`|𝐚(\sigma )|=𝐚(\sigma )`$ on $`(\alpha ,\beta )`$, from (18) we have
$$\frac{d}{d\sigma }\left(𝐚(\sigma )𝐅(\sigma )^\epsilon \right)|_{\sigma =\sigma _1}0.$$
Hence
$`0`$ $``$ $`𝐚^{}(\sigma _1)+\epsilon 𝐅^{}(\sigma _1)𝐅(\sigma _1)^{\epsilon 1}`$
$`=`$ $`𝐛(\sigma _1)+\epsilon 𝐚(\sigma _1)𝐅(\sigma _1)^\epsilon `$
$`=`$ $`𝐛(\sigma _1)\epsilon 𝐚(\sigma _1)^2.`$
By (17) $`𝐛(\sigma _1)<\epsilon 𝐚(\sigma _1)^2`$. This is a contradiction, proving Claim 2.
From Claim 2 we immediately have $`𝐚(\sigma )=\mathrm{O}\left(𝐅(\sigma )^{\epsilon /2}\right)`$, and thus by Corollary 11(b) $`𝐚(\sigma )=\mathrm{o}\left(𝐅(\sigma )^\epsilon \right)`$. Then from Corollary 8(b)
$$𝐛(\sigma )=\mathrm{o}\left(𝐚(\sigma )^2\right)=\mathrm{o}\left(𝐅(\sigma )^\epsilon \right)\text{ as }\sigma \alpha +.$$
This finishes the proof of (a).
Part (b) is now a trivial consequence of part (a) and Corollary 8(c).
###### Remark 13.
Corollary 12(b) readily shows many Dirichlet series satisfying (A1)–(A3) are not admissible.
* $`\zeta (s)^k`$, $`k=1,2,\mathrm{}`$ , is not admissible as it has a pole at its abscissa $`\alpha =1`$.
* The zeta function
$$\underset{j=1}{\overset{k}{}}\left(1n_j^s\right)^{m_j}$$
of a finitely generated multiplicative number system is not admissible as it has a pole at its abscissa $`\alpha =0`$.
###### Corollary 14.
The function $`\widehat{F}(x)`$ grows much faster than $`x^{\alpha +1}`$, namely
$$\underset{x\mathrm{}}{lim}\frac{\widehat{F}(x)}{x^{\alpha +1}}=\mathrm{}.$$
###### Proof.
This is clear from Corollary 10, Corollary 11(a) and (12). ∎
###### Definition 15.
For $`\alpha `$, a real-valued function $`g(x)`$ that is eventually defined on the reals and eventually positive is said to have regular variation at infinity with index $`\alpha `$, written simply as $`g(x)\mathrm{𝖱𝖵}_\alpha `$, if for any $`y>0`$
(19)
$$\underset{x\mathrm{}}{lim}\frac{g(xy)}{g(x)}=y^\alpha .$$
###### Corollary 16.
$`\widehat{F}(x)\mathrm{𝖱𝖵}_{\alpha +1}`$.
###### Proof.
We assume $`x,y>0`$. In the expressions for $`\widehat{F}(xy)`$ and $`\widehat{F}(x)`$ given by Theorem 7 let $`\sigma =\sigma _x`$ (for $`x`$ sufficiently large) and divide to obtain
$`{\displaystyle \frac{\widehat{F}(xy)}{\widehat{F}(x)}}`$ $`=`$ $`(y^{\sigma _x+1}){\displaystyle \frac{\mathrm{exp}\left((\mathrm{log}y)^2/\left(2𝐛(\sigma _x)\right)\right)+R(xy,\sigma _x)}{1+R(x,\sigma _x)}}`$
$``$ $`y^{\alpha +1}\text{as }x\mathrm{}`$
since both $`R(x,\sigma _x)0`$ and $`R(xy,\sigma _x)0`$ as $`x\mathrm{}`$ by Theorem 7 and (12); and since $`𝐛(\sigma _x)\mathrm{}`$ as $`x\mathrm{}`$ by (6) and (12). ∎
###### Lemma 17.
Let $`𝐆(s)=_{n1}g(n)/n^s`$ be a Dirichlet series with nonnegative real coefficients and abscissa $`\alpha 0`$, and let $`G(x)={\displaystyle \underset{1nx}{}}g(n)`$, $`\widehat{G}(x)={\displaystyle _1^x}G(u)𝑑u`$.
If $`\widehat{G}(x)\mathrm{𝖱𝖵}_{\alpha +1}`$ then
* $`G(x)\mathrm{𝖱𝖵}_\alpha `$, and
* $`G(x){\displaystyle \frac{\alpha +1}{x}}\widehat{G}(x).`$
###### Proof.
This is an immediate consequence of Lemma 11.21 from . ∎
###### Corollary 18.
$`F(x)\mathrm{𝖱𝖵}_\alpha `$ and
$$F(x)\frac{\alpha +1}{x}\widehat{F}(x)\frac{x^{\sigma _x}𝐅(\sigma _x)}{\sigma _x\sqrt{2\pi 𝐛(\sigma _x)}}$$
as $`x\mathrm{}`$.
###### Proof.
By Corollary 10, Corollary 16, Lemma 17 and (12). ∎
###### Corollary 19.
The function $`F(x)`$ grows much faster than $`x^\alpha `$, namely
$$\underset{x\mathrm{}}{lim}\frac{F(x)}{x^\alpha }=\mathrm{}.$$
###### Proof.
By Corollary 11(a), Corollary 18 and (12). ∎
From this Corollary it is immediate that $`\zeta (s)`$ is not admissible (a fact already noted in Remark 13).
## 4. Tenenbaum’s Condtions
A version of admissibility conditions for Dirichlet series due to Tenenbaum , 1988, is given in the following.<sup>2</sup><sup>2</sup>2Tenenbaum actually uses $`T(\sigma )=\sigma 𝐛(\sigma )/\epsilon `$, which makes our (T4) unnecessary, and he uses the saddlepoint $`\sigma _x`$ instead of $`\sigma `$ in (T1)–(T3), which he labels as (H2)–(H4).
###### Definition 20.
Suppose $`𝐅(s)`$ satisfies conditions (A1)–(A3) and there is a function $`T:(\alpha ,\beta )(0,\mathrm{})`$, for some $`\beta >\alpha `$, such that as $`\sigma \alpha +`$
1. $`\sigma ^2𝐛(\sigma )\mathrm{}`$
2. $`{\displaystyle \frac{𝐛(\sigma )^3}{𝐜(\sigma )^2}}\mathrm{}`$
3. $`T(\sigma ){\displaystyle \frac{\left|𝐅(\sigma +it)\right|}{𝐅(\sigma )}}1`$, for $`\sigma (\alpha ,\beta )`$ and for $`\delta (\sigma )|t|T(\sigma )`$,
where $`\delta (\sigma )=\left|𝐛(\sigma )𝐜(\sigma )\right|^{1/5}`$
4. $`{\displaystyle \frac{\sqrt{𝐛(\sigma )}}{T(\sigma )}}0`$
5. $`\left|𝐜(\sigma +it)\right|\left|𝐜(\sigma )\right|`$ for $`\sigma (\alpha ,\beta )`$ and $`t`$
6. $`\underset{\sigma \alpha +}{lim\; inf}\left|𝐜(\sigma )\right|>0`$.
Then we say that $`𝐅(s)`$ is T-admissible, as witnessed by $`T(\sigma )`$.
Tenenbaum uses $`T(\sigma )=\sigma 𝐛(\sigma )/\epsilon (\sigma )`$ where $`\epsilon (\sigma )0`$ as $`\sigma \alpha +`$. This choice of $`T(\sigma )`$ makes condition (T4) unnecessary. Furthermore he gives an error term that is important to his applications in number theory, especially to the function $`\psi (x,y)`$.
###### Theorem 21.
If $`𝐅(s)`$ is T-admissible then it is admissible.
###### Proof.
Let $`𝐅(s)=\mathrm{exp}\left(𝐇(s)\right)`$ be a T-admissible Dirichlet series as witnessed by $`T(\sigma ):(\alpha ,\beta )(0,\mathrm{})`$. (T1) shows that
(21)
$$𝐛(\sigma )\mathrm{}\text{ as }\sigma \alpha +,$$
so we can assume that $`𝐛(\sigma )`$ is positive on $`(\alpha ,\beta )`$. From (21) and (T4) it is clear that
$$T(\sigma )\mathrm{}\text{ as }\sigma \alpha +.$$
By (21) and (T6) one has
$$\delta (\sigma )0\text{ as }\sigma \alpha +,$$
so (A4) holds. As (T1) is (A5) we only need to verify that (A6)–(A8) hold.
For (A7) we have for $`\sigma (\alpha ,\beta )`$ and $`|t|\delta (\sigma )`$
$`\left|𝖱(\sigma +it)\right|`$ $``$ $`\left|𝐜(\sigma )t^3\right|\text{by (}\text{3}\text{)}`$
$``$ $`\left|𝐜(\sigma )\delta (\sigma )^3\right|`$
$`=`$ $`|𝐜(\sigma )|\left|𝐛(\sigma )𝐜(\sigma )\right|^{3/5}`$
$`=`$ $`\left({\displaystyle \frac{𝐜(\sigma )^2}{𝐛(\sigma )^3}}\right)^{1/5}`$
$``$ $`0\text{as }\sigma \alpha +\text{by (T2)}.`$
For (A6) we have from (T3) for $`\sigma (\alpha ,\beta )`$
$$T(\sigma )\frac{\left|𝐅\left(\sigma +i\delta (\sigma )\right)\right|}{𝐅(\sigma )}1.$$
Multiplying this by (T4) gives
$$\sqrt{𝐛(\sigma )}\frac{\left|𝐅\left(\sigma +i\delta (\sigma )\right)\right|}{𝐅(\sigma )}0,$$
which, in view of (5) and the fact that (A7) holds, gives (A6).
Finally (A8) is verified as follows, where $`\sigma (\alpha ,\beta )`$:
$`\sigma \sqrt{𝐛(\sigma )}{\displaystyle _{|t|\delta (\sigma )}}{\displaystyle \frac{\left|𝐅(\sigma +it)\right|}{𝐅(\sigma )}}{\displaystyle \frac{dt}{\sigma ^2+t^2}}`$
$``$ $`\sigma \sqrt{𝐛(\sigma )}{\displaystyle _{\delta (\sigma )|t|T(\sigma )}}{\displaystyle \frac{\left|𝐅(\sigma +it)\right|}{𝐅(\sigma )}}{\displaystyle \frac{dt}{\sigma ^2+t^2}}+2\sigma \sqrt{𝐛(\sigma )}{\displaystyle _{T(\sigma )}^{\mathrm{}}}{\displaystyle \frac{dt}{t^2}}`$
$``$ $`{\displaystyle \frac{\sqrt{𝐛(\sigma )}}{T(\sigma )}}\left(\sigma {\displaystyle _{\delta (\sigma )|t|T(\sigma )}}{\displaystyle \frac{dt}{\sigma ^2+t^2}}\right)+{\displaystyle \frac{2\sigma \sqrt{𝐛(\sigma )}}{T(\sigma )}}\text{by (T3)}`$
$`=`$ $`\mathrm{o}(1)\text{by (T4)}.`$
The conditions of Tenenbaum have proved to be very practical, giving the asymptotics for many naturally occurring examples of Dirichlet series to which the saddlepoint method applies.
###### Example 22.
The function
$$𝐅(s):=e^{\zeta (s)}$$
is readily proved to be T-admissible, witnessed by $`T(\sigma )=𝐛(\sigma )`$, after noting
* $`\zeta (s)={\displaystyle \frac{1}{s1}}+g(s)`$, where $`g(s)`$ is holomorphic
* there is a constant $`C>0`$ such that for $`\sigma [1,2]`$ and $`|t|1`$ we have $`\left|\zeta (\sigma +it)\right|C\mathrm{log}|t|`$.
From the T-admissibility of $`\mathrm{exp}\left(\zeta (s)\right)`$ one easily has the T-admissibility of
$$𝐅_\lambda (s):=\mathrm{exp}\left(\zeta (s\lambda )\right)=\mathrm{exp}\left(\underset{n=1}{\overset{\mathrm{}}{}}n^\lambda n^s\right)$$
for $`\lambda 0`$.<sup>3</sup><sup>3</sup>3The asymptotics for $`F_\lambda (x)`$ are also analyzed in $`\mathrm{\S }`$11.5 of by the saddlepoint method, after changing the path of the Perron integral. (See Footnote 1 for errata to $`\mathrm{\S }`$11.5.) R. Warlimont first pointed out the example of $`\mathrm{exp}(\zeta (s))`$ to us. Later he found related examples of admissible functions, such as $`\mathrm{exp}\left(\zeta (s)^k\right)`$, that subsequently turned out to be T-admissible as well.
###### Example 23.
Tenenbaum studies the counting functions $`\psi (x,y)`$ for the zeta functions
$$\zeta (s,y):=\underset{py}{}\left(1p^s\right)^1.$$
As noted in Remark 13, the functions $`\zeta (s,y)`$ are not admissible. These functions satisfy all the conditions for being T-admissible except (T3), and for $`y`$ in a suitable range (depending on $`x`$) they satisfy (T3) provided $`\sigma =\sigma _x`$. This leads to asymptotics for $`\psi (x,y)`$ as $`x`$ and $`y`$ tend to infinity with $`y`$ suitably constrained.
###### Example 24.
The function
$$𝐅_k(s):=\mathrm{exp}\left(\frac{1}{1k^s}\right)$$
is admissible for $`k=2,\mathrm{}`$ , but not T-admissible. It is clear that each $`𝐅_k(s)`$ satisfies (A1)–(A3), and has a Dirichlet series expansion with abscissa of convergence $`\alpha =0`$. To see that $`𝐅_k(s)`$ is not T-admissible note that
$$\frac{\left|𝐅_k(\sigma +it)\right|}{𝐅_k(\sigma )}$$
is, for each $`\sigma >0`$, positive and periodic as a function of $`t`$, and thus does not uniformly go to 0 on $`[\delta (\sigma ),T(\sigma )]`$ as $`\sigma \alpha +`$. Consequently $`𝐅_k(s)`$ does not satisfy condition (T3).
To show that $`𝐅_k(s)`$ is admissible let
$`\delta _k(\sigma )`$ $`=`$ $`(k^\sigma 1)^{7/5}`$
$`T_k(\sigma )`$ $`=`$ $`(k^\sigma 1)^3.`$
One has
$`𝐚_k(\sigma )`$ $`=`$ $`{\displaystyle \frac{k^\sigma }{\left(k^\sigma 1\right)^2}}\mathrm{log}k`$
$`𝐛_k(\sigma )`$ $`=`$ $`{\displaystyle \frac{k^{2\sigma }+k^\sigma }{\left(k^\sigma 1\right)^3}}\left(\mathrm{log}k\right)^2`$
$`𝐜_k(\sigma )`$ $`=`$ $`{\displaystyle \frac{k^{3\sigma }+4k^{2\sigma }+k^\sigma }{\left(k^\sigma 1\right)^4}}\left(\mathrm{log}k\right)^3.`$
Verifying (A4)–(A6) is routine. For (A7) we proceed as in the proof of Theorem 21, namely for $`|t|\delta (\sigma )`$ one has
$$\left|𝖱(\sigma +it)\right|\left|𝐜(\sigma )\delta (\sigma )^3\right|0\sigma \alpha +.$$
This leaves (A8), which is usually the challenging part of the verification of admissibility. First note that
$`{\displaystyle \frac{\sigma \sqrt{𝐛_k(\sigma )}}{𝐅_k(\sigma )}}{\displaystyle _{|t|T_k(\sigma )}}\left|𝐅_k(\sigma +it)\right|{\displaystyle \frac{dt}{\sigma ^2+t^2}}`$ $``$ $`\sigma \sqrt{𝐛_k(\sigma )}{\displaystyle _{|t|T_k(\sigma )}}{\displaystyle \frac{dt}{t^2}}`$
$`=`$ $`{\displaystyle \frac{\sigma \sqrt{𝐛_k(\sigma )}}{T_k(\sigma )}}0\text{ as }\sigma \alpha +.`$
Thus we only need to show that
$$\frac{\sigma \sqrt{𝐛_k(\sigma )}}{𝐅_k(\sigma )}_{\delta (\sigma )|t|T_k(\sigma )}\left|𝐅_k(\sigma +it)\right|\frac{dt}{\sigma ^2+t^2}0\text{ as }\sigma \alpha +.$$
Substituting $`\tau =t\mathrm{log}k`$, we need to show
$$\frac{\sigma \sqrt{𝐛_k(\sigma )}}{𝐅_k(\sigma )}_{\delta (\sigma )\mathrm{log}k|\tau |T_k(\sigma )\mathrm{log}k}\left|𝐅_k\left(\sigma +i\frac{\tau }{\mathrm{log}k}\right)\right|\frac{(\mathrm{log}k)d\tau }{\left(\sigma \mathrm{log}k\right)^2+\tau ^2}0$$
as $`u1+`$. Letting $`u=k^\sigma `$ it suffices to show
$$_{(u1)^{7/5}\mathrm{log}k}^{(u1)^3\mathrm{log}k}\frac{\mathrm{log}u}{(u1)^{3/2}}\mathrm{exp}\left(\frac{2\left(\mathrm{cos}(\tau )1\right)}{(u1)\left(u^22u\mathrm{cos}(\tau )+1\right)}\right)\frac{d\tau }{(\mathrm{log}u)^2+\tau ^2}0$$
as $`u1+`$.
One can do this by noting that as $`u1+`$ the integrand rapidly and uniformly approaches 0 outside neighborhoods of radius $`(u1)^{7/5}`$ about the points $`\tau =2m\pi `$, indeed much faster than $`(u1)^3`$. Thus it suffices to show that
(22)
$$_U\frac{\mathrm{log}u}{(u1)^{3/2}}\mathrm{exp}\left(\frac{2\left(\mathrm{cos}(\tau )1\right)}{(u1)\left(u^22u\mathrm{cos}(\tau )+1\right)}\right)\frac{d\tau }{(\mathrm{log}u)^2+\tau ^2}0$$
as $`u1+`$, where $`U`$ is the union of the intervals
$$[2m\pi (u1)^{7/5},(2m\pi +(u1)^{7/5}]$$
about the points $`2m\pi `$, $`m1`$, such that $`2m\pi (u1)^{7/5}<(u1)^3`$. The integral in (22) is bounded by
(23)
$$2\zeta (2)_0^{(u1)^{7/5}}\frac{\mathrm{log}u}{(u1)^{3/2}}\mathrm{exp}\left(\frac{2\left(\mathrm{cos}(\tau )1\right)}{(u1)\left(u^22u\mathrm{cos}(\tau )+1\right)}\right)𝑑\tau .$$
Let $`J(u,\tau )`$ be the integrand in (23). Then
$`{\displaystyle _0^{(u1)^{7/5}}}J(u,\tau )𝑑\tau `$ $`=`$ $`{\displaystyle _0^{(u1)^{3/2}}}J(u,\tau )𝑑\tau +{\displaystyle _{(u1)^{3/2}}^{(u1)^{7/5}}}J(u,\tau )𝑑\tau `$
$``$ $`J(u,0)(u1)^{3/2}+J(u,(u1)^{3/2})(u1)^{7/5}`$
$``$ $`0\text{as }u1+.`$
This proves $`𝐅_k(s)`$ is admissible, and thus the class of admissible functions is wider than the class of T-admissible functions.
## 5. Closure under Product
The goal of this section is to prove that the product of two admissible functions $`𝐅_1(s)`$ and $`𝐅_2(s)`$ with the same abscissa of convergence is again admissible.
###### Theorem 25.
Suppose $`𝐅_1(s)`$ and $`𝐅_2(s)`$ are admissible with the same abscissa of convergence $`\alpha `$. Then $`𝐅_1(s)𝐅_2(s)`$ is admissible.
###### Proof.
We assume $`𝐅_j(s),\delta _j(\sigma ),𝐛_j(s)`$ satisfy (A1)–(A8) for $`j=1,2`$, and we assume $`\beta _j>\alpha `$ chosen such that $`𝐛_j(\sigma )>0`$ for $`\sigma (\alpha ,\beta _j)`$.
Let
$`𝐅(s)`$ $`:=`$ $`𝐅_1(s)𝐅_2(s)`$
$`\beta `$ $`:=`$ $`\mathrm{min}(\beta _1,\beta _2)`$
$`\delta (\sigma )`$ $`:=`$ $`\mathrm{min}(\delta _1(\sigma ),\delta _2(\sigma ))\text{for }\sigma (\alpha ,\beta ).`$
We have
$`𝐅_j(s)`$ $`=`$ $`e^{𝐇_j(s)}(j=1,2)`$
$`𝐅(s)`$ $`=`$ $`e^{𝐇(s)}`$
$`𝐇(s)`$ $`=`$ $`𝐇_1(s)+𝐇_2(s)`$
$`𝐚(s)`$ $`=`$ $`𝐚_1(s)+𝐚_2(s)`$
$`𝐛(s)`$ $`=`$ $`𝐛_1(s)+𝐛_2(s)`$
$`𝖱(s)`$ $`=`$ $`𝖱_1(s)+𝖱_2(s).`$
It is easy to check that (A1)–(A5) hold for $`𝐅(s)`$.
Next,
$`𝖱(\sigma +it)`$ $`=`$ $`𝖱_1(\sigma +it)+𝖱_2(\sigma +it)`$
$``$ $`0\text{uniformly for }|t|\delta (\sigma )`$
since each of the $`𝖱_j`$ satisfy (A7) and since $`\delta (\sigma )\delta _j(\sigma )`$ for $`j=1,2`$. So (A7) also holds for $`𝐅`$.
To prove (A6) and (A8) for $`𝐅`$ we first observe that for $`\sigma >\alpha `$, for $`t`$ and for $`j=1,2`$
(24)
$$\left|𝐅_j(\sigma +it)\right|𝐅_j(\sigma ),$$
and for $`\sigma (\alpha ,\beta )`$ and $`j=1,2`$
(25) $`𝐛_j(\sigma )`$ $`>`$ $`0`$
(26) $`𝐛(\sigma )`$ $`=`$ $`𝐛_1(\sigma )+𝐛_2(\sigma )2\mathrm{max}(𝐛_1(\sigma ),𝐛_2(\sigma )).`$
From (24) we have for $`\sigma >\alpha `$, for $`t`$ and for $`j=1,2`$
(27)
$$\frac{\left|𝐅(\sigma +it)\right|}{𝐅(\sigma )}=\frac{\left|𝐅_1(\sigma +it)\right|}{𝐅_1(\sigma )}\frac{\left|𝐅_2(\sigma +it)\right|}{𝐅_2(\sigma )}\frac{\left|𝐅_j(\sigma +it)\right|}{𝐅_j(\sigma )}.$$
Choose $`\gamma _1(\alpha ,\beta )`$ such that for $`\sigma (\alpha ,\gamma _1)`$ and $`j=1,2`$
(28)
$$𝐛_j(\sigma )\delta _j(\sigma )^2>1.$$
This is possible by (8).
Now suppose that we are given $`\epsilon (0,1)`$.
Choose $`\gamma _2(\alpha ,\gamma _1)`$ such that for $`\sigma (\alpha ,\gamma _2)`$ and $`j=1,2`$
(29) $`𝐛_j(\sigma )\mathrm{exp}\left(𝐛_j(\sigma )\delta _j(\sigma )^2\right)<\epsilon ^2`$
(30) $`\sigma \sqrt{𝐛_j(\sigma )}{\displaystyle _{|t|\delta _j(\sigma )}}{\displaystyle \frac{\left|𝐅_j(\sigma +it)\right|}{𝐅_j(\sigma )}}{\displaystyle \frac{dt}{\sigma ^2+t^2}}<\epsilon .`$
We can do this because the $`𝐅_j`$ satisfy (A6) and (A8).
Choose $`\gamma (\alpha ,\gamma _2)`$ such that for $`\sigma (\alpha ,\gamma )`$ and $`j=1,2`$
(31)
$$\frac{\left|𝐅_j(\sigma +it)\right|}{𝐅_j(\sigma )}<2\mathrm{exp}\left(𝐛_j(\sigma )t^2/2\right)\text{for }|t|\delta _j(\sigma ).$$
In view of (5) we can do this because the $`𝐅_j`$ satisfy (A7).
Claim: For $`\sigma (\alpha ,\gamma )`$
(32) $`𝐛(\sigma )\mathrm{exp}\left(𝐛(\sigma )\delta (\sigma )^2\right)<2\epsilon ^2`$
(33) $`\sigma \sqrt{𝐛(\sigma )}{\displaystyle _{|t|\delta (\sigma )}}{\displaystyle \frac{\left|𝐅(\sigma +it)\right|}{𝐅(\sigma )}}{\displaystyle \frac{dt}{\sigma ^2+t^2}}<12\epsilon .`$
This will prove that (A6) and (A8) hold for $`𝐅`$.
We start by fixing $`\sigma (\alpha ,\gamma )`$.
* Case (i): $`\delta _2(\sigma )\delta _1(\sigma )`$.
Then $`\delta (\sigma )=\delta _2(\sigma )`$.
* Subcase (ia): $`𝐛_1(\sigma )𝐛_2(\sigma )`$.
Then by (26) and (29)
$`𝐛(\sigma )\mathrm{exp}\left(𝐛(\sigma )\delta (\sigma )^2\right)`$ $`<`$ $`2𝐛_2(\sigma )\mathrm{exp}\left(𝐛_2(\sigma )\delta _2(\sigma )^2\right)`$
$`<`$ $`2\epsilon ^2.`$
Also by (26), (27) for $`j=2`$ and (30) for $`j=2`$ we have
$`\sigma \sqrt{𝐛(\sigma )}{\displaystyle _{|t|\delta (\sigma )}}{\displaystyle \frac{\left|𝐅(\sigma +it)\right|}{𝐅(\sigma )}}{\displaystyle \frac{dt}{\sigma ^2+t^2}}`$
$``$ $`\sigma \sqrt{2𝐛_2(\sigma )}{\displaystyle _{|t|\delta _2(\sigma _2)}}{\displaystyle \frac{\left|𝐅_2(\sigma +it)\right|}{𝐅_2(\sigma )}}{\displaystyle \frac{dt}{\sigma ^2+t^2}}`$
$`<`$ $`\sqrt{2}\epsilon <12\epsilon .`$
* Subcase (ib): $`𝐛_2(\sigma )𝐛_1(\sigma )`$.
By (29) for $`j=2`$
$$\sqrt{𝐛_2(\sigma )}\mathrm{exp}\left(𝐛_2(\sigma )\delta _2(\sigma )^2/2\right)<\epsilon ,$$
so
(34)
$$\sqrt{𝐛_1(\sigma )}\mathrm{exp}\left(𝐛_1(\sigma )\delta _2(\sigma )^2/2\right)<\epsilon $$
since
$$\sqrt{x}\mathrm{exp}\left(x\delta _2(\sigma )^2/2\right)$$
is decreasing for $`x>1/\delta _2(\sigma )^2`$, and since by Subcase (ib) and (28) for $`j=2`$
$$𝐛_1(\sigma )𝐛_2(\sigma )>\frac{1}{\delta _2(\sigma )^2}.$$
Then by (26) and (34)
$`𝐛(\sigma )\mathrm{exp}\left(𝐛(\sigma )\delta (\sigma )^2\right)`$ $`<`$ $`2𝐛_1(\sigma )\mathrm{exp}\left(𝐛_1(\sigma )\delta _2(\sigma )^2\right)`$
$`<`$ $`2\epsilon ^2.`$
Also from (34) we have
$$\sqrt{𝐛_1(\sigma )}\mathrm{exp}\left(𝐛_1(\sigma )t^2/2\right)<\epsilon \text{for }\delta _2(\sigma )|t|.$$
Combined with (31) for $`j=1`$ this gives
(35)
$$\sqrt{𝐛_1(\sigma )}\frac{\left|𝐅_1(\sigma +it)\right|}{𝐅_1(\sigma )}<2\epsilon \text{for }\delta _2(\sigma )|t|\delta _1(\sigma ).$$
By (26) and (27) for $`j=1`$
$`\sigma \sqrt{𝐛(\sigma )}{\displaystyle _{|t|\delta (\sigma )}}{\displaystyle \frac{\left|𝐅(\sigma +it)\right|}{𝐅(\sigma )}}{\displaystyle \frac{dt}{\sigma ^2+t^2}}`$
$``$ $`\sigma \sqrt{2𝐛_1(\sigma )}{\displaystyle _{|t|\delta _2(\sigma )}}{\displaystyle \frac{\left|𝐅_1(\sigma +it)\right|}{𝐅_1(\sigma )}}{\displaystyle \frac{dt}{\sigma ^2+t^2}}`$
$`=`$ $`\sqrt{2}\left(J_1(\sigma )+J_2(\sigma )\right),`$
where
$`J_1(\sigma )`$ $`=`$ $`\sigma \sqrt{𝐛_1(\sigma )}{\displaystyle _{\delta _2(\sigma )|t|\delta _1(\sigma )}}{\displaystyle \frac{\left|𝐅_1(\sigma +it)\right|}{𝐅_1(\sigma )}}{\displaystyle \frac{dt}{\sigma ^2+t^2}}`$
$`J_2(\sigma )`$ $`=`$ $`\sigma \sqrt{𝐛_1(\sigma )}{\displaystyle _{|t|\delta _1(\sigma )}}{\displaystyle \frac{\left|𝐅_1(\sigma +it)\right|}{𝐅_1(\sigma )}}{\displaystyle \frac{dt}{\sigma ^2+t^2}}.`$
By (35)
$$J_1(\sigma )2\epsilon \sigma _{\delta _2(\sigma )|t|\delta _1(\sigma )}\frac{dt}{\sigma ^2+t^2}<2\pi \epsilon ,$$
and by (30)
$$J_2(\sigma )<\epsilon .$$
Thus
$$\sigma \sqrt{𝐛(\sigma )}_{|t|\delta (\sigma )}\frac{\left|𝐅(\sigma +it)\right|}{𝐅(\sigma )}\frac{dt}{\sigma ^2+t^2}<\sqrt{2}\left(2\pi +1\right)\epsilon <12\epsilon ,$$
and the claim is proved in Case (i).
Case (ii), where $`\delta _1(\sigma )\delta _2(\sigma )`$, is handled likewise. So (A6) and (A8) hold for $`𝐅`$, and the theorem is proved.
## 6. Open questions
###### Problem 1.
Is the sum of two admissible functions also admissible?
###### Problem 2.
Is the product of any two admissible functions also admissible?
###### Problem 3.
Given two admissible functions $`𝐅_j(x)=\mathrm{exp}\left(𝐇_j(s)\right)`$ is the function $`\mathrm{exp}\left(𝐇_1(s)𝐇_2(s)\right)`$ admissible?
###### Problem 4.
If $`𝐅(s)`$ is admissible, does it follow that $`e^{𝐅(s)}`$ is also admissible?
We suspect, by analogy with Hayman’s work, that this is true.
###### Problem 5.
Can the notion of admissible be extended to include $`\zeta (s)`$? |
warning/0507/astro-ph0507090.html | ar5iv | text | # Numerical Relativistic Hydrodynamics Based on the Total Variation Diminishing Scheme
## 1 Introduction
Many high-energy astrophysical problems involve relativistic flows, and thus understanding relativistic flows is important for correctly interpreting astrophysical phenomena. For instance, intrinsic beam velocities larger than $`0.9c`$ are typically required to explain the apparent superluminal motions observed in relativistic jets in microquasars in the Galaxy (Mirabel & Rodríguez, 1999) as well as in extragalactic radio sources associated with active galactic nuclei (Zensus, 1997). In some powerful extragalactic radio sources, ejections from galactic nuclei produce true beam velocities of more than $`0.98c`$. Relativistic descriptions are also inevitable in other situations of rapid expansion such as the early stages of supernova explosions (Burrows, 2000) and the production of energetic gamma-ray bursts (Mészáros, 2002). In the general fireball model of gamma-ray bursts, the internal energy of gas is converted into the bulk kinetic energy during expansion and this expansion leads to relativistic outflows with high bulk Lorentz factors $`100`$. Since such relativistic flows are highly nonlinear and intrinsically complex, in addition to possessing large Lorentz factors, often studying them numerically is the only possible approach.
For numerical study of non-relativistic hydrodynamics, explicit finite difference upwind schemes have been developed and implemented successfully. The schemes which have been used for astrophysical researches include the Roe scheme (Roe, 1981), the total variation diminishing (TVD) scheme (Harten, 1983), the piecewise parabolic method (PPM) scheme (Colella & Woodward, 1984), and the essentially non-oscillatory (ENO) scheme (Harten et al., 1987). These schemes are based on exact or approximate Riemann solvers using the characteristic decomposition of the hyperbolic system of hydrodynamic conservation equations. They all are able to capture sharp discontinuities robustly in the complex flows, and to describe the physical solution accurately.
Although the upwind schemes were originally developed for non-relativistic hydrodynamics, some have been extended to special relativistic hydrodynamics. For instance, Dolezal & Wong (1995) adapted the ENO scheme to one-dimensional relativistic hydrodynamics. They fulfilled the ENO scheme using the local characteristic approach which depends on the local linearizion of the system of conservation equations. Martí & Müller (1996) adapted the PPM scheme to one-dimensional relativistic hydrodynamics using an exact relativistic Riemann solver to calculate numerical fluxes at cell interfaces. Donat et al. (1998) and Aloy et al. (1999) constructed multidimensional relativistic hydrodynamic codes based on the ENO scheme and the PPM scheme, respectively. Reviews of various numerical approaches and test problems can be found in Martí & Müller (2003) and Wilson & Mathews (2003). These works showed that the advantage of the upwind schemes, high accuracy and robustness, are carried over to relativistic hydrodynamics.
In this paper we describe a multidimensional code for special relativistic hydrodynamics based on the total variation diminishing (TVD) scheme (Harten, 1983). The TVD scheme is an explicit Eulerian finite difference upwind scheme and an extension of the Roe scheme to second-order accuracy in space and time. The advantage of the TVD scheme is that a code based on it is simple and fast, and yet performs well. A non-relativistic hydrodynamic code based the TVD scheme was built and applied to astrophysical problems such as the large scale structure formation of the universe by one of authors (Ryu et al., 1993). The special relativistic hydrodynamic code in this paper was built by extending the non-relativistic code. All the components of the the non-relativistic code was kept, so the relativistic code has the structure parallel to the non-relativistic counterpart. It makes the relativistic code comprehensible and easily usable. Through tests, we demonstrate that the newly developed code for special relativistic hydrodynamics can handle interesting astrophysical problems involving large Lorentz factors or ultrarelativistic regimes where energy densities greatly exceed rest mass densities.
This paper is organized as follows. In Section 2 we describe the step by step procedures for building the code including the basic equations, characteristic decomposition, TVD scheme, multidimensional extension, and Lorentz transformation. Tests are presented in Section 3. A summary and discussion follows in Section 4.
## 2 Numerical Relativistic Hydrodynamics
### 2.1 Basic Equations
The ideal relativistic hydrodynamic equations can be written as a hyperbolic system of conservation equations
$$\frac{D}{t}+\frac{}{x_j}\left(Dv_j\right)=0,$$
(1)
$$\frac{M_i}{t}+\frac{}{x_j}\left(M_iv_j+p\delta _{ij}\right)=0,$$
(2)
$$\frac{E}{t}+\frac{}{x_j}\left[\left(E+p\right)v_j\right]=0,$$
(3)
where the equation of state is given by
$$p=\left(\gamma 1\right)\left(e\rho \right).$$
(4)
Here, $`D`$, $`M_i`$, and $`E`$ are the mass density, momentum density, and total energy density in the reference frame, and $`\rho `$, $`v_j`$, and $`e`$ are the mass density, velocity, and internal plus mass energy density in the local rest frame, respectively. In general, the adiabatic index $`\gamma `$ is taken as $`5/3`$ for mildly relativistic cases and as $`4/3`$ for ultrarelativistic cases where $`e\rho `$. In equations (1)–(3), the indices $`i`$ and $`j`$ run over $`x`$, $`y`$, and $`z`$ and the conventional Einstein summation is used. The speed of light is set to unity ($`c1`$) throughout this paper.
The quantities in the reference frame are related to those in the local rest frame via Lorentz transformation
$$D=\mathrm{\Gamma }\rho ,$$
(5)
$$M_i=\mathrm{\Gamma }^2\left(e+p\right)v_i,$$
(6)
$$E=\mathrm{\Gamma }^2\left(e+p\right)p,$$
(7)
where the Lorentz factor is given by
$$\mathrm{\Gamma }=\frac{1}{\sqrt{1v^2}}$$
(8)
with $`v^2=v_x^2+v_y^2+v_z^2`$.
In the non-relativistic limit, the quantities $`D`$, $`M_i`$, and $`E`$ approach their non-relativistic counterparts $`\rho ^N`$, $`\rho ^Nv_i^N`$, and $`E^N+\rho ^Nc^2`$ and equations (1)–(3) reduce to the non-relativistic hydrodynamic equations
$$\frac{\rho ^N}{t}+\frac{}{x_j}\left(\rho ^Nv_j^N\right)=0,$$
(9)
$$\frac{\rho ^Nv_i^N}{t}+\frac{}{x_j}\left(\rho ^Nv_i^Nv_j^N+p^N\delta _{ij}\right)=0,$$
(10)
$$\frac{E^N}{t}+\frac{}{x_j}\left[\left(E^N+p^N\right)v_j^N\right]=0,$$
(11)
where the pressure is given by
$$p^N=\left(\gamma 1\right)\left(E^N\frac{1}{2}\rho ^Nv_{}^{N}{}_{}{}^{2}\right).$$
(12)
### 2.2 Characteristic Decomposition
Equations (1)–(3) can be written as
$$\frac{\stackrel{}{q}}{t}+\frac{\stackrel{}{F}_j}{x_j}=0$$
(13)
with the state and flux vectors
$$\stackrel{}{q}=\left[\begin{array}{c}D\\ M_i\\ E\end{array}\right],\stackrel{}{F}_j=\left[\begin{array}{c}Dv_j\\ M_iv_j+p\delta _{ij}\\ \left(E+p\right)v_j\end{array}\right],$$
(14)
or as
$$\frac{\stackrel{}{q}}{t}+A_j\frac{\stackrel{}{q}}{x_j}=0,A_j=\frac{\stackrel{}{F}_j}{\stackrel{}{q}}.$$
(15)
Here, $`A_j`$ is the $`5\times 5`$ Jacobian matrix composed with the state and flux vectors. The construction of the matrix $`A_j`$ can be simplified by introducing a parameter vector, $`\stackrel{}{u}`$, as
$$A_j=\frac{\stackrel{}{F}_j}{\stackrel{}{u}}\frac{\stackrel{}{u}}{\stackrel{}{q}}.$$
(16)
We choose the parameter vector which consists of the physical quantities in the local rest frame,
$$\stackrel{}{u}=\left[\begin{array}{c}\rho \\ v_i\\ e\end{array}\right].$$
(17)
In building an upwind code to solve a hyperbolic system of conservation equations, the eigen-structure (eigenvalues and eigenvectors) of the Jacobian matrix is required. Eigen-structures for relativistic hydrodynamics in multidimensions were previously described, for instance, in Donat et al. (1998). However, the state vector in this paper is different from that of Donat et al. (1998), so the eigen-structure is different. In the following, our eigen-structure of equation (16) is presented. We first define the specific enthalpy, $`h`$, and the the sound speed, $`c_s`$, respectively as
$$h=\frac{e+p}{\rho },c_s^2=\frac{\gamma p}{\rho h}.$$
(18)
Then the eigenvalues of $`A_x`$ for $`j=x`$ are
$$a_1=\frac{\left(1c_s^2\right)v_x\sqrt{\left(1v^2\right)c_s^2\left[1v^2c_s^2\left(1c_s^2\right)v_x^2\right]}}{1v^2c_s^2},$$
(19)
$$a_2=v_x,$$
(20)
$$a_3=v_x,$$
(21)
$$a_4=v_x,$$
(22)
$$a_5=\frac{\left(1c_s^2\right)v_x+\sqrt{\left(1v^2\right)c_s^2\left[1v^2c_s^2\left(1c_s^2\right)v_x^2\right]}}{1v^2c_s^2}.$$
(23)
The eigenvalues $`a_{15}`$ represent the five characteristic speeds associated with two sound wave modes ($`a_{1,5}`$) and three entropy modes ($`a_{24}`$).
The complete set of the corresponding right eigenvectors ($`A_x\stackrel{}{R}=a\stackrel{}{R}`$) is
$$\stackrel{}{R}_1=[\frac{1v_xa_1}{\mathrm{\Gamma }h\left(1v_x^2\right)},a_1,\frac{\left(1v_xa_1\right)v_y}{1v_x^2},\frac{\left(1v_xa_1\right)v_z}{1v_x^2},1]^\mathrm{T},$$
(24)
$$\stackrel{}{R}_2=[\frac{\mathrm{\Gamma }\left(2h1\right)v_y}{h},0,1,0,0]^\mathrm{T},$$
(25)
$$\stackrel{}{R}_3=[\frac{\mathrm{\Gamma }^2\left(2h1\right)\left(v^2v_x^2\right)+h}{\mathrm{\Gamma }h},v_x,0,0,1]^\mathrm{T},$$
(26)
$$\stackrel{}{R}_4=[\frac{\mathrm{\Gamma }\left(2h1\right)v_z}{h},0,0,1,0]^\mathrm{T},$$
(27)
$$\stackrel{}{R}_5=[\frac{1v_xa_5}{\mathrm{\Gamma }h\left(1v_x^2\right)},a_5,\frac{\left(1v_xa_5\right)v_y}{1v_x^2},\frac{\left(1v_xa_5\right)v_z}{1v_x^2},1]^\mathrm{T}.$$
(28)
The complete set of the left eigenvectors ($`\stackrel{}{L}A_x=a\stackrel{}{L}`$), which are orthonormal to the right eigenvectors, is
$$\stackrel{}{L}_1=[\frac{\mathrm{\Gamma }h\left(v_xa_5\right)}{\left(h1\right)\left(a_1a_5\right)},\mathrm{\Delta }_{12},\frac{\mathrm{\Gamma }^2\left(2h1\right)\left(v_xa_5\right)v_y}{\left(h1\right)\left(a_1a_5\right)},\frac{\mathrm{\Gamma }^2\left(2h1\right)\left(v_xa_5\right)v_z}{\left(h1\right)\left(a_1a_5\right)},\mathrm{\Delta }_{15}],$$
(29)
$$\stackrel{}{L}_2=[\frac{\mathrm{\Gamma }hv_y}{h1},\frac{\left[\mathrm{\Gamma }^2\left(2h1\right)\left(v^2v_x^2\right)+h\right]v_xv_y}{\left(h1\right)\left(1v_x^2\right)},\frac{\mathrm{\Gamma }^2\left(2h1\right)v_y^2}{h1}+1,\frac{\mathrm{\Gamma }^2\left(2h1\right)v_yv_z}{h1},$$
$$\frac{\left[\mathrm{\Gamma }^2\left(2h1\right)\left(v^2v_x^2\right)+h\right]v_y}{\left(h1\right)\left(1v_x^2\right)}],$$
(30)
$$\stackrel{}{L}_3=[\frac{\mathrm{\Gamma }h}{h1},\frac{\left[\mathrm{\Gamma }^2\left(2h1\right)\left(v^2v_x^2\right)+1\right]v_x}{\left(h1\right)\left(1v_x^2\right)},\frac{\mathrm{\Gamma }^2\left(2h1\right)v_y}{h1},\frac{\mathrm{\Gamma }^2\left(2h1\right)v_z}{h1},$$
$$\frac{\mathrm{\Gamma }^2\left(2h1\right)\left(v^2v_x^2\right)1}{\left(h1\right)\left(1v_x^2\right)}],$$
(31)
$$\stackrel{}{L}_4=[\frac{\mathrm{\Gamma }hv_z}{h1},\frac{\left[\mathrm{\Gamma }^2\left(2h1\right)\left(v^2v_x^2\right)+h\right]v_xv_z}{\left(h1\right)\left(1v_x^2\right)},\frac{\mathrm{\Gamma }^2\left(2h1\right)v_yv_z}{h1},\frac{\mathrm{\Gamma }^2\left(2h1\right)v_z^2}{h1}+1,$$
$$\frac{\left[\mathrm{\Gamma }^2\left(2h1\right)\left(v^2v_x^2\right)+h\right]v_z}{\left(h1\right)\left(1v_x^2\right)}],$$
(32)
$$\stackrel{}{L}_5=[\frac{\mathrm{\Gamma }h\left(v_xa_1\right)}{\left(h1\right)\left(a_5a_1\right)},\mathrm{\Delta }_{52},\frac{\mathrm{\Gamma }^2\left(2h1\right)\left(v_xa_1\right)v_y}{\left(h1\right)\left(a_5a_1\right)},\frac{\mathrm{\Gamma }^2\left(2h1\right)\left(v_xa_1\right)v_z}{\left(h1\right)\left(a_5a_1\right)},\mathrm{\Delta }_{55}],$$
(33)
where the auxiliary variables are defined as
$$\mathrm{\Delta }_{12}=\frac{\left[\mathrm{\Gamma }^2\left(2h1\right)\left(v^2v_x^2\right)+1\right]\left(v_xa_5\right)v_x}{\left(h1\right)\left(1v_x^2\right)\left(a_1a_5\right)}+\frac{1}{a_1a_5},$$
(34)
$$\mathrm{\Delta }_{15}=\frac{\left[\mathrm{\Gamma }^2\left(2h1\right)\left(v^2v_x^2\right)+1\right]\left(v_xa_5\right)}{\left(h1\right)\left(1v_x^2\right)\left(a_1a_5\right)}\frac{a_5}{a_1a_5},$$
(35)
$$\mathrm{\Delta }_{52}=\frac{\left[\mathrm{\Gamma }^2\left(2h1\right)\left(v^2v_x^2\right)+1\right]\left(v_xa_1\right)v_x}{\left(h1\right)\left(1v_x^2\right)\left(a_5a_1\right)}+\frac{1}{a_5a_1},$$
(36)
$$\mathrm{\Delta }_{55}=\frac{\left[\mathrm{\Gamma }^2\left(2h1\right)\left(v^2v_x^2\right)+1\right]\left(v_xa_1\right)}{\left(h1\right)\left(1v_x^2\right)\left(a_5a_1\right)}\frac{a_1}{a_5a_1}.$$
(37)
The eigenvalues and eigenvectors of $`A_y`$ and $`A_z`$ can be obtained by properly redefining indices. We note that the eigenvalues are same regardless of the choice of state or parameter vectors. But the right and left eigenvectors are different or can be presented in different forms.
### 2.3 One-Dimensional Functioning Code Based on the TVD Scheme
The TVD scheme we employ to build one-dimensional functioning code is practically identical to that in Harten (1983) and Ryu et al. (1993). But for completeness, the procedure is shown here. The state vector $`\stackrel{}{q}_i^n`$ at the cell center $`i`$ at the time step $`n`$ is updated by calculating the modified flux vector $`\overline{\stackrel{}{f}}_{x,i\pm 1/2}`$ along the $`x`$-direction at the cell interface $`i\pm 1/2`$ as follows:
$$L_x\stackrel{}{q}_i^n=\stackrel{}{q}_i^n\frac{\mathrm{\Delta }t^n}{\mathrm{\Delta }x}\left(\overline{\stackrel{}{f}}_{x,i+1/2}\overline{\stackrel{}{f}}_{x,i1/2}\right),$$
(38)
$$\overline{\stackrel{}{f}}_{x,i+1/2}=\frac{1}{2}\left[\stackrel{}{F}_x(\stackrel{}{q}_i^n)+\stackrel{}{F}_x(\stackrel{}{q}_{i+1}^n)\right]\frac{\mathrm{\Delta }x}{2\mathrm{\Delta }t^n}\underset{k=1}{\overset{5}{}}\beta _{k,i+1/2}\stackrel{}{R}_{k,i+1/2}^n,$$
(39)
$$\beta _{k,i+1/2}=Q_k(\frac{\mathrm{\Delta }t^n}{\mathrm{\Delta }x}a_{k,i+1/2}^n+\gamma _{k,i+1/2})\alpha _{k,i+1/2}\left(g_{k,i}+g_{k,i+1}\right),$$
(40)
$$\gamma _{k,i+1/2}=\{\begin{array}{ccc}\left(g_{k,i+1}g_{k,i}\right)/\alpha _{k,i+1/2}\hfill & \mathrm{for}& \alpha _{k,i+1/2}0,\hfill \\ 0\hfill & \mathrm{for}& \alpha _{k,i+1/2}=0,\hfill \end{array}$$
(41)
$$g_{k,i}=\mathrm{sign}(\stackrel{~}{g}_{k,i+1/2})\mathrm{max}\{0,\mathrm{min}[|\stackrel{~}{g}_{k,i+1/2}|,\mathrm{sign}(\stackrel{~}{g}_{k,i+1/2})\stackrel{~}{g}_{k,i1/2}]\},$$
(42)
$$\stackrel{~}{g}_{k,i+1/2}=\frac{1}{2}\left[Q_k(\frac{\mathrm{\Delta }t^n}{\mathrm{\Delta }x}a_{k,i+1/2}^n)\left(\frac{\mathrm{\Delta }t^n}{\mathrm{\Delta }x}a_{k,i+1/2}^n\right)^2\right]\alpha _{k,i+1/2},$$
(43)
$$\alpha _{k,i+1/2}=\stackrel{}{L}_{k,i+1/2}^n\left(\stackrel{}{q}_{i+1}^n\stackrel{}{q}_i^n\right),$$
(44)
$$Q_k(x)=\{\begin{array}{ccc}x^2/4\epsilon _k+\epsilon _k\hfill & \mathrm{for}& |x|<2\epsilon _k,\hfill \\ |x|\hfill & \mathrm{for}& |x|2\epsilon _k.\hfill \end{array}$$
(45)
Here, $`k=1`$ to 5 stand for the five characteristic modes. The internal parameters $`\epsilon _k`$’s are associated with numerical viscosity, and defined for $`0\epsilon _k0.5`$; $`\epsilon _{1,5}=0.10.3`$ for the sound wave modes and $`\epsilon _{24}=00.1`$ for the entropy modes are reasonable choices.
We note that the flux limiter in equation (42) is the min-mod limiter. The min-mod limiter is known to be very stable but has the cost of additional diffusion. To reproduce sharper structures with less diffusion, other flux limiters, such as the monotonized central difference limiter (MC limiter)
$$g_{k,i}=\mathrm{sign}(\stackrel{~}{g}_{k,i+1/2})\mathrm{max}\{0,\mathrm{min}[\frac{1}{2}(|\stackrel{~}{g}_{k,i+1/2}|+\mathrm{sign}(\stackrel{~}{g}_{k,i+1/2})\stackrel{~}{g}_{k,i1/2}),2|\stackrel{~}{g}_{k,i+1/2}|,$$
$$2\mathrm{s}\mathrm{i}\mathrm{g}\mathrm{n}(\stackrel{~}{g}_{k,i+1/2})\stackrel{~}{g}_{k,i1/2}]\},$$
(46)
or the superbee limiter
$$g_{k,i}=\mathrm{sign}(\stackrel{~}{g}_{k,i+1/2})\mathrm{max}\{0,\mathrm{min}[|\stackrel{~}{g}_{k,i+1/2}|,2\mathrm{s}\mathrm{i}\mathrm{g}\mathrm{n}(\stackrel{~}{g}_{k,i+1/2})\stackrel{~}{g}_{k,i1/2}],\mathrm{min}[2|\stackrel{~}{g}_{k,i+1/2}|,$$
$$\mathrm{sign}(\stackrel{~}{g}_{k,i+1/2})\stackrel{~}{g}_{k,i1/2}]\},$$
(47)
may be used; however, these limiters are more susceptible to oscillations at discontinuities. In the tests described in §3, the min-mod limiter was used.
In order to define the physical quantities at the cell interfaces, the TVD scheme originally used the Roe’s linearizion technique (Harten, 1983). Although it is possible to implement this linearizion technique in the relativistic domain in a computationally feasible way (see Eulderink & Mellema, 1995), there is unlikely to be a significant advantage from the computational point of view. Instead, we simply calculate the algebraic averages of quantities at two adjacent cell centers to define the physical quantities at the cell interfaces;
$$v_{x,i+1/2}=\frac{v_{x,i}+v_{x,i+1}}{2},v_{y,i+1/2}=\frac{v_{y,i}+v_{y,i+1}}{2},v_{z,i+1/2}=\frac{v_{z,i}+v_{z,i+1}}{2},$$
(48)
$$h_{i+1/2}=\frac{h_i+h_{i+1}}{2},$$
(49)
$$c_{s,i+1/2}=\left[\frac{\left(\gamma 1\right)\left(h_{i+1/2}1\right)}{h_{i+1/2}}\right]^{1/2}.$$
(50)
### 2.4 Multidimensional Extension
To extend the one-dimensional code to multidimensions, the procedure described in the previous subsection is applied separately to the $`y`$ and $`z`$-directions. Multiple spatial dimensions are treated through the Strang-type dimensional splitting (Strang, 1968). Then, the state vector is updated by
$$\stackrel{}{q}^{n+1}=L_zL_yL_x\stackrel{}{q}^n.$$
(51)
In order to maintain second-order accuracy in time, the order of the dimensional splitting is permuted as follows
$$L_zL_yL_x,L_xL_yL_z,L_xL_zL_y,L_yL_zL_x,L_yL_xL_z,L_zL_xL_y.$$
(52)
The time step $`\mathrm{\Delta }t^n`$ is restricted by the usual Courant stability condition
$$\mathrm{\Delta }t^n=\mathrm{min}[\frac{C_{\mathrm{Cour}}\mathrm{\Delta }x}{\mathrm{max}(a_{k,i+1/2}^n)_x},\frac{C_{\mathrm{Cour}}\mathrm{\Delta }y}{\mathrm{max}(a_{k,i+1/2}^n)_y},\frac{C_{\mathrm{Cour}}\mathrm{\Delta }z}{\mathrm{max}(a_{k,i+1/2}^n)_z}].$$
(53)
The Courant constant should be $`C_{\mathrm{Cour}}<1`$. We typically use $`C_{\mathrm{Cour}}0.9`$. The time step is calculated at the beginning of a permutation sequence and used through the complete sequence.
### 2.5 Lorentz Transformation
In the code, the conserved quantities $`D`$, $`M_i`$, and $`E`$ in the reference frame are evolved in time, but the physical quantities $`\rho `$, $`v_j`$, and $`e`$ in the local rest frame are needed for fluxes to be estimated. The quantities $`\rho `$, $`v_j`$, and $`e`$ can be obtained through Lorentz transformation of equations (5)–(7) at each time step. Schneider et al. (1993) showed that the transformation is reduced to a single quartic equation for $`v`$
$$f(v)=\left[\gamma v\left(EMv\right)M\left(1v^2\right)\right]^2\left(1v^2\right)v^2\left(\gamma 1\right)^2D^2=0,$$
(54)
where $`M^2=M_x^2+M_y^2+M_z^2`$. They also showed that the physically meaningful solution for $`v`$ is between the lower limit, $`v_1`$, and the upper limit, $`v_2`$,
$$v_1=\frac{\gamma E\sqrt{\left(\gamma E\right)^24\left(\gamma 1\right)M^2}}{2\left(\gamma 1\right)M},v_2=\frac{M}{E},$$
(55)
and that the solution is unique. Once $`v`$ is known, the quantities $`\rho `$, $`v_j`$, and $`e`$ can be straightforwardly calculated from the following relations
$$\rho =\frac{D}{\mathrm{\Gamma }},$$
(56)
$$v_x=\frac{M_x}{M}v,v_y=\frac{M_y}{M}v,v_z=\frac{M_z}{M}v,$$
(57)
$$e=EM_xv_xM_yv_yM_zv_z.$$
(58)
Equation (54) could be solved using a numerical procedure such as the Newton-Raphson root-finding method, as suggested in Schneider et al. (1993). A problem with the numerical approach is, however, that iterations can fail to converge. For instance, convergence can fail if one of the relativistic conditions is violated due to numerical errors, e.g., $`M>E`$, in a cell. This occurs mostly in extreme regimes. In addition, we found that convergence is often slow or sometimes fails in the limit $`ME`$. On the other hand, quartic equations have analytic solutions. The general form of roots can be found in standard books such as Abramowitz & Stegun (1972) or on webs such as “http://mathworld.wolfram.com/QuarticEquation.html”. Although it is too complicated to prove analytically, we found numerically that for the physical meaningful values of $`v`$ and $`c_s`$, $`v<1`$ and $`c_s<\sqrt{\gamma 1}`$, among the four roots of equation (54), two are complex and the other two are real. While the smaller real root is smaller than the lower limit $`v_1`$, the larger real root is between the two limits $`v_1`$ and $`v_2`$. So the larger real root is the one we are looking for, and we use its analytic formula in our code. The advantages of the analytic approach are obvious. It always gives a solution we are looking for, and it is easier to predict and deal with unphysical situations if one of the relativistic conditions is violated due to numerical errors.
## 3 Numerical Tests
### 3.1 Relativistic Shock Tube
We have performed two sets of relativistic shock tube tests in the one, two, and three-dimensional computational boxes with $`x=[0,1]`$, $`y=[0,1]`$, and $`z=[0,1]`$. Initially two different physical states are set up perpendicular to the direction along which waves propagate; along the $`x`$-axis in the one-dimensional calculation, along the diagonal line connecting $`(0,0)`$ and $`(1,1)`$ in the two-dimensional calculation, and along the diagonal line connecting $`(0,0,0)`$ and $`(1,1,1)`$ in the three-dimensional calculation. The initial states of the first test are
$$(\rho ,v_x,v_y,v_z,p)=\{\begin{array}{cc}(10,0,0,0,13.3)\hfill & 0x,\left(x+y\right)/2,\left(x+y+z\right)/31/2,\hfill \\ (1,0,0,0,10^6)\hfill & 1/2<x,\left(x+y\right)/2,\left(x+y+z\right)/31.\hfill \end{array}$$
(59)
The initial states of the second test are
$$(\rho ,v_x,v_y,v_z,p)=\{\begin{array}{cc}(1,0,0,0,10^3)\hfill & 0x,\left(x+y\right)/2,\left(x+y+z\right)/31/2,\hfill \\ (1,0,0,0,10^2)\hfill & 1/2<x,\left(x+y\right)/2,\left(x+y+z\right)/31.\hfill \end{array}$$
(60)
In equations (59) and (60), the expressions within inequalities are for one, two, and three dimensions, respectively. The first test involves a mildly relativistic flow and the second test involves a highly relativistic flow. In both tests, we assume the adiabatic index $`\gamma =5/3`$ and the outflow condition is used for the $`x`$, $`y`$, and $`z`$-boundaries. Both tests were previously considered by several authors (e.g., Martí & Müller, 1996). The estimation of accuracy was done by comparing the numerical solutions with the exact solutions described in Thompson (1986) and Martí & Müller (1994). In Figures 1(a) and (b), our numerical solutions are shown as open circles and the exact solutions are represented by solid lines.
Figure 1(a) shows the mildly relativistic shock tube test done using $`256`$, $`256^2`$, and $`256^3`$ cells with a Courant constant $`C_{\mathrm{Cour}}=0.9`$ and the parameters $`\epsilon _{1,5}=0.1`$ and $`\epsilon _{24}=0`$. The plots of one, two, and three-dimensions correspond to times $`t=0.4`$, $`0.4\sqrt{2}`$, and $`0.4\sqrt{3}`$, respectively. Structures such as the shock front, contact discontinuity and rarefaction wave are accurately produced. There are actually slight improvements in accuracy in the multidimensional calculations. Figure 1(b) shows the highly relativistic shock tube test done again using $`256`$, $`256^2`$, and $`256^3`$ cells with a Courant constant $`C_{\mathrm{Cour}}=0.6`$ and the parameters $`\epsilon _{1,5}=0.1`$ and $`\epsilon _{24}=0`$. The plots of one, two, and three-dimensions correspond to times $`t=0.4`$, $`0.4\sqrt{2}`$, and $`0.4\sqrt{3}`$, respectively. The flow is more extreme, but the structure is correctly reproduced without spurious oscillations. But in the rest mass density profile the peak does not reach the value of the exact solution due to the coarseness of computational cells. According to our tests, in a one-dimensional calculation, the peak can be accurately reproduced when $`2048`$ numerical cells are used. There are also improvements in accuracy in the multidimensional calculations.
For a more quantitative comparison, we have calculated the norm errors of the rest mass density, velocity, and pressure for different dimensions. The errors shown in Table 1 are calculated at the same times as in Figure 1. The errors are gradually reduced as the dimensionality increases and demonstrate a good agreement between the numerical and exact solutions. Note that the values of $`E(p)`$ exceeding unity are still acceptable because these are from the initial large value of pressure.
### 3.2 Relativistic Wall Shock
A one-dimensional relativistic wall shock test has been performed in the computational box of $`x=[0,1]`$. Initially a gas with extreme velocity occupying all numerical cells propagates along the $`x`$-axis against a reflecting wall placed at $`x=1`$. As the gas hits the wall, it is compressed and heated and eventually a reverse shock is generated. The initial condition of this test is
$$(\rho ,v_x,v_y,v_z,p)=(1,0.999999,0,0,10^4)0x1.$$
(61)
The adiabatic index $`\gamma =5/3`$ is assumed and the inflow boundary condition is used at $`x=0`$. It is another test which was widely used by several authors (e.g., Donat et al., 1998).
The relativistic jump condition for strong shocks with negligible preshock pressure is given by Blandford & McKee (1976)
$$v_s=\frac{\left(\gamma 1\right)\mathrm{\Gamma }v}{\mathrm{\Gamma }+1},$$
(62)
$$\rho ^{}=\rho \frac{\gamma \mathrm{\Gamma }+1}{\gamma 1},$$
(63)
$$v^{}=0,$$
(64)
$$p^{}=\rho \left(\mathrm{\Gamma }1\right)\left(\gamma \mathrm{\Gamma }+1\right).$$
(65)
Here, $`v_s`$ is the shock velocity and the superscript $``$ represents the postshock quantities, while the quantities without any superscript refer to the preshock gas.
Figure 2 shows the structure at $`t=0.75`$ when the reverse shock is located at $`x=0.5`$. The calculation has been done using $`512`$ computational cells with a Courant constant $`C_{\mathrm{Cour}}=0.9`$ and the parameters $`\epsilon _{1,5}=0.3`$ and $`\epsilon _{24}=0.1`$. The numerical solution is drawn with open circles and the exact solution is represented by solid lines. The numerical and exact solutions match exactly without any oscillation or overshoot in the rest mass density, velocity, and pressure profiles.
With different inflow velocities, we have calculated the mean errors in the rest mass density, velocity, and pressure. The errors are calculated for the same time as in Figure 2 and given in Table 2. Note that the order of the mean errors is $`10^3`$, and that the accuracy does not depend systematically on the investigated Lorentz factor. The mean error in the rest mass density is $`0.5\%`$ for all the Lorentz factors and about $`0.25\%`$ for the maximum Lorentz factor. This accuracy is comparable to or better than that of other published upwind scheme codes.
### 3.3 Relativistic Blast Wave
The propagation of a relativistic blast wave has been tested in the two-dimensional computational box with $`x=[0,1]`$ and $`y=[0,1]`$. A gas of high density and pressure is initially confined in a spherical region and the subsequent explosion is allowed to evolve. This makes a spherical blast wave propagate outward. The initial condition of this test is
$$(\rho ,v_x,v_y,v_z,p)=\{\begin{array}{cc}(10,0,0,0,10^3)\hfill & 0\sqrt{x^2+y^2}1/2,\hfill \\ (1,0,0,0,1)\hfill & \mathrm{outside}.\hfill \end{array}$$
(66)
The adiabatic index is taken to be $`\gamma =4/3`$ and the reflecting and outflow boundary conditions are used.
The calculation has been done using $`512^2`$ cells with a Courant constant $`C_{\mathrm{Cour}}=0.6`$ and the parameters $`\epsilon _{1,5}=0.1`$ and $`\epsilon _{24}=0`$. To test the symmetry properties of the code, the calculation has been stopped before a reverse shock reaches the inner reflecting boundary. Figure 3 shows the profiles of the rest mass density, velocity, and pressure measured along the diagonal line connecting $`(0,0)`$ and $`(1,1)`$ at $`t=0.7`$. The spherical blast wave successfully propagates to a larger radius, and we have found that all structures in it preserve the initial symmetry.
### 3.4 Relativistic Hawley-Zabusky Shock
In order to test the applicability of the code to complex relativistic flows, we have performed a two-dimensional test simulation of the relativistic version of the Hawley-Zabusky shock. The test was originally suggested by Hawley & Zabusky (1989) for non-relativistic hydrodynamics. Almost the same physical values as in the original paper are used here. Initially a plane-parallel shock with a Mach number $`1.2`$ propagates along the $`x`$-axis into two regions of different density. The regions are separated by oblique discontinuity whose inclination is $`30^{}`$ with respect to the $`x`$-axis. The density jumps three times across the discontinuity. The initial configuration is summarized as
$$(\rho ,v_x,v_y,v_z,p)=\{\begin{array}{cc}(1,0.6,0,0,0.48)\hfill & 0x1/16,0y1,\hfill \\ (1,0,0,0,0.48)\hfill & 1/16<x\sqrt{3}y+1/4,0y1,\hfill \\ (3,0,0,0,0.48)\hfill & \mathrm{outside}.\hfill \end{array}$$
(67)
The adiabatic index $`\gamma =1.4`$ is used. The inflow and outflow conditions are used at the $`x`$-boundaries and the reflecting condition is used at the $`y`$-boundaries.
The simulation has been done in the two-dimensional computational box with $`x=[0,8]`$ and $`y=[0,1]`$ using a uniform numerical grid of $`2048\times 256`$ cells. A Courant constant $`C_{\mathrm{Cour}}=0.9`$ and the parameters $`\epsilon _{1,5}=0.1`$ and $`\epsilon _{24}=0`$ were used. We have simulated this test until $`t=20`$ in order to see the long term evolution. The passage of the planar shock through the discontinuity causes the Kelvin-Helmholtz instability to occur along the discontinuity and end up with formation of vortices. The vortices roll up, interact, and merge during the simulation; the detailed morphology and the number of vortices formed are somewhat sensitive to numerical resolution. Figure 4 shows the gray-scale images of the rest mass density at different times ($`t=2`$, $`11`$, and $`20`$). Because all the structures are dragged to the right boundary as time goes on, only the left, middle, and right half of the computational box are shown at $`t=2`$, $`11`$, and $`20`$, respectively. The vortices along the discontinuity are clearly formed and overall the morphology is similar to that of the non-relativistic simulation.
### 3.5 Relativistic Extragalactic Jets
Finally, in order to test the applicability of the code to realistic relativistic flows, we have simulated a two-dimensional relativistic extragalactic jet propagating into homogeneous medium. The relativistic jet inflows with a velocity $`0.99`$ to the computational box of $`x=[0,4]`$ and $`y=[0,1]`$. The jet has initially radius $`1/8`$ ($`32`$ cells) and Mach number $`8.76`$. The density ratio of the jet to the ambient medium is $`0.1`$ and the pressure of the jet is in equilibrium with that of the ambient medium. The initial condition for jet inflow and ambient medium is summarized as
$$(\rho ,v_x,v_y,v_z,p)=\{\begin{array}{cc}(1,0.99,0,0,0.1)\hfill & 0x1/32,0y1/8,\hfill \\ (10,0,0,0,0.1)\hfill & \mathrm{outside}.\hfill \end{array}$$
(68)
The adiabatic index $`\gamma =4/3`$ is used. The inflow and outflow conditions are used at the $`x`$-boundaries and the reflecting and outflow conditions are used at the $`y`$-boundaries.
The simulation has been done using a uniform numerical grid of $`1024\times 256`$ cells with a Courant constant $`C_{\mathrm{Cour}}=0.3`$ and the parameters $`\epsilon _{1,5}=0.3`$ and $`\epsilon _{24}=0.1`$. Figure 5 shows the gray-scale images of logarithm of the rest mass density, pressure, and Lorentz factor at $`t=5`$ when the bow shock reaches the right boundary. We can clearly see the dominant structures of bow shock, working surface, contact discontinuity, and cocoon. It is clear that the internal structure of the relativistic jet is less complex compared to that of a non-relativistic jet due to the effects of high Lorentz factor. The overall morphology and dynamics of our simulation match roughly with those of previous works, e.g., Duncan & Hughes (1994), although the initial conditions and the plotted epoch are different.
## 4 Summary and Discussion
A multidimensional code for special relativistic hydrodynamics was described. It differs from previous codes in the following aspects: 1) It is based on the total variation diminishing (TVD) scheme (Harten, 1983), which is an explicit Eulerian finite difference upwind scheme and an extension of the Roe scheme to second-order accuracy in space and time. 2) It employs a new set of conserved quantities, and so the paper describes a new eigen-structure for special relativistic hydrodynamics. 3) For the Lorentz transformation from the conserved quantities in the reference frame to the physical quantities in the local rest frame, an analytic formula is used.
To demonstrate the performance of the code, several tests were presented, including relativistic shock tubes, a relativistic wall shock, a relativistic blast wave, the relativistic version of the Hawley-Zabusky shock, and a relativistic extragalactic jet. The relativistic shock tube tests showed that the code clearly resolves mildly relativistic and highly relativistic shocks within $`24`$ numerical cells, although it requires more cells for resolving contact discontinuities. The relativistic wall shock test showed that the code correctly captures very strong shocks with very high Lorentz factors. The relativistic blast wave test showed that blast waves propagate through ambient medium while preserving the symmetry. The test simulations of the relativistic version of the Hawley-Zabusky shock and a relativistic extragalactic jet proved the robustness and flexibility of the code, and that the code can be applied to studies of practical astrophysical problems.
The strong points of the new code include the following: 1) Based on the TVD scheme, the code is simple and fast. The core routine of the TVD relativistic hydrodynamics is only about 300 lines long in the three-dimensional version. It runs only about $`1.52`$ times slower than the non-relativistic counterpart (per time step). Yet, tests have shown that the code is accurate and reliable enough to be suited for astrophysical applications. In addition, the use of an analytic formula for Lorentz transformation makes the code robust, so it ran for all the tests we have performed without failing to converge. 2) The code has been built in a way to be completely parallel to the non-relativistic counterpart. So it can be easily understood and used, once one is familiar with the non-relativistic code. In addition, the techniques developed for the non-relativistic code such as parallelization can be imported transparently.
Finally, the code is currently being applied for studies of relativistic jet interactions with inhomogeneous external media and turbulence of relativistic flows. The results will be reported in separate papers.
EC is grateful to Paul Wiita for his advice and comments on this work. EC was supported by the GSU College of Arts and Sciences, and by Research Program Enhancement funds to the Program in Extragalactic Astronomy. DR was supported by the KOSEF grant R01-2004-000-10005-0. |
warning/0507/astro-ph0507491.html | ar5iv | text | # Cloud Structure and Physical Conditions in Star-forming Regions from Optical Observations. II. Analysis
## 1 INTRODUCTION
The physical and chemical state in the interstellar medium (ISM) continuously changes through the combination of star formation and stellar death (Cox & Smith 1974; Jenkins & Meloy 1974; McKee 1990; Thornton et al. 1998; Marri & White 2004). The evolution of a galaxy is governed to a large extent by these processes. The final stages of star formation involve the disruption of interstellar material from which the stars formed. If enough material is available in interstellar clouds forming O and B stars, one generation of stars can lead to the formation of a second generation, and so forth, but eventually the cloud remnants can no longer sustain star formation. The propagation of star formation occurs in regions of enhanced gas density behind shocks (e.g. Elmegreen & Lada 1977; McKee & Tan 2003 ). Early type (O and B) stars are associated with stellar winds, expanding H II regions, and supernova explosions that lead to shocks in the surrounding medium. Since these stars tend to form in clusters (Clark et al. 2005), the effects are magnified as seen in the Magellanic Clouds. Clearly, knowledge of the physical conditions and chemical composition of the ISM in star-forming regions will help us understand the above processes and their effects on the evolution of the ISM and the Galaxy.
High-resolution optical observations of interstellar absorption (e.g., Welty & Hobbs 2001; Pan et al. 2004, hereafter Paper I) have revealed complex velocity structure on many Galactic lines of sight, especially on sight lines through star-forming clouds. Determination of the physical properties of individual interstellar clouds from observations of absorption lines requires high spectral resolution to distinguish the individual components contributing to the generally complex line profiles. Furtheremore, ultraviolet (UV) spectra provide considerable information on abundances and general physical conditions in the ISM (e.g., Spitzer & Jenkins 1975). In Paper I, we presented high-resolution optical spectra of interstellar CN, CH, CH<sup>+</sup>, Ca I, K I, and Ca II absorption along 29 lines of sight in three star-forming regions, $`\rho `$ Oph, Cep OB2, and Cep OB3. We obtained velocity component structure by simultaneously analyzing the spectra of the six species along a given sightline. Column densities and Doppler parameters of the individual velocity components were derived through profile fitting. To complement the results of the optical survey, presented in Paper I, we here present an analysis of far ultraviolet (FUV) spectra of 15 of the stars in that sample. The FUV data were either acquired in dedicated observing programs or extracted from the existing data archives of the Far Ultraviolet Spectroscopic Explorer (FUSE) and Hubble Space Telescope/Space Telescope Imaging Spectrograph (HST/STIS). The main contribution of the present paper is an analysis of the results from Paper I and those for H<sub>2</sub> and CO in terms of cloud structure for gas with densities of 10 to 1500 cm<sup>-3</sup>. We emphasize that our data for the diffuse molecular gas associated with the three star-forming regions probe the large-scale structure of the parent molecular cloud as it disperses. We investigate whether a connection exists between the levels of star formation and conditions in the surrounding diffuse gas.
The remainder of the paper has the following outline. The next section gives a brief background of the three star-forming regions, and $`\mathrm{\S }`$ 3 describes the FUV observations and their analysis. A comparison of Doppler parameters is the focus of $`\mathrm{\S }`$ 4. In $`\mathrm{\S }`$ 5, we examine correlations between column densities of different species. A chemical analysis is performed in $`\mathrm{\S }`$ 6 to derive gas densities for components with CH and CN. The analysis utilizes available results for C<sub>2</sub> from the literature. The amount of H<sub>2</sub> excitation and its relationship to the chemical results is described here also. This is followed by a discussion on the distribution of species within clouds, on large- and small-scale cloud structures, and a comparison of general properties for the molecular gas in the three star-forming regions. The final section presents a summary of our work.
## 2 Background
The divergence in the modes of star formation is perhaps best understood in terms of the physical properties of the molecular cloud and environmental effects such as nearby H II regions. $`\rho `$ Oph, Cep OB2, and Cep OB3 are three regions that provide fruitful laboratories for the investigation of cloud dispersal as a result of star formation. Our data on diffuse molecular gas, whose $`LSR`$ velocities are similar to those of the radio-emitting clouds, probe the later stages of dispersal. Extensive studies (e.g., Sargent 1979; de Geus, Bronfman, & Thaddeus 1990; Patel et al. 1998) on these regions outlined basic properties for the cloud complexes.
### 2.1 The $`\rho `$ Oph Region
The region around the star $`\rho `$ Oph, also called the $`\rho `$ Oph Molecular Cloud, is a well-known concentration of dark nebulae and molecular clouds with a size of about 6$`\times `$5 pc<sup>2</sup>, a total molecular mass of 80 $`M_{}`$ (de Geus et al. 1990), and at a distance of 140 pc. One of its prominent characteristics is the extremely active star formation in the cloud core, where stars are heavily clustered with an estimated star-forming efficiency $``$ 20% (cf. Green & Young 1992; Doppmann et al. 2003; Phelps & Barsony 2004). Radio observations (Loren 1989; de Geus et al. 1990) showed that CO molecules mainly exist within a velocity range of $`V_{LSR}`$ 1.5–5.0 km s<sup>-1</sup>, and that the cloud forms the boundary between the bulk of the Ophiuchus molecular gas and the Upper-Scorpius stellar groups.
The $`\rho `$ Oph Molecular Cloud is part of the Ophiuchus molecular clouds, a filamentary system of clouds with a total molecular mass of 10<sup>4</sup> $`M_{}`$ (de Geus 1992). The Ophiuchus molecular clouds, in turn, belong to a larger cloud complex that is associated with the Scorpio-Centaurus OB association. This OB association consists of three subgroups, Lower-Centaurus Crux (LCC), Upper-Centaurus Lupus (UCL), and Upper-Scorpius (US). Blaauw (1958) showed that the three subgroups are separated in position and age, and that the youngest subgroup lies in Upper-Scorpius, which is adjacent to the $`\rho `$ Oph Molecular Cloud. Based on photometric measurements, de Geus, de Zeeuw, & Lub (1989) derived ages of 11–12, 14–15, and 5–6 Myr for the three subgroups, respectively. Somewhat surprisingly, they found that the middle subgroup, UCL, is the oldest one. They argued that either the classic picture of sequential star formation (Blaauw 1964; Elmegreen & Lada 1977) is not valid for the Sco OB2 Association as a whole because the oldest subgroup is inbetween the two younger ones, or for some unknown reason, massive star formation was initiated near the middle of the original giant molecular cloud. However, Pan (2002) conjectured that LCC and UCL may actually have a similar age. More recently, Sartori, Lepine, & Dias (2003) found that the US subgroup is 4−8 Myr old and that UCL and LCC have the same age, 16-20 Myr.
### 2.2 The Cepheus Bubble
A region of enhanced infrared emission associated with Cep OB2 shows the location of a giant ring-shaped cloud system, the so called Cepheus Bubble. The distance of the Bubble is somewhat uncertain. Earlier measurements indicated distances from 700 to 900 pc (Garrison & Kormendy 1976; Kun, Balazs, & Toth 1987). Recently, de Zeeuw et al. (1999) determined the distances of nearby OB associations using proper motion and parallaxes measured by Hipparcos. They obtained a distance of 615 pc for Cep OB2. As de Zeeuw et al. noticed, their distances are systematically smaller than the previous photometric determinations. If a distance of 750 pc is adopted, the ring-shaped cloud system has a diameter of about 120 pc and the total molecular mass for the system is about 10<sup>5</sup> $`M_{}`$ (Patel et al. 1998)<sup>4</sup><sup>4</sup>4If the distance of 615 pc, as suggested by de Zeeuw et al. (1999), is adopted, the size and mass of the cloud system will be slightly smaller. Our discussion on large- and small-scale structure will not be affected by the choice of the distance.. A dozen star-forming regions have been identified in the rim of the Bubble. Each star-forming region has a size of a few pc, similar to that of the $`\rho `$ Oph region, and total molecular mass of about 10<sup>3</sup> $`M_{}`$. Some large star-forming regions consist of several star-forming globules. For instance, IC 1396 includes at least four star-forming globules with a total mass of about 2200 $`M_{}`$ and a size of 12 pc (Patel et al. 1995); Weikard et al. (1996) reached similar conclusions.
Based on measurements of CO emission, Patel et al. (1998) proposed that the Bubble was created by stellar winds, and likely supernova explosions, from now deceased stars in the cluster NGC 7160 as well as the evolved stars, VV Cep and $`\mu `$ Cep, in the region. Patel et al. (1998) suggested that these stars were the first generation formed. The first generation stars have an age of about 13-18 Myr. After about $``$ 7 Myr, the expanding shell reached a size of about 30 pc and became unstable. The instability led to the formation of a second generation of stars, current numbers of Cep OB2, which are about $``$ 7 Myr old. The IRAS point sources in the globules and molecular clouds associated with the Bubble are the third generation, whose formation was triggered by stellar winds from Cep OB2 stars. From H I 21 cm data, Abraham, Balazs & Kun (2000) obtained a similar picture for the formation of the Bubble. The latter model showed that a supernova explosion might have occurred as recently as about 2 Myr ago. The explosion expanded the pre-existing Bubble further.
### 2.3 Cloud System Associated with Cep OB3
The molecular clouds producing the Cep OB3 Association represent a filamentary system, showing clumpy and irregular structure (Sargent 1977, 1979). The cloud system is located at a distance of about 700–800 pc (Garrison 1970; Moreno-Corral et al. 1993), with a size of $``$ 60$`\times `$20 pc<sup>2</sup> and a total molecular mass of $``$ 10<sup>4</sup> $`M_{}`$ (Sargent 1979). A few identified star-forming regions are embedded in the cloud system. The star-forming regions in this cloud system have a size similar to those in the Cepheus Bubble (a few pc), but have smaller total molecular mass (100–500 $`M_{}`$). Our lines of sight toward Cep OB3 probe two star-forming regions in the cloud system, Cepheus B and Cepheus F, with total molecular mass of 100 and 300, respectively (Yu et al. 1996).
The Cep OB3 Association consists of two subgroups separated by $``$ 13 pc on the plane of the sky. Based on photometric measurements, Blaauw (1964) assigned ages of 8 and 4 Myr to the two subgroups. However, ages derived from such isochrone fitting must be treated with caution in young OB associations because masses and ages remain model dependent. Consequently, the shape of both the evolutionary tracks and of the isochrones change from one model to another, leading to discrepancies in the age estimates. Much lower ages for the younger subgroup can be found in the literature. For instance, Garmany (1973) and Assousa, Herbst & Turner (1977) obtained an age of $``$ 1 Myr for it. Sargent (1979) suggested that the formation of the younger group was triggered by the interaction between stellar winds of the older subgroup and the original molecular cloud.
## 3 FUV Observations and Analysis
In Paper I, we presented observational data for interstellar CN, CH, CH<sup>+</sup>, Ca I, K I, and Ca II absorption on 29 directions in $`\rho `$ Oph, Cep OB2, and Cep OB3. To complement these results, we obtained and analyzed FUV data for H<sub>2</sub> and CO. Eight stars from the sample in Paper I were observed with FUSE under our programs A051 and B030. For each star, reduction and calibration were performed using CALFUSE V2.4. We then re-binned the data by a factor of 4, yielding a two-pixel resolution element of $``$ 0.06 Å. FUSE archival data for two more stars, HD 206773 and HD 217312, were downloaded from programs B071 and P193, respectively. Four additional stars have been studied before; we used their published H<sub>2</sub> column density, $`N`$(H<sub>2</sub>), and $`T_{1,0}`$ results from Rachford et al. (2003) (under FUSE program P116). Although no published H<sub>2</sub> results are available for HD 208266, we included this star in our analysis because there were CH and CO observations. Moreover, these data can be utilized in tandem to predict $`N`$(H<sub>2</sub>) toward HD 208266. Table 1 lists all FUV datasets that we utilized, specifying which space telescope was the source, what S/N was obtained, and which molecules were analyzed.
Three $`BX`$ bands of H<sub>2</sub> (2$``$0, 3$``$0, and 4$``$0, between 1042 and 1083 Å) were chosen to model the molecule’s column density. This range is covered by the four overlapping FUSE detector segments, LiF-1A, LiF-2B, SiC-1A, and SiC-2B, although 2$``$0 is not covered by LiF-2B. Each segment was fitted independently with the code ISMOD, which uses the simplex method to minimize the rms of the residuals down to about 10<sup>-4</sup> in relative parameter steps. The fit was based on the CH cloud components in Paper I, preserving the velocity separations and relative fractions in column density among components as fixed input. We, however, allowed the $`b`$-values to vary during the fits, constraining them by the mass ratio of CH to H<sub>2</sub> and by the kinetic temperature ($`T`$) in the gas as given by the (fitted) $`T_{1,0}`$ rotational temperature of H<sub>2</sub> (see $`\mathrm{\S }`$ 6.2). Both the thermal and non-thermal components of the velocity field were thus determined. Free parameters included the total column density, $`N`$(H<sub>2</sub>), radial velocity of H<sub>2</sub>, six rotational temperatures relative to the ground state, $`T_J^{}`$$`_{{}_{}{}^{},0}`$, and the placement of the stellar continuum. Our final H<sub>2</sub> model and the associated uncertainties were computed by averaging the results from fits of the four independent segments. In a few cases of high-J levels, where one of the four segment yielded a deviant result, that value was not included in the final average. The values of total $`N`$(H<sub>2</sub>) are very robust and practically independent of the cloud structure because of the damping wings for the $`J^{}`$ = 0 and 1 lines.
For determinations of CO column density, we first searched the HST/STIS archive and downloaded high-resolution observations of $`AX`$ bands of CO. Six stars happen to have E140H data with R $``$ 100,000, covering $`AX`$ bands blueward of $``$ 1360 Å (transitions 7$``$0 and higher). For the ISMOD fits, we used $`f`$-values from Chan, Cooper & Brion (1993), see also Morton and Noreau (1994). These values are within a few percent of those recommended by Eidelsberg et al. (1999), except for the $`f`$-value for 11$``$0 where Chan et al.’s value is 11% smaller. The higher HST/STIS resolution allowed us to derive relative cloud strengths as well. When the CO along a line of sight is relatively weak, not enough fitting leverage is available from the $`AX`$ bands alone. In such cases we expanded the analysis to include simultaneous fits of Rydberg bands from the FUSE data with R $``$ 17,500. For CO toward HD 203374A, our simultaneous fits of $`AX`$ and Rydberg bands were published in Sheffer, Federman, & Andersson (2003). An additional three stars have medium-resolution data from the E140M grating (R = 46,000). Despite the lower resolution, robust fits are reliably obtainable in these cases because the spectral coverage includes 14 $`AX`$ bands down to the lowest S/N region around Ly-$`\alpha `$. Nevertheless, in our fits we employed at most six bands for simultaneous solutions, making sure to include bands spanning the largest possible range in $`f`$-value. Adding more bands to the fit did not change the results significantly. Finally, for the six stars without any archival HST/STIS data, our CO modeling was based solely on the three Rydberg bands, $`BX`$ 0$``$0, $`CX`$ 0$``$0, and $`EX`$ 0$``$0. Understandably, the inferred CO component structure for these stars are less reliable than for those from $`AX`$ bands, due to the lower resolution of FUSE and the limited range in $`f`$-values sampled, which were based on the results of Federman et al. (2001). However, since the $`BX`$ 0$``$0 band is relatively optically thin, the total $`N`$(CO) is approximately independent of the cloud structure. The derived total H<sub>2</sub> and CO column densities along the fifteen lines of sight are listed in Table 2, along with the kinetic temperature ($`T_{1,0}`$) of H<sub>2</sub>. For comparison, we also tabulate CH column densities from Paper I and abundance ratios (CH/H<sub>2</sub> and CO/H<sub>2</sub>) along these sight lines. Table 3 provides the column densities for individual H<sub>2</sub> rotational levels.
## 4 Doppler Parameter and CN Excitation Temperature
The Doppler parameters, $`b`$-values, of individual absorption components seen for any individual species provide upper limits on the temperature and internal turbulent velocity ($`v_t`$) in the interstellar gas because $`b=(2kT/m+v_t^2)^{1/2}`$. If two species with significantly different atomic weight $`m`$ coexist in the same volume of gas, comparison of their $`b`$-values provides estimates of the relative contributions of thermal and turbulent broadening, if both lines are fully resolved. Alternately, known temperatures and turbulent velocities for two species can allow estimates of their relative volumetric distributions. The extracted $`b`$-values in the present paper, except for CN $`b`$-values from the doublet ratio method in this section, are based on Gaussian instrumental widths.
Earlier studies suggested that CN resides in denser regions whereas CH<sup>+</sup> is found mainly in regions of lower density (Cardelli et al. 1990; Federman et al. 1994). The CH molecule, on the other hand, can be present in both low- and high-density gas (100 vs. 600 cm<sup>-3</sup>). It can be synthesized through the non-equilibrium CH<sup>+</sup> chemistry in low density diffuse clouds (Draine & Katz 1986; Zsargó & Federman 2003), and in moderately dense gas it is produced via C<sup>+</sup> \+ H$`{}_{2}{}^{}`$ CH$`{}_{2}{}^{+}+h\nu `$ (Federman et al. 1994). However, it is not a trivial task to disentangle the amount of CH formed with CH<sup>+</sup> or associated with CN. Examining CH<sup>+</sup>–like CH and CN–like CH components may elucidate the chemical schemes occurring in the gas (Lambert, Sheffer, & Crane 1990). Of our 125 CH components, there were 12 CN–like CH components, which are CH components with corresponding CN, but no CH<sup>+</sup> \[$`N`$(CH<sup>+</sup>) $`<3.0\times 10^{11}`$ cm<sup>-2</sup>\], at the same $`V_{LSR}`$, and 78 CH<sup>+</sup>–like CH components, CH components with detected CH<sup>+</sup> at the same V<sub>LSR</sub> but no CN \[$`N`$(CN) $`<2.0\times 10^{11}`$ cm<sup>-2</sup>\]. Typically, CH<sup>+</sup> and CN column densities are 7 times greater than these limits. Table 4 presents average $`b`$-values of the CH<sup>+</sup>-like CH and the CN-like CH components, along with mean Doppler parameters of all defined components for other species. The analysis reveals that CN-like CH components have an average $`b`$-value indistinguishable from that of CN components (0.83 $`\pm `$ 0.11 km s<sup>-1</sup> vs. 0.90 $`\pm `$ 0.11 km s<sup>-1</sup>), whereas CH<sup>+</sup>-like CH components possess a larger average $`b`$-value (1.10 $`\pm `$ 0.16 km s<sup>-1</sup>). Simple statistical tests show that the average $`b`$-values for a given species agree within their mutual uncertainties in the diffuse gas surronding different star-forming regions, although Cep OB3 seems to have consistently larger $`b`$-values than $`\rho `$ Oph and Cep OB2. Since thermal broadening plays a minor role, our results indicate that cloud turbulence in the three regions does not differ significantly.
Compared with earlier studies, our average $`b`$-values of atomic species (Table 4) are slightly greater than those obtained from very high-resolution spectra of interstellar atomic lines. Based on spectra with resolution of 0.3-0.6 km s<sup>-1</sup>, Welty and colleagues obtained median $`b`$-values of 0.66 km s<sup>-1</sup> for Ca I components (Welty, Hobbs, & Morton 2003), 0.67 km s<sup>-1</sup> for K I (Welty & Hobbs 2001), and 1.33 km s<sup>-1</sup> for Ca II components (Welty, Morton, & Hobbs 1996). Three causes may lead to the differences. Their spectra had much higher resolution than ours so that they could discern more closely blended components which we did not, different lines of sight were probed, and the median values are slightly smaller than corresponding average values. (For example, the median $`b`$ is 0.66 km s<sup>-1</sup> while the corresponding average is 0.70 km s<sup>-1</sup> for well defined Ca I components.) With a resolution of $``$ 0.50 km s<sup>-1</sup>, Andersson, Wannier, & Crawford (2002) obtained spectra of interstellar CH absorption toward 18 stars in southern molecular cloud envelopes. Based on F-tests (Lupton 1993), they fitted all but one of these CH profiles with a single-velocity component. (However, they noticed that additional velocity components may be present). Their $`b`$-values range from 0.3 to 3.6 km s<sup>-1</sup>, with an average of 1.7 km s<sup>-1</sup>. This mean $`b`$-value is greater than ours, 1.04 km s<sup>-1</sup> (Table 4). Because we simultaneously analyzed the spectra of six species along a given sightline (Paper I), in general, we found more CH velocity components per line of sight. This is one of the reasons, if not main one, that our average $`b`$ value is smaller than theirs. Based on ultra-high-resolution ($`\delta `$v $``$ 0.35 km s<sup>-1</sup>) observations along five lines of sight, Crawford (1995) found average $`b`$-values of 0.63, 1.5, 2.3, and 0.8 km s<sup>-1</sup> for CN, CH, CH<sup>+</sup>, and K I components. Overall, our average Doppler parameters are in agreement with those determined by previous studies using spectra with higher resolution.
The $`b`$-value of a CN component can also be derived by using the doublet ratio (Strömgren 1948) — see also Gredel, van Dishoeck, & Black (1991). Because CN $`R`$(1) and $`P`$(1) absorption lines arise from the same rotational level, $`N^{\prime \prime }`$=1, the same column density should be inferred from them. Therefore, the $`b`$-value of a CN component can be determined by performing a curve of growth analysis on the $`R`$(1) and $`P`$(1) values of $`W_\lambda `$ and requiring they give the same column density and $`b`$-value. To check $`b`$-values of CN components derived from the profile fitting in Paper I, we applied the doublet ratio method to eleven components where both CN $`R`$(1) and $`P`$(1) lines were detected and obtained an average $`b`$-value of 0.70 $`\pm `$ 0.14 km s<sup>-1</sup>, which is closer to that of Crawford (1995) based on ultra-high-resolution spectra. Considering the uncertainties, it also agrees with the mean Doppler parameter in Table 4, which we determined from profile fitting.
Once the $`b`$-value of a CN component is determined, column densities of rotational levels of $`N^{\prime \prime }`$=0 \[from CN $`R`$(0)\] and $`N^{\prime \prime }`$=1 \[from CN $`R`$(1) and $`P`$(1)\] can be derived by using a curve of growth. Knowing columns of these rotational levels allows us to obtain the CN excitation temperature, $`T_{ex}`$, of CN components. $`T_{ex}`$ is governed by the Boltzmann equation,
$$T_{ex}=\frac{h\nu }{k_B}\left(ln\left[\frac{g_1N(0)}{g_0N(1)}\right]\right)^1=\frac{5.442}{ln\left[3N(0)/N(1)\right]},$$
(1)
where $`h\nu `$ is the energy separation between rotational levels $`N^{\prime \prime }`$=0 and $`N^{\prime \prime }`$=1, while $`g`$ and $`N`$ are the statistical weight and column density for each level. The derived CN excitation temperatures for these eleven components range from 2.55 K to 2.89 K with an average of 2.75 $`\pm `$ 0.10 K. The derived excitation temperatures indicate that no significant excitation in addition to that due to the CMB is observed in these cloudlets, including three components toward HD 204827 where very strong CN absorption is detected and $`T_{ex}`$ is between 2.72 and 2.78 K.
The $`b`$-values in Table 4 for CO toward Cep OB2 and OB3 are most like those for CN, suggesting that these species coexist in the diffuse molecular gas probed by optical absorption. There is also a hint that the $`b`$-values for Cep OB3 are larger, as seen in other species, but in the case of CO, the difference could be the result of coarser resolution in $`FUSE`$ spectra for the stars in Cep OB3. Other evidence makes the connection between CO and CN stronger. First, while the input for the synthesis of CO bands was based on the component structure seen in CH, the best fit revealed a component structure very similar to that for CN. This is illustrated in Table 5, where the $`rms`$ deviations between CO and CH relative column density fractions on the one hand and between CO and CN on the other are presented for syntheses based on STIS spectra. Along each line of sight, the $`rms`$ deviations are smaller for CN than for CH. The correspondence is not perfect, however, because CO and CH are detected in cases where only upper limits for CN are available (e.g., Crenny & Federman 2004). Second, there is an approximate linear relationship between $`N`$(CO) and $`N`$(CN) for individual velocity components. Since upper limits exist for both quantities, correlations were sought through the use of the package ASURV, revision 1.1 (see Isobe, Feigelson, & Nelson 1986; Isobe & Feigelson 1990; La Valley, Isobe, & Feigelson 1992). Using Schmitt’s Method, we obtain log\[$`N`$(CO)\] $`=`$ ($`1.16\pm 0.13`$) log\[$`N`$(CN)\] $`+`$ ($`0.75\pm 1.59`$). This analysis suggests that all but 3 of 24 limits are consistent with detections. When all data are treated as detections, Schmitt’s Method gives a slope of $`1.35\pm 0.17`$ and an intercept of $``$($`1.56\pm 2.02`$) and a linear least-squares fit indicates respective values of $`1.63\pm 0.16`$ and $``$($`4.89\pm 1.85`$). The slopes are steeper when all data are considered detections because the upper limits on CN are less confining. Each of these observational effects show that columns of CO and CN are tightly coupled in the denser portions of diffuse molecular clouds, a conclusion consistent with our earlier findings (Zsargó & Federman 2003; Crenny & Federman 2004).
## 5 Correlations among Column Densities
Comparisons of column densities for different species provide constraints on chemical models, elemental depletions, species distributions, and the processes contributing to the ionization in diffuse clouds (e.g., Hobbs 1976; Chaffee & White 1982; Federman et al. 1994; Welty & Hobbs 2001). Based on high resolution and high S/N spectra, Paper I derived column densities for CN, CH, CH<sup>+</sup>, Ca I, K I, and Ca II for individual velocity components along 29 lines of sight. The data set has good internal self-consistency and is also large enough to allow us to discuss correlations based on individual components rather than total line of sight columns as usually done in previous studies (e.g., Danks, Federman, & Lambert 1984; Welty et al. 2003). Therefore, presumably, this data set will provide a good view of relationships among species. We examine some possible correlations between column densities of observed species. In each case, we first perform a least-squares fit to log\[$`N`$(Y)\] versus log\[$`N`$(X)\] to obtain the linear correlation coefficient, $`r`$, and the corresponding probability of no correlation, $`p`$, — i.e., the probability for a random distribution of N measurements to result in a correlation coefficient $`r^{}r`$, according to Bevington & Robinson (1992). Then, a regression analysis<sup>5</sup><sup>5</sup>5We used the subroutine regrwt.f, obtained from the Penn State statistical software archive at http://www.astro.psu.edu/statcodes, with slight modifications, to perform the regressions. is performed to get the slope, intercept, and an estimate of the scatter (the root mean square distance of points from the fit line). In all the fits, we made one 3$`\sigma `$ pass through the data to eliminate outliers. It turns out that three points are excluded from the fit for the relationship of log\[$`N`$(CN)\] versus log\[$`N`$(CH<sup>+</sup>-corrected CH)\] and one from fits of CN vs. CH, Ca I vs. K I, and Ca I vs. Ca II. Table 6 provides a summary of the correlations. The analysis shows that all correlations listed in the table have confidence levels, 1 $``$ $`p`$, exceeding 99.95%. In the following subsections, we describe some of these relationships.
### 5.1 CH versus CH<sup>+</sup>
When all CH velocity components are considered, the data set suggests a weak, at most, trend of increasing CH abundance with increasing CH<sup>+</sup> column. However, the correlation becomes much stronger when we include only CH<sup>+</sup>-like CH components, as shown in the bottom panel of Figure 1. A closer look reveals that the correction is weak for components with higher CH<sup>+</sup> columns. If only components with log\[$`N`$(CH<sup>+</sup>)\] $`>`$ 12.6 are considered, column densities of the two species are not well correlated. This is clearly indicated by the greater amount of dispersion seen at higher column densities for a linear fit to all the data. Therefore, the correlation is sought for components with log\[$`N`$(CH<sup>+</sup>)\] $``$ 12.6.
In previous studies, based on total column densities along lines of sight, CH column densities were sometimes found well correlated with $`N`$(CH<sup>+</sup>). For instance, Crane, Lambert, & Sheffer (1995) found that lines of sight with detectable amounts of both CH and CH<sup>+</sup> show a large range in CH to CH<sup>+</sup> column density ratios, whereas Federman, Welty, & Cardelli (1998) showed that, for lines of sight with no CN detections, CH was well correlated with CH<sup>+</sup>, and that the CH predicted by CH<sup>+</sup> synthesis agreed well with observations. Our findings explain why different correlations were seen in these studies. Because many of the sight lines in the study of Crane et al. have large amounts of CN, some CH on these sight lines is associated with CN instead of CH<sup>+</sup>. Therefore, $`N`$(CH) should not be well correlated with $`N`$(CH<sup>+</sup>). On the other hand, Federman et al. (1998) focused on lines of sight with detected CH and CH<sup>+</sup> but no corresponding CN. All CH in this study was CH<sup>+</sup>-like CH and a good correlation is expected.
### 5.2 CN Versus CH
The top panel of Figure 1 reveals that column densities of CN components are correlated with columns of their corresponding CH components (one component of HD 217312 was eliminated by a 3$`\sigma `$ pass). This component-based correlation is much tighter than the sight-line-based relationship (cf. Figure 7 of Federman et al. 1994). Moreover, the correlation is even stronger when only CN-like CH components are included, as shown in the bottom panel of Figure 2. This study strengthens our previous statement, based on $`b`$-values, that CN-like CH components are closely associated with the CN components.
We have discussed CN-like CH and CH<sup>+</sup>-like CH. How about CH components that have both corresponding CN and CH<sup>+</sup> components? Is a part of their CH associated with CN, and another part with CH<sup>+</sup>? To seek answers to these questions, we calculated the amount of CH corresponding to CH<sup>+</sup> for these components by applying the $`N`$(CH<sup>+</sup>-like CH) vs. $`N`$(CH) relation given in Table 6, and subtracted the result from the total CH column density, yeilding $`N`$(CH<sup>+</sup>-corrected CH). \[If log\[$`N`$(CH<sup>+</sup>)\] $`>`$ 12.6, $`N`$(CH) = 4.0 $`\times 10^{12}`$ cm<sup>-2</sup> is taken from the total CH column as CH<sup>+</sup>-like CH\]. The top panel of Figure 2 presents the correlation between $`N`$(CN) and $`N`$(CH<sup>+</sup>-corrected CH), where three components are excluded from the final fitting by the 3$`\sigma `$ pass. All three components have small amounts of CN, (0.4−0.7) $`\times 10^{12}`$ cm<sup>-2</sup>, but are still too much for their corresponding CH<sup>+</sup>-corrected CH in the relationship. It is possible that their amounts of CN are overestimated because columns of about 0.4 $`\times 10^{12}`$ cm <sup>-2</sup> represent marginal detections. As one can see from Figure 1, the $`N`$(CH<sup>+</sup>-like CH) vs. $`N`$(CH<sup>+</sup>) relationship shows quite a bit of scatter. This scatter may cause the three components to be outliers as well.
Use of CH<sup>+</sup>-corrected CH interestingly does improve the correlation compared to the original $`N`$(CN) vs. $`N`$(CH) plot. The correlation coefficient slightly increases from 0.80 to 0.86. As one can see from Figure 2 and Table 6, the regression is now very similar to that for $`N`$(CN) vs. $`N`$(CN-like CH). This analysis indicates that our assumption is a reasonable one. In other words, we conclude that there are two kinds of CH in diffuse molecular clouds: CN-like CH and CH<sup>+</sup>-like CH. Disentangling the amount of CH in each appears possible when utilizing the relation between columns of CH<sup>+</sup>-like CH and CH<sup>+</sup>. In the past, only ultra-high resolution observations toward $`\zeta `$ Oph (Lambert et al. 1990; Crawford 1997) distinguished between these two possibilities.
### 5.3 K I Versus CH
Welty & Hobbs (2001) found an essentially linear relationship between $`N`$(K I) and $`N`$(CH) by using total column densities along lines of sight. Based on column densities for individual velocity components, we obtain a very similar correlation, with a slope of 0.96 $`\pm `$ 0.04, as shown in the upper panel of Figure 3. Three facts suggest that these two species respond very similarly to changes in physical conditions in diffuse molecular clouds: 1) the strong linear correlation between $`N`$(K I) and $`N`$(CH), 2) the similar average $`b`$-values (Table 4) and 3) the similar column density relationship for components and for lines of sight. This also implies that they generally coexist, at least in the density range sampled by these data. However, we also note that there are exceptions. A few K I components with reasonable strengths do not have corresponding CH components. For instance, the velocity component at $`V_{LSR}`$ 4.3 km s<sup>-1</sup> toward HD 216532 has a K I column density of 3.2 $`\times `$ 10<sup>11</sup> cm<sup>-2</sup>, but no detected CH. However, there are relatively strong Ca I and Ca II lines with column densities of 5.2 $`\times `$ 10<sup>9</sup> cm<sup>-2</sup> and 10.0 $`\times `$ 10<sup>11</sup> cm<sup>-2</sup>, respectively. Considering the $`b`$-value of the K I component, 1.5 km s<sup>-1</sup>, is greater than the average value of 0.92 km s<sup>-1</sup> (as is the case for the Ca I and Ca II component), we think this K I component may arise from a slightly lower density region where no significant amount of CH is present, but one with a relatively large volume so that a reasonably strong K I line is observed. In other words, this may indicate that K I could be distributed in less dense regions where no observable amount of CH resides.
### 5.4 Ca I Versus K I
The lower panel of Figure 3 presents the relationship between $`N`$(Ca I) and $`N`$(K I). The figure shows that the column densities of the two species are well correlated. However, the best fit slope of 0.60 $`\pm `$ 0.04, indicated by the solid line in the plot, is significantly smaller than 1.0. Compared with the relationship of Welty et al. (2003), which was based on total column densities, our best fit slope is the same as theirs (0.60 $`\pm `$ 0.08) while the scatter and the intercept (2.88 $`\pm `$ 0.47 versus 3.11 $`\pm `$ 0.68) are slightly smaller. Welty et al. (2003) discussed possibilities that could lead to a slope less than unity for the comparison between log\[$`N`$(Ca I)\] and log\[$`N`$(K I)\]. They concluded that it is most likely that calcium depletion has a steeper dependence on local density than does potassium depletion. This is not all that unexpected because calcium is a refractory element while potassium is more volatile, having lower condensation temperature. Once again, three facts, the good correspondence in column density, similar average $`b`$-values (see Table 4), and the similar correlations between relationships based on individual components and on lines of sight, suggest that these two species generally share comparable volumes.
In addition to above the correlations, Table 6 also lists relationships for log\[$`N`$(Ca I)\] vs. log\[$`N`$(Ca II)\] and log\[$`N`$(K I)\] vs. log\[$`N`$(Ca II)\]. Column densities of these species are correlated although not as tightly as the relationships above, with about twice the scatter. The column density ratios can differ by a factor of 10 for individual components. Such a large range in column density ratios may indicate that these species do not track each other well. In other words, they seem to be probing different portions of a cloud. There are actually many Ca II components without corresponding Ca I or K I. This latter fact arises because Ca II is believed to be distributed more broadly (Welty et al. 2003).
## 6 Chemical Analysis
### 6.1 CN Chemistry
The route to CN in diffuse molecular clouds involves CH to C<sub>2</sub> to CN (Federman et al. 1984, 1994). We use the steady-state analytical expressions in Federman et al. (1994), with updated rate coefficients (Knauth et al. 2001; Pan, Federman, & Welty 2001), to extract estimates for gas density in the material containing CN by matching the observed column densities for CN and C<sub>2</sub> (when available). Steady state is appropriate for this scheme because the photochemical time scales are less than 100 yrs., much shorter than the sound-crossing time. $`Observed`$ CH column densities are adopted for use in estimating the CN and C<sub>2</sub> column densities. In our earlier papers (i.e., Federman et al. 1984, 1994) the analysis is based on line-of-sight column densities, but here we obtain results for individual velocity components where CN is detected. The CH and CN column densities mainly come from the results presented in Paper I. The exception is CH toward $`\lambda `$ Cep; we adopt the results from the ultra-high resolution measurements of Crane et al. (1995). (In passing, we note that the 2 CH components have a separation consistent with what we found for CN. Of the 2 CH<sup>+</sup> components, only the bluer one is associated with gas containing CN; the redder component is seen in K I absorption but not in absorption from other molecules $``$ see Paper I.) Along several sight lines, results for C<sub>2</sub> are also available (Danks & Lambert 1983; Federman & Lambert 1988; Federman et al. 1994). When more than one CN component is present for a given direction, the relative amounts of C<sub>2</sub> in each component is taken to be the same as CN because these molecules likely coexist (Federman et al. 1994).
To extract estimates of gas densities from the chemical model (Federman et al. 1994, Pan et al. 2001), in addition to column densities of CN and CH (C<sub>2</sub> column provides a further constraint), O, C<sup>+</sup>, and N abundances, temperature, and UV radiation field are needed. Our sample contains some of the most molecular-rich diffuse clouds studied to date. Special care is necessary when choosing the appropriate atomic abundances in these circumstances. The O and C<sup>+</sup> abundances are now available along sight lines in our sample. Cartledge et al. (2001) found that the O abundance toward $`\rho `$ Oph D and HD 207198 differs by about 1-$`\sigma `$ from the diffuse cloud average, while that for C<sup>+</sup> is comparable to the average (Sofia et al. 2004). André et al. (2003) and Cartledge et al. (2004), respectively, provided results on O toward HD 206773 and $`\lambda `$ Cep that are not atypical. Thus, we continue to use the average abundance for diffuse clouds in the chemical analysis. For C<sup>+</sup> this is not unreasonable because the largest CO columns encountered in our survey are about 10<sup>16</sup> cm<sup>-2</sup>, which represents less than 10% of the carbon budget, and the neutral carbon column densities (Jenkins & Tripp 2001) are smaller still. Most calculations are based on the average interstellar radiation field of Draine (1978) and on $`T`$ $`=`$ 50 K, the temperature commonly inferred from excitation of low-lying rotational levels in H<sub>2</sub> and C<sub>2</sub> (e.g., van Dishoeck & Black 1986; Federman et al. 1994). The latter value is not critical because the results for $`n`$ are not very sensitive to $`T`$ — the estimated gas densities change by less than 10% when the temperature is changed by a factor of 2 (Pan et al. 2001). For especially molecular-rich clouds, lower values for $`T`$ are adopted, while slightly larger temperatures (65 K) are used for sight lines toward Cep OB3 where fewer CN components are present. These adopted temperatures and the UV radiation field are justified in $`\mathrm{\S }`$6.2.
The chemical model also involves the optical depths of grains at 1000 Å, $`\tau _{uv}`$, and threshold grain optical depths for ultraviolet photons, $`\tau _{uv}^\mathrm{o}`$, in somewhat indirect ways. These optical depths were used to determine the amount that photodissociation rates are attenuated and a coarse treatment of the C<sup>+</sup> to CO transition, respectively. For most sight lines, we previously adopted $`\tau _{uv}`$ $`=`$ 2$`A_V`$ (e.g., Federman et al. 1994), but such a value does not appear to be appropriate for our sample of sight lines. As noted by Federman et al. (1994), larger grains are present toward the stars comprising $`\rho `$ Oph, and a value of 1.4$`A_V`$ is used. Though the extinction law for directions in Cep OB3 is similar to the law representing the average for diffuse sight lines, the molecular gas is more clumpy (Sargent 1979) and Federman et al. (1994) adopted 1.4$`A_V`$ here as well. We do the same in light of the numerous components seen in K I absorption. For Cep OB2, extinction laws for several stars are available (Massa, Savage, & Fitzpatrick 1983; Savage et al. 1985; Aiello et al. 1988; Fitzpartick & Massa 1990; Megier et al. 1997), although sometimes only in tabular form. For the most part, the laws do not differ much from the average law, except for the sight line very rich in molecules $``$ HD 204827, where the ultraviolet extinction rises quite rapidly at the shortest wavelengths (Fitzpatrick & Massa 1990). Our starting point for the directions in Cep OB2 for $`\tau _{uv}`$ was 2$`A_V`$. The number of components toward stars in this association as well as in Cep OB3 created a further complication, which is compounded by the fact that foreground material along the $``$ 750-pc pathlengths is likely. First, we assumed that all the molecular gas is located in the star-forming regions. Extinction only from the molecular components was considered in the chemical analysis; it was estimated from the fraction of K I column associated with CH components. The fraction always exceeded 0.50 and was usually 0.80 or more. Second, the line-of-sight extinction was used because shadowing of one diffuse molecular cloud on another is more than likely. Since photodissociation is the dominant destruction pathway for most of the clouds in our study, uncertainties in $`\tau _{uv}`$ lead to uncertainties of about 30% in the inferred gas densities. Since the directions under study have a significant amount of reddening \[$`E`$($`BV`$) $``$ 0.3 mag\], but diffuse cloud abundances for C<sup>+</sup>, the perscription of Federman & Huntress (1989) for treating the C<sup>+</sup> to CO transition had to be revisited. The cause of the apparent inconsistency between reddening and C<sup>+</sup> abundance is the presence of numerous molecular components along each direction in Cep OB2 and OB3. For these associations, the threshold grain optical depth for ultraviolet photons, $`\tau _{uv}^\mathrm{o}`$, was increased to 3.75 from 2.0. As a result, only toward HD 204827, where $`N`$(CN) reaches the highest values, is there likely to be much conversion of C<sup>+</sup> into CO.
The results of our analysis, which are based on the rate equations from Federman et al. (1994) with updated rate data from Knauth et al. (2001) and Pan et al. (2001), appear in Table 7. The simplified model cloud has constant density and temperature. For each cloud along a specific direction, we list the factor giving the enhancement over Draine’s (1978) interstellar UV radiation field ($`I_{uv}`$), $`\tau _{uv}`$, the kinetic temperature ($`T`$), the gas density \[$`n`$ $`=`$ $`n`$(H) $`+`$ 2$`n`$(H<sub>2</sub>)\], the observed column densities for CH, C<sub>2</sub>, and CN \[$`N_o`$(CH), $`N_o`$(C<sub>2</sub>), and $`N_o`$(CN)\], and the predicted columns for C<sub>2</sub> and CN that best match the observations \[$`N_p`$(C<sub>2</sub>) and $`N_p`$(CN)\]. The column densities are given in units of 10<sup>12</sup> cm<sup>-2</sup>. The predicted column densities are always within 30% of the observed values; thereby providing another measure of the uncertainty in inferred gas density.
The inferred gas densities range from about 100 to slightly higher than 1000 cm<sup>-3</sup>, not untypical for diffuse molecular gas, and consistent with those by Weikard et al. (1996) — based on CO observations, they estimated densities for CO cloudlets of about 500 to 2000 cm<sup>-3</sup>. Our inferred densities probably represent lower limits to the true densities. First, we assumed that all the CH in a component participated in the production of CN although we find CN-like and CH<sup>+</sup>-like CH (see Section 5.2). Use of the relationship between CH<sup>+</sup>-like CH and CH<sup>+</sup> to account for CH in the lower density gas along the line of sight and not involved in CN chemistry is possible. We chose not to pursue this because the dispersion in the relationship (Fig. 1) could lead to unphysical results (negative columns) in some of our sample. Instead, we note that the clouds with the lowest densities are the ones most likely affected by this assumption. Second, we obtained the extinction for the molecular gas along the line of sight from the fraction of K I in CH components. The extinction would be smaller if only the CN components were used, resulting in higher densities to offset more photodissociation.
We analyzed the CN chemistry for several of these sight lines in the past (Federman & Lambert 1988; Federman et al. 1994; Pan et al. 2001). For the most part, the inferred gas densities are similar, usually within 30%. Much of the difference is likely caused by our current emphasis on individual components along a line of sight. The other factor, highlighted by the comparison with the results of Pan et al. (2001) based on the same chemical rates, is the use of updated column densities. For example, the difference in density for the $`0.8`$ km s<sup>-1</sup> component toward HD 206267C arises from the decrease in $`N`$(CN) from $`0.9\times 10^{12}`$ to $`0.4\times 10^{12}`$ cm<sup>-2</sup>.
Two sight lines, HD 207198 and HD 204827, yield surprisingly low gas densities for the large CN columns present. For HD 207198, the three CN components contain substantial amounts of CH<sup>+</sup>. If the CH associated with CH<sup>+</sup> does not take part in CN production and therefore must be removed from the CH column used in the chemical model, then densities a factor of 2 larger would be inferred. This modification is less effective for the clouds toward HD 204827. Instead, the low values inferred for $`n`$ result from the significant extinction. The large extinction, consistent with HD 204827 lying deepest within the CO contours of any sight line in our sample (see Fig. 8), greatly lessens the importance of destruction through photodissociation. This contrasts with the results for HD 62542 where high densities are required to reproduce the large CN column on a sight line with $`A_V`$ of 1 mag (Cardelli et al. 1990; Federman et al. 1994). These points are illustrated in Table 7 for the $`+1.9`$ km s<sup>-1</sup> component toward HD 207538. The second entry shows that $`n`$ increases five-fold for $`I_{uv}`$ $`=`$ 5, highlighting the importance of photochemistry for most of our sample. The third entry indicates that when about half the CH is associated with CN, $`n`$ increases a factor of $``$ 2.
### 6.2 H<sub>2</sub> Excitation
Column densities for H<sub>2</sub> rotational levels in the ground vibrational level carry information on physical conditions that provides a justification for our inputs in the chemical analysis. The relative populations of $`J`$ $`=`$ 0 and 1 give a reliable estimate for kinetic temperature (Savage et al. 1977). Since rotational transitions between ortho-H<sub>2</sub> (odd $`J`$) and para-H<sub>2</sub> (even $`J`$) are strictly forbidden, collisions involving protons and H<sub>2</sub> produce a thermal distribution between $`J`$ $`=`$ 0 and 1. High-J levels may be populated by photon pumping, collisions in shock-heated gas, and in the formation of H<sub>2</sub> molecules on dust grains. Column density ratios such as $`N`$(4) and $`N`$(2) yield upper limits to the UV flux causing the photon pumping when the other processes are neglected (but see Gry et al. 2005).
In Table 2, we listed kinetic temperatures from the H<sub>2</sub> analysis ($`\mathrm{\S }`$ 3) for lines of sight toward Cep OB2 and Cep OB3. The listed values of $`T_{1,0}`$ suggest that directions toward Cep OB3 have slightly higher kinetic temperatures (80 K on average) than those toward Cep OB2 (74 K). The temperatures used in the chemical analysis ($`\mathrm{\S }`$ 6.1) are somewhat lower, 65 K for Cep OB3 and 50 K for others, as would be expected for the less extensive, denser portions of a cloud where CN is detected. As for higher J levels, the observed ratios, log\[$`N`$(4)/$`N`$(2)\], presented in Table 3 lie between $``$2.68 and $``$0.95. Placing our results for H<sub>2</sub> on Fig. 3 of Browning, Tumlinson, & Shull (2003), which is based on $`I_{uv}`$ $``$ 15, allows us to infer the relevant UV flux to adopt. The data for HD 203374A and HD 208440 are to the left of the points in the figure in Browning et al., indicating fluxes consistent with the average interstellar value. On the other hand, the data for 9 Cep, HD 209339, and HD 217035A lie to the right, suggesting fluxes along these lines of sight are higher than the average interstellar value. An UV enhancement factor as large as 100 may be needed to explain our H<sub>2</sub> results. Because other processes leading to high-J level populations and because the presence of less shielded foreground gas that is mainly atomic were not considered, $`I_{uv}`$ values inferred from H<sub>2</sub> data should be regarded as upper limits for the chemical analysis ($`\mathrm{\S }`$6.1). Furthermore, such large fluxes are usually associated with correspondingly high densities and temperatures in photon dominated regions (e.g., Knauth et al. 2001), but our CO data indicate subthermal excitation in diffuse molecular gas. Therefore, our use of the average interstellar radiation field ($`I_{uv}`$ $``$ 1) seems reasonable under these circumstances along lines of sight with CN absorption. In other words, it appears that the material studied here is sufficiently far from the background star ($``$ 3 pc) that only the average interstellar field is important.
For the 14 directions with H<sub>2</sub> measurements, knowledge of both the $`b`$-value for each component and the mean kinetic temperature along the line of sight yields an estimate for the turbulent velocity in each component. The turbulent velocity for all but one of the 72 components is $``$ 1 km s<sup>-1</sup>. The results for gas toward Cep OB2 and OB3 are indistinguishable. Since the material along these sight lines is predominantly atomic, a comparison based on sound speeds derived from a ratio of specific heats of 5/3 is appropriate. For temperatures of 70 to 80 K, the sound speed is also about 1 km s<sup>-1</sup>. Thus, the turbulence appears to be sonic. One component toward $`\lambda `$ Cep is the exception: its turbulent velocity is $``$ 2.5 km s<sup>-1</sup>. The component resembles the CH<sup>+</sup>-like CH component toward $`\zeta `$ Oph (Lambert et al. 1990; Crawford 1997).
## 7 Discussion
### 7.1 Distribution of Species
In previous sections, similarities in average $`b`$-values and correlations between column densities allowed us to suggest that some species coexist, while others probe different portions of a diffuse molecular cloud with typical visual extinctions of 1 to 2 mag. We here investigate the distribution of observed species from other aspects. If two species coexist, their spectral profiles along a given line of sight should appear similar, although their line strengths may be quite different. For easy comparison, we construct apparent optical depth (AOD) profiles (Savage & Sembach 1991) for the absorption lines, and overplot scaled AOD profiles for different species on a given line of sight. Figure 4 shows the comparison of these profiles along the line of sight toward HD 207308. The scaling factors differ from species to species and from panel to panel. For instance, CH and K I AOD profiles are scaled by factors of 1/2.84 and 1/15.1 in the middle panel, respectively.
By comparing AOD profiles for six species on each of our sight lines, we find, along a given sight line, (1) that the CN profile is the narrowest one, whereas the Ca II profile is the widest one, (2) that CH and K I profiles are usually similar in both width and shape, and (3) that in most cases, CH<sup>+</sup> and Ca I profiles are wider than CH, in roughly 80% of the cases for CH<sup>+</sup>, and 60% for Ca I. The similarity among the profiles for the three species is not all that clear. In some cases, Ca I profiles are similar to CH<sup>+</sup> ones. In other cases, they may be more similar to those of CH than to those of CH<sup>+</sup>. It may depend upon the local gas density. Here we are assuming line width, or equivalently $`b`$-value, is a measure of the distribution along the line of sight.
The comparison of AOD profiles suggests that Ca II is the most widely distributed among the six species, followed by CH<sup>+</sup> and Ca I, then K I and CH, and finally CN. Distributions of velocity components reinforces this statement. Figure 5 is a histogram plot for distributions of velocity components with respect to $`V_{LSR}`$ along lines of sight in Cep OB2 (excluding components toward HD 206267A and $`\lambda `$ Cep because component structures are not available for some species). The Figure shows that Ca II components are much more widely distributed than those of other species, and that distributions of CH, K I, and Ca I have very similar shape. This again indicates correlations among the latter three species.
Previous studies of diffuse molecular clouds (Federman et al. 1984) showed that the CN molecule is produced in observable quantities only after significant amounts of precursor molecules, such as CH and C<sub>2</sub>, are available. In particular, the relationship between $`N`$(CN) and $`N`$(H<sub>2</sub>) has a slope much steeper than that seen for $`N`$(CH) vs. $`N`$(H<sub>2</sub>) (e.g., Federman 1982; Danks et al. 1984). Federman et al. (1984) attributed the steepness in slope to the number of chemical steps before the molecule can be detected. Higher densities facilitate the transformation along the chemical sequence. Our observations reveal that CN has the narrowest profile, smallest $`b`$-values, and smallest number of velocity components among the species. In other words, CN is the least widely distributed and occupies the smallest volume among the species. Therefore, we believe that CN mainly resides only in denser regions of diffuse molecular clouds, and that no observable amount of CN is present in low density clouds or in cloud envelopes. The correspondence between CO and CN noted above suggests, for the most part, that the same applies to CO. This is consistent with the slope found for $`N`$(CO) vs. $`N`$(H<sub>2</sub>) — see $`\mathrm{\S }`$ 7.4; CO can exist in regions with somewhat lower densities because faster (ion-molecule) reactions dominate its production (e.g., van Dishoeck & Black 1986).
In $`\mathrm{\S }\mathrm{\S }`$ 4 and 5, we showed that there are two kinds of CH-containing material: CN-like CH, which is associated with CN production, and CH<sup>+</sup>-like CH, which is related to the formation of CH<sup>+</sup>. This suggests that CH could coexist with either CN and CH<sup>+</sup>. Since in all velocity components where CN is detected, CH is observed, we believe that CH exists in the whole volume where CN resides. On the other hand, generally wider profiles, more components, and larger $`b`$-values indicate that CH occupies a larger volume than does CN.
Similarities in $`b`$-values, AOD profiles, and a good correlation between their column densities all suggest that CH and K I generally coexist, or share a large fraction of their volumes in diffuse clouds. The combined set of facts leads us to believe that these two species can also exist in less dense gas, or less dense regions of clouds (such as envelopes) where CN is not present. A relevant question is to what extent do the CH and K I distributions extend into the region of lower density. In other words, what is the lowest gas density that can be probed by CH and K I observations? The relatively low ionization potential (4.341 eV) for K I suggests that it cannot exist in a very low density environment because it must be shielded from photoionization. In addition, the $`b`$-values of CH and K I are only slightly larger than that of CN. This indicates that their distributions cannot extend much beyond the region containing CN. Therefore, CH and K I appear to be distributed in high- and moderately high-density regions of diffuse clouds ($`n`$ 30 cm<sup>-3</sup>). We also note that K may be depleted onto grains in very high density gas, and that there are some K I components without corresponding CH components. The facts suggest that K I may be absent in extremely dense regions of clouds, but slightly extend to less dense regions compared to CH (also see $`\mathrm{\S }`$ 5.3).
Because the AOD profiles for CH<sup>+</sup> and Ca I are generally wider than those for CH and K I, these species may occupy bigger volumes of a cloud. On the other hand, similar $`b`$-values and strong correlations between $`N`$(K I) and $`N`$(Ca I) suggest that K I and Ca I share a large fraction of their volumes. According to the chemical models of Federman (1982) and Danks et al. (1984) and the depletion characteristics of calcium (see next section), both CH<sup>+</sup> and Ca I cannot have a significant abundance in dense gas where CN and CH reside. CH<sup>+</sup> and Ca I only exist in less dense gas, and so their distributions extend throughout the envelope of a cloud beyond those of CH and K I. This is also consistent with the fact that Ca I has a greater first ionization potential than does K I (6.1 eV vs. 4.3 eV). On the other hand, neither CH<sup>+</sup> nor Ca I can exist in a very low density enviroment where insignificant amounts of H<sub>2</sub> (e.g., Federman 1982) are available for CH<sup>+</sup> production and where the main forms of Ca will be Ca II or Ca III. Thus, we conclude that CH<sup>+</sup> and Ca I mainly exist in moderately high- to intermediate-density gas ($`n`$ 10–300 cm<sup>-3</sup>).
The Ca II ion is the most widely distributed species among the species we observed. Its absorption can arise in components not seen in neutral atoms or molecules (see Figs. 4 and 5), and its $`b`$-value is the largest among our observed species. Considering that Ca readily depletes onto grains, Ca II should exist in a moderately high to low-density enviroment. Figure 6 is a schematic showing the distributions of species we discussed. In some low density clouds, the density is so low that molecules are not present even in the densest regions. Ca II may be the only species among those studied here in such clouds.
### 7.2 Ca and K Depletions
Cardelli, Federman & Smith (1991) noticed that the Ca I absorption line profiles were more similar to those of CH<sup>+</sup> than to those of CH on some lines of sight and suggested that calcium depletion depends on the density. For total column density along a line of sight, they found an inverse linear relationship between the ratios $`N`$(CN)/$`N`$(CH) and $`N`$(Ca I)/$`N`$(CH). Using a simple model, assuming photoionization equilibrium in predominantly neutral gas, where carbon is the primary source of electrons and where the space densities and column densities have the same functional form, they concluded that the fractional density of calcium varies roughly as $`n^3`$, due to the depletion. However, Welty et al. (2003) recently found a steeper relationship between the column density ratios with a slope of $``$2.5, implying that the density dependence of Ca column density is not steeper than $`n^{1.8}`$.
Comparing our spectra, we found some cases where the Ca I profiles are similar to those of CH<sup>+</sup> and in other cases the similarity is with CH profiles. A closer examination shows that the similarities may depend on the local gas density where the absorption arises. In most (roughly 75%) of the high density cases \[relatively large ratio of $`W_\lambda `$(CN)/$`W_\lambda `$(CH) or strong CN absorption\], Ca I profiles resemble those of CH<sup>+</sup>. For intermediate density clouds \[small ratio of $`W_\lambda `$(CN)/$`W_\lambda `$(CH) or no CN detection\], we did not find a clear trend, but we note that in some cases, the Ca I, CH<sup>+</sup>, and CH profiles are quite similar. These findings on profiles are consistent with the picture for the distribution of species drawn in $`\mathrm{\S }`$ 7.1. In a high $`n`$ enviroment, a large fraction of CH absorbers come from a region where no detectable amount of Ca I and CH<sup>+</sup> reside, whereas along the line of sight a large fraction of Ca I and CH<sup>+</sup> absorbers originate from a common volume. Then, Ca I profiles are more similar to those of CH<sup>+</sup> than those of CH. For a intermediate density cloud, the three species may coexist over a large fraction of the volume, and their profiles would be similar. Therefore, qualitatively, the calcium depletion depends upon the local gas density.
As already noted, using total column densities, Cardelli et al. (1991) and Welty et al. (2003) obtained different slopes for the relationship between $`N`$(CN)/$`N`$(CH) and $`N`$(Ca I)/$`N`$(CH). Figure 7 presents relations between the ratios $`N`$(CN)/$`N`$(CH) and both $`N`$(Ca I)/$`N`$(CH) and $`N`$(K I)/$`N`$(CH) for individual components. Fair correlations are found between the ratios with correlation coefficients of $``$0.68 and $``$0.40 and corresponding correlation confidences of 99.98% and 99.50%, respectively. Regression fits yield slopes of $``$1.03 $`\pm `$ 0.15 and $``$2.03 $`\pm `$ 0.20 for the respective relationships. Within the framework of Cardelli et al., these slopes imply that the calcium and potassium column densities are proportional to $`n^{2.9}`$ and $`n^{2.0}`$, respectively, due to their depletions. However, caution should be exercised in applying the framework because assumptions may not be valid. These species do not occupy exactly the same volume in a cloud, as stated in $`\mathrm{\S }`$ 7.1, and Welty et al. (2003) argued that the photoionization equilibrium may not strictly hold in diffuse clouds. While the inferred density dependence due to depletion may not be very accurate, it appears that calcium depletion varies more steeply on local density than potassium does, as suggested by both this analysis and by the slope in Figure 3.
### 7.3 Cloud Structure
#### 7.3.1 Large-scale Structure
The Cepheus bubble, at a distance of $``$ 750 pc from the Sun, has been the subject of many investigations. Radio H I and CO observations revealed that the bubble has structures on a scale of a few pc. Expanding shells of dense gas around H II regions have been inferred in previous studies of CO emission from the molecular clouds around OB stars. The large-scale CO molecular cloud shows a clumpy and irregular structure, with small cloudlets having sizes of a few parsecs (e.g., Patel et al. 1995, 1998; Weikard et al. 1996). Figure 8 is the map of CO emission for the Cep OB2 region adapted from Patel et al. (1995; 1998), with our stars projected onto the map. According to their locations on the CO map, we divided our stars into several groups (that do not reflect coeval stellar groups). For example, we consider HD 206267A, C, D, and HD 206773 as one group. They are associated with the same CO cloud, cloud number 18 in Table 2 of Patel et al. (1998). HD 207538, HD 208266, and $`\nu `$ Cep are in another group associated with cloud number 21. The AOD profiles for absorption lines are very similar for a given species in a group. The lines of sight in each group have most velocity components in common, especially within the $`V_{LSR}`$ range in which CO emission was detected from a cloudlet. Spectra for sight lines passing near edges of CO clouds show strong CN absorption. However, there are some exceptions. For example, the line of sight toward 14 Cep passes near the edge of one CO cloud, but only one weak CN component with $`N`$(CN) = 0.45 $`\times `$ 10<sup>12</sup> cm<sup>-2</sup> was detected. There is no nearby CO cloud toward HD 207198, but strong CN and CO absorption was detected on the line of sight. Since 14 Cep has the smallest extinction \[$`E`$($`BV`$)=0.35\] and the smallest estimated distance (650 pc) among our lines of sight in Cep OB2, 14 Cep most likely is located in front of its nearby CO cloudlet. In general, our data are consistent with the large-scale structure suggested by maps of CO radio emission. The $`V_{LSR}`$ of CN , CO, and CH components on a line of sight (including those toward $`\rho `$ Oph and Cep OB3) are in a range within which CO emission is detected from a cloudlet. For IC 1396 in Cep OB2, this correspondence is consistent with the conclusion of Weikard et al. (1996) that much of the material seen in CO emission is in front of HD 206267. The sensitivity of absorption measurements at visible and UV wavelengths allows us to detect smaller molecular columns than is possible through emission line maps, providing further detail of cloud dispersal in regions of star formation.
#### 7.3.2 Small-scale Structure
While it is well established that the column densities of some atomic (H I, Na I, K I) and molecular (H<sub>2</sub>CO, OH) species vary on scales of 10 to 10<sup>5</sup> AU along diffuse sight lines, detailed information on the physical conditions causing the variations is generally not available. A key question posed by the observations is whether the subparsec structure is caused by density variations (e.g., Frail et al. 1994; Lauroesch & Meyer 1999; Crawford et al. 2000), by fluctuations in ionization equilibrium (Lauroesch & Meyer 1999; Welty & Fitzpartrick 2001), by the geometric structure of clouds (Heiles 1997), or by something else. A crucial point is to determine accurate physical conditions for the gas showing variations in column density. With gas densities inferred from a chemical model for CH, C<sub>2</sub>, and CN (see $`\mathrm{\S }`$ 6.1) for velocity components on lines of sight toward members of multiple star systems, we now discuss small scale structure in diffuse molecular gas.
There are three multiple star systems, $`\rho `$ Oph A/B/C/D, HD 206267 A/C/D, and HD 217035 A/B, in our data set. Separations between members of systems range from 450 to 38,320 AU (0.02 to 0.18 pc). Substantial variations in CN absorption are observed among sight lines of $`\rho `$ Oph. There are striking differences in the CN, CH, CH<sup>+</sup>, Ca I, and K I profiles among the three sight lines of HD 206267. No CN is detected toward either line of sight in the HD 217035 system, but a significant difference in the CH<sup>+</sup> profiles for the two members of the system is observed. The corresponding column densities differ by a factor of $``$2 in individual velocity components.
Pan et al. (2001) discussed gas density variations on sight lines of the multiple star systems HD 206267 and HD 217035. Because Paper I employed a different method to determine column densities of individual components and a slightly different data reduction procedure, velocity component structures and their column densities might be slightly different from those in Pan et al. However, the conclusion of Pan et al. that the component gas densities can differ by factors of 5 between adjacent lines of sight is not altered. From Table 7, we see that the derived gas densities of the components toward members of HD 206267 at the same or similar $`V_{LSR}`$ differ by factors of 2−7. (If the upper limits for gas density are considered, the difference can be as high as 10.) In the case of the $`\rho `$ Oph multiple star system, the diffuse molecular cloud toward $`\rho `$ Oph C has a density about twice that inferred along the other three nearby sight lines. While less dramatic than the results for the system HD 206267, the result for $`\rho `$ Oph reinforces the conclusions of Pan et al. Density contrasts up to factors of 10 over scale lengths of 10,000 to 20,000 AU are present in the densest portions of diffuse clouds, those sampled by CN absorption; variations of the order of 10<sup>4</sup> suggested by 21 cm observations (e.g., Frail et al. 1994) must have a different origin. Possibilities include fluctuations in ionization equilibrium (Lauroesch & Meyer 1999; Welty & Fitzpatrick 2001) and the geometric structure of clouds (Heiles 1997). The smallest scale length probed by our measurements, 450 AU between $`\rho `$ Oph A and B, shows no variation in modeled density greater than about 35%.
Cloud thicknesses for individual components (cloudlets) along the line of sight toward the background star can be estimated as well. The total column density of protons, $`N_{tot}`$(H), for a component can be estimated from our measurements of $`N`$(K I) (from Paper I) and the relationship between $`N`$(K I) and $`N_{tot}`$(H), log\[$`N`$(K I)\] $`=(26.9\pm 2.7)+(1.8\pm 0.1)`$ log\[$`N_{tot}`$(H)\] derived by Welty & Hobbs (2001), under the assumption that it applies to components although they used total column densities. Then $`n`$ and $`N_{tot}`$(H) are used to estimate the thickness of a cloudlet. For the multiple star systems, $`\rho `$ Oph and HD 206267, the thickness is within 50% of 1 pc, except for the two low-density components toward HD 206267A where the thickness is about 5 pc. Since $`n`$ refers to the dense region where CN resides, the mean $`n`$ for the whole cloudlet should be lower. Our estimates should then be regarded as lower limits. If the lower limits are not far below their true values, they suggest that some of the cloudlets (in Cep OB2) are sheet-like with aspect ratios of 5 to 10 because maps of CO millimeter-wave emission show that a typical size of a cloudlet on the sky is a few parsecs (Patel et al. 1995, 1998; Weikard et al. 1996). Our results, therefore, provide evidence for the need to consider non-spherical geometries (Heiles 1997) when analyzing small scale structure in diffuse clouds.
### 7.4 Individual Cloud Systems
The diffuse gas associated with the three star-forming regions, Cep OB2, Cep OB3, and $`\rho `$ Oph (belonging to Sco OB2), were observed. Table 8 lists the average number, $``$M$``$, of velocity components per sight line for each species and the average column density of each component for each species.<sup>6</sup><sup>6</sup>6CO column densities for individual components were obtained along 15 lines of sight toward Cep OB2 and Cep OB3 by analyzing HST and FUSE spectra — see $`\mathrm{\S }`$ 3. Because different spectral resolutions comprise this sample, it may be not fully appropriate to include the CO results in the comparison. Still, our analyses show that CO has fewer components than CH but more than CN, and that, on average, sight lines toward Cep OB2 and Cep OB3 have about the same number of CO components per line of sight, but components in Cep OB2 have much larger mean column density. The table shows that lines of sight toward Cep OB2 and Cep OB3 have many more Ca II and K I components than those toward $`\rho `$ Oph. This is consistent with radio observations (Sargent 1977; Weikard et al. 1996; Patel et al. 1998), which revealed that clouds in Cep OB2 and Cep OB3 have complicated structures. Table 8 also reveals that Cep OB3 has the greatest mean number of Ca II and CH<sup>+</sup> components per sight line, the largest average column density per Ca II component, but the smallest corresponding values for CN among the three star-forming regions. This indicates that a larger fraction of low density material is probed by the lines of sight toward Cep OB3. Because CN resides only in the denser region of diffuse molecular clouds, while Ca II is widely distributed, the ratio, $``$M$`{}_{\mathrm{CN}}{}^{}\times N`$(CN)$``$/ ($``$M$`{}_{\mathrm{CaII}}{}^{}\times N(`$Ca II)$``$), may reflect the fraction of denser material in the clouds. The ratio for $`\rho `$ Oph is about 2 and 10 times larger than for Cep OB2 and Cep OB3, respectively.
Lines of sight toward Cep OB3 have larger mean extinction, $`E(BV)`$, (0.78 vs. 0.53 mag), greater $``$M$`{}_{\mathrm{CaII}}{}^{}`$, and larger $`N`$(Ca II)$``$ compared to those toward Cep OB2. This suggests that more material is intercepted along lines of sight in Cep OB3. However, Cep OB3 clouds have a lower mean density. This indicates that Cep OB2 clouds are geometrically thinner, along the line of sight, than Cep OB3 clouds.
Using total H<sub>2</sub> column densities from Copernicus observations, Federman (1982) and Federman et al. (1980) found that total CH and CO column densities along lines of sight were well correlated with total H<sub>2</sub> column densities. In the present study, we obtained total H<sub>2</sub> column densities along 15 sight lines toward Cep OB2 and Cep OB3 based on $`FUSE`$ observations. Although $`FUSE`$ provides higher quality spectra, we still could not extract reliable column densities for individual components due to very large line optical depths, thus preventing us from seeking correlations between $`N`$(CH) and $`N`$(CO) with $`N`$(H<sub>2</sub>) on a component basis as we did for other relationships in $`\mathrm{\S }`$ 5. Figure 9 shows plots of total column density for CH and CO versus H<sub>2</sub> along the 15 lines of sight. In general, the Figure reveals the same trends, as previous studies (Federman et al. 1980; Federman 1982; Danks et al. 1984): $`N`$(CH) and $`N`$(CO) increase with increasing $`N`$(H<sub>2</sub>) with a slope of 0.95 $`\pm `$ 0.10 for CH (Cep OB2) and 3.16 $`\pm `$ 0.34 (Cep OB2) and 2.93 $`\pm `$ 0.58 (Cep OB3) for CO. For comparison, Danks et al. (1984) obtained a slope of 0.90 $`\pm `$ 0.10 for $`N`$(CH) vs. $`N`$(H<sub>2</sub>), and Federman et al. (1980) found a slope of about 2 for the CO trend. However, Cep OB2 sight lines are clearly distinguished from those toward Cep OB3. Lines of sight toward Cep OB2 have higher $`N`$(CH) and $`N`$(CO) column densities for a given $`N`$(H<sub>2</sub>). This implies that the physical conditions in these star-forming regions, such as average gas density, may be quite different. Since there are a larger fraction of CN components toward stars in Cep OB2, the diffuse molecular clouds toward Cep OB2 have higher average gas density. The direction toward HD 217312 in Cep OB3 represents an intermediate case: there is a large enough fraction of CN components to place it among the CH data for Cep OB2, but not enough for the CO data because CO traces the densest gas with CN.
Absorption from the diffuse molecular gas for a specific star-formation region shows similar overall structure, and the length scales probed by our measurements are similar for the three star-forming regions. However, the diffuse molecular clouds associated with Cep OB3 exhibit different characteristics compared to those associated with $`\rho `$ Oph and Cep OB2. The cloudlets have lower gas density and lower total molecular mass (see $`\mathrm{\S }`$ 2). Compared with Cep OB2 and Cep OB3, material in the $`\rho `$ Oph star-forming cloud is much more compact. This is consistent with our finding that lines of sight toward Cep OB2 and Cep OB3 have many more Ca II and K I components. It is not clear, however, if the differences seen in the diffuse gas surrounding the parent molecular clouds are related to the star formation histories in the clouds. The molecular cloud system associated with Cep OB3 is currently forming its second generation of stars, whereas a third generation of stars is forming in clouds associated with Cep OB2 and possibly in the $`\rho `$ Oph region. The older stellar subgroup in Cep OB3 has an age of $``$ 7 Myr, whereas the oldest stars in Cep OB2 and in Sco OB2 (to which $`\rho `$ Oph belongs) are about 15 Myr, with the second older groups $``$ 7 Myr (see $`\mathrm{\S }`$ 2). It would be useful to clarify possible connections between physical conditions and differences in star-forming history. It seems that star formation in Cep OB2 and Cep OB3 has proceeded via sequential triggering (Patel et al. 1995; Sargent 1979), whereas star formation scenarios for $`\rho `$ Oph are still a matter of debate (de Geus 1989; Sartori et al. 2003). de Geus et al. (1989) derived respective ages of 11–12, 14–15, and 5–6 Myr for three stellar subgroups, LCC, UCL, and US, of Sco OB2. They argued that either the classic picture of sequential star formation (Blaauw 1964; Elmegreen & Lada 1977) is not valid for the Sco OB2 Association as a whole because the oldest subgroup, UCL, is in between the two younger ones, or for some unknown reason massive star formation was initiated near the middle of the original giant molecular cloud. However, Pan (2002) found that uncertainties in the above derived ages are large, and that LCC and UCL might actually have a similar age and thus star formation in Sco OB2 may not contradict the model for sequential star formation. Based on spatial positions and ages of stellar subgroups in Sco OB2, Pan (2002) conjected that LCC and UCL could represent the first generation stars, the members of US the second, and the newly formed stars in the $`\rho `$ Oph cloud the third generation. If all three star-forming regions follow the classical picture of sequential star formation (Elmegreen & Lada 1977), and assuming stellar generations in Sco OB2 as that suggested by Pan (2002), then the age periods between the 1<sup>st</sup> and 2<sup>nd</sup> and between the 2<sup>nd</sup> and 3<sup>rd</sup> are about 7–9 and 5–6 Myr in all three star-forming regions. Support for this picture comes from Sartori et al. (2003). They found that the US subgroup is 4–8 Myr old and that LCC and UCL have the same age, 16–20 Myr, by adopting Hipparcos distances, newer isochrones, and temperatures derived from the spectral types. However, from the spatial distribution, the space velocity, and the age distribution of stars in Sco OB2, Sartori et al (2003) argued that star formation in $`\rho `$ Oph, most likely, is connected to star formation in nearby spiral arms rather than sequential triggering. This may explain why material in $`\rho `$ Oph is more compact than in Cep OB2 and Cep OB2.
## 8 Summary
We obtained and analyzed FUV data on H<sub>2</sub> and CO along 15 lines of sight toward Cep OB2 and Cep OB3 to complement the optical survey results of Paper I. Average $`b`$-values of individual components are found to be consistent with those obtained from ultra-high resolution spectra. The $`b`$-values for CN components are also consistent with those derived from the doublet ratio method (Strömgren 1948; Gredel et al. 1991). The inferred CN excitation temperatures range from 2.55 to 2.89 K, with an average of 2.75 $`\pm `$ 0.10 K, indicating that no significant excitation in addition to that due to the CMB is present in the diffuse molecular clouds in our sample.
With our large high-resolution data set, we examined some possible correlations between column densities of different species based on individual velocity components. In general, relationships are found to be tighter compared with those based on total column densities. The main correlations are the following. (1) There are two kinds of CH in diffuse molecular gas: CN-like CH and CH<sup>+</sup>-like CH. Disentangling the amount of CH in each appears possible when utilizing the relation between column densities of CH<sup>+</sup>-like CH and of CH<sup>+</sup>. (2) Column densities of CN and CH components are well correlated. The relationship becomes even stronger if the corresponding column density of CH<sup>+</sup>-like CH is subtracted from the total CH column. A close correspondence between CN and CO is also seen. (3) The trends between CH and K I, and Ca I and K I show tight correlatations. The slope of the relationship for $`N`$(Ca I) vs. $`N`$(K I) is 0.60 $`\pm `$ 0.04, significantly smaller than 1.0, which may suggest that calcium depletion depends more strongly on local density than does potassium depletion.
We investigated the spatial distributions of species by analyzing apparent optical depth profiles, distributions of components with respect to $`V_{LSR}`$, $`b`$-values, ionization potentials of elements, and correlations between column densities. These analyses show that different species are restricted to specific density regions. The CN and CO molecules mainly reside in denser regions of diffuse molecular clouds. No observable amount of CN is present in low density clouds or in cloud envelopes. The species CH and K I are distributed in high- and moderately high-density gas ($`n30`$ cm<sup>-3</sup>). (K I may not be present in very dense regions of clouds because there may be enhanced K depletion.) CH<sup>+</sup> and Ca I are mainly distributed in moderately high- and intermediate-density regions ($`n`$ 10–300 cm<sup>-3</sup>). The Ca II ion is the most widely distributed among the observed species; it can exist in moderately high-density regions but favors a relatively low density environment.
Several correlations and analyses show that both Ca and K may be depleted onto grains in denser gas, and that the Ca depletion has a steeper dependence on local density. Within the framework of Cardelli et al (1991), our data suggest that Ca column density varies roughly as $`n^{2.9}`$, whereas K column is proportional to $`n^{2.0}`$ due to their depletions. However, we note that their assumptions involving space densities and column densities may not be applicable.
Gas densities for individual velocity components where CN is detected were inferred from a chemical model. The inferred gas densities are independent of assumptions about cloud shape, on which some previous calculations relied. Based on these derived gas densities, we examined the large- and small-scale structure of clouds. In general, our data are consistent with the large-scale structure suggested by maps of CO radio emission. Our analysis reveals the presence of variations in gas density among sight lines in two multiple star systems, $`\rho `$ Oph and HD 206267. The gas density is seen to vary by factors of 5−10 over scales of $``$ 10,000 AU. This indicates that observed column density variations are due in part to a change in gas density. However, variations of the order of 10<sup>4</sup> suggested by 21 cm observations (e.g., Frail et al. 1994) must have a different origin. Cloud thicknesses for individual components (cloudlets) are estimated to be $``$ 1pc, which may suggest that some of the cloudlets are sheet-like with aspect ratios of 5 to 10. The estimation of cloud thickness provides evidence for the need to consider non-spherical geometries (Heiles 1997) when analyzing small-scale structure in diffuse clouds.
Comparisons show that there are both similarities and differences in general characteristics of diffuse gas in the three star-forming regions, $`\rho `$ Oph, Cep OB2 and Cep OB3. For example, cloud turbulence in the three regions does not differ significantly (but note that Cep OB3 seems to have consistently larger $`b`$-values than $`\rho `$ Oph and Cep OB2). Clouds in Cep OB2 and Cep OB3 have more complex (clumpy) structures than those in $`\rho `$ Oph; in other words, material in $`\rho `$ Oph is more compact than in Cep OB2 and Cep OB2. The molecular cloudlets associated with Cep OB3 have lower gas density and lower total molecular mass, compared to those associated with $`\rho `$ Oph and Cep OB2. The differences may be related to the star formation histories for the regions.
This paper presents analyses of the physical and chemical structure of the diffuse gas in star-forming regions based on absorption at visible and UV wavelengths. As in Federman et al. (1997) and Knauth et al. (2001), such analyses provide results for foreground diffuse molecular gas that are both consistent and complement results of the denser gas probed at longer wavelength. The combined set of data yields a comprehensive view of the effects of star formation on molecular clouds.
This research made use of the Simbad database operated at CDS Strasbourg, France. K. P. acknowledges KPNO for providing board and lodging in Tucson during observing runs. We thank the anonymous referee for constructive suggestions that improved the paper. The research is based in part on observations made with the NASA–CNES–CSA Far Ultravoilet Spectroscopic Explorer (FUSE), which is operated for NASA by the Johns Hopkins University under NASA contract NAS5–32985. Additional observations made with the NASA/ESA Hubble Space Telescope were obtained from the Multiwavelength Archive at the Space Telescope Science Institute; STScI is operated by the Association of Universities for Research in Astronomy, Inc. under NASA contract NAS5–26555. This work was supported by NASA grants NAG5–4957, NAG5–8961, and NAG5–10305 and grants GO–08693.03–A and AR–09921.01–A from the Space Telescope Science Institute. We acknowledge use of the regression subroutine obtained from the Penn State statistical software archive. |
warning/0507/math0507086.html | ar5iv | text | # Conformal blocks revisited
## 1. Canonical construction of the Virasoro algebra
In this section we fix a $``$-algebra $`R`$ and a $`R`$-algebra $`𝒪`$ isomorphic to the formal power series ring $`R[[t]]`$. In other words, $`𝒪`$ comes with a principal ideal $`𝔪`$ so that $`𝒪`$ is complete for the $`𝔪`$-adic topology and $`𝒪/𝔪^j`$ is for $`j=1,2,\mathrm{}`$ a free $`R`$-module of rank $`j`$. A generator $`t`$ of $`𝔪`$ will then identify $`𝒪`$ with $`R[[t]]`$. We denote by $`L`$ the localization of $`𝒪`$ obtained by inverting a generator of $`𝔪`$. For $`N`$, $`𝔪^NL`$ has the obvious meaning as a $`𝒪`$-submodule of $`L`$. The *$`𝔪`$-adic topology* on $`L`$ is the topology that has the collection of cosets $`\{f+𝔪^N\}_{fL,N}`$ as a basis of open subsets. We sometimes write $`F^NL`$ for $`𝔪^N`$. We further denote by $`\theta `$ the $`L`$-module of continuous $`R`$-derivations from $`L`$ into $`L`$ and by $`\omega `$ the $`L`$-dual of $`\theta `$. These $`L`$-modules come with filtrations (making them principal filtered $`L`$-modules): $`F_N\theta `$ consists of the derivations that take $`𝔪`$ to $`𝔪^{N+1}`$ and $`F^N\omega `$ consists of the $`L`$-homomorphisms $`\theta L`$ that take $`F^0\theta `$ to $`𝔪^N`$. So in terms of the generator $`t`$ above, $`L=R((t))`$, $`\theta =R((t))\frac{d}{dt}`$, $`F^N\theta =R[[t]]t^{N+1}\frac{d}{dt}`$, $`\omega =R((t))dt`$ and $`F^N\omega =R[[t]]t^{N1}dt`$.
The residue map $`\mathrm{Res}:\omega R`$ which assigns to an element of $`R((t))dt`$ the coefficient of $`t^1dt`$ is independent of the choice of $`t`$. The $`R`$-bilinear map
$$r:L\times \omega R,(f,\alpha )\mathrm{Res}(f\alpha )$$
is a topologically perfect pairing of filtered $`R`$-modules: we have $`r(t^k,t^{l1}dt)=\delta _{k,l}`$ so that any $`R`$-linear $`\varphi :LR`$ which is continuous (i.e., $`\varphi `$ zero on $`𝔪^N`$ for some $`N`$) is definable by an element of $`\omega `$ (namely by $`_{k<N}\varphi (t^k)t^{k1}dt`$) and likewise for a $`R`$-linear continuous map $`\omega R`$.
### A trivial Lie algebra
If we regard $`L^\times `$ as an algebraic group over $`R`$ (or rather as a group object in a category of ind schemes over $`R`$), then its Lie algebra, denoted here by $`𝔩`$, is $`L`$, regarded as a $`R`$-module with trivial Lie bracket. It comes with a decreasing filtration $`F^{}𝔩`$ (as a Lie algebra) defined by the valuation. The universal enveloping algebra $`U𝔩`$ is clearly $`\mathrm{Sym}_R^{}(𝔩)`$. The ideal $`U_+𝔩U𝔩`$ generated by $`𝔩`$ is also a right $`𝒪`$-module (since $`𝔩`$ is). We complete it $`𝔪`$-adically: given an integer $`N0`$, then an $`R`$-basis of $`U_+𝔩/(U𝔩F^N𝔩)`$ is the collection $`t^{k_1}\mathrm{}t^{k_r}`$ with $`k_1k_2\mathrm{}k_r<N`$. So elements of the completion
$$U_+𝔩\overline{U}_+𝔩:=\underset{}{\mathrm{lim}}_NU_+𝔩/U𝔩F^N𝔩$$
are series of the form $`_{i=1}^{\mathrm{}}r_it^{k_{i,1}}\mathrm{}t^{k_{i,r_i}}`$ with $`r_iR`$, $`k_{1,i}k_{2,i}\mathrm{}k_{i,r_i}`$, $`\{k_{i,1}\}_i`$ bounded from below and $`lim_i\mathrm{}k_{i,r_i}=\mathrm{}`$. We put $`\overline{U}𝔩:=R\overline{U}_+𝔩`$, which could of course have been defined directly as
$$U𝔩\overline{U}𝔩:=\underset{}{\mathrm{lim}}_NU𝔩/U𝔩F^N𝔩.$$
We will refer to this construction as the *$`𝔪`$-adic completion on the right*. (In the present case, this is no different from $`𝔪`$-adic completion on the left, because $`𝔩`$ is commutative.)
Any $`D\theta `$ defines an $`R`$-linear map $`\omega L`$ which is selfadjoint relative our topological pairing: $`r(D,\alpha ,\beta )=r(\alpha ,D,\beta )`$. We use that pairing to identify $`D`$ with an element of the closure of $`\mathrm{Sym}^2𝔩`$ in $`\overline{U}𝔩`$. Let $`C(D)`$ be half this element, so that in terms of the above topological basis:
$$C(D)=\frac{1}{2}\underset{l}{}D,t^{l1}dtt^l.$$
In particular we have for $`D=D_k=t^{k+1}\frac{d}{dt}`$
$$C(D_k)=\frac{1}{2}\underset{i+j=k}{}:t^it^j:.$$
We here adhered to the normal ordering convention (the factor with the highest index comes last and hence acts first), to make the righthand side look like an element of $`\overline{U}𝔩`$, although there is no need for this as $`t^it^j=t^jt^i`$. Observe that the map $`C:\theta \overline{U}𝔩`$ is continuous.
### Oscillator algebra
The residue map defines a central extension of $`𝔩`$, the *oscillator algebra* $`\widehat{𝔩}`$, which as a $`R`$-module is simply $`𝔩\mathrm{}R`$ and has Lie bracket
$$[f+\mathrm{}r,g+\mathrm{}s]:=\mathrm{}\mathrm{Res}(gdf).$$
So $`[t^k,t^l]=\mathrm{}k\delta _{k,l}`$ and $`R\mathrm{}`$ is central. We filter $`\widehat{𝔩}`$ by letting $`F^N\widehat{𝔩}`$ be $`F^N𝔩`$ for $`N>0`$ and $`F^N𝔩+R\mathrm{}`$ for $`N0`$. We complete $`U\widehat{𝔩}`$ $`𝔪`$-adically on the right:
$$U\widehat{𝔩}\overline{U}\widehat{𝔩}:=\underset{}{\mathrm{lim}}_NU\widehat{𝔩}/U\widehat{𝔩}F^N𝔩.$$
Since $`\mathrm{}`$ is in the center of $`𝔩`$, this is a $`R[\mathrm{}]`$-algebra (for a similar reason it is even a $`R[e,\mathrm{}]`$-algebra if $`e=t^0`$ denotes the unit element of $`L`$ viewed as an element of $`𝔩`$). As a $`R[\mathrm{}]`$-algebra it is a quotient of the tensor algebra of $`𝔩`$ (over $`R`$) tensored with $`R[\mathrm{}]`$, $`_R^{}𝔩_RR[\mathrm{}]`$, by the two-sided ideal generated by the elements $`fggf\mathrm{}\mathrm{Res}(gdf)`$. As a $`R[\mathrm{}]`$-module it has for topological basis the collection $`t^{k_1}\mathrm{}t^{k_r}`$ with $`r0`$, $`k_1k_2\mathrm{}k_r`$. Since $`\widehat{𝔩}`$ is not abelian, the left and right $`𝔪`$-adic topologies differ. For instance, $`_{k1}t^kt^k`$ does not converge in $`\overline{U}\widehat{𝔩}`$, whereas $`_{k1}t^kt^k`$ does. Notice that the obvious surjection $`\pi :U\widehat{𝔩}U𝔩`$ is simply the reduction modulo $`\mathrm{}`$ of $`U\widehat{𝔩}`$ and likewise for their completions. Observe that the filtrations of $`𝔩`$ and $`\widehat{𝔩}`$ determine decreasing filtrations of the their (completed) universal enveloping algebras. For instance, $`F^NU\widehat{𝔩}=_{r0}_{n_1+\mathrm{}+n_rN}F^{n_1}\widehat{𝔩}\mathrm{}F^{n_r}\widehat{𝔩}`$.
Denote by $`𝔩_2`$ the image of $`𝔩^2\widehat{𝔩}^2U\widehat{𝔩}`$. Under the reduction modulo $`\mathrm{}`$, $`𝔩_2`$ maps onto $`\mathrm{Sym}^2(𝔩)U𝔩`$ with kernel $`R\mathrm{}`$. Its closure $`\overline{𝔩}_2`$ in $`\overline{U}\widehat{𝔩}`$ maps onto the closure of $`\mathrm{Sym}^2(𝔩)`$ in $`\overline{U}𝔩`$ with the same kernel. We denote by $`\widehat{\theta }`$ the set of pairs $`(D,u)\theta \times \overline{𝔩}_2`$ for which $`C(D)`$ is the mod $`\mathrm{}`$ reduction of $`u`$ so that we have an exact sequence
$$0\mathrm{}R\widehat{\theta }\theta 0$$
of $`R`$-modules and a natural $`R`$-homomorphism $`\widehat{C}:\widehat{\theta }\overline{U}\widehat{𝔩}`$. The generator $`t`$ defines a (noncanonical) section of $`\widehat{\theta }\theta `$:
$$D\theta \widehat{D}:=(D,\frac{1}{2}\underset{j}{}:D,t^{j1}dtt^j:)\theta \times \widehat{𝔩}.$$
###### Lemma 1.1.
We have
1. $`[\widehat{C}(\widehat{D}),f]=\mathrm{}D(f)`$ as an identity in $`\overline{U}\widehat{𝔩}`$ (where $`f𝔩\widehat{𝔩}`$) and
2. $`[\widehat{C}(\widehat{D}_k),\widehat{C}(\widehat{D}_l)]=\mathrm{}(lk)\widehat{C}(\widehat{D}_{k+l})+\mathrm{}^2\frac{1}{12}(k^3k)\delta _{k+l,0}`$.
###### Proof.
For the first statement we compute $`[\widehat{C}(\widehat{D}_k),t^l]`$. The only terms in the expansion of $`_{i+j=k}:t^it^j:`$ that can contribute are for the form $`[t^{k+l}t^l,t^l]`$ or $`[t^lt^{k+l},t^l]`$, (depending on whether $`k+2l0`$ or $`k+2l0`$) and with coefficient $`\frac{1}{2}`$ if $`k+2l=0`$ and $`1`$ otherwise. In all cases the result is $`\mathrm{}lt^{k+l}=\mathrm{}D_k(t^l)`$.
Formula (i) implies that
$$\begin{array}{c}[\widehat{C}(\widehat{D}_k),\widehat{C}(\widehat{D}_l)]=\underset{N\mathrm{}}{lim}\underset{|i|N}{}\frac{1}{2}\left(D_k(t^i)t^{li}+t^iD_k(t^{li})\right)\hfill \\ \hfill =\mathrm{}\underset{N\mathrm{}}{lim}\underset{|i|N}{}\left(it^{k+i}t^{li}+t^i(li)t^{k+li}\right)\end{array}$$
This is up to a reordering equal to $`\mathrm{}(lk)\widehat{C}(\widehat{D}_{k+l})`$. The terms which do not commute and are in the wrong order are those for which $`0<k+i=(li)`$ (with coefficient $`i`$) and for which $`0<i=(k+li)`$ (with coefficient $`(li)`$). This accounts for the extra term $`\mathrm{}^2\frac{1}{12}(k^3k)\delta _{k+l,0}`$
This lemma suggests we rescale $`\widehat{C}`$ as
$$T:=\frac{1}{\mathrm{}}\widehat{C}:\widehat{\theta }\overline{U}\widehat{𝔩}[\frac{1}{\mathrm{}}]$$
and write $`c_0`$ for $`(0,\mathrm{})\widehat{\theta }`$ since then
1. $`T`$ is injective and maps $`\widehat{\theta }`$ onto a Lie subalgebra of $`\overline{U}\widehat{𝔩}[\frac{1}{\mathrm{}}]`$ with $`c_0\widehat{\theta }1`$,
2. if we transfer the Lie bracket to $`\widehat{\theta }`$, we find that
$$[\widehat{D}_k,\widehat{D}_l]=(lk)\widehat{D}_{k+l}+\frac{k^3k}{12}\delta _{k+l,0}c_0,$$
3. and $`\mathrm{ad}_{T(\widehat{D})}`$ leaves $`𝔩`$ invariant (as a subspace of $`\overline{U}\widehat{𝔩}`$) and acts on that subspace by derivation with respect to $`D`$,
Thus we get a central extension of $`\theta `$ with a one-dimensional center canonically isomorphic to $`R`$ ($`c_0`$ corresponds to $`1R`$), the *Virasoro algebra* (of the $`R`$-algebra $`L`$).
###### Remark 1.2.
An alternative coordinate free definition of the Virasoro algebra, based on the algebra of pseudodifferential operators on $`L`$, can be found in .
### Fock representation
It is clear that $`F^𝔩=F^1\widehat{𝔩}`$ is an abelian subalgebra of $`\widehat{𝔩}`$. We put
$$𝔽:=(U\widehat{𝔩}/U\widehat{𝔩}F^1𝔩)[\frac{1}{\mathrm{}}].$$
This is a representation of $`𝔩`$ over $`R`$; at the same time it is a $`R[e,\mathrm{},\mathrm{}^1]`$-module and as such it is free with basis the collection $`t^{k_r}\mathrm{}t^{k_1}v_o`$, where $`v_o`$ denotes the image of $`1`$, $`r0`$ and $`1k_1k_2\mathrm{}k_r`$ (for $`r=0`$, read $`v_o`$). This shows that we may identify $`𝔽`$ with $`(\overline{U}\widehat{𝔩}/\overline{U}\widehat{𝔩}F^1\widehat{𝔩})[\frac{1}{\mathrm{}}]`$, which makes it a representation for $`\widehat{\theta }`$ over $`R`$. This is the *Fock representation* of $`\widehat{\theta }`$. The $`R[e,\mathrm{},\mathrm{}^1]`$-subalgebra of $`U\widehat{𝔩}`$ generated by the elements $`t^k`$, $`k=1,2,\mathrm{}`$ is a polynomial algebra (in infinitely many variables) which projects isomorphically onto $`𝔽`$. But this subalgebra depends on $`t`$; $`𝔽`$ itself does not seem to be an $`R`$-algebra in a canonical way. On the other hand, it is naturally filtered with $`F^N𝔽=F^N\overline{U}\widehat{𝔩}(v_0)`$.
It follows from Lemma 1.1 that when $`D\theta `$, then
$$\begin{array}{c}T(\widehat{D})t^{k_r}\mathrm{}t^{k_1}v_0=\hfill \\ \hfill =\left(\underset{i=1}{\overset{r}{}}X_{\alpha _r}t^{k_r}\mathrm{}D(t^{k_i})\mathrm{}t^{k_1}\right)v_0+t^{k_r}\mathrm{}t^{k_1}T_𝔤(\widehat{D})v_0.\end{array}$$
For $`DF^0\theta `$ we have $`T_𝔤(\widehat{D})v_0=0`$ and so $`F^0\theta `$ acts on $`𝔽`$ by coefficientwise derivation.
## 2. The Sugawara construction
In this section, we fix a base field $`k`$ of caracteristic zero and a simple Lie algebra $`𝔤`$ over $`k`$ of finite dimension. We retain the data and the notation of Section 1, except that we now assume $`R`$ to be a $`k`$-algebra.
### Loop algebras
The space of symmetric invariant bilinear forms $`𝔤\times 𝔤k`$ is of dimension one and has a canonical generator $`(|)`$ (namely the form which takes the value $`2`$ on the small coroots). We choose an orthonormal basis $`\{X_\kappa \}_\kappa `$ of $`𝔤`$ relative to this form. The dual of this line is the space $`𝔤`$-invariants in $`\mathrm{Sym}^2𝔤`$; we shall denote this line by $`𝔠`$. The form $`(|)`$ singles out a generator of $`𝔠`$, namely the Casimir element $`c=_\kappa X_\kappa X_\kappa `$. It is well-known and easy to prove that $`𝔠`$ maps to the center of $`U𝔤`$. In particular, $`c`$ acts in any finite dimensional irreducible representation of $`𝔤`$ by a scalar. In the case of the adjoint representation we denote half this scalar by $`\stackrel{ˇ}{h}`$ (for this happens to be the *dual Coxeter number* of $`𝔤`$) so that
$$\underset{\kappa }{}[X_\kappa [X_\kappa ,Y]]=2\stackrel{ˇ}{h}Y\text{for all }Y𝔤\text{.}$$
Let $`L𝔤`$ stand for $`𝔤_kL`$, but considered as a filtered $`R`$-Lie algebra (so we restrict the scalars to $`R`$): $`F^NL𝔤=𝔤_k𝔪^N`$. An argument similar as for $`r`$ shows that the pairing
$$r_𝔤:(𝔤_kL)(𝔤_k\omega )R,(Xf,Y\kappa )(X|Y)\mathrm{Res}(f\kappa )$$
is topologically perfect; the basis dual to $`(X_\kappa t^l)_{\kappa ,l}`$ is $`(X_\kappa t^{l1}dt)_{\kappa ,l}`$.
For an integer $`N0`$, the quotient $`U_+L𝔤/UL𝔤F^NL𝔤`$ is a free $`R`$-module with generators $`X_{\kappa _1}t^{k_1}\mathrm{}X_{\kappa _r}t^{k_r}`$, $`k_1\mathrm{}k_r<N`$. We complete $`UL𝔤`$ $`𝔪`$-adically on the right:
$$\overline{U}L𝔤:=\underset{}{\mathrm{lim}}_NUL𝔤/UL𝔤F^NL𝔤.$$
A central extension $`\widehat{L𝔤}`$ of $`L𝔤`$ by $`𝔠`$ is defined by endowing $`L𝔤𝔠`$ with the Lie bracket
$$[Xf+cr,Yg+cs]:=[X,Y]fg+c\mathrm{Res}(gdf)(X|Y).$$
We filter $`\widehat{L𝔤}`$ by letting for $`N>0`$, $`F^N\widehat{L𝔤}=F^NL𝔤`$ and for $`N0`$, $`F^N\widehat{L𝔤}=F^NL𝔤+𝔠`$. Then $`U\widehat{L𝔤}`$ is a filtered $`R[c]`$-algebra whose reduction modulo $`c`$ is $`UL𝔤`$. Since the residue is zero on $`𝒪`$, the inclusion of $`F^0L𝔤`$ in $`\widehat{L𝔤}`$ is a homomorphism of Lie algebras. The $`𝔪`$-adic completion on the right
$$\overline{U}\widehat{L𝔤}:=\underset{}{\mathrm{lim}}_NU\widehat{L𝔤}/(U\widehat{L𝔤}F^NL𝔤)$$
is still a $`R[c]`$-algebra and the obvious surjection $`\overline{U}\widehat{L𝔤}\overline{U}L𝔤`$ is the reduction modulo $`c`$. These (completed) enveloping algebras inherit a decreasing filtration from $`L`$.
### Sugawara representation
If $`c`$ is regarded as an element of $`𝔤_k𝔤`$, then tensoring with it defines the $`R`$-linear map
$$𝔩_R𝔩L𝔤_RL𝔤,fg\underset{\kappa }{}X_\kappa fX_\kappa g,$$
which, composed with $`L𝔤_RL𝔤\widehat{L𝔤}_R\widehat{L𝔤}U\widehat{L𝔤}`$, yields $`\gamma :𝔩_R𝔩U\widehat{L𝔤}`$. Since $`\gamma (fggf)=_\kappa [X_\kappa f,X_\kappa g]=cdim𝔤\mathrm{Res}(gdf)`$, $`\gamma `$ drops and extends naturally to an $`R`$-module homomorphism $`\widehat{\gamma }:𝔩_2U\widehat{L𝔤}`$ which sends $`\mathrm{}`$ to $`cdim𝔤`$. It extends continuously to a map from the closure $`\overline{𝔩}_2`$ of $`𝔩_2`$ in $`\overline{U}\widehat{𝔩}`$ to $`\overline{U}\widehat{L𝔤}`$. As $`\overline{𝔩}_2`$ contains the image of $`C:\widehat{\theta }\overline{U}\widehat{𝔩}`$, we get a natural $`R`$-homomorphism
$$\widehat{C}_𝔤:=\widehat{\gamma }\widehat{C}:\widehat{\theta }\overline{U}\widehat{L𝔤}.$$
We may also describe $`\widehat{C}_𝔤`$ in the spirit of Section 1: given $`D\theta `$, then the $`R`$-linear map
$$1D:𝔤_k\omega 𝔤_kL$$
is continuous and selfadjoint relative to $`r_𝔤`$ and the perfect pairing $`r_𝔤`$ allows us to identify it with an element of $`\overline{U}L𝔤`$; this element is our $`\widehat{C}_𝔤(D)`$. The choice of $`t`$ yields the lift
$$\widehat{C}_𝔤(\widehat{D})=\frac{1}{2}\underset{\kappa }{}\underset{l}{}:X_\kappa D,t^{l1}dtX_\kappa t^l:\overline{U}\widehat{L𝔤}$$
(so that in particular, $`\widehat{C}_𝔤(\widehat{D}_k)=\frac{1}{2}_{\kappa ,l}:X_\kappa t^{kl}X_\kappa t^l:`$). This formula can be used to define $`\widehat{C}_𝔤`$, but this approach does not exhibit its naturality.
###### Lemma 2.1.
We have
1. $`[C_𝔤(\widehat{D}),Xf]=(c+\stackrel{ˇ}{h})XD(f)`$, for every $`D\theta `$, $`X𝔤`$ and $`fL`$ and
2. $`[C_𝔤(\widehat{D}_k),C_𝔤(\widehat{D}_l)]=(c+\stackrel{ˇ}{h})(kl)C_𝔤(\widehat{D}_{k+l})+c(c+\stackrel{ˇ}{h})\delta _{k+l,0}\frac{1}{12}dim𝔤(k^3k)`$.
For the proof (which is a bit tricky, but not very deep), we refer to Lecture 10 of (our $`C_𝔤(\widehat{D}_k)`$ is their $`T_k`$).
###### Corollary 2.2 (Sugawara representation).
The $`R`$-linear map
$$T_𝔤:=\frac{1}{c+\stackrel{ˇ}{h}}\widehat{C}_𝔤:\widehat{\theta }\overline{U}\widehat{L𝔤}[\frac{1}{c+\stackrel{ˇ}{h}}]$$
which sends the central $`c_0\widehat{\theta }`$ to $`c(c+\stackrel{ˇ}{h})^1dim𝔤`$ is a homomorphism of $`R`$-Lie algebras. Moreover, if $`D\theta `$, then $`\mathrm{ad}_{T_𝔤(\widehat{D})}`$ leaves $`L𝔤`$ invariant (as a subspace of $`\overline{U}\widehat{L𝔤}`$) and acts on that subspace by derivation with respect to $`D`$.
### Fock type representation for $`𝔤`$
Consider the $`U\widehat{L𝔤}`$-module
$$𝔽(L𝔤):=(U\widehat{L𝔤}/UF^1L𝔤)[\frac{1}{c+\stackrel{ˇ}{h}}]$$
If $`v_0𝔽(L𝔤)`$ denotes the image of $`1`$, then as a $`R[\frac{1}{c+\stackrel{ˇ}{h}}]`$-module it has for basis is the collection $`\{X_{\kappa _r}t^{k_r}\mathrm{}X_{\kappa _1}t^{k_1}(v_0):r0,0k_1k_2\mathrm{}k_r\}`$. It is easy to see that we can identify $`𝔽(L𝔤)`$ with $`\overline{U}\widehat{L𝔤}/\overline{U}F^1L𝔤[\frac{1}{c+\stackrel{ˇ}{h}}]`$, so that $`\widehat{\theta }`$ acts on $`𝔽(L𝔤)`$. It follows from Corollary 2.2 that when $`D\theta `$, then
$$\begin{array}{c}T_𝔤(\widehat{D})X_{\kappa _r}t^{k_r}\mathrm{}X_{\kappa _1}t^{k_1}(v_0)=\hfill \\ \hfill =\left(\underset{i=1}{\overset{r}{}}X_{\kappa _r}t^{k_r}\mathrm{}X_{\kappa _i}D(t^{k_i})\mathrm{}X_{\kappa _1}t^{k_1}\right)v_0+\\ \hfill +X_{\kappa _r}t^{k_r}\mathrm{}X_{\kappa _1}t^{k_1}T_𝔤(\widehat{D})v_0.\end{array}$$
Thus $`\widehat{\theta }`$ is faithfully represented as a Lie algebra of $`R[\frac{1}{c+\stackrel{ˇ}{h}}]`$-linear endomorphisms of $`𝔽(L𝔤)`$. If $`DF^0\theta `$, then clearly $`T_𝔤(\widehat{D})v_0=0`$ and hence:
###### Lemma 2.3.
The Lie subalgebra $`F^0\theta `$ of $`\widehat{\theta }`$ acts on $`𝔽(L𝔤)`$ by coefficientwise derivation.
This lemma has an interesting corollary. In order to state it, consider the module of $`k`$-derivations $`RR`$ (denoted here simply by $`\theta _R`$ instead of the more correct $`\theta _{R/k}`$) and the module of $`k`$-derivations of $`L`$ which preserve $`RL`$ (denoted by $`\theta _{L,R}`$). Since $`LR((t))`$ as an $`R`$-algebra, every $`k`$-derivation $`RR`$ extends to one from $`L`$ to $`L`$. So $`\theta _{L,R}/\theta `$ can be identified with $`\theta _R`$.
###### Corollary 2.4.
The central extension $`\widehat{\theta }`$ of $`\theta `$ by $`Rc_0`$ naturally extends to a central extension of $`R`$-Lie algebras $`\widehat{\theta }_{L,R}`$ of $`\theta _{L,R}`$ by $`Rc_0`$ (so with $`\widehat{\theta }_{L,R}/\widehat{\theta }=\theta _{L,R}/\theta \theta _R`$) and the Sugawara representation $`T_𝔤`$ of $`\widehat{\theta }`$ on $`𝔽(L𝔤)`$ extends to $`\widehat{\theta }_{L,R}`$ in such a manner that it preserves any $`U\widehat{L𝔤}`$-submodule of $`𝔽(L𝔤)`$ and for every $`D\theta _{L,R}`$, the following relations hold in $`\mathrm{End}(𝔽(L𝔤))`$:
1. $`[T_𝔤(\widehat{D}),Xf]=X(Df)`$ for $`X𝔤`$, $`fL`$, and
2. (Leibniz rule) $`[T_𝔤(\widehat{D}),r]=Dr`$ for $`rR`$.
These assertions also hold for the oscillator representation $`T`$ on $`𝔽`$.
###### Proof.
The generator $`t`$ can be used to define a section of $`\theta _{L,R}\theta _R`$: the set of elements of $`\theta _{L,R}`$ which kill $`t`$ is a $`k`$-Lie subalgebra of $`\theta _{L,R}`$ which projects isomorphically onto $`\theta _R`$. Now if $`D\theta _{L,R}`$, write $`D=D_{\mathrm{vert}}+D_{\mathrm{hor}}`$ with $`D_{\mathrm{vert}}\theta `$ and $`D_{\mathrm{hor}}(t)=0`$ and define a $`k[\frac{1}{c+\stackrel{ˇ}{h}}]`$-linear operator $`\widehat{D}`$ in $`𝔽(L𝔤)`$ as the sum of $`T_𝔤(\widehat{D}_{\mathrm{vert}})`$ and coefficientwise derivation by $`D_{\mathrm{hor}}`$. This map clearly satisfies the two properties and preserves any $`U\widehat{L𝔤}`$-submodule of $`𝔽(L𝔤)`$. As to its dependence on $`t`$: another choice yields a decomposition of the form $`D=(D_{\mathrm{hor}}+D_0)+(D_{\mathrm{vert}}D_0)`$ with $`D_0F^0\theta `$ and in view of Lemma 2.3 $`\widehat{D}_0`$ acts in $`𝔽(L𝔤)`$ by coefficientwise derivation.
The same argument works for the oscillator representation. ∎
### Semilocal case
This refers to the situation where we allow the $`R`$-algebra $`L`$ to be a finite direct sum of $`R`$-algebras isomorphic to $`R((t))`$: $`L=_{iI}L_i`$, where $`I`$ is a nonempty finite index set and $`L_i`$ as before. We use the obvious convention, for instance, $`𝒪`$, $`𝔪`$, $`\omega `$, $`𝔩`$ are now the direct sums over $`I`$ (as filtered objects) of the items suggested by the notation. If $`r:L\times \omega R`$ denotes the sum of the residue pairings of the summands, then $`r`$ is still topologically perfect. In this setting, the oscillator algebra $`\widehat{𝔩}`$ is of course not the direct sum of the $`\widehat{𝔩}_i`$, but rather the quotient of $`_i\widehat{𝔩}_i`$ that identifies the central generators $`c_{0,i}`$ of the summands with a single $`c_0`$. We thus get a Virasoro extension $`\widehat{\theta }`$ of $`\theta `$ by $`c_0R`$ and a (faithful) oscillator representation of $`\widehat{\theta }`$ in $`\overline{U}\widehat{𝔩}`$. The decreasing filtrations are the obvious ones. In likewise manner we define $`\widehat{L𝔤}`$ (a central extension of $`_{iI}L𝔤_i`$) and construct $`\widehat{\theta }_{L,R}`$ and the associated Sugawara representation: Corollaries 2.2 and 2.4 continue to hold.
## 3. The WZW connection
From now on we assume that our base field $`k`$ is algebraically closed and $`R`$ is a noetherian $`k`$-algebra. We place ourselves in the semi-local case. In what follows, we often allow $`𝔤`$ to be replaced by $`k`$, viewed as an abelian Lie algebra, where the substitutions are the obvious ones; for instance, $`\widehat{L𝔤}`$, $`T_𝔤`$ and $`𝔽(L𝔤)`$ become $`\widehat{𝔩}`$, $`T`$ and $`𝔽`$.
### Abstract conformal blocks
Let $`A`$ be a $`R`$-subalgebra of $`L`$ and let $`\theta _{A/R}`$ have the usual meaning as the Lie algebra of $`R`$-derivations $`AA`$. We assume that:
1. as a $`R`$-algebra, $`A`$ is flat and of finite type,
2. $`A𝒪=R`$ and the $`R`$-module $`L/(A+𝒪)`$ is locally free of finite rank,
3. the annihilator $`\mathrm{ann}(A)`$ of $`A`$ in $`\omega `$ contains the image of $`dA`$ and the resulting map $`\mathrm{Hom}_A(\mathrm{ann}(A),A)\theta _{A/R}`$ is an isomorphism.
We denote $`\theta _{A,R}`$ the Lie algebra of $`k`$-derivations $`AA`$ which preserve $`R`$. So the quotient $`\theta _R^A:=\theta _{A,R}/\theta _{A/R}`$ consists of the $`k`$-derivations $`RR`$ that extend to one of $`A`$ and hence is a submodule of $`\theta _R`$. It is clear that $`\theta _{A,R}\theta _{L,R}`$. We denote by $`\widehat{\theta }_{A,R}`$ the preimage of $`\theta _{A,R}`$ in $`\widehat{\theta }_{L,R}`$ and by $`\widehat{\theta }_R^A`$ the quotient $`\widehat{\theta }_{A,R}/\theta _{A/R}`$. These are central extensions of $`\theta _{A,R}`$ resp. $`\theta _R^A`$ by $`𝔠_kR=cR`$.
We put
$$A𝔤:=𝔤_kA$$
and view it as a Lie subalgebra of $`L𝔤`$. Since the residue vanishes on $`\mathrm{ann}(A)`$, the inclusion $`A𝔤\widehat{L𝔤}`$ is a Lie algebra homomorphism. We define the *universal covacuum space* as the space of $`A𝔤`$-covariants in $`𝔽(L𝔤)`$, $`𝔽(L𝔤)_{A𝔤}:=𝔽(L𝔤)/A𝔤𝔽(L𝔤)`$.
###### Proposition 3.1.
For $`D\theta _{A/R}`$, $`T_𝔤(\widehat{D})`$ lies in the closure of $`A𝔤\widehat{L𝔤}`$ in $`\overline{U}\widehat{L𝔤}`$.
The Sugawara representation of the Lie algebra $`\widehat{\theta }_{A,R}`$ on $`𝔽(L𝔤)`$ preserves the space of $`A𝔤`$-covariants in $`𝔽(L𝔤)`$, $`𝔽(L𝔤)_{A𝔤}:=𝔽(L𝔤)/A𝔤𝔽(L𝔤)`$, and acts on $`𝔽(L𝔤)_{A𝔤}`$ via the central extension $`\widehat{\theta }_R^A`$ of $`\theta _R^A`$ (with the central $`c_0`$ acting as multiplication by $`(c+\stackrel{ˇ}{h})^1cdim𝔤`$).
These assertions also hold if we replace $`𝔤`$ by the abelian Lie algebra $`k`$ (in particular, $`T(\widehat{D})`$ lies in the closure of $`𝔞\widehat{𝔩}`$ in $`\overline{U}\widehat{𝔩}`$).
###### Proof.
We only do this for $`𝔤`$. Since $`D`$ maps $`\mathrm{ann}(A)`$ to $`AL`$, $`1D`$ maps the submodule $`𝔤\mathrm{ann}(A)`$ of $`𝔤\omega `$ to the submodule $`𝔤A=A𝔤`$ of $`𝔤L=L𝔤`$. It is clear that $`𝔤\mathrm{ann}(A)`$ and $`A𝔤`$ are each others annihilator relative to the pairing $`r_𝔤`$. This implies that $`\widehat{C}(\widehat{D})`$ lies in the closure of the image of $`A𝔤_kL𝔤+L𝔤_kA𝔤`$ in $`\overline{U}\widehat{L𝔤}`$. It follows that $`\widehat{C}(\widehat{D})`$ has the form $`cr+_\kappa _{n1}X_\kappa f_{\kappa ,n}X_\kappa g_{\kappa ,n}`$ with $`rR`$, one of $`f_{\kappa _n},g_{\kappa ,n}L`$ being in $`A`$ and the order of $`f_{\kappa _n}`$ smaller than that of $`g_{\kappa ,n}`$ for almost all $`\kappa ,n`$. Since the elements of $`A`$ have order $`0`$ and $`X_\kappa f_{\kappa ,n}X_\kappa g_{\kappa ,n}X_\kappa g_{\kappa ,n}X_\kappa g_{\kappa ,n}(modcR)`$, we can assume that all $`f_{\kappa ,n}`$ lie in $`A`$ and so the first assertion follows.
If $`D\theta _{A,R}`$, then for $`X𝔤`$ and $`fA`$, we have $`[D,Xf]=X(Df)`$, which is an element of $`A𝔤`$ (since $`DfA`$). This shows that $`T_𝔤(\widehat{D})`$ preserves $`A𝔤𝔽(L𝔤)`$. If $`D\theta _{A/R}`$, then it follows from the proven part that $`T_𝔤(\widehat{D})`$ maps $`𝔽(L𝔤)`$ to $`A𝔤𝔽(L𝔤)`$ and hence induces the zero map in $`𝔽(L𝔤)_{A𝔤}`$. So $`\widehat{\theta }_{A,R}`$ acts on $`𝔽(L𝔤)_{A𝔤}`$ via $`\widehat{\theta }_R^A`$. ∎
More relevant than $`𝔽(L𝔤)_{A𝔤}`$ will be certain finite dimensional quotients thereof obtained as follows. Let $`\mathrm{}`$ be a positive integer and let $`V:iV_i`$ be a map which assigns to every $`iI`$ a finite dimensional irreducible representation $`V_i`$ of $`𝔤`$. Make $`_{iI}V_i`$ a $`k`$-representation of $`F^0L𝔤`$ by letting $`c`$ act as multiplication by $`\mathrm{}`$ and $`𝔤_k𝒪_i`$ on the $`i`$th factor via its projection onto $`𝔤`$. If we induce this up to $`\widehat{L𝔤}`$ we get a representation $`\stackrel{~}{}_{\mathrm{}}(V)`$ of $`\widehat{L𝔤}`$ which clearly is a quotient of $`𝔽(L𝔤)`$. Its irreducible quotient is denoted by $`_{\mathrm{}}(V)`$. The following is known (see the book of Kac ). First of all, $`_{\mathrm{}}(V)`$ is nonzero precisely when each $`𝔤`$-representation $`V_i`$ *is of level $`\mathrm{}`$*, i.e., has the property that for every nilpotent $`X𝔤`$, $`X^{\mathrm{}+1}`$ yields the zero map in $`V_i`$. Assume this is the case. Then $`_{\mathrm{}}(V)`$ is integrable as a $`\widehat{L𝔤}`$-module: if $`Y𝔤`$ is nilpotent and $`fL`$, then $`Yf`$ acts locally nilpotently in $`_{\mathrm{}}(V)`$.
The set of isomorphism classes of finite dimensional irreducible $`𝔤`$-representations of level $`\mathrm{}`$ is finite. It is clear that this set, which we denote by $`P_{\mathrm{}}`$, is invariant under dualization (and more generally, under all outer automorphisms of $`𝔤`$).
Adopting the physicists terminology, we might call the $`R`$-module of $`R`$-linear forms on $`_{\mathrm{}}(V)_{A𝔤}`$ the *conformal block* attached to $`A`$. The following proposition says that it is of finite rank and describes in essence the WZW-connection.
###### Proposition 3.2 (Finiteness).
The space $`_{\mathrm{}}(V)`$ is finitely generated as a $`UA𝔤`$-module (so that $`_{\mathrm{}}(V)_{A𝔤}`$ is a finitely generated $`R`$-module). The Lie algebra $`\widehat{\theta }_R^A`$ acts on $`_{\mathrm{}}(V)_{A𝔤}`$ via the Sugawara representation with the central $`c_0`$ acting as multiplication by $`\frac{c}{c+\stackrel{ˇ}{h}}dim𝔤`$.
###### Proof.
Pick a Cartan subalgebra $`𝔥𝔤`$, let $`\mathrm{\Delta }𝔥^{}`$ be its root system and $`𝔤=𝔥_{\alpha \mathrm{\Delta }}𝔤_\alpha `$ the associated decomposition. Recall that $`𝔤_\alpha `$ has a nilpotent generator and $`[𝔤_\alpha ,𝔤_\beta ]`$ is contained in $`𝔤_{\alpha +\beta }`$ when $`\alpha +\beta \mathrm{\Delta }`$ and is contained in $`𝔥`$ otherwise. Recall also that $`𝔤`$ is spanned by $`_{\alpha \mathrm{\Delta }}𝔤_\alpha `$ and that $`𝔥`$ normalizes each $`𝔤_\alpha `$. Choose a generator $`t_i`$ of $`𝔪_i`$. Using a simple (PBW-type) argument we see that $`_{\mathrm{}}(V)`$ is spanned over $`R`$ by the subspaces
$$(𝔤_{\alpha _k}q_k)^{r_k}\mathrm{}(𝔤_{\alpha _1}q_1)^{r_1}V$$
with $`k0`$, each $`q_\rho `$ a negative power of some $`t_i`$ and such that $`r_1\mathrm{}r_k`$. Choose a finite set $`Q`$ of powers of $`t_i`$’s that spans an $`R`$-supplement of $`𝒪+A`$ in $`L`$ and let $`M`$ be the $`R`$-span of the above subspaces for which all $`q_\rho `$ lie in $`Q`$. It is clear that then $`M`$ generates $`_{\mathrm{}}(V)`$ as a $`UA𝔤`$-module. Since each line $`𝔤_\alpha q`$ acts locally nilpotently in $`_{\mathrm{}}(V)`$, there exists a positive integer $`N`$ such that $`(𝔤_\alpha q)^NV=0`$ for all $`\alpha \mathrm{\Delta }`$, $`qQ`$. So the only generating lines that contribute to $`M`$ have $`r_1<N`$ and are hence finite in number. This shows that $`M`$ is a finitely generated $`R`$-module. The rest follows from 3.1. ∎
###### Remark 3.3.
We expect the $`R`$-module $`_{\mathrm{}}(V)_{A𝔤}`$ is also flat and that this is a consequence from a related property for the $`UA𝔤`$-module $`_{\mathrm{}}(V)`$). Such a result, or rather an algebraic proof of it, might simplify the argument in (see Section 4 for our version) which shows that the sheaf of conformal blocks attached to a degenerating family of pointed curves is locally free.
### Propagation principle
The following proposition is a bare version of what is known as *propagation of vacua*, or rather the generalization of this fact that can be found in Beauville . It often allows us to reduce the discussion to the case where $`I`$ is a singleton.
###### Proposition 3.4.
Let $`JI`$ be such that $`A`$ maps onto $`_{jJ}L_j/𝒪_j`$. Denote by $`BA`$ the kernel, so that we have a surjective Lie homomorphism $`B𝔤𝔤^J`$ via which $`B𝔤`$ acts on $`_{jJ}V_j`$. Then the map of $`B𝔤`$-modules $`_{\mathrm{}}(V|IJ)_{jJ}V_j_{\mathrm{}}(V)`$ induces an isomorphism on covariants:
$$\begin{array}{ccc}\left(_{\mathrm{}}(V|IJ)_{jJ}V_j\right)_{B𝔤}& \stackrel{}{}& _{\mathrm{}}(V)_{A𝔤}.\end{array}$$
###### Proof.
It suffices to do the case when $`J`$ is a singleton $`\{o\}`$. The hypotheses clearly imply that $`_{\mathrm{}}(V|I\{o\})V_o_{\mathrm{}}(V)_{A𝔤}`$ is onto. The kernel is easily shown to be $`B𝔤(_{\mathrm{}}(V|I\{o\})V_o)`$. ∎
###### Remark 3.5.
This proposition is sometimes used in the opposite direction: if $`𝔪_oA`$ is a principal ideal with the property that for a generator $`t𝔪_o`$, the $`𝔪_o`$-adic completion of $`A`$ gets identified with $`R((t))`$, then let $`\stackrel{~}{I}`$ be the disjoint union of $`I`$ and $`\{o\}`$, $`\stackrel{~}{V}`$ the extension of $`V`$ to $`\stackrel{~}{I}`$ which assigns to $`o`$ the trivial representation and $`\stackrel{~}{A}:=A[t^1]`$. With $`(\stackrel{~}{I},\{o\})`$ taking the role of $`(I,J)`$, we then find that $`_{\mathrm{}}(V)_{A𝔤}_{\mathrm{}}(\stackrel{~}{V})_{\stackrel{~}{A}𝔤}`$.
### Conformal blocks in families
We transcribe the preceding in more geometric terms. This leads us to sheafify many of the notions we introduced earlier; we shall modify our notation (or its meaning) accordingly. Suppose given a proper and flat morphism between $`k`$-varieties $`\pi :𝒞S`$ whose base $`S`$ is smooth and fibers are reduced connected curves which have complete intersection singularities only (we are here not assuming that $`𝒞`$ is smooth over $`k`$). Since the family is flat, the arithmetic genus of the fibers is locally constant; we simply assume it is globally so and denote this constant genus by $`g`$. We also suppose given disjoint sections $`S_i𝒞`$, indexed by the finite nonempty set $`I`$ whose union $`_{iI}S_i`$ lies in the smooth part of $`𝒞`$ and meets every irreducible component of a fiber. The last condition ensures that if $`j:𝒞^{}:=𝒞_{iI}S_i𝒞`$ is the inclusion, then $`\pi j`$ is an affine morphism.
We denote by $`(𝒪_i,𝔪_i)`$ the formal completion of $`𝒪_𝒞`$ along $`S_i`$, by $`_i`$ the subsheaf of fractions of $`𝒪_i`$ with denominator a local generator of $`𝔪_i`$ and by $`𝒪`$, $`𝔪`$ and $``$ the corresponding direct sums. But we keep on using $`\omega `$, $`\theta `$, $`\widehat{\theta }`$ etc. for their sheafified counterparts. So these are now all $`𝒪_S`$-modules and the residue pairing is also one of $`𝒪_S`$-modules: $`r:\times \omega 𝒪_S`$. We write $`𝒜`$ for $`\pi _{}j_{}j^{}𝒪_𝒞`$ (a sheaf of $`𝒪_S`$-algebras) and often identify this with its image in $``$. We denote by $`\theta _{𝒜/S}`$ the sheaf of $`𝒪_S`$-derivations $`𝒜𝒜`$ and by $`\omega _{𝒜/S}`$ for the direct image on $`S`$ of the the relative dualizing sheaf of $`𝒞^{}/S`$, (in the present situation the relative dualizing sheaf of $`\pi `$, $`\omega _{𝒞/S}`$, is simply the direct image on $`𝒞`$ of the sheaf of relative differentials on the open subset of $`𝒞`$ where $`\pi `$ is smooth). So $`\omega _{𝒜/S}`$ is torsion free an hence embeds in $`\omega `$.
###### Lemma 3.6.
The properties $`A_1`$, $`A_2`$ and $`A_3`$ hold for the sheaf $`𝒜`$. Precisely,
1. $`𝒜`$ is as a sheaf of $`𝒪_S`$-algebras flat and of finite type,
2. $`𝒜𝒪=𝒪_S`$ and $`R^1\pi _{}𝒪_𝒞=/(𝒜+𝒪)`$ is locally free of rank $`g`$,
3. we have $`\theta _{𝒜/S}=\mathrm{Hom}_𝒜(\omega _{𝒜/S},𝒜)`$ and $`\omega _{𝒜/S}`$ is the annihilator of $`𝒜`$ with respect to the residue pairing.
###### Proof.
Property $`𝒜_1`$ is clear. It is also clear that $`𝒪_S=\pi _{}𝒪_𝒞𝒜𝒪`$ is an isomorphism. The long exact sequence defined by the functor $`\pi _{}`$ applied to the the short exact sequence
$$0𝒪_𝒞j_{}j^{}𝒪_𝒞/𝒪0$$
tells us that $`R^1\pi _{}𝒪_𝒞=/(𝒜+𝒪)`$; in particular, the latter is locally free of rank $`g`$. Hence $`𝒜_2`$ holds as well.
In order to verify $`𝒜_3`$, we note that $`\pi _{}\omega _{𝒞/S}`$ is the $`𝒪_S`$-dual of $`R^1\pi _{}𝒪_S`$, and hence is locally free of rank $`g`$. The first part of $`𝒜_3`$ from the correponding local property $`\theta _{𝒞/S}=\mathrm{Hom}_{𝒪_𝒞}(\omega _{𝒞/S},𝒪_𝒞)`$ by apply $`\pi _{}j^{}`$ to either side. This local property is known to hold for families of curves with complete intersection singularities. (A proof under the assumption that $`𝒞`$ is smooth—which is does not affect the generality, since $`\pi `$ is locally the restriction of that case and both sides are compatible with base change—runs as follows: if $`j^{}:𝒞^{}𝒞`$ denotes the locus where $`\pi `$ is smooth, then its complement is of codimension $`2`$ everywhere. Clearly, $`\theta _{𝒞/S}`$ is the $`𝒪_𝒞`$-dual of $`\omega _{𝒞/S}`$ on $`𝒞^{}`$ and since both are inert under $`j_{}^{}j^{}^{}`$, they are equal everywhere.)
The last assertion essentially restates the well-known fact that the polar part of a rational section of $`\omega _{𝒞/S}`$ must have zero residue sum, but can otherwise be arbitrary. More precisely, the image of $`\omega _{𝒜/S}`$ in $`\omega /F^1\omega `$ is the kernel of the residue map $`\omega /F^1\omega 𝒪_S`$. The intersection $`\omega _{𝒜/S}F^1\omega `$ is $`\pi _{}\omega _{𝒞/S}`$ and is hence locally free of rank $`g`$. Since $`\mathrm{ann}(F^1\omega )=𝒪`$, it follows that $`\mathrm{ann}(\omega _{𝒜/S})𝒪`$ and $`/(\mathrm{ann}(\omega _{𝒜/S})+𝒪)`$ are locally free of rank 1 and $`g`$ respectively. Since $`𝒜`$ has these properties also and is contained in $`\mathrm{ann}(\omega _{𝒜/S})`$, we must have $`𝒜=\mathrm{ann}(\omega _{𝒜/S})`$. ∎
For what follows one usually supposes that the fibers are stable $`I`$-pointed curves (meaning that every fiber of $`\pi j`$ has only ordinary double points as singularities and has finite automorphism group) and is versal (so that in the discriminant $`D_\pi `$ of $`\pi `$ is a reduced normal crossing divisor), but we shall not make these assumptions yet. Instead, we assume the considerable weaker property that the sections of the sheaf $`\theta _S(\mathrm{log}D_\pi )`$ of vector fields on $`S`$ tangent to $`D_\pi `$ lift locally on $`S`$ to vector fields on $`𝒞`$. This is for instance the case if $`𝒞`$ is smooth and $`\pi `$ is multi-transversal with respect to the (Thom) stratification of $`\mathrm{Hom}(T𝒞,\pi ^{}TS)`$ by rank . Notice that the restriction homomorphism $`\theta _S(\mathrm{log}D_\pi )𝒪_{D_\pi }\theta _{D_\pi }`$ is an isomorphism.
Let $`\theta _{𝒞,S}\theta _𝒞`$ denote the sheaf of derivations which preserve $`\pi ^{}𝒪_S`$. Applying $`\pi _{}j_{}j^{}`$ to the exact sequence $`0\theta _{𝒞/S}\theta _{𝒞,S}\theta _{𝒞,S}/\theta _{𝒞/S}0`$ yields the exact sequence
$$0\theta _𝒜\theta _{𝒜,S}\theta _S(\mathrm{log}D_\pi )0.$$
(using our litfability assumption and the fact that $`\pi j`$ is affine).
We have defined $`\widehat{\theta }_{𝒜,S}`$ as the preimage of $`\theta _{𝒜,S}`$ in $`\widehat{\theta }_{,S}`$ and $`\widehat{\theta }_S(\mathrm{log}D_\pi )`$ as the quotient $`\widehat{\theta }_{,S}/\theta _𝒜`$. These centrally extend $`\theta _{𝒜,S}`$ and $`\theta _S`$ by $`c_0𝒪_S`$.
Observe that $`𝔤=𝔤_{𝒪_S}`$ is now a sheaf of Lie algebras over $`𝒪_S`$. The same applies to $`\widehat{𝔩}`$ and so we have a Virasoro extension $`\widehat{\theta }_S`$ of $`\theta _S`$ by $`c_0𝒪_S`$. We have also defined $`𝒜𝔤=𝔤_k𝒜`$, which is a Lie subsheaf of $`𝔤`$ as well as of $`\widehat{𝔤}`$ and the Fock type $`\widehat{𝔤}`$-module $`(\widehat{𝔤})`$. The *universal covacuum module* is the sheaf of $`𝒜𝔤`$-covariants in the latter,
$$(\widehat{𝔤})_𝒞:=(\widehat{𝔤})_{𝒜𝔤}:=(\widehat{𝔤})/𝒜𝔤(\widehat{𝔤}).$$
From Proposition 3.1 we get:
###### Corollary 3.7.
The representation of the Lie algebra $`\widehat{\theta }_{𝒜,S}`$ on $`(\widehat{𝔤})`$ preserves $`𝒜𝔤(\widehat{𝔤})`$ and acts on $`(\widehat{𝔤})_𝒞`$ via the central extension $`\widehat{\theta }_S(\mathrm{log}D_\pi )`$ of $`\theta _S(\mathrm{log}D_\pi )`$ (with the central $`c_0`$ acting as multiplication by $`(c+\stackrel{ˇ}{h})^1cdim𝔤`$). This construction has a base change property along the smooth part $`D_\pi ^{}`$ of the discriminant: we may identify $`(\widehat{𝔤})_𝒞𝒪_{D_\pi ^{}}`$ with $`(\widehat{|_{D_\pi ^{}}𝔤})_{𝒞|_{D_\pi ^{}}}`$ and the action of $`\widehat{\theta }_S(\mathrm{log}D_\pi )𝒪_{D_\pi ^{}}`$ on it factors through $`\widehat{\theta }_{D_\pi ^{}}`$.
The bundle of integrable representations $`_{\mathrm{}}(V)`$ over $`S`$ is defined in the expected manner: it is obtained as a quotient of $`(\widehat{𝔤})`$ in the way $`_{\mathrm{}}(V)`$ is obtained from $`𝔽(L𝔤)`$. According to Corollary 3.7, $`_{\mathrm{}}(V)`$ comes with a $`k`$-linear representation of the Lie algebra $`\widehat{\theta }_S(\mathrm{log}D_\pi )`$ that satisfies the Leibniz rule and on which $`c_0`$ acts as multiplication by $`\mathrm{}(\mathrm{}+\stackrel{ˇ}{h})^1dim𝔤`$. We write $`_{\mathrm{}}(V)_𝒞`$ for $`_{\mathrm{}}(V)_{𝒜𝔤}`$.
###### Corollary 3.8 (WZW-connection).
The $`𝒪_S`$-module $`_{\mathrm{}}(V)_𝒞`$ is of finite rank; It is also locally free over $`SD_\pi `$ and the Lie action of $`\widehat{\theta }_S(\mathrm{log}D_\pi )`$ on $`_{\mathrm{}}(V)_𝒞`$ defines a flat connection on the associated bundle of projective spaces, $`_S(_{\mathrm{}}(V)_𝒞)S`$ with a logarithmic singularity along $`D_\pi `$. The same base change property holds along the smooth part $`D_\pi ^{}`$ of the discriminant as in Corollary 3.7.
###### Proof.
The first assertion follows from 3.2. The rest is clear except perhaps the local freeness of $`_{\mathrm{}}(V)_𝒞`$ on $`SD_\pi `$. But this follows from the local existence of a connection in the $`𝒪_S`$-module $`_{\mathrm{}}(V)_𝒞`$. ∎
###### Remark 3.9.
In case the base field $`k`$ is $``$ we could equally well work in the complex-analytic category. The preceding corollary then leads to an interesting topological functor: if $`\mathrm{\Sigma }`$ is a compact smooth oriented surface, $`I\mathrm{\Sigma }`$ a finite nonempty subset such that for every $`iI`$ is given a $`𝔤`$-representation $`V_i`$, then there is canonically associated to these data a projective space $`(\mathrm{\Sigma },I,V)`$: for every complex structure on $`\mathrm{\Sigma }`$ compatible with the given smooth structure and orientation we have defined a conformal block as above. The complex structures with these property form a simply connected (even contractible) space $`𝒞(\mathrm{\Sigma })`$. To see this: such a structure amounts to giving a conformal structure on $`\mathrm{\Sigma }`$, that is, a Riemann metric up to multiplication by the exponential of a real valued function. The space of Riemann metrics is a convex cone in some linear space (hence contractible) and the linear space of real valued functions on $`\mathrm{\Sigma }`$ acts sufficiently nice on it to ensure that the orbit space remains simply connected. So any two such complex structures can be joined a path whose homotopy class is unique. Corollary 3.8 then implies that the WZW-connection enables us to identify their associated conformal blocks in a unique manner. Perhaps the best way to describe $`(\mathrm{\Sigma },I,V)`$ is as the space of flat sections of the flat projective space bundle of conformal blocks over $`𝒞(\mathrm{\Sigma })`$. The naturality implies that the group of orientation preserving diffeomorphisms of the pair $`\mathrm{\Sigma }`$ which fix $`I`$ pointwise acts on $`(\mathrm{\Sigma },I,V)`$. Since its identity component is contractible, the action will be through the mapping class group $`\mathrm{\Gamma }(\mathrm{\Sigma },I)`$. It is conjectured that this action leaves invariant a Fubini-Study metric on $`(\mathrm{\Sigma },I,V)`$.
### Propagation principle continued
In the preceding subsection we made the assumption throughout that a union $`\stackrel{~}{S}`$ of sections of $`𝒞S`$ is given to ensure that $`𝒞\stackrel{~}{S}S`$ is affine. However, the propagation principle permits us to abandon that assumption. In fact, this leads us to let $`V`$ stand for any map which assigns to every closed point $`p`$ of $`𝒞`$ an irreducible $`𝔤`$-representation $`V_p`$ of level $`\mathrm{}`$, subject to the condition that its *support*, $`\mathrm{Supp}(V)`$ (i.e., the set of $`p`$ for which $`V_p`$ is not the trivial representation), is a trivial finite cover over $`S`$ and contained in the locus where $`\pi :𝒞S`$ is smooth. Our earlier defined $`_{\mathrm{}}(V)`$ is now better written as $`_{\mathrm{}}(V|_{\mathrm{Supp}(V)})`$. Since $`𝒞\mathrm{Supp}(V)`$ need not be affine over $`S`$, this does not yield the right notion of conformal block. We can find however, at least locally over $`S`$, additional sections $`S_i^{}`$ of $`𝒞\mathrm{Supp}(V)S`$ so that $`𝒞^{}:=𝒞\mathrm{Supp}(V)_iS_i^{}`$ is affine over $`S`$. Then we can form $`_{\mathrm{}}(V|𝒞𝒞^{})`$ and Remark 3.5 shows that the resulting conformal block $`_{\mathrm{}}(V|𝒞𝒞^{})_{𝔤[\pi _{}𝒪_𝒞^{}]}`$ is independent of the choices made. This suggests that we let $`_{\mathrm{}}(V)`$ stand for the sheaf associated to the presheaf
$$SU\underset{}{\mathrm{lim}}_{\stackrel{~}{S}}_{\mathrm{}}(V|_{\stackrel{~}{S}}),$$
where $`\stackrel{~}{S}`$ runs over the closed subvarieties of $`\pi ^1U`$ that trivial finite covers of $`S`$, contain $`\mathrm{Supp}(V)`$ and have affine complement in $`\pi ^1U`$ and perhaps also justifies our custom of writing $`_{\mathrm{}}(V)_𝒞`$ for the associated sheaf of conformal blocks. It is clear that in this set-up there is also no need anymore to insist that the fibers of $`\pi `$ be connected.
### The genus zero case and the KZ-connection
We take $`C^1`$ and let $`\mathrm{}`$ and $`zCV_z`$ be as usual. Choose an affine coordinate $`z`$ on $`C`$ such that $`\mathrm{}\mathrm{Supp}(V)`$ and let $`z_1,\mathrm{},z_N`$ enumerate the distinct points of $`\mathrm{Supp}(V)`$. Write $`V_i`$ for $`V_{z_i}`$. The local field attached to $`\mathrm{}`$ has parameter $`t_{\mathrm{}}=z^1`$. If $`_{\mathrm{}}`$ denotes the representation of $`\widehat{𝔤((z^1))}`$ attached to the trivial representation (with generator $`v_0`$), then by the propagation principle 3.4 we have $`_{\mathrm{}}(V)_C=(_{\mathrm{}}V_1\mathrm{}V_N)_{𝔤[[z]]}`$, where $`𝔤[[z]]`$ acts on $`V_i`$ via its evaluation at $`z_i`$. According to , the $`𝔤[[z]]`$-homomorphism $`U𝔤[[z]]_{\mathrm{}}`$, $`uuv_0`$ is surjective and its kernel is the left ideal generated by the constants $`𝔤`$ and $`(Xz)^{1+\mathrm{}}`$, where $`X𝔤`$ is a regular nilpotent element (as these form a single orbit under the adjoint representation, it does not matter which one we take). This implies that $`_{\mathrm{}}(V)_^1`$ can be identified with a quotient of the space of $`𝔤`$-covariants $`(V_1\mathrm{}V_N)_𝔤`$, namely its biggest quotient on which $`(_{i=1}^Nz_iX_{(i)})^{1+\mathrm{}}`$ acts trivially (where $`X_{(i)}`$ acts on $`V_i`$ as $`X`$ and on the other tensor factors $`V_j`$, $`ji`$, as the identity).
Now regard $`z_1,\mathrm{},z_N`$ as variables. If $`S`$ denotes the open subset of $`(z_1,\mathrm{},z_N)^N`$ of pairwise distinct $`N`$-tuples, then we have in an evident manner a family $`𝒞`$ over $`S`$, $`\pi :C\times SS`$, with $`N+1`$ sections (including the one at infinity) so that we also have defined $`𝒞^{}S`$. This determines a sheaf of conformal blocks $`_{\mathrm{}}(V)_𝒞`$ over $`S`$. According to the preceding, we have an exact sequence
$$\begin{array}{ccc}(V_1\mathrm{}V_N)_𝔤_k𝒪_S& \stackrel{(_iz_iX_{(i)})^{1+\mathrm{}}}{}& (V_1\mathrm{}V_N)_𝔤_k𝒪_S_{\mathrm{}}(V)_𝒞0.\end{array}$$
We identify its WZW connection, or rather, a natural lift of that connection to $`V_1\mathrm{}V_N_k𝒪_S`$. In order to compute the covariant derivative with respect to the vector field $`_i:=\frac{}{z_i}`$ on $`S`$, we follow our receipe and lift it to $`𝒞`$ in the obvious way (with zero component along $`C`$). We continue to denote that lift by $`_i`$ and determine its (Sugawara) action on $`_{\mathrm{}}(V)`$. We first observe that $`_i`$ is tangent to all the sections, except the $`i`$th. Near that section we decompose it as $`(\frac{}{z}+_i)\frac{}{z}`$, where the first term is tangent to the $`i`$-th section and the second term is vertical. The action of the former is easily understood (in terms of an obvious trivialization of $`_{\mathrm{}}(V)`$, it acts as derivation with respect to $`z_i`$). The vertical term, $`\frac{}{z}`$, acts via the Sugawara representation, that is, it acts on the $`i`$th slot as $`\frac{1}{\mathrm{}+\stackrel{ˇ}{h}}_\kappa X_\kappa (zz_i)^1X_\kappa `$ and as the identity on the others. This action does not induce one in $`V_1\mathrm{}V_N_k𝒪_S`$. To make it so, we add to this the action of the element
$$\left(\frac{1}{\mathrm{}+\stackrel{ˇ}{h}}\underset{\kappa }{}X_\kappa (zz_i)^1\right)X_\kappa ^{(i)}𝔤[𝒞^{}]U\widehat{𝔤},$$
for this sum then acts in $`V_1\mathrm{}V_N_k𝒪_S`$ as
$$\frac{1}{\mathrm{}+\stackrel{ˇ}{h}}\underset{ji}{}\frac{1}{z_jz_i}X_\kappa ^{(j)}X_\kappa ^{(i)}.$$
Let us regard the Casimir element $`c`$ as a tensor of $`𝔤_k𝔤`$, and denote by $`c^{(i,j)}`$ its action in $`V_1\mathrm{}V_N`$ on the $`i`$th and $`j`$th factor (since $`c`$ is symmetric, we have $`c^{(i,j)}=c^{(j,i)}`$, so that we need not worry about the order here). We conclude that the WZW-connection is induced by the connection on $`V_1\mathrm{}V_N_k𝒪_S`$ defined by the connection form
$$\frac{1}{\mathrm{}+\stackrel{ˇ}{h}}\underset{i=1}{\overset{N}{}}\underset{ji}{}\frac{dz_i}{z_jz_i}c^{(i,j)}=\frac{1}{\mathrm{}+\stackrel{ˇ}{h}}\underset{1i<jN}{}\frac{d(z_iz_j)}{z_iz_j}c^{(i,j)}$$
This lift of the WZW-connection is known as the Knizhnik-Zamolodchikov connection. It is not difficult to verify that it is still flat (see for instance ).
## 4. Double point degeneration
### Factorization
Let $`\mathrm{}`$, $`\pi :𝒞S`$ and $`V`$ be as in the preceding section so that we have defined the sheaf of conformal blocks $`_{\mathrm{}}(V)_𝒞`$. We assume in addition that we are given a section $`S_0`$ along which $`\pi `$ has an ordinary double point and a partial normalization $`\nu :\stackrel{~}{𝒞}𝒞`$ which only separates the branches. So $`\nu `$ is an isomorphism over the complement of $`S_0`$ and $`S_0`$ has two disjoint lifts to $`𝒞`$, which we shall denote by $`S_+`$ and $`S_{}`$. The *factorization principle* expresses $`_{\mathrm{}}(V)_𝒞`$ in terms of conformal blocks attached to the partial normalization of $`\stackrel{~}{𝒞}`$.
Recall that $`P_{\mathrm{}}`$ denotes the set of isomorphism classes of irreducible $`𝔤`$-representations of level $`\mathrm{}`$ and is invariant under dualization: if $`\lambda P_{\mathrm{}}`$, then $`\lambda ^{}P_{\mathrm{}}`$. Let $`V_\lambda `$ be a $`𝔤`$-representation in the equivalence class $`\lambda P_{\mathrm{}}`$ and choose $`𝔤`$-equivariant dualities $`b_\lambda :V_\lambda V_\lambda ^{}k`$, $`\lambda P_{\mathrm{}}`$.
###### Proposition 4.1.
Let $`\stackrel{~}{V}_{\lambda ,\lambda ^{}}`$ be the representation valued map on $`\stackrel{~}{𝒞}`$ which is constant equal to $`V_\lambda `$ resp. $`V_\lambda ^{}`$ on $`S_+`$ resp. $`S_{}`$ and on any other closed point of the representation already assigned to its image in $`𝒞`$. Then the contractions $`b_\lambda `$ define an isomorphism
$$\begin{array}{ccc}_{\lambda P_{\mathrm{}}}_{\mathrm{}}(\stackrel{~}{V}_{\lambda ,\lambda ^{}})_{\stackrel{~}{𝒞}}& \stackrel{}{}& _{\mathrm{}}(V)_𝒞.\end{array}$$
This is almost a formal consequence of:
###### Lemma 4.2.
Let $`M`$ be a finite dimensional representation of $`𝔤\times 𝔤`$ which is of level $`\mathrm{}`$ relative to both factors. If $`M^\delta `$ denotes that same space viewed as $`𝔤`$-module with respect to the diagonal embedding $`\delta :𝔤𝔤\times 𝔤`$, then the contraction $`_{\lambda P_{\mathrm{}}}M(V_\lambda V_\lambda ^{})M`$ (each component is defined by $`b^\lambda `$; the symbol $``$ stands for the exterior tensor product of representations) induces an isomorphism between covariants:
$$\begin{array}{ccc}_{\lambda P_{\mathrm{}}}\left(M(V_\lambda V_\lambda ^{})\right)_{𝔤\times 𝔤}& \stackrel{}{}& M_𝔤^\delta .\end{array}$$
###### Proof.
Without loss of generality we may assume that $`M`$ is irreducible, or more precisely, equal to $`V_\mu V_\mu ^{}`$ for some $`\mu ,\mu ^{}P_{\mathrm{}}`$. Then $`M^\delta =V_\mu V_\mu ^{}`$. By Schur’s lemma, $`M_𝔤^\delta `$ is one-dimensional if $`\mu ^{}=\mu ^{}`$ and trivial otherwise. That same lemma applied to $`𝔤\times 𝔤`$ shows that $`(M(V_\lambda V_\lambda ^{}))_{𝔤\times 𝔤}`$ is zero unless $`(\mu ,\mu ^{})=(\lambda ^{},\lambda )`$, in which case it is one-dimensional. The lemma follows. ∎
###### Proof of 4.1.
The issue is local on $`S`$ and so we may assume that there exists an an affine open-dense subset $`𝒞^{}`$ of $`𝒞`$ which contains its singular locus, is disjoint with $`\mathrm{Supp}(V)`$ and is such that its complement $`\stackrel{~}{S}:=𝒞𝒞^{}`$ is a trivial finite cover over $`S`$. Then $`_{\mathrm{}}(V)_𝒞=_{\mathrm{}}(V|\stackrel{~}{S})_{\pi _{}𝒪_𝒞^{}[𝔤]}`$. Let $`\stackrel{~}{𝒞}^{}:=\nu ^1𝒞^{}`$. Evaluation in $`S_0`$ resp. $`S_+,S_{}`$ defines epimorphisms $`\pi _{}𝒪_𝒞^{}𝒪_S`$ resp. $`\pi _{}\nu _{}𝒪_{\stackrel{~}{𝒞}^{}}𝒪_S𝒪_S`$ whose kernels may be identified by means of $`\nu `$. If we denote that common kernel by $``$ and $`𝔤`$ has the evident meaning, then the argument used to prove Proposition 3.2 shows that $`M:=_{\mathrm{}}(V|_{\stackrel{~}{S}})_𝔤`$ is a $`𝒪_S`$-module of finite rank. It underlies a representation of $`𝒪_S𝔤𝒪_S𝔤`$ of level $`\mathrm{}`$ relative to both factors and is such that $`M_𝔤^\delta =_{\mathrm{}}(V)_𝒞`$. Now apply Lemma 4.2. ∎
### Local freeness
We continue with the situation of Proposition 4.1, except that we now assume (for simplicity only, actually) that the base $`S`$ is $`\mathrm{Spec}(k)`$: $`C`$ is a reduced complete curve with complete intersection singularities only which has at $`S_0C`$ an ordinary double point. Choose generators $`t_\pm `$ of the maximal ideals of the completed local rings of $`C`$ at $`S_\pm `$. There is then a canonical *smoothing* of $`C`$, that is, a way of making $`C`$ the closed fiber of a flat morphism $`\tau :𝒞\mathrm{\Delta }`$, with $`\mathrm{\Delta }`$ the spectrum of the discrete valuation ring $`k[[\tau ]]`$, such that the generic fiber is smooth: in the product $`\stackrel{~}{C}\times \mathrm{\Delta }`$, blow up $`(S_\pm ,o)`$ and let $`\stackrel{~}{𝒞}`$ be the formal neighborhood of the strict transform of $`\stackrel{~}{C}\times \{o\}`$. So at the preimage of $`(S_\pm ,o)`$ on the strict transform of $`\stackrel{~}{C}\times \{o\}`$ we have the formal coordinate chart $`(t_\pm ,\tau /t_\pm )`$. Now let $`𝒞`$ be the quotient of $`\stackrel{~}{𝒞}`$ obtained by identifying these formal charts up to order: $`(t_+,\tau /t_+)=(\tau /t_{},t_{})`$, so that $`(t_+,t_{})`$ is now a formal chart of $`𝒞`$ on which we have $`\tau =t_+t_{}`$ (in either domain $`\tau `$ represents the same regular function). We thus have defined a flat morphism $`\tau :𝒞\mathrm{\Delta }`$ whose closed fiber may be identified with $`C`$. The domain $`𝒞`$ is smooth over $`k`$ and that the generic fiber of $`\tau `$ is smooth over $`k((\tau ))`$. This is our canonical smoothing of $`C`$. Let us also notice that $`𝒞`$ is at every $`pC\{S_0\}`$ canonically identified with $`(C,p)\times \mathrm{\Delta }`$ with $`\tau `$ given as the projection on the second factor.
Let $`I`$ be a finite subset of the smooth part of $`C`$ such that $`CI`$ is affine and $`V`$ is a $`𝔤`$-representation valued map on $`I`$. We extend $`V`$ canonically to $`𝒞`$ by letting it be zero at $`S_0`$ and constant on $`\{p\}\times \mathrm{\Delta }`$ when $`pS_0`$; we denote that extension by $`V/\mathrm{\Delta }`$. Then we have defined the $`k[[\tau ]]`$-module $`_{\mathrm{}}(V/\mathrm{\Delta })`$. It is clear that $`_{\mathrm{}}(V/\mathrm{\Delta })`$ is naturally identified with $`_{\mathrm{}}(V)[[\tau ]]`$. We assume that the complement of the support of $`V`$ is affine and we denote by $`A`$ its $`k`$-algebra of regular functions. Then the complement of the support of $`V/\mathrm{\Delta }`$ is affine over $`\mathrm{\Delta }`$ and if $`𝒜`$ denotes the corresponding $`k[[\tau ]]`$-algebra of regular functions, then $`A=𝒜/\tau 𝒜`$. According to Proposition 3.2, the conformal block $`_{\mathrm{}}(V/\mathrm{\Delta })_𝒞`$ is a finitely generated $`k[[\tau ]]`$-module; its reduction modulo $`\tau `$ is clearly $`_{\mathrm{}}(V)_{A𝔤}`$. If we denote by $`B`$ the algebra of regular functions on $`C\mathrm{Supp}(V)S_0=\stackrel{~}{C}\mathrm{Supp}(V)(S_+S_{})`$, then Proposition 4.1 identifies the latter with $`_{\lambda P_{\mathrm{}}}_{\mathrm{}}(\stackrel{~}{V}_{\lambda ,\lambda ^{}})_{B𝔤}`$. It is our goal to extend this identification to one of $`_{\mathrm{}}(V/\mathrm{\Delta })_𝒞`$ with $`_{\lambda P_{\mathrm{}}}_{\mathrm{}}(\stackrel{~}{V}_{\lambda ,\lambda ^{}})_{B𝔤}[[\tau ]]`$ (which will imply that $`_{\mathrm{}}(V/\mathrm{\Delta })_𝒞`$ is a free $`k[[\tau ]]`$-module).
Put $`𝒪_{}:=k[[t_+]]k[[t_{}]]`$ and $`L_{}:=k((t_+))k((t_{}))`$.
###### Lemma 4.3.
The rational map $`\stackrel{~}{C}\times \mathrm{\Delta }\stackrel{~}{𝒞}𝒞`$ identifies $`k[[t_+,t_{}]]`$ with the subalgebra of $`L_{}[[\tau ]]`$ of elements of the form $`_{n0,m0}a_{m,n}(t_+^{mn}\tau ^n,t_{}^{nm}\tau ^m)`$. Furthermore, any continuous $`k`$-derivation $`k[[t_+,t_{}]]k[[t_+,t_{}]]`$ which fixes $`\tau =t_+t_{}`$ defines in $`L_{}`$ an operator of the form
$$(t_+\frac{}{t_+},0)+\underset{m,n0}{}a_{m,n}(t_+^{mn+1}\tau ^n\frac{}{t_+},t_{}^{nm+1}\tau ^m\frac{}{t_{}}).$$
###### Proof.
Let $`f=_{n,m0}a_{m,n}t_+^mt_{}^n`$. If we substitute $`t_{}=\tau /t_+`$, this becomes at $`(S_+,o)`$: $`_{n0}(_{m0}a_{m,n}t_+^{mn})\tau ^n`$. So the coefficient of $`\tau ^n`$ has polar part $`_{m=0}^{n1}a_{m,n}t_+^{mn}`$ and constant term $`a_{n,n}`$. Similarly at $`(S_{},o)`$, $`f`$ is written as $`_{m0}(_{n0}a_{m,n}t_+^{nm})\tau ^m`$ with the coefficient of $`\tau ^m`$ having polar part $`_{n=0}^{m1}a_{m,n}t_+^{nm}`$ and constant term $`a_{m,m}`$. The polar parts and constant terms of these expressions determine all the $`a_{n,m}`$.
The second assertion is left as an exercise. ∎
Given $`\lambda P_{\mathrm{}}`$, then the Casimir element $`c=_\kappa X_\kappa X_\kappa `$ acts in $`V_\lambda `$ as a scalar, a scalar we shall denote by $`c(\lambda )`$. Observe that $`c(\lambda ^{})=c(\lambda )`$.
Let $`_{\mathrm{}}^\pm (V_\lambda )`$ denote the representation attached to $`V_\lambda `$ of the central extension of $`𝔤((t_\pm ))`$, so that $`_{\mathrm{}}^+(V_\lambda )_{\mathrm{}}^{}(V_\lambda ^{})`$ is a representation of the central extension $`\widehat{L_{}𝔤}`$ of $`L𝔤`$.
###### Lemma 4.4.
There exists a tensor valued Laurent series
$$\epsilon ^\lambda =\underset{d=0}{\overset{\mathrm{}}{}}\epsilon _d^\lambda \tau ^d(_{\mathrm{}}^+(V_\lambda )_{\mathrm{}}^{}(V_\lambda ^{}))[[\tau ]]\text{ with }\epsilon _d^\lambda F^d_{\mathrm{}}^+(V_\lambda )F^d_{\mathrm{}}^{}(V_\lambda ^{})$$
whose constant term $`\epsilon _0^\lambda V_\lambda V_\lambda ^{}`$ is the dual of $`b_\lambda `$ and which is annihilated by the image of $`𝔤[[t_+,t_{}]]`$ in $`\widehat{L_{}𝔤}`$. Moreover, for any continuous $`k`$-derivation $`D:k[[t_+,t_{}]]k[[t_+,t_{}]]`$ which fixes $`\tau `$, $`\epsilon ^\lambda `$ is an eigenvector of $`T_𝔤(D)`$ with eigenvalue $`\frac{c(\lambda )}{2(\mathrm{}+\stackrel{ˇ}{h})}`$.
###### Proof.
We first observe that the choice of the coordinates $`t_+`$ and $`t_{}`$ defines a grading on all the relevant objects on which we have defined the associated filtration $`F`$ (e.g., the degree zero summand of $`_{\mathrm{}}(V_\lambda )`$ is $`V_\lambda `$). It is known that the pairing $`b_\lambda :V_\lambda \times V_\lambda ^{}k`$ extends (in a unique manner) to a perfect pairing
$$b_\lambda :_{\mathrm{}}^+(V_\lambda )\times _{\mathrm{}}^{}(V_\lambda ^{})k$$
characterized by the property that $`b_\lambda (Xt_+^ku,u^{})+b_\lambda (u,Xt_{}^ku^{})=0`$ for all $`X𝔤`$ and $`k`$. It follows from that the restriction of $`b_\lambda `$ to $`_{\mathrm{}}^+(V_\lambda )_d\times _{\mathrm{}}^{}(V_\lambda ^{})_d^{}`$ is zero when $`dd^{}`$ and is perfect when $`d=d^{}`$. So if $`\epsilon _d^\lambda _{\mathrm{}}^+(V_\lambda )_d_{\mathrm{}}^{}(V_\lambda ^{})_d`$ denotes the latter’s inverse, then we have for all $`k`$, $`X𝔤`$
$$(Xt_+^k,0)\epsilon _{d+k}^\lambda +(0,Xt_{}^k)\epsilon _d^\lambda =0.$$
This amounts to the property that $`(Xt_+^k,\tau ^kXt_{}^k)`$ annihilates $`_{d0}\epsilon _d^\lambda \tau ^d`$. Since Lemma 4.3 says that the elements of $`𝔤[[t_+,t_{}]]`$ have series expansions in $`(Xt_+^k,Xt_{}^k)\widehat{L_{}𝔤}`$ with coefficients in $`k[[\tau ]]`$, the first statement of the lemma follows.
The second will be a straightforward computation. Using Lemma 4.3, we write $`D`$ as an operator in $`L_{}`$. We then find that we need to verify the following two assertions:
1. For every pair $`m,n0`$, $`\tau ^nT_𝔤(t_+^{mn+1}\frac{}{t_+})\tau ^mT_𝔤(t_{}^{nm+1}\frac{}{t_{}})`$ kills $`\epsilon ^\lambda `$ and
2. $`T_𝔤(t_+\frac{}{t_+})(\epsilon ^\lambda )=\frac{c(\lambda )}{2(\mathrm{}+\stackrel{ˇ}{h})}\epsilon ^\lambda `$.
As to (i), if we substitute
$$T_𝔤(t_+^{mn+1}\frac{}{t_+})=\frac{1}{2(\mathrm{}+\stackrel{ˇ}{h})}\underset{j}{}\underset{\kappa }{}:X_\kappa t_+^{mnj}X_\kappa t_+^j:$$
and do likewise for $`T_𝔤(t_{}^{nm+1}\frac{}{t_{}})`$, this assertion follows easily.
For (ii) we first observe that $`T_𝔤(t_+\frac{}{t_+})`$ preserves the grading of $`_{\mathrm{}}^+(V_\lambda )`$ and acts on $`_{\mathrm{}}^+(V_\lambda )_0=V_\lambda `$ as $`(2\mathrm{}+2\stackrel{ˇ}{h})^1_\kappa X_\kappa X_\kappa `$. This is just multiplication by $`\frac{c(\lambda )}{2(\mathrm{}+\stackrel{ˇ}{h})}`$. If $`u_{\mathrm{}}^+(V_\lambda )_d`$ is written $`u=Y_rt_+^{k_r}\mathrm{}Y_1t_+^{k_1}v`$ with $`vV_\lambda `$, $`Y_\rho 𝔤`$, $`d=k_r+\mathrm{}+k_1`$, then $`T_𝔤(t_+\frac{}{t_+})(u)=du+Y_rt_+^{k_r}\mathrm{}Y_1t_+^{k_1}T_𝔤(t_+\frac{}{t_+})(v)=(d\frac{c(\lambda )}{2(\mathrm{}+\stackrel{ˇ}{h})})u`$. Since $`t_+\frac{}{t_+}\tau ^d=d\tau ^d`$, it follows that $`\epsilon _d^\lambda \tau ^d`$ is an eigenvector of $`T_𝔤(t_+\frac{}{t_+})`$ with eigenvalue $`\frac{c(\lambda )}{2(\mathrm{}+\stackrel{ˇ}{h})}`$. ∎
###### Theorem 4.5.
The $`k[[\tau ]]`$-homomorphism
$`E=(E_\lambda )_\lambda :_{\mathrm{}}(V/\mathrm{\Delta })=_{\mathrm{}}(V)[[\tau ]]`$ $`_{\lambda P_l}_{\mathrm{}}(\stackrel{~}{V}_{\lambda ,\lambda ^{}})[[\tau ]],`$
$`U={\displaystyle \underset{k0}{}}u_k\tau ^k`$ $`{\displaystyle \underset{\lambda P_{\mathrm{}}}{}}U\epsilon ^\lambda :={\displaystyle \underset{\lambda P_{\mathrm{}}}{}}{\displaystyle \underset{k,d0}{}}u_k\epsilon _d^\lambda \tau ^{k+d}`$
is also a map of $`𝒜𝔤`$-representations if we let $`𝒜𝔤`$ act on the left hand side via the inclusion $`𝒜B[[\tau ]]`$. The resulting $`k[[\tau ]]`$-homomorphism of covariants,
$$_{\mathrm{}}(V/\mathrm{\Delta })_𝒞_{\lambda P_l}_{\mathrm{}}(\stackrel{~}{V}_{\lambda ,\lambda ^{}})_{B𝔤}[[\tau ]]=_{\lambda P_l}_{\mathrm{}}(\stackrel{~}{V}_{\lambda ,\lambda ^{}})_{\stackrel{~}{C}}[[\tau ]]$$
is an isomorphism. In particular, $`_{\mathrm{}}(V/\mathrm{\Delta })_𝒞`$ is a free $`k[[\tau ]]`$-module. Moreover, covariant derivation in $`_{\mathrm{}}(V/\mathrm{\Delta })_𝒞`$ by $`\tau \frac{d}{d\tau }`$ respects each summand $`_{\mathrm{}}(\stackrel{~}{V}_{\lambda ,\lambda ^{}})_{\stackrel{~}{C}}[[\tau ]]`$ and acts there as the first order differential operator $`\tau \frac{d}{d\tau }+\frac{c(\lambda )}{2(\mathrm{}+\stackrel{ˇ}{h})}`$.
###### Proof.
The first statement is immediate from Lemma 4.4. So the map on covariants is defined and is $`k[[\tau ]]`$-linear. If we reduce modulo $`\tau `$ we get the map
$$u_{\mathrm{}}(V)_{A𝔤}\underset{\lambda P_{\mathrm{}}}{}u\epsilon _0^\lambda _{\lambda P_l}_{\mathrm{}}(\stackrel{~}{V}_{\lambda ,\lambda ^{}})_{B𝔤}.$$
We recognize its domain and range as $`_{\mathrm{}}(V/\mathrm{\Delta })_C`$ and $`_{\lambda P_l}_{\mathrm{}}(\stackrel{~}{V}_{\lambda ,\lambda ^{}})_{\stackrel{~}{C}}`$ and we observe that the map itself is just the inverse of the isomorphism of Proposition 4.1. Since the range is a free $`k[[\tau ]]`$-module, this implies that the $`k[[\tau ]]`$-homomorphism of covariants is an isomorphism.
According to Corollary 2.4 covariant derivation with respect to $`\tau \frac{d}{d\tau }`$ in $`_{\mathrm{}}(V/\mathrm{\Delta })_𝒞`$ is defined by means of a $`k`$-derivation $`D`$ of $`𝒜`$ which lifts $`\tau \frac{d}{d\tau }`$: if we write $`D=\tau \frac{d}{d\tau }+_{n0}\tau ^nD_n`$, where $`D_n`$ is a vector field on the smooth part of $`C`$, then the covariant derivative is induced by $`T_𝔤(D)=\tau \frac{d}{d\tau }+_{n0}\tau ^nT_𝔤(D_n)`$ acting on $`_{\mathrm{}}(V)[[\tau ]]=_{\mathrm{}}(V/\mathrm{\Delta })_𝒞`$. From the last clause of Lemma 4.4 we get that when $`U_{\mathrm{}}(V)[[\tau ]]`$,
$$\begin{array}{c}T_𝔤(D)E_\lambda (U)=T_𝔤(D)(U\epsilon ^\lambda )=\hfill \\ \hfill =T_𝔤(D)(U)\epsilon ^\lambda \frac{c(\lambda )}{2(\mathrm{}+\stackrel{ˇ}{h})}U\epsilon ^\lambda =E_\lambda T_𝔤(D)(U)\frac{c(\lambda )}{2(\mathrm{}+\stackrel{ˇ}{h})}E_\lambda (U).\end{array}$$
Since $`T_𝔤(D)`$ acts on $`_{\mathrm{}}(\stackrel{~}{V}_{\lambda ,\lambda ^{}})_{\stackrel{~}{C}}[[\tau ]]`$ simply as derivation by $`\tau \frac{d}{d\tau }`$, the last clause follows. ∎
###### Corollary 4.6.
The monodromy of the WZW connection acting on $`_{\mathrm{}}(V/\mathrm{\Delta })_𝒞`$ has finite order and acts in the summand $`_{\mathrm{}}(\stackrel{~}{V}_{\lambda ,\lambda ^{}})_{\stackrel{~}{C}}`$ as multiplication by $`\mathrm{exp}(\pi \sqrt{1}\frac{c(\lambda )}{\mathrm{}+\stackrel{ˇ}{h}})`$.
###### Proof.
The multivalued flat sections of $`_{\mathrm{}}(V/\mathrm{\Delta })_𝒞`$ decompose under $`E`$ as a direct sum labeled by $`P_{\mathrm{}}`$. The summand corresponding to $`P_{\mathrm{}}`$ is the set of solutions of the differential equation $`\tau \frac{d}{d\tau }U+\frac{c(\lambda )}{2(\mathrm{}+\stackrel{ˇ}{h})}U=0`$. These are clearly of the form $`u\tau ^{c(\lambda )/2(\mathrm{}+\stackrel{ˇ}{h})}`$ with $`u_{\mathrm{}}(\stackrel{~}{V}_{\lambda ,\lambda ^{}})_{\stackrel{~}{C}}`$. If we let $`\tau `$ run over the unit circle, then we see that the monodromy is as asserted. Finally, we observe that $`\frac{c(\lambda )}{\mathrm{}+\stackrel{ˇ}{h}}`$. ∎
### Verlinde formula
Theorem 4.5 and its Corollary 4.6 extend to the case where the base $`S`$ is an arbitrary $`k`$-variety as in Proposition 4.1 and the smoothing is arbitrary. This is based on a versality argument, which shows that our smoothing construction is not so special as it may appear. To be concrete, suppose that we are given a family of smooth curves $`\stackrel{~}{𝒞}S`$ with pairwise disjoint sections $`\{S_i\}_{iI}\{S_{},S_+\}`$ and a generator $`t_\pm `$ of the ideal defining $`S_\pm `$ in the formal completion along $`𝒞`$. Assume that the complement of the union of these sections is affine over $`S`$ and that this family is *versal* as a family of pointed curves. Then $`\stackrel{~}{𝒞}S`$ factors through a family whose fibers have a single double point: $`\stackrel{~}{𝒞}𝒞S`$, where $`\stackrel{~}{𝒞}`$ is obtained by identifying the sections $`S_+`$ and $`S_{}`$. We regard the latter as endowed with the sections $`\{S_i\}_{iI}`$ so that $`(𝒞,\{S_i\}_{iI})S`$ is a family of pointed curves. Then the smoothing of $`𝒞`$ over $`S\times \mathrm{\Delta }`$ defined by $`t_\pm `$ with its sections $`\{S_i\times \mathrm{\Delta }\}_{iI}`$ will be a versal (as a family of pointed curves) so that any deformation of $`(𝒞,\{S_i\}_{iI})S`$ is obtained from this one by means of a base change. As any two versal deformations are isomorphic, it follows that 4.5 and 4.6 apply to any versal deformation of $`(𝒞,\{S_i\}_{iI})S`$.
The preceding leads to a formula of the dimension of a conformal block. By theorem 4.5, or rather the generalization discussed above, the dimension of a conformal block stays the same under a degeneration. Since every pointed curve degenerates into one with ordinary double points whose normalization consists of curves of genus zero, it suffices to do the computation for such a degenerate curve. But then we may invoke 4.1 to reduce to the case of a smooth rational curve, which can be dealt with using our discussion of the genus zero case. We can even arrange that the support of the representation map meets every component in at most three points. A more refined approach involves the notion of a fusion ring .
###### Remark 4.7.
In case $`k=`$, one can work in the complex-analytic category. Then the fact that the singularity is formally regular singular ensures that the flat multivalued sections converge on simply connected sectors based at $`o`$ (see for instance ) so that the monodromy has the expected interpretation. If $`C\{S_0\}`$ is smooth, then the we have an associated fibration over the punctured unit disk whose monodromy is given by a Dehn twist along a vanishing circle on the general fiber. In terms of the setting of Remark 3.9: an isotopy class $`\alpha `$ of embedded circles in $`\mathrm{\Sigma }`$ defines the class of a Dehn twist $`D_\alpha \mathrm{\Gamma }(\mathrm{\Sigma },I)`$ and the action of $`D_\alpha `$ on $`(\mathrm{\Sigma },I,V)`$ is given by the above formula (hence is of finite order). |
warning/0507/astro-ph0507002.html | ar5iv | text | # STEllar Content and Kinematics from high resolution galactic spectra via Maximum A Posteriori
## 1 Introduction
For decades now, the spectral indices from the Lick group have been used to study the properties of stellar populations (Faber et al., 1985; Worthey, 1994; Trager et al., 1998). Since the profile and depth of the lines involved in these spectral indices are affected by the Line Of Sight Velocity Distribution (hereafter LOSVD) of the stars, it is necessary to correct the measured depths by a factor depending on the moments of the velocity distribution (Davies et al., 1993; Kuntschner, 2000, 2004). The latter moments must be determined using specialized code (Bender, 1990; Saha & Williams, 1994; Pinkney et al., 2003; Merritt, 1997; Kuijken & Merrifield, 1993; van der Marel & Franx, 1993). These kinematical deconvolution routines have been used for some time and have undergone 2 major mutations. First, thanks to the increasing power of computers, it became affordable to swap back and forth from direct space to Fourier space, so that many disturbances such as border effects and saturation could be avoided. It became straightforward to mask problematic regions of the data, such as dead pixels, emission lines, etc… The second evolution of these codes allowed the use of multiple superimposed stellar templates to best match the observed spectrum (Rix & White, 1992; Cappellari & Emsellem, 2004). It has also been proposed to use single stellar populations as synthetic templates, and this approach has proved to be useful in addressing the template mismatch problem (Falcón-Barroso et al., 2003). Moreover, this technique can save precious telescope time since it circumvents the need for observing template stars.
In Ocvirk et al. 2005 (hereafter Paper I), we introduced STECMAP, a method for recovering non-parametrically the stellar content of a given galaxy from its integrated light spectrum. Using STECMAP requires, as a preliminary, convolving the data or models with the proper Point Spread Function (PSF), which can be of both physical (i.e. the stellar LOSVD) and instrumental (the instrument’s PSF) origin. Adjusting the LOSVD to fit the data does not only constrain the kinematics of the observed galaxy but will also reduce the mismatch due to errors in the determination of the redshift or anomalous PSF, which is ultimately a necessary step when fitting galaxy spectra.
Here we propose to constrain the velocity distribution simultaneously with the stellar content, by merging the kinematic deconvolution and the stellar content reconstruction in one global Maximum A Posteriori likelihood inversion method. Hence, STECMAP becomes STECKMAP (STEllar Content and Kinematics via Maximum A Posteriori likelihood). In this respect, STECKMAP resembles the method proposed by e.g. Falcón-Barroso et al. (2003), except that it takes advantage of the treatment of the stellar content by STECMAP. Together with the stellar age distribution and the age-metallicity relation, the LOSVD is described non-parametrically and the only a priori we use are smoothness and positivity.
We also tentatively address the case of age-dependent kinematics, i.e. we try to recover the individual LOSVDs and ages of several superimposed kinematical subcomponents. This approach is motivated by the fact that galaxies often display several kinematical components. Ellipticals and dwarf ellipticals are for instance known to often harbor kinematically decoupled cores (De Rijcke et al., 2004; Balcells & Quinn, 1990; Bender & Surma, 1992), and spiral galaxies are usually assumed to be constituted of a thin and a thick disk, a bulge and a halo (Freeman & Bland-Hawthorn, 2002). The variety of the dynamical properties of the components has a counterpart in their stellar content, as a signature of the formation and evolution of the galaxy. For instance, the halo of the Milky Way is believed to consist mainly of old, metal poor stars, while the bulge is more metal rich, and the thin disk is mainly younger than the bulge (Freeman & Bland-Hawthorn, 2002). It is thus natural to let any stellar sub population have its own LOSVD. This possibility has been recently addressed by De Bruyne et al. (2004a, b), in a slightly different framework: they use individual stars as templates for the different components, while we propose to use synthetic SSP models. Such a method would allow us to separate the several kinematical components of galaxies from integrated light spectra, and constrain for example their age-velocity dispersion and age-metallicity relation. The highly detailed stellar content and kinematical information that can be obtained for the Milky Way or for nearby galaxies that can be resolved into stars, such as M31 (Ferguson et al., 2002; Ibata et al., 2004), could be extended to a larger sample of more distant galaxies. This technique could also be useful in detecting and characterizing major stellar streams in age and velocity from integral field spectroscopy of galaxies.
In the whole paper we use the Pégase-HR SSP models (Le Borgne et al., 2004) in order to illustrate and investigate the behaviour of the problems through simulations and inversions of mock data. Indeed, Pégase-HR, with its high spectral resolution ($`R=\mathrm{10\hspace{0.17em}000}`$), is an adequate choice for testing the recovery of detailed kinematical information in the form of non-parametric LOSVDs. The problems and methods we describe are however by no means specific to Pégase-HR (and its wavelength coverage) and STECKMAP could be used with any possible SSP model, depending on the type of data that is being analyzed.
We will start with the modeling of the kinematics. Then, we will address the idealized linear problem of recovering the LOSVD when the stellar content is known, i.e. the template is assumed to be perfect. Section 4 deals with simultaneous age and LOSVD reconstruction of composite populations. Finally, section 5 investigates the case of age-dependent kinematics in a simplified context where the metallicity and extinction are known a priori.
## 2 Models of galaxy spectra
In this section we present the modeling of galaxy spectra, taking into account the composite nature of the stellar population, in age, metallicity and extinction, and finally its kinematics.
### 2.1 The composite reddened population at rest
We model the SED of the composite reddened population at rest using the ingredients defined in Paper I:
$$F_{\mathrm{rest}}(\lambda )=f_{\mathrm{ext}}(E,\lambda )_{t_{\mathrm{min}}}^{t_{\mathrm{max}}}\mathrm{\Lambda }(t)B(\lambda ,t,Z(t))dt,$$
(1)
where $`\mathrm{\Lambda }(t)`$ is the luminosity weighted stellar age distribution, $`Z(t)`$ is the age-metallicity relation, and $`B(\lambda ,t,Z)`$ is the flux-averaged single stellar population basis of an isochrone population of age $`t`$, $`f_{\mathrm{ext}}`$ the extinction law, and metallicity $`Z`$. We recall briefly the main properties of the Pégase-HR SSP basis we used in this paper. As mentioned earlier, spectral resolution is $`R=\mathrm{10\hspace{0.17em}000}`$ over the full optical domain $`\lambda \lambda =[4000,6800]`$ Å, sampled in steps of $`0.2`$ Å. The models are available for metallicities $`Z[0.0001,\mathrm{\hspace{0.17em}0.1}]`$ and considered reliable between $`t_{\mathrm{min}}=10`$ Myr and $`t_{\mathrm{max}}=15`$ Gyr. The IMF used is described in Kroupa et al. (1993) and the stellar masses range from $`0.1`$$`\mathrm{M}_{}`$ to $`120`$$`\mathrm{M}_{}`$. The extinction law $`f_{\mathrm{ext}}`$ was taken from Calzetti (2001).
### 2.2 Model kinematics
Stellar motions in galaxies define a LOSVD corresponding to projected local velocity distributions integrated along the line of sight and across one resolved spatial element.
#### 2.2.1 Global kinematics
The motion of the stars can to first approximation be accounted for by assuming that the velocities of all stars of all ages along the line of sight are taken from the same velocity distribution (hence “global”). The model spectral energy distribution, $`\varphi (\lambda )`$, is the convolution of the assumed normalized LOSVD, $`g(\mathrm{v})`$, defined for $`\mathrm{v}[\mathrm{v}_{\mathrm{min}},\mathrm{v}_{\mathrm{max}}]`$ with the model spectrum at rest $`F_{\mathrm{rest}}(\lambda )`$. The convolved spectrum $`\varphi (\lambda )`$ reads:
$$\varphi (\lambda )=_{\mathrm{v}_{\mathrm{min}}}^{\mathrm{v}_{\mathrm{max}}}F_{\mathrm{rest}}\left(\frac{\lambda }{1+\mathrm{v}/c}\right)g(\mathrm{v})\frac{\mathrm{dv}}{1+\mathrm{v}/c},$$
(2)
where $`c`$ is the velocity of light. The above expression reads as a standard convolution
$$\stackrel{~}{\varphi }(w)=c_{u_{\mathrm{min}}}^{u_{\mathrm{max}}}\stackrel{~}{F}(wu)\stackrel{~}{g}(u)du,$$
(3)
with the following reparameterization:
$`w`$ $``$ $`\mathrm{ln}(\lambda ),u\mathrm{ln}(1+{\displaystyle \frac{\mathrm{v}}{c}}),`$ (4)
$`\stackrel{~}{F}(w)`$ $``$ $`F_{\mathrm{rest}}(\mathrm{e}^w)=F_{\mathrm{rest}}(\lambda ),`$ (5)
$`\stackrel{~}{g}(u)`$ $``$ $`g(c(\mathrm{e}^u1))=g(\mathrm{v}),\stackrel{~}{\varphi }(w)\varphi (\mathrm{e}^w)=\varphi (\lambda ),`$ (6)
$`u_{\mathrm{min}}`$ $`=`$ $`\mathrm{ln}(1+{\displaystyle \frac{\mathrm{v}_{\mathrm{min}}}{c}}),u_{\mathrm{max}}=\mathrm{ln}(1+{\displaystyle \frac{\mathrm{v}_{\mathrm{max}}}{c}}).`$ (7)
#### 2.2.2 Age-dependent kinematics
We now allow the LOSVD to depend on the age of the stars. For simplicity, we consider here only unreddened mono-metallic populations, i.e. $`f_{\mathrm{ext}}(E,\lambda )=1`$ and $`Z(t)=Z_0`$. We introduce the age-velocity distribution, $`\mathrm{\Xi }(\mathrm{v},t)`$, defined in $`[\mathrm{v}_{\mathrm{min}},\mathrm{v}_{\mathrm{max}}]\times [t_{\mathrm{min}},t_{\mathrm{max}}]`$, which gives the contribution of stars of velocity and age in $`[\mathrm{v},\mathrm{v}+\mathrm{dv}]\times [t,t+\mathrm{d}t]`$ to the total observed light. Thus, for a given age $`t`$, $`\mathrm{\Xi }(\mathrm{v},t)`$ is the LOSVD of the single stellar population of age $`t`$. The age-velocity distribution $`\mathrm{\Xi }(\mathrm{v},t)`$ is related to the stellar age distribution, $`\mathrm{\Lambda }(t)`$, by:
$$_{\mathrm{v}_{\mathrm{min}}}^{\mathrm{v}_{\mathrm{max}}}\mathrm{\Xi }(\mathrm{v},t)\mathrm{dv}=\mathrm{\Lambda }(t),$$
(8)
The model spectrum of such a population thus reads:
$`\varphi (\lambda )={\displaystyle _{t_{\mathrm{min}}}^{t_{\mathrm{max}}}}{\displaystyle _{\mathrm{v}_{\mathrm{min}}}^{\mathrm{v}_{\mathrm{max}}}}B({\displaystyle \frac{\lambda }{1+\mathrm{v}/c}},t,Z_0)\mathrm{\Xi }(\mathrm{v},t){\displaystyle \frac{\mathrm{dv}\mathrm{d}t}{1+\mathrm{v}/c}},`$ (9)
The above expression can be rewritten more conveniently
$$\stackrel{~}{\varphi }(w)=c_{u_{\mathrm{min}}}^{u_{\mathrm{max}}}_{t_{\mathrm{min}}}^{t_{\mathrm{max}}}\stackrel{~}{B}(wu,t)\stackrel{~}{\mathrm{\Xi }}(u,t)dtdu,$$
(10)
using the same reparameterization as in Sect. 2.2.1 and
$`\stackrel{~}{B}(w,t)`$ $``$ $`B(\mathrm{e}^w,t,Z_0)=B(\lambda ,t,Z_0),`$ (11)
$`\stackrel{~}{\mathrm{\Xi }}(u,t)`$ $``$ $`\mathrm{\Xi }(c(\mathrm{e}^u1),t)=\mathrm{\Xi }(\mathrm{v},t).`$ (12)
In the rest of the paper, we will use exclusively the standard (i.e. reparameterized) convolutions as in Eq. (3) and Eq. (10). For readability, we will drop the superscript $`\stackrel{~}{}`$ and set the speed of light to one.
## 3 Kinematical deconvolution
Sect. 2.2.1 shows that with proper reparameterization, the convolution of a model spectrum at rest, $`F(w)`$, with the stellar LOSVD, $`g(u)`$, reads as a standard convolution, given by Eq. (3). Finding the LOSVD when the observed spectrum, $`\varphi (w)`$, and the template spectrum, $`F(w)`$, are given is a classical deconvolution problem. Our goal here is not to discuss the respective qualities of the many different methods available in the literature to solve this problem. Most rely on fitting the data while imposing some a priori on the LOSVD, i.e. they provide Maximum A Posteriori (MAP) estimates of the LOSVD. Let us describe briefly our method to obtain such a solution with the purpose of coupling it in a later step with STECMAP.
### 3.1 The convolution kernel
Here we discretize Eq. (3) to obtain a matrix form defining the convolution kernel. We use an evenly spaced set
$$\left\{u_j=u_{\mathrm{min}}+(j1/2)\delta u;j=1,2,\mathrm{},p\right\}$$
spanning $`[u_{\mathrm{min}},u_{\mathrm{max}}]`$ with constant step $`\delta u(u_{\mathrm{max}}u_{\mathrm{min}})/p`$. We expand the LOSVD as a sum of $`p`$ gate functions:
$$g(u)=\frac{1}{\delta u}\underset{j}{}g_j\theta \left(\frac{uu_j}{\delta u}\right)$$
where
$$\theta (x)=\{\begin{array}{cc}1\hfill & \text{if }1/2<x1/2\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}.$$
Injecting this expansion into Eq. (3) leads to:
$`\varphi (w)`$ $`=`$ $`{\displaystyle \frac{1}{\delta u}}{\displaystyle \underset{j=1}{\overset{j=p}{}}}g_j{\displaystyle \underset{u_{\mathrm{min}}}{\overset{u_{\mathrm{max}}}{}}}F(wu)\theta \left({\displaystyle \frac{uu_j}{\delta u}}\right)du,`$ (13)
$``$ $`{\displaystyle \underset{j=1}{\overset{j=p}{}}}g_jF(wu_j).`$
Similarly, we now sample along the wavelengths by integrating over a small $`\delta w`$:
$`\varphi _i`$ $``$ $`{\displaystyle \frac{1}{\delta w}}{\displaystyle \varphi (w)\theta \left(\frac{ww_i}{\delta w}\right)dw},`$ (14)
$``$ $`{\displaystyle \underset{j=1}{\overset{j=p}{}}}g_jF(w_iu_j),`$
where $`\{w_j;j=1,2,\mathrm{},m\}`$ is a set of *logarithmic* wavelengths spanning the spectral range with a constant step.
Using matrix notation and accounting for data noise, the observed SED reads:
$$𝐲=𝐊𝐠+𝐞,$$
(15)
where $`𝐲=(\varphi _1,\varphi _2,\mathrm{},\varphi _m)^{}`$ is the measured spectrum, and $`𝐞=(e_1,e_2,\mathrm{},e_m)^{}`$ accounts for modelling errors and noise. The vector of sought parameters $`𝐠=(g_1,g_2,\mathrm{},g_p)^{}`$ is the discretized LOSVD. The vector $`𝐬=𝐊𝐠`$ is the model of the observed spectrum and the matrix $`𝐊`$
$$K_{ij}=F(w_iu_j),(i,j)\{1,\mathrm{},m\}\times \{1,\mathrm{},p\},$$
(16)
is called the convolution kernel.
The convolution theorem (Press, 2002) states that the Fourier transform of the convolution of two functions is equal to the frequency-wise product of the individual Fourier transforms of the two functions. Applying this theorem yields another equivalent expression for the model spectrum $`𝐬`$:
$$𝐬=^1\mathrm{diag}(𝐅)𝐠,$$
(17)
where $``$ is the discrete Fourier operator defined in Press (2002) as:
$`_{ij}`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \frac{2\mathrm{i}\pi }{m}}(i1)(j1)\right),(i,j)[1,\mathrm{}m]^2,`$ (18)
$`^1`$ $`=`$ $`{\displaystyle \frac{1}{m}}^{}.`$ (19)
Note that since $`m`$ is the size of the template spectrum $`𝐅`$, the discretized LOSVD $`𝐠`$, which is initially of size $`p`$ needs to be symmetrically padded with zeros to the size $`m`$ in order to transform the Toeplitz matrix into a circulant one. The diagonal matrix $`\mathrm{diag}(𝐅)`$ carries the coefficients of the Fourier transform of the model spectrum at rest $`𝐅`$. This notation involving the Fourier operator, $``$, will be very useful for a number of algebraic derivations in the rest of the paper. In practice, from a computational point of view, it is more efficient to implement any forward or inverse Fourier transform through FFT. Similarly, the product $`\mathrm{diag}(𝐅)𝐠`$ is in practice implemented as a frequency-wise product of the individual FFTs.
### 3.2 Regularization and MAP
A number of earlier publications have shown that the maximum likelihood solution to Eq. (15) is very sensitive to the noise in the data $`𝐞`$. Hence, in the spirit of Paper I, we choose to regularize the problem by requiring the LOSVD to be smooth. To do so, we use the quadratic penalization $`P(𝐠)`$ as defined by Eq. (29) in Paper I:
$$P(𝐠)=𝐠^{}𝐋^{}𝐋𝐠.$$
(20)
In the rest of the paper, the penalization is Laplacian, i.e. $`𝐋=𝐃_2`$, where $`𝐃_2`$ is the discrete second order difference operator, as defined in Pichon et al. (2002). The objective function, $`Q_\mu `$, to be minimized is given by:
$$Q_\mu (𝐠)=\chi ^2(𝐲|𝐠)+\mu P(𝐠),$$
(21)
where the $`\chi ^2`$ is defined by
$$\chi ^2(𝐲|𝐠)=\left(𝐲𝐬(𝐠)\right)^{}𝐖\left(𝐲𝐬(𝐠)\right).$$
(22)
The vector $`𝐲`$ is the observed spectrum and the weight matrix is the inverse of the covariance matrix of the noise: $`𝐖=\mathrm{Cov}(𝐞)^1`$. The parameter $`\mu `$ controls the smoothness of the LOSVD through its coefficients, $`𝐠`$. It can be set on the basis of simulations (as described in Paper I) or automatically by GCV (Wahba, 1990), according to the SNR of the data. In the latter case, the properties of the convolution kernel can be used to speed up the computation of the GCV function. Further regularization is provided by the requirement of positivity, implemented through quadratic reparameterization. Minimizing $`Q_\mu `$ yields the regularized solution $`𝐠_\mu `$. Efficient minimization procedures require the analytical expression of the gradients of $`Q_\mu `$, given in Sect. A.1.
### 3.3 Simulations
We applied this deconvolution technique to mock data, created from Pégase-HR SSPs of several ages and metallicities, with $`R=\mathrm{10\hspace{0.17em}000}`$ at $`40006800`$ Å. In a first set of experiments, the model spectrum at rest was a solar metallicity $`10`$ Gyr single stellar population. It was convolved with various LOSVDs, both Gaussian and non Gaussian, with velocity dispersions ranging from $`30`$ to $`500`$ km/s. It was then perturbed with Gaussian noise at levels ranging from SNR$`=5`$ to $`100`$ per pixel, and deconvolved using the model spectrum at rest as template (i.e. no template mismatch). In all cases, the LOSVDs are adequately recovered. Figure 1 shows the reconstruction of a Gaussian LOSVD, for SNR$`=10`$ per pixel. However, there are necessarily some biases in the reconstruction of the sharp features of the LOSVD. This is expected since we introduced regularization via smoothing. To illustrate the relationship between regularization and bias, we performed a new set of similar simulations for a non-Gaussian LOSVD (sum of 2 Gaussians) with SNR=20 per pixel and varied the smoothing parameter $`\mu `$. The results are shown in figure 2. The panels a and b correspond to $`\mu =10`$ while the panels c and d correspond to $`\mu =1000`$. The model, median and interquartiles of 500 reconstructions are displayed. We also plotted the whole set of 500 recovered solutions, in order to show the locus of the solutions. One can see that the biases of the median recontruction are reduced when lowering $`\mu `$. The highest bump is correctly reproduced for $`\mu =10`$ while it is not for $`\mu =1000`$. But on the other hand the solutions are much more widely spread when $`\mu =10`$. This means that most solutions taken from the set of low $`\mu `$ simulations can be very far from the model, while all the large $`\mu `$ solutions lie reasonably close to the model.
The regularization acts as a Wiener filter in the sense that it damps the high frequency components of the solution. Regularization improves the significance of an individual reconstruction (it will nearly always lie reasonably close to the model), at the cost of introducing a bias.
### 3.4 Age and metallicity mismatch
We take advantage of the large range of ages and metallicities of single stellar populations covered by Pégase-HR to shortly illustrate the effects of template mismatch on LOSVD determinations. In this section we show the results of a large number of simulations aiming at characterizing the error made when a wrong template is chosen for the kinematical inversion of data. For this purpose, mock data were created by convolving a single stellar population of age, $`a_0`$, and metallicity, $`Z_0`$, with a centered Gaussian LOSVD of dispersion $`\sigma _v=50`$ km/s. It was perturbed by Gaussian noise corresponding to SNR$`=100`$ per pixel and then deconvolved, using as template a single stellar population of age $`a_1`$, and metallicity $`Z_1`$. The spectral resolution and wavelength range are the same as in Sect. 3.3. Fig. 3 shows the error on the measured velocity dispersion. The latter is measured as the r.m.s of the reconstructed LOSVD. If the parameters of the template are different from those of the model, the velocity dispersion error increases very quickly. The age metallicity degeneracy is visible as a valley of smaller error, following the upper-left to bottom-right diagonal of the figures. Of course, the $`\chi ^2`$ distance between the model and the mock data follows a similar 2D distribution, and will lead to the rejection of highly mismatched LOSVD estimates. However, in practice, it is usually not straightforward to quantify all the sources of error. It is thus somewhat arbitrary to set an upper limit of $`\chi ^2`$ for the admissible solutions, and the error on the kinematics is thus hard to quantify. This experiment illustrates in this context the long known issue that when the correct model is not available, large errors on the determination of kinematics are expected. In order to reduce the error in the estimates of the kinematical properties of a stellar assembly, it is necessary to allow for a wide range of modulations of the template. This is naturally achieved by making the non parametric stellar content account for the changes of the template, as discussed in the next section.
## 4 Recovering stellar content and global kinematics
The mixed inversion described in this section couples the recovery of both the stellar content and the kinematics, thereby turning STECMAP into STECKMAP. Proper application of this method provides an interpretation of the observed object in terms of stellar content and kinematics.
### 4.1 Inverse problem
For a given model spectrum at rest, $`F_{\mathrm{rest}}(\lambda )`$, and a LOSVD, $`g(v)`$, the emitted SED, $`\varphi (\lambda )`$ is given by Eq. (2). We now wish to account also for the additional variables involved in $`F_{\mathrm{rest}}`$, given by Eq. (1), namely the stellar age distribution, $`\mathrm{\Lambda }(t)`$, the age-metallicity relation $`Z(t)`$, and the color excess $`E(BV)=E`$. Injecting Eq. (1) into the convolution Eq. (3) yields the emitted SED:
$$\varphi (w)=f_{\mathrm{ext}}(E,wu)\mathrm{\Lambda }(t)B(wu,t,Z(t))g(u)dtdu,$$
(23)
Solving Eq. (23) for $`\mathrm{\Lambda }`$, $`Z`$, $`E`$ and $`g`$ when $`\varphi `$, $`f_{\mathrm{ext}}`$ and $`B`$ are given is the inverse problem we are tackling here.
### 4.2 Discretization and parameters
Expanding the two time-dependent unknowns $`\mathrm{\Lambda }(t)`$ and $`Z(t)`$ as a sum of $`n`$ gate functions and injecting into Eq. (1) yields the discrete model spectrum at rest:
$$𝐅=\mathrm{diag}(𝐟_{\mathrm{ext}}(E))𝐁𝐱,$$
(24)
This discretization is explained in details in Sec.5 of Paper I. Similarly, we develop the LOSVD $`g(u)`$ as a sum of $`p`$ gate funtions as in Sect. 3. Note that the reddened model at rest plays the role of the stellar template in a classical kinematic convolution. Injecting Eq. (24) into Eq. (17) thus allows us to express the model spectrum, $`𝐬`$, as
$$𝐬=^{}\mathrm{diag}(\mathrm{diag}(𝐟_{\mathrm{ext}}(E))𝐁𝐱)𝐠,$$
(25)
However, here, the template is this time modulated by the unknowns describing the stellar content.
### 4.3 Smoothness and metallicity constraints
The discrete problem of finding the stellar age distribution $`𝐱`$, the age-metallicity relation $`𝐙`$, the extinction $`E`$ and the LOSVD $`𝐠`$ for an observed spectrum $`𝐲`$ and given an extinction law $`f_{\mathrm{ext}}`$ and a SSP basis $`B`$ is of course likely to be very ill-conditioned since it arises as the combination of several ill-conditioned problems. It therefore requires regularization. We also want the metallicity of the components to remain in the model range. We use the standard penalization $`P`$ and the binding function $`C`$ defined in Paper I to build the penalization $`P_𝝁`$ for this problem:
$`P_𝝁=\mu _𝐱P(𝐱)+\mu _𝐙P(𝐙)+\mu _CC(𝐙)+\mu _vP(𝐠),`$ (26)
where $`𝝁(\mu _𝐱,\mu _𝐙,\mu _C,\mu _v)`$. Again, we choose $`𝐋=𝐃_2`$ as defined in Pichon et al. (2002), so that the penalization $`P`$ is actually Laplacian. The objective function, $`Q_𝝁`$, is now defined as:
$$Q_𝝁=\chi ^2(𝐬(𝐱,𝐙,E,𝐠))+P_𝝁(𝐱,𝐙,E,𝐠).$$
(27)
and its partial derivatives are given in Sect. A.2. Note that there is in principle an additionnal formal degeneracy for this inverse problem. If the set $`(𝐱,𝐙,E,𝐠)`$ is a solution to (23), then $`(\alpha 𝐱,𝐙,E,𝐠/\alpha )`$ is also a solution for any scalar $`\alpha `$, because the age distribution $`𝐱`$ and the LOSVD $`𝐠`$ are not explicitly normalized in this formulation. However, the adopted regularization lifts this degeneracy. The penalization function $`P`$ is quadratic ($`P(\alpha 𝐱)=\alpha ^2P(𝐱)`$). Thus, if $`𝐱`$ or $`𝐠`$ is too large in norm, the solution is unattractive. Practically, the algorithm reaches a solution where $`𝐱`$ and $`𝐠`$ are similar in norm. In any case, this degeneracy would easily be remedied by adding a normalizing term to the penalization $`P_𝝁`$ of the form $`𝐱1`$, which would force the discretized stellar age distribution $`𝐱`$ to have unitary norm. Following the same principle, one could equivalently choose to normalize the LOSVD rather than the stellar age distribution.
### 4.4 Simulations
Let us now test the behaviour of STECKMAP by applying it to mock data. The latter were produced using an arbitrary stellar age distribution $`𝐱`$, an age-metallicity relation $`𝐙`$, a LOSVD $`𝐠`$ and an extinction parameter $`E`$. Several simulations were performed with various input models: bumpy age distributions, increasing or decreasing age-metallicity relation and extinctions, Gaussian and non Gaussian wide or narrow LOSVDs, in various pseudo-observational contexts. Figure 4 shows the results of two of these experiments. In the top line, the model is a young metal-poor population superimposed to an older metal-rich population. In the bottom panels, the model has a rather constant stellar age distribution, a non-monotonic age-metallicity relation and a strongly non Gaussian LOSVD. In both cases the 3 unknowns are correctly recovered. In these examples, the data quality mimics that of the best Sloan Digital Sky Survey galaxies: the resolution is $`R2000`$ and SNR$`=30`$ per $``$1 Å pixel. The wavelength domain of Pégase-HR is however narrower than the SDSS’s. These simulations simply aim at demonstrating the generally good behaviour of the method, and show that accounting for the kinematics does not fundamentally weaken the constraints on the stellar content. For a more thorough study of the informational content of the Pégase-HR wavelength range, the reader can refer to the systematic double burst simulations with variable spectral resolution and SNR per Å performed in Paper I.
## 5 Recovery of age-dependent kinematics
In this section we present an implementation of the recovery of age-dependent kinematics, i.e the situation when each sub-population has its own LOSVD. In this experiment, we restrict ourselves to the case where the stellar populations have a known metallicity and are seen without extinction. This choice is mainly motivated by the numerical cost of such a large inversion procedure. The modeling is given by Eq. (10). Finding the age-velocity distribution $`\mathrm{\Xi }(u,t)`$ when the mono-metallic basis $`B`$ and the observed spectrum $`\varphi `$ are given is the inverse problem. It arises as the combination of a linear age inversion and a kinematical deconvolution.
### 5.1 A sum of convolutions
The age-velocity distribution, $`\mathrm{\Xi }(u,t)`$, is expanded as a linear combination of normalized 2D gate functions $`\theta _{ij}(u,t)`$:
$$\theta _{ij}(u,t)\frac{1}{\delta u\delta t}\theta \left(\frac{uu_i}{\delta u}\right)\theta \left(\frac{tt_j}{\delta t}\right).$$
In other words, $`\mathrm{\Xi }(u,t)`$ is represented by a 2D array $`𝐯`$ of size $`(p,n)`$, i.e. $`p`$ is the size of each LOSVD and $`n`$ is the number of age bins. The linear step in $`u`$ is $`\delta u`$ and the step in $`t`$ is $`\delta t`$.
By injecting the expansion into Eq. (10) we obtain
$`\varphi (w)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{p}{}}\underset{j=1}{\overset{n}{}}v_{ij}\theta _{ij}(u,t)B(wu,t)\mathrm{d}t\mathrm{d}u},`$ (28)
$``$ $`{\displaystyle \underset{i=1}{\overset{p}{}}}{\displaystyle \underset{j=1}{\overset{n}{}}}v_{ij}B_j(wu_i).`$
As in the previous sections, $`B_j(u)`$ is a time-averaged single stellar population of age $`t_j\pm \frac{1}{2}\delta t`$. We then discretize along wavelengths by averaging over small $`\delta w`$:
$`\varphi _k`$ $`=`$ $`{\displaystyle \frac{1}{\delta w}}{\displaystyle \underset{i=1}{\overset{p}{}}}{\displaystyle \underset{j=1}{\overset{n}{}}}v_{ij}{\displaystyle B_j(wu_i)\theta \left(\frac{ww_k}{\delta w}\right)dw},`$ (29)
$``$ $`{\displaystyle \underset{i=1}{\overset{p}{}}}{\displaystyle \underset{j=1}{\overset{n}{}}}v_{ij}B_j(w_ku_i),`$
where $`(w_j)_{j\{0,\mathrm{},m\}}`$ is a set of constant step logarithmic wavelengths. The above expression also reads in matrix form as a sum of kernel convolutions. Finally, the model SED of the emitted light reads:
$$𝐬=\underset{j=1}{\overset{n}{}}𝐊_j𝐯_j,$$
(30)
where $`𝐬=(\varphi _1,\varphi _2,\mathrm{},\varphi _m)`$, $`𝐯_j=(v_{1j},v_{2j},\mathrm{},v_{pj})`$ and
$$𝐊_j=\left[\begin{array}{cccc}\hfill K_{11j}& \hfill K_{12j}& \hfill \mathrm{}& \hfill K_{1pj}\\ \hfill K_{21j}& \hfill K_{22j}& \hfill \mathrm{}& \hfill K_{2pj}\\ \hfill \mathrm{}& \hfill \mathrm{}& \hfill \mathrm{}& \hfill \mathrm{}\\ \hfill K_{m1j}& \hfill K_{m2j}& \hfill \mathrm{}& \hfill K_{mpj}\end{array}\right],$$
(31)
with
$$K_{ikj}B_j(w_ku_i).$$
(32)
With this notation, $`𝐊_j`$ and $`𝐯_j`$ are respectively the convolution kernel and the LOSVD of the sub-population of age $`t_j`$, and the model spectrum $`𝐲`$ is the sum of the convolution of the kernel of each sub-population by its own LOSVD.
### 5.2 2D age-velocity smoothness constraints
In the previous sections, the unknowns were mono-dimensional functions of time or velocity. Here, the unknown is a 2D distribution, and we thus have to implement a 2D smoothing constraint. We wish to allow the smoothness in age to be distinct from the smoothness in velocity. We thus construct two penalizing functions, $`P_a`$ and $`P_v`$, relying on the standard function $`P`$. $`P_a`$ computes the sum of the Laplacians of the columns of $`𝐯`$ while $`P_v`$ computes the sum of the Laplacians of the lines of $`𝐯`$. The smoothness in the direction of the velocities (respectively ages) is set by $`\mu _v`$ (respectively $`\mu _a`$). We define the vectors $`𝐯_j=(v_{1j},v_{2j},\mathrm{},v_{pj})`$ as the columns of $`𝐯`$, i.e. the LOSVD s of the subpopulations. We similarly define the $`𝐯^i=(v_{i1},v_{i2},\mathrm{},v_{in})`$ as the lines of $`𝐯`$. With this notation, the penalization $`P_𝝁`$ reads:
$`P_𝝁(𝐯)`$ $``$ $`\mu _aP_a(𝐯)+\mu _vP_v(𝐯),`$ (33)
$``$ $`\mu _a{\displaystyle \underset{i=1}{\overset{p}{}}}P(𝐯^i)+\mu _v{\displaystyle \underset{j=1}{\overset{n}{}}}P(𝐯_j).`$
The objective function, $`Q_𝝁`$, is now fully specified as $`Q_𝝁=\chi ^2+P_𝝁`$. Its gradients are given in appendix A.3. We choose here the smoothing parameters, $`𝝁(\mu _a,\mu _v)`$, on the basis of simulations.
### 5.3 Simulations of a bulge-disk system
We studied the feasibility of separating two age-dynamically distinct populations, i.e. two components which do not overlap in an age-velocity distribution diagram, in a regime of very high quality model and data. We performed simulations in the idealized case of a very simplified spiral galaxy consisting of a bulge-disk system of solar metallicity seen without extinction at some intermediate inclination, in two observational contexts. The corresponding ages and projected kinematical parameters are given in Table 1. The resolution of the pseudo-data is $`R=\mathrm{10\hspace{0.17em}000}`$ at $`40006800`$ Å, and the SNR is $`100`$ per $`0.2`$ Å pixel.
* Case 1: The galaxy is resolved, and the fiber aperture is small compared to the angular size of the galaxy. The line of sight is offset by a couple kpc from the center along the major axis. The projected model age-velocity distribution involves 2 superimposed components: an old, non-rotating kinematically hot population representing the bulge, and a young, rotating, kinematically cold component. The model and the median of 30 reconstructions are shown in Fig. 5. The separation of the components is clear and their parameters can be recovered with good accuracy, considering the difficulty of the task.
* Case 2: The galaxy is unresolved. The difference with the former situation is that because of the spatial integration, both age-velocity distributions are centered. For a given dynamical model, the projected dispersion of the disk component depends on its inclination. Fig. 6 shows that the separation is successful and that the ages and integrated kinematical properties of both components can be measured.
## 6 Conclusions
The non-parametric kinematical deconvolution of a galaxy spectrum is efficiently performed using a MAP formalism (Sect. 3). Regularization through smoothness requirements and positivity improve significantly the behaviour of the inversion with respect to noise in the data. This improvement occurs at the cost of introducing some bias in the reconstructed LOSVD, but this bias remains reasonable. Strong non-Gaussianities of LOSVDs are reliably detected from mock data generated using Pégase-HR SSPs for SNR down to $`20`$ per $`0.2`$ Å pixel.
When the template does not exactly match the model spectrum at rest, i.e. there is some template mismatch, the error on the velocity dispersion increases very quickly (Sect. 3.4). For example, in our experiments, where $`\sigma _\mathrm{V}=50`$ km/s with $`R=\mathrm{10\hspace{0.17em}000}`$ data, the error on the measured velocity dispersion amounts up to 10-20 per cent if the template differs from the model by more than 0.3 dex in age and metallicity, perpendicular to the age-metallicity degeneracy.
The formal similarity between the non-parametric kinematical deconvolution and the recovery of the stellar content allows us to merge both processes in a “mixed” inversion where the observed spectrum is fitted by determining the stellar content and the kinematics simultaneously (Sect. 4). This circumvents the need for iterations where kinematical and stellar content analyses are carried out one after the other, until convergence is reached; this provides an efficient method to analyze large sets of data.
Satisfactory reconstructions of the stellar age distribution, the age-metallicity relation, the extinction and the global LOSVD were obtained from mock data down to $`R=2000`$, SNR$`=30`$ per 1 Å pixel in the $`\mathrm{4\hspace{0.17em}000}\mathrm{6\hspace{0.17em}800}`$ Å range (simulating SDSS data in the Pégase-HR range), indicating the good behaviour of the method. Since in our simulations, the introduction of the kinematics into STECMAP did not affect the recovery of the stellar content, we consider that the error estimates and separability analysis given in Paper I remain valid.
In a more exploratory part of this work, we showed the feasibility of recovering age-dependent kinematics in a simplified mono-metallic unreddened context (Sect. 5). We were able to separate the bulge and disc component of a simplified model spiral galaxy in integrated light provided very high quality data (SNR=100 per $`0.2`$ Å pixel in the optical domain) and models are available, i.e. we constrain both components in velocity dispersion and age. This separation was also carried out successfully in the setup corresponding to an unresolved galaxy.
Further investigations are needed to extend this technique to a regime where the metallicity and extinction are unknown. We expect that letting the metallicity be a free parameter would certainly lead to a more degenerate problem, as shows the degradation of the resolution in age found in Paper I compared to fixed metallicity problems. On the contrary, we do not expect the addition of the extinction as a free parameter or a more complex form of extiction law or flux calibration correction, possibly non parametric, to deteriorate the conditioning of the problem. The results are encouraging, and the feasibility of such age-dependent kinematics reconstructions deserves to be tackled in realistic specific pseudo-observational regimes in the future.
As mentioned in Paper I, the SSP models were considered to be perfect and noiseless. It still has to be investigated how instrumental error sources such as flux and wavelength calibration error, additive noise, contamination by adjacent objects, and, equally important, model errors, can affect the robustness of such sophisticated interpretations.
## Perspectives
STECKMAP will be very useful to interpret data of large spectroscopic surveys, complete or in progress, such as 2DFGRS<sup>1</sup><sup>1</sup>1http://www.mso.anu.edu.au/2dFGRS/, SDSS<sup>2</sup><sup>2</sup>2http://www.sdss.org/, DEEP2<sup>3</sup><sup>3</sup>3http://www.deep.berkeley.edu/, or VVDS<sup>4</sup><sup>4</sup>4http://www.oamp.fr/virmos/vvds.htm, especially where both constraints on the stellar content and the dynamics are required. STECKMAP’s analysis of the spectroscopic survey data or of a SNR selected subsample, combined with survey photometry could provide estimates of the stellar and dynamical masses (which must be corrected for fiber aperture though), thereby allowing astronomers the prospect of investigating the dark matter content in galaxies on a statistically significant sample, in the spirit of Padmanabhan et al. (2004).
The application of age-dependent kinematics to integral field spectroscopy data from, for example SAURON (Bacon et al., 2001; de Zeeuw et al., 2002), OASIS (McDermid et al., 2004), MUSE (Henault et al., 2003) or MPFS (Chilingarian et al., 2004) could significantly boost the amount of information extracted from this data.
The inner parts of elliptical or dwarf elliptical galaxies have shown via adaptive optics new kinematically decoupled structures (cores or central disks), which were precedently unresolved (McDermid et al., 2004; Bacon et al., 2001). Similarly, if decoupled structures are unresolved and remain so, even with adaptive optics, it may still be possible to separate components in age-velocity space. Hence, the technique presented in Sect. 5 extends the range of investigation for the inner components of galaxies even further in redshift and distance with the current generation of instruments. The faint, generalized counterparts of kinematically decoupled cores, i.e. stellar streams generated by minor merging and accretion of satellites, may also be detectable by an age-dependent kinematics reconstruction in systems which can not be resolved into stars, provided that they are sufficiently distinct from the bulk stars of the galaxy in the age-velocity space. This will enlarge the sample of galaxies for which such detailed information is available, and may make it statistically significant.
Acknowledgments
We are grateful to A. Siebert for useful comments and helpful suggestions. We would like to thank D. Munro for freely distributing his Yorick programming language (available at http://www.maumae.net/yorick/doc/index.html), together with its MPI interface, which we used to implement our algorithm in parallel. PO thanks the MPA for their hospitality and funding from a Marie Curie studentship.
## Appendix A Gradient computations
### A.1 Kinematic deconvolution
In this section we derive the gradient of $`Q_\mu `$ with respect to the LOSVD $`𝐠`$. First, we rewrite the $`\chi ^2`$ term as:
$$\chi ^2=𝐫^{}𝐖𝐫,$$
(34)
where the residuals vector $`𝐫`$ is defined by
$$𝐫=𝐲^1\mathrm{diag}(𝐅)𝐠,$$
(35)
The derivative of the $`\chi ^2`$ then reads:
$$\frac{\chi ^2}{𝐠}=2^{}\mathrm{diag}(𝐅)^{}𝐖𝐫,$$
(36)
where the asterisk $``$ denotes the complex conjugate. Since the stellar template and the LOSVD can play symmetrical roles in Eq. (17) we can also write the derivative of $`\chi ^2`$ relatively to the stellar template:
$$\frac{\chi ^2}{𝐅}=2^{}\mathrm{diag}(𝐠)^{}𝐖𝐫,$$
(37)
This expression will be useful for later derivations of gradients for more complex problems in the following appendices.
### A.2 Gradients of the mixed inversion
Here we show how to obtain the partial derivatives of $`Q_\mu =\chi ^2+P_\mu `$ as defined in Sect. 4. Given that writing the derivatives of the penalizing functions $`P_\mu `$ is straightforward, we will in this appendix focus on the gradients of the $`\chi ^2`$. In the mixed inversion, the reddened model spectrum at rest plays the role of the stellar template $`𝐅`$ in the classical kinematic deconvolution of equation (15). $`\chi ^2/𝐠`$ can thus be obtained by replacing $`𝐅\mathrm{diag}(𝐟_{\mathrm{ext}}(E))𝐁𝐱`$ in equation (36).
$$\frac{\chi ^2}{𝐠}=2^{}\mathrm{diag}(\mathrm{diag}(𝐟_{\mathrm{ext}}(E))𝐁𝐱)^{}𝐖𝐫,$$
(38)
where $`𝐫=𝐲𝐬`$ is the residuals vector, with $`𝐬`$ as given by Eq. (25). To obtain the other partial derivatives we use the following relation. For any parameter $`\alpha `$ we have
$$\frac{\chi ^2}{\alpha }=\left(\frac{\chi ^2}{𝐅}\right)^{}\frac{𝐅}{\alpha }.$$
(39)
The first term $`\chi ^2/𝐅`$ is given by Eq. (37) while the second term reads, considering each unknown:
$`{\displaystyle \frac{𝐅}{𝐱}}`$ $`=`$ $`\mathrm{diag}(𝐟_{\mathrm{ext}})𝐀,`$ (40)
$`{\displaystyle \frac{𝐅}{𝐙}}`$ $`=`$ $`\mathrm{diag}(𝐱){\displaystyle \frac{𝐁}{𝐙}}\mathrm{diag}(𝐟_{\mathrm{ext}}),`$ (41)
$`{\displaystyle \frac{𝐅}{E}}`$ $`=`$ $`\mathrm{diag}({\displaystyle \frac{𝐟_{\mathrm{ext}}}{E}})𝐀𝐱,`$ (42)
with the same notation as in the appendix of the STECMAP paper.
### A.3 Gradients for the age-dependent kinematics recovery
Again, we focus on the partial derivatives of the $`\chi ^2`$. Using Eq. (17), the model can be rewritten using the Fourier operator
$$𝐬=\underset{j=1}{\overset{n}{}}^{}\mathrm{diag}(𝐁_j)𝐯_j,$$
(43)
where $`𝐁_𝐣`$ is the discretized time-averaged single stellar population of age $`[t_{j1},t_j]`$, The derivatives of $`\chi ^2`$ relatively to $`𝐯`$ can be derived directly from Eq. (36) since the model is just a sum of convolutions. Replacing $`𝐅𝐁_j`$ and $`𝐠𝐯_j`$ yields the gradient of the $`\chi ^2`$:
$$\frac{\chi ^2}{𝐯_j}=2^{}\mathrm{diag}(𝐅)^{}𝐖𝐫,$$
(44)
with the residuals vector $`𝐫=𝐲𝐬`$. Finally, the derivative of $`Q_\mu `$ relatively to $`𝐯`$ is the matrix defined by
$$\frac{Q_\mu }{𝐯}=(\frac{Q_\mu }{𝐯_1},\frac{Q_\mu }{𝐯_2}\mathrm{}\frac{Q_\mu }{𝐯_n}).$$
(45) |
warning/0507/math0507102.html | ar5iv | text | # On the strong consistency of asymptotic M-estimators
## 1 Introduction
After the seminal work<sup>1</sup><sup>1</sup>1The interested reader may find a quite recent account in \[Ald97\] and references therein. of Fisher, the asymptotic properties of maximum likelihood estimators, and in particular their consistency, were studied by various authors, including Doob \[Doo34\], Cramér \[Cra46\], and Huzurbazar \[Huz48\]. Nowadays, one of the best known result regarding consistency goes back to Wald, who gave in \[Wal49\] a short and elegant proof of strong consistency of parametric maximum likelihood estimators. Since that time, several authors studied various versions of such consistency problems, including among others, Le Cam \[LC53\], Kiefer and Wolfowitz \[KW56\], Bahadur \[Bah67, Bah71\], Huber \[Hub67\], Perlman \[Per72\], Wang \[Wan85\], and Pfanzagl \[Pfa88, Pfa90\].
Wald’s original proof relies roughly on local compactness of the parameter space, on continuity and coercivity<sup>2</sup><sup>2</sup>2By coercivity we mean that the log-likelihood tends to $`\mathrm{}`$ when the parameter tends to $`\mathrm{}`$. of the log-likelihood, on the law of large numbers, and last but not least on local uniform integrability of the log-likelihood. It does not require differentiability, and makes extensive use of likelihood ratios. The integrability assumption has been weakened by many authors, including for instance Kiefer and Wolfowitz in \[KW56\] and Perlman in \[Per72\], see also \[Bah71\]. One can find a modern presentation of Wald’s method for $`M`$-estimators in van der Vaart’s monograph \[vdV98\].
Pfanzagl gave in \[Pfa88, Pfa90\] a proof of strong consistency of asymptotic maximum likelihood estimators for nonparametric “concave models” with respect to the estimated parameter, including nonparametric mixtures. His approach relies in particular on a simplification of an earlier work of Wang in \[Wan85\] based on uniform local bound of the likelihood ratio.
The present work was initially motivated by the inverse problems considered in \[CL06\]. Our aim is to simplify Pfanzagl’s approach, and to extend the framework from asymptotic maximum likelihood to more general asymptotic M-estimators. In particular, log-likelihood ratios are replaced by contrast differences. The hypotheses appearing in our main Theorem are unnecessarily strong. However, they allow a simple and short presentation. We emphasize the role played by a sort of contraction map $`a^{}`$ defined on the parameter space. We do not assume any coercivity of the contrast as in \[Wal49\]. However, we require the compactness of the space of the estimated parameter, as in \[KW56\] and \[vdV98\] for example. This compactness comes usually for free in the case of fully nonparametric models. We do not make use of any Uniform Law of Large Numbers. Our method does not belong to the Glivenko-Cantelli approaches of consistency, as in \[Dud98\], \[Fio00\], \[AK94\], \[vdV98\] and \[vdG03, vdG00\] and references therein.
Let $`\mathrm{\Theta }`$ be a separable Hausdorff topological space with countable base. Let $`(P_\theta )_{\theta \mathrm{\Theta }}`$ be a known family of Borel measures on a measurable space $`𝒳`$. Let $`\theta ^{}\mathrm{\Theta }`$ be some unknown point of $`\mathrm{\Theta }`$ such that $`P^{}:=P_\theta ^{}`$ is a probability measure. Let $`(X_n)_n`$ be an i.i.d. sequence of observed random variables defined on a probability space $`(\mathrm{\Omega },,)`$ and taking their values in $`𝒳`$, with common law $`P^{}`$. Let $`(\widehat{\theta }_n)_n`$ be a sequence of random variables defined on $`(\mathrm{\Omega },,)`$, taking their values in $`\mathrm{\Theta }`$, and such that $`(\widehat{\theta }_n)_n`$ is $`_n`$-measurable for any $`n`$, where $`_n:=\sigma (X_0,\mathrm{},X_n)`$. We say that $`(\widehat{\theta }_n)_n`$ is *strongly consistent* if and only if
$$\text{a.s. }\underset{n+\mathrm{}}{lim}\widehat{\theta }_n=\theta ^{}.$$
(1)
We use in the sequel the abbreviations “a.s.” for *almost sure*, “a.a.” for *almost all*, and “a.e.” for *almost everywhere*. Let $`\mathrm{\Theta }\times 𝒳(\theta ,x)m(\theta ,x)`$ be a known function such that $`m_\theta :=m(\theta ,)`$ is measurable for any $`\theta \mathrm{\Theta }`$. For any $`n`$, we define the random function $`M_n:\mathrm{\Theta }`$ by
$$M_n(\theta ):=\frac{1}{n}\underset{i=1}{\overset{n}{}}m(\theta ,X_i).$$
This can be written also $`M_n(\theta )=_nm_\theta `$ where $`_n:=\frac{1}{n}(\delta _{X_1}+\mathrm{}+\delta _{X_n})`$ is the empirical measure. We say that $`(\widehat{\theta }_n)_n`$ is a *sequence of asymptotic M-estimators* if and only if
$$\text{a.s. }\overline{\mathrm{lim}}_{n+\mathrm{}}\left(\underset{\mathrm{\Theta }}{sup}M_nM_n(\widehat{\theta }_n)\right)=0.$$
(2)
The term *asymptotic* is used for the same notion (with the likelihood) by Pfanzagl in \[Pfa88\]. In the literature, some authors, including Wald and Perlman, use the term *approximate* rather than *asymptotic*. However, the term *approximate* has been used by Bahadur in a different sense in \[Bah71, page 34\].
For example, if for large enough $`n`$, there exists an $`_n`$-measurable $`\widehat{\theta }_n`$ in $`\mathrm{\Theta }`$ such that $`M_n(\widehat{\theta _n})=sup_\mathrm{\Theta }M_n`$, then such a random sequence $`(\widehat{\theta }_n)_n`$ fulfils (2).
For any probability measure $`P`$ on $`𝒳`$, let $`\mathrm{L}_+^1(𝒳,P)`$ (resp. $`\mathrm{L}_{}^1(𝒳,P)`$) be the set of random variables $`Z:𝒳`$ such that $`Z^+:=\mathrm{max}(+Z,0)`$ (resp. $`Z^{}:=\mathrm{max}(Z,0)`$) is in $`\mathrm{L}^1(𝒳,P)`$. On $`E(𝒳,P):=\mathrm{L}_{}^1(𝒳,P)\mathrm{L}_+^1(𝒳,P)`$, the expectation $`P(Z)=P(Z^+)P(Z^{})`$ makes sense and takes its values in $`\overline{}:=\{\pm \mathrm{}\}`$. For any $`\theta \mathrm{\Theta }`$ such that $`m_\theta E(𝒳,P^{})`$, we define the *contrast* $`M^{}(\theta )\overline{}`$ by
$$M^{}(\theta ):=P^{}m_\theta .$$
(3)
In the sequel, we say that the model is *identifiable* when for any $`\theta \mathrm{\Theta }`$, the condition $`P_\theta =P^{}`$ implies that $`\theta =\theta ^{}`$.
###### Example 1.1 (Log-Likelihood).
Assume that for some fixed Borel measure $`Q`$ on $`𝒳`$, one has $`P_\theta Q`$ for any $`\theta \mathrm{\Theta }`$. Let $`f_\theta :=dP_\theta /dQ`$ and assume that $`f_\theta >0`$ on $`𝒳`$ for any $`\theta \mathrm{\Theta }`$. Define $`m(\theta ,x):=\mathrm{log}(f_\theta (x))`$. Then $`M_n:\mathrm{\Theta }`$ is the *log-likelihood* random functional given by $`M_n(\theta )=_nm_\theta =_n\mathrm{log}(f_\theta )`$. We will speak about sequences of “asymptotic maximum likelihood estimators”. The log-likelihood ratio is
$$M_n(\theta _1)M_n(\theta _2)=_n\mathrm{log}(f_{\theta _1}/f_{\theta _2}).$$
As usual for the log-likelihood, when $`M^{}(\theta ^{})`$ is finite, one can write for any $`\theta `$
$$M^{}(\theta )M^{}(\theta ^{})=\mathrm{𝐄𝐧𝐭}\left(P_\theta ^{}|P_\theta \right),$$
where $`\mathrm{𝐄𝐧𝐭}\left(P_{\theta _1}|P_{\theta _2}\right)`$ is the Kullback-Leibler relative entropy of $`P_{\theta _1}`$ with respect to $`P_{\theta _2}`$. In particular, $`M^{}(\theta )M^{}(\theta ^{})`$ with equality if and only if $`P_\theta =P_\theta ^{}`$, which implies $`\theta =\theta ^{}`$ if the model is identifiable. Notice that when $`Q`$ is the Lebesgue measure on $`𝒳=^n`$, then $`M^{}(\theta ^{})=_𝒳f_\theta ^{}(x)\mathrm{log}(f_\theta ^{}(x))𝑑x`$ is the Shannon entropy of $`f_\theta ^{}`$.
###### Example 1.2 (Beyond the log-likelihood).
Assume that for some fixed Borel measure $`Q`$ on $`𝒳`$, one has $`P_\theta Q`$ for any $`\theta \mathrm{\Theta }`$, with $`P_\theta (𝒳)1`$ and $`f_\theta :=dP_\theta /dQ`$. Let $`\mathrm{\Phi },\mathrm{\Psi }:(0,+\mathrm{})`$ be two smooth functions. Assume that $`\mathrm{\Psi }(f_\theta )\mathrm{L}^1(𝒳,Q)`$ for any $`\theta \mathrm{\Theta }`$. Define $`m_\theta `$ by
$$m_\theta =\mathrm{\Phi }(f_\theta )_𝒳\mathrm{\Psi }(f_\theta )𝑑Q+P_\theta (𝒳).$$
This gives rise the the following empirical contrast
$$M_n(\theta )=_n(\mathrm{\Phi }(f_\theta ))_𝒳\mathrm{\Psi }(f_\theta )𝑑Q+P_\theta (𝒳).$$
In particular, if $`\theta \mathrm{\Theta }`$ is such that $`\mathrm{\Phi }(f_\theta )\mathrm{L}^1(𝒳,P^{})`$ where here again $`P^{}:=P_\theta ^{}`$,
$$M^{}(\theta )=P^{}(\mathrm{\Phi }(f_\theta ))_𝒳\mathrm{\Psi }(f_\theta )𝑑Q+P_\theta (𝒳).$$
Assume now that $`uu\mathrm{\Phi }^{}(u)`$ is locally integrable on $`_+`$, and consider the case where $`\mathrm{\Psi }`$ is the $`\mathrm{\Phi }`$-transform given for any $`u(0,+\mathrm{})`$ by
$$\mathrm{\Psi }(u)=_0^uv\mathrm{\Phi }^{}(v)𝑑v.$$
For $`\mathrm{\Phi }:u\mathrm{log}(u)`$, one has $`\mathrm{\Psi }:uu`$ and we recover the log-likelihood contrast
$$M^{}(\theta )=P^{}(\mathrm{log}(f_\theta )).$$
For $`\mathrm{\Phi }:uu`$, one has $`\mathrm{\Psi }:u\frac{1}{2}u^2`$, and we get the quadratic contrast
$$M^{}(\theta )=\frac{1}{2}f_\theta f_\theta ^{}_{\mathrm{L}^2(𝒳,Q)}^2+\frac{1}{2}f_\theta ^{}_{\mathrm{L}^2(𝒳,Q)}^2+P_\theta (𝒳).$$
In both cases, the map $`\theta M^{}(\theta )`$ admits $`\theta ^{}`$ as unique maximum provided that the model is identifiable. More generally, define the $`\mathrm{\Phi }`$-transform $`\mathrm{\Theta }:(0,+\mathrm{})^2`$ by
$`\mathrm{\Theta }(u,v):`$ $`=u\mathrm{\Phi }(v)\mathrm{\Psi }(v)`$
$`=u\mathrm{\Phi }(v){\displaystyle _0^v}w\mathrm{\Phi }^{}(w)𝑑w.`$
When $`\theta `$ and $`\theta ^{}`$ are such that both $`\mathrm{\Theta }(f_\theta ^{},f_\theta ^{})`$ and $`\mathrm{\Theta }(f_\theta ^{},f_\theta )`$ belong to $`\mathrm{L}^1(𝒳,Q)`$,
$$M^{}(\theta )=_𝒳(\mathrm{\Theta }(f_\theta ^{},f_\theta )\mathrm{\Theta }(f_\theta ^{},f_\theta ^{}))𝑑Q+_𝒳\mathrm{\Theta }(f_\theta ^{},f_\theta ^{})𝑑Q+P_\theta (𝒳).$$
Notice that $`\mathrm{\Theta }`$ is linear in $`\mathrm{\Phi }`$. One can consider useful examples for which the function $`\mathrm{\Phi }`$ is bounded, in such a way that $`m_\theta `$ is bounded for any $`\theta \mathrm{\Theta }`$. For instance, let us examine the case where $`\mathrm{\Phi }:u(1+u)^2`$. Then, $`\mathrm{\Psi }:uu^2(1+u)^2`$, and the map $`\theta M^{}(\theta )`$ admits $`\theta ^{}`$ as unique maximum, provided identifiability holds, since for any $`(u,v)_+^2`$,
$$\mathrm{\Theta }(u,v)=\frac{u+v^2}{(1+v)^2}\text{ and }\mathrm{\Theta }(u,v)\mathrm{\Theta }(u,u)=\frac{(vu)^2}{(1+u)(1+v)^2}.$$
The function $`\mathrm{\Psi }`$ is additionally bounded here. The similar case $`\mathrm{\Phi }:u(1+u^2)^1`$ is also quite interesting. Notice that $`\mathrm{\Theta }(u,)`$ is concave on $`(0,+\mathrm{})`$ as soon as $`\mathrm{\Phi }`$ is concave, non decreasing, with $`\mathrm{\Phi }^{}(v)+v\mathrm{\Phi }^{\prime \prime }(v)0`$ for any $`v>0`$. Observe that this is not the approach of Pfanzagl in \[Pfa90\], which is more related to the log-likelihood ratio. Notice that in the case of the log-likelihood, one has $`\mathrm{\Phi }:u\mathrm{log}(u)`$, which gives $`\mathrm{\Psi }:uu`$ and $`\mathrm{\Theta }:(u,v)u\mathrm{log}(v)v`$, and thus $`\mathrm{\Theta }(u,v)\mathrm{\Theta }(u,u)=u\mathrm{log}(u/v)+uv`$. It might be possible to extensively study such “$`\mathrm{\Phi }`$-estimators”, in the spirit of the “$`\mathrm{\Phi }`$-calculus” developed in \[Cha04, Cha06\]. This is however outside the scope of this short article.
One can notice that the observation of Lindsay in \[Lin83a, Lin83b\] regarding the nature of maximum likelihood for nonparametric mixture models remains valid for more general models provided that $`m`$ is concave.
## 2 Main result and Corollaries
With the settings given in the Introduction, the following Theorem holds.
###### Theorem 2.1.
Assume that $`\mathrm{\Theta }`$ is compact and that the following assumptions hold.
1. For $`P^{}`$-a.a. $`x𝒳`$, the map $`m(,x)`$ is continuous on $`\mathrm{\Theta }`$;
2. There exists a continuous map $`a^{}:\mathrm{\Theta }\mathrm{\Theta }`$ which may depend on $`\theta ^{}`$ such that for any $`\theta \theta ^{}`$, there exists a neighborhood $`V\mathrm{\Theta }`$ of $`\theta `$ for which $`sup_V\left(mm_a^{}\right)\mathrm{L}_+^1(𝒳,P^{})`$ and $`P^{}(m_\theta m_{a^{}(\theta )})<0`$.
Then any sequence $`(\widehat{\theta }_n)_n`$ of asymptotic M-estimators is strongly consistent.
###### Proof.
Postponed to section 4. ∎
The quantity $`P^{}(m_\theta m_{a^{}(\theta )})`$ in (A2) has a meaning in $`\overline{}`$ since the first part of (A2) ensures that $`m_\theta m_{a^{}(\theta )}\mathrm{L}_+^1(𝒳,P^{})`$. Moreover, $`P^{}(m_\theta m_{a^{}(\theta )})`$ reads $`M^{}(\theta )M^{}(a^{}(\theta ))`$ when the couple $`(m_\theta ,m_{a^{}(\theta )})`$ is in $`\mathrm{L}_{}^1(𝒳,P^{})\times \mathrm{L}_+^1(𝒳,P^{})`$ or in $`\mathrm{L}_+^1(𝒳,P^{})\times \mathrm{L}_{}^1(𝒳,P^{})`$.
Since $`\theta ^{}`$ is unknown in practice, each assumption in Theorem 2.1 must hold for any $`\theta ^{}\mathrm{\Theta }`$ such that $`P_\theta ^{}`$ is a probability measure, in order to make the result useful.
###### Remark 2.2 (Assumptions).
The first part of (A2) is in a way an $`M`$-estimator version of the integrability condition considered by Kiefer and Wolfowitz for the log-likelihood in \[KW56\]. The assumptions (A1) and (A2) required by Theorem 2.1 can be weakened. However, they permit a streamlined presentation. In particular, only lower semi-continuity is needed in (A1), see for instance \[Pfa88\]. Additionally, and following for example \[Per72, page 266\], the uniform integrability assumption (A2) can be weakened, by considering blocks of $`k>1`$ observations instead of one observation, see also \[vdV98, comments following Theorem 5.14\].
As stated in the following Corollary, Theorem 2.1 implies a version of Wald consistency Theorem for asymptotic $`M`$-estimators, see \[Wal49\], \[Per72, Section 2 page 269\], and \[vdV98, Theorem 5.14\].
###### Corollary 2.3 (Perlman-Wald).
Assume that $`\mathrm{\Theta }`$ is compact, and that for $`P^{}`$-a.a. $`x𝒳`$, the map $`m(,x)`$ is continuous on $`\mathrm{\Theta }`$. Assume that for any $`\theta `$ in $`\mathrm{\Theta }`$, there exists a neighborhood $`V`$ such that $`sup_Vm\mathrm{L}^1(𝒳,P^{})`$. Assume in addition that $`M^{}`$ achieves its supremum over $`\mathrm{\Theta }`$ at $`\theta ^{}`$, and only at $`\theta ^{}`$. Then, any sequence of asymptotic $`M`$-estimators is strongly consistent.
###### Proof.
One has $`m_\theta \mathrm{L}^1(𝒳,P^{})`$ for any $`\theta `$ in $`\mathrm{\Theta }`$, and thus $`M^{}:\mathrm{\Theta }`$ is well defined. Moreover, (A2) holds with a constant map $`a^{}\theta ^{}`$. Namely, for any $`\theta \theta ^{}`$, one has on one hand $`P^{}(m_\theta m_\theta ^{})<0`$ since $`M^{}(\theta )<M^{}(\theta ^{})`$, and on the other hand
$$\underset{V}{sup}(mm_a^{})=m_\theta ^{}+\underset{V}{sup}m\mathrm{L}^1(𝒳,P^{}).$$
As stated in the following Corollary, Theorem 2.1 implies the main result of Pfanzagl in \[Pfa88\] for concave models, itself based on an earlier result of Wang in \[Wan85\]. This is typically the case for mixtures models, for which $`\mathrm{\Theta }`$ is a convex set of probability measures on some measurable space, cf. section 3.
###### Corollary 2.4 (Pfanzagl-Wang).
Let $`Q`$ be a reference Borel measure on $`𝒳`$. Consider the case where $`\mathrm{\Theta }`$ is a convex compact subset of a linear space such that for any $`\theta \mathrm{\Theta }`$, $`P_\theta (𝒳)1`$ and $`P_\theta Q`$ with $`f_\theta :=dP_\theta /dQ>0`$ on $`𝒳`$. Suppose that $`Q`$-a.e. on $`𝒳`$, the map $`\theta f_\theta (x)`$ is concave and continuous on $`\mathrm{\Theta }`$. Assume that the model is identifiable. Consider $`m_\theta :=\mathrm{log}(f_\theta )`$ and the related log-likelihood $`M_n`$. Then any sequence of asymptotic log-likelihood estimators is strongly consistent.
###### Proof.
First of all, we notice that it is not possible to take $`a^{}\theta ^{}`$ since we cannot ensure that the condition $`m_\theta ^{}m_\theta =\mathrm{log}(f_\theta ^{}/f_\theta )\mathrm{L}_+^1(𝒳,P^{})`$ of (A2) is true. However, the concavity of the model allows to take a map $`a^{}`$ which is a strict contraction around $`\theta ^{}`$. Namely, for an arbitrary $`\lambda (0,1)`$, let us take
$$a^{}(\theta ):=\lambda \theta ^{}+(1\lambda )\theta .$$
The concavity of the model yields
$$m_{a^{}(\theta )}m_\theta =\mathrm{log}\left(\frac{f_{\lambda \theta ^{}+(1\lambda )\theta }}{f_\theta }\right)\mathrm{log}\left(\frac{\lambda f_\theta ^{}+(1\lambda )f_\theta }{f_\theta }\right)\mathrm{log}(1\lambda ).$$
Now, we have $`\mathrm{log}(1\lambda )\mathrm{L}^1(𝒳,P^{})`$ since $`\lambda <1`$. Define the function $`\mathrm{\Phi }:_+`$ by $`\mathrm{\Phi }(u):=u\mathrm{log}(\lambda u+(1\lambda ))`$. The concavity of the model yields
$$P^{}(m_{a^{}(\theta )}m_\theta )_𝒳f_\theta ^{}\mathrm{log}\left(\frac{\lambda f_\theta ^{}+(1\lambda )f_\theta }{f_\theta }\right)𝑑Q=_𝒳\mathrm{\Phi }\left(\frac{f_\theta ^{}}{f_\theta }\right)f_\theta 𝑑Q.$$
Let us show that the right hand side of the inequality above is strictly positive when $`\theta \theta ^{}`$. One has $`P_\theta (𝒳)>0`$ since $`f_\theta >0`$. Define $`\mathrm{\Psi }(u):=u\mathrm{\Phi }(1/u)`$. Jensen’s inequality for the probability measure $`P_\theta (𝒳)^1P_\theta `$ and the convex function $`\mathrm{\Phi }`$ yields
$$_𝒳\mathrm{\Phi }\left(\frac{f_\theta ^{}}{f_\theta }\right)f_\theta 𝑑Q\mathrm{\Psi }(P_\theta (𝒳)).$$
(4)
It is enough to show that either (4) is strict or the right hand side of (4) is strictly positive. Since $`\lambda >0`$, the function $`\mathrm{\Phi }`$ is strictly convex. Thus equality holds in (4) if and only if $`P_\theta (f_\theta ^{}=\alpha f_\theta )=1`$ for some $`\alpha _+`$. The only admissible case is $`\alpha =P_\theta (𝒳)^1>1`$ since $`P_\theta ^{}(𝒳)=1`$ and since identifiability forbids $`P_\theta (f_\theta ^{}=f_\theta )=1`$. Therefore, if $`P_\theta (𝒳)=1`$, inequality (4) is necessarily strict. On the other hand, $`\mathrm{\Psi }(1)=0`$ and $`\mathrm{\Psi }(u)>0`$ when $`u<1`$. Thus the right hand side of (4) is always non negative, and is strictly positive as soon as $`P_\theta (𝒳)<1`$. We conclude that $`P^{}(m_{a^{}(\theta )}m_\theta )>0`$ as soon as $`\theta \theta ^{}`$. This shows that (A2) holds with $`V=\mathrm{\Theta }`$, and the proof is thus complete. ∎
###### Remark 2.5 (About the map $`a^{}`$).
Let $`a^{}:\mathrm{\Theta }\mathrm{\Theta }`$ be a map which satisfies the condition $`P^{}(m_\theta m_{a(\theta )})<0`$ for any $`\theta \theta ^{}`$ of (A2). Then, the impossibility of $`P^{}(m_\theta m_\theta )<0`$ for any $`\theta `$ yields that
* $`a^{}(\theta )\theta `$ for any $`\theta \theta ^{}`$. In particular,
+ the map $`a^{}`$ cannot be the identity map ;
+ if $`a^{}`$ is constant, then $`a^{}\theta ^{}`$ ;
+ the point $`\theta ^{}`$ is the only possible fixed point for $`a^{}`$.
The proof of Corollary 2.3 gives an example where $`a^{}\theta ^{}`$ works and fulfills (A2). In contrast, Corollary 2.4 provides a situation where a constant $`a^{}`$ does not fulfill (A2). However, we have shown in the proof of Corollary 2.4 that an $`a^{}`$ map which is a strict contraction around $`\theta ^{}`$ fulfills (A2). Actually, when $`\mathrm{\Theta }`$ has the structure of a convex subset of a vector space, any strict contraction around $`\theta ^{}`$ fulfills the properties of $`a^{}`$ listed above. The existence of a fixed point can be related to Brouwer-like fixed point Theorems. For instance, any continuous mapping of a non-empty compact convex subset of $`^d`$ into itself contains at least one fixed point. Consequently, when $`\mathrm{\Theta }`$ is a non-empty compact and convex subset of $`^d`$, any continuous $`a^{}`$ map admits $`\theta ^{}`$ as a unique fixed point. There exists numerous dimension free Brouwer-like fixed points theorems, due to Schauder, Tikhonov, Kakutani, …, see for instance \[Zei86\] and \[Goe02\].
###### Remark 2.6 (Infinite values of $`m`$).
Theorem 2.1 does not allow $`m`$ to take the value $`\mathrm{}`$. This limitation is due to the fact that differences of the form $`m_\theta m_\theta ^{}`$ do not make sense if $`m`$ is allowed to take the value $`\mathrm{}`$. The consistency proof of Wald does not suffer from such a limitation since it does not rely on $`m`$ differences, but it requires however strong uniform integrability assumptions. A careful reading of the proof of Theorem 2.1 shows that only differences of the form $`m_\theta m_{a^{}(\theta )}`$ are involved. On the other hand, according to Remark 2.5, $`a^{}(\theta )\theta `$ for any $`\theta \theta ^{}`$. Consequently, one may allow, in Theorem 2.1, the map $`m(\theta ,x)`$ to take the value $`\mathrm{}`$ for at most one value of $`\theta `$. For the log-likelihood, $`m_\theta =\mathrm{log}(f_\theta )`$ and one has $`m_\theta (x)=\mathrm{}`$ if and only if $`f_\theta (x)=0`$. One may allow $`f_\theta 0`$ for at most one value of $`\theta `$ in Corollary 2.4.
###### Remark 2.7.
Let $`\theta \mathrm{\Theta }`$ such that $`m_\theta E(𝒳,P^{})`$. Then, the law of large numbers applies and gives that $`P^{}`$-a.s., $`lim_nM_n(\theta )=M^{}(\theta )\overline{}`$, and the a.s. subset of $`𝒳`$ may depend on $`\theta `$. In particular $`M_n(\theta )=M^{}(\theta )+o_P(1)`$. For a sequence $`(\widehat{\theta }_n)_n`$ satisfying (2), one can write for any $`\theta \mathrm{\Theta }`$ with finite $`M_n(\theta )`$
$`M_n(\widehat{\theta }_n)`$ $`=M_n(\widehat{\theta }_n)M_n(\theta )+M_n(\theta )`$
$`\left(\underset{\mathrm{\Theta }}{sup}M_nM_n(\widehat{\theta }_n)\right)+M_n(\theta )`$
$`=o_P(1)+M(\theta )`$
where the last step follows by (2) and the law of large numbers.
## 3 Log-Likelihood and mixtures models
For any topological space $`𝒵`$ equipped with its Borel $`\sigma `$-field, we denote by $`_1(𝒵)`$ the set of probability measures on $`𝒵`$, and by $`𝒞_b(𝒵)`$ the set of bounded real valued continuous functions on $`𝒵`$. The Prohorov topology on $`_1(𝒵)`$ is defined as follows: $`\theta _n\theta `$ in $`_1(𝒵)`$ if and only if $`_𝒵f𝑑\theta _n_𝒵f𝑑\theta `$ for any $`f𝒞_b(𝒵)`$. It is known that a subset of $`_1(𝒵)`$ is compact if and only if it is tight. As a consequence, $`_1(𝒵)`$ is not compact in general. Following \[Pfa88, section 5 page 149\], the set sub-probabilities provides a compactification which allows the following consistency result for asymptotic log-likelihood estimators of nonparametric mixture models.
###### Corollary 3.1 (Pfanzagl).
Let $`𝒵`$ be a locally compact Hausdorff topological space with countable base. Let $`Q`$ be a measure on a measurable space $`𝒳`$. Let $`k:𝒳\times 𝒵(0,+\mathrm{})`$ be such that $`k(x,z)𝑑Q(x)=1`$ for any $`z𝒵`$ and $`k(x,)𝒞_b(𝒵)`$ for any $`x𝒳`$. Let $`\mathrm{\Theta }:=_1(𝒵)`$ and consider the family $`(P_\theta )_{\theta \mathrm{\Theta }}`$ of probability measures on $`𝒳`$ defined by $`dP_\theta =f_\theta dQ`$ with $`f_\theta (x):=k(x,z)𝑑\theta (z)`$. Assume that the model is identifiable. Let $`m:\mathrm{\Theta }\times 𝒳`$ be the map defined by $`m(\theta ,x):=\mathrm{log}f_\theta (x)`$, and $`M_n`$ be the corresponding log-likelihood. Then any sequence of asymptotic maximum likelihood estimators is strongly consistent for the Prohorov topology.
###### Proof.
As explained above, $`\mathrm{\Theta }=_1(𝒵)`$ is not compact for the Prohorov topology, and one must consider a suitable compactification, as in \[Bah71\] for instance. Let $`𝒞_0(𝒵)`$ be the set of real valued continuous functions on $`𝒵`$ which vanish at infinity. Let $`\overline{\mathrm{\Theta }}`$ be the set of Borel measures $`\theta `$ on $`𝒵`$ such that $`\theta (𝒵)1`$ (i.e. sub-probabilities), equipped with the vague topology related to $`𝒞_0(𝒵)`$. Namely, $`\theta _n\theta `$ in $`\overline{\mathrm{\Theta }}`$ if and only if $`_𝒵f𝑑\theta _n_𝒵f𝑑\theta `$ for any $`f𝒞_0(𝒵)`$. The injection $`\mathrm{\Theta }\overline{\mathrm{\Theta }}`$ is continuous; $`\overline{\mathrm{\Theta }}`$ is a compact metrizable topological space, and thus has a countable base. Moreover, $`\overline{\mathrm{\Theta }}`$ is convex, and for any $`\theta \overline{\mathrm{\Theta }}`$, there exists $`\theta ^{}\mathrm{\Theta }`$ and $`\alpha [0,1]`$ such that $`\theta =\alpha \theta ^{}`$.
We extend the set of probability measures $`(P_\theta )_{\theta \mathrm{\Theta }}`$ on $`𝒳`$ to the set of sub-probability measures $`(P_\theta )_{\theta \overline{\mathrm{\Theta }}}`$ on $`𝒳`$, where $`dP_\theta =f_\theta dQ`$ and $`f_\theta (x):=k(x,z)𝑑\theta (z)`$. One has by virtue of Fubini-Tonelli Theorem that $`P_\theta (𝒳)=\theta (𝒵)`$, and thus $`P_\theta _1(𝒳)`$ if and only if $`\theta \mathrm{\Theta }:=_1(𝒵)`$. Notice that $`\theta ^{}`$ is taken in $`\mathrm{\Theta }`$.
Let $`\theta \overline{\mathrm{\Theta }}`$ such that $`P_\theta =P_\theta ^{}`$. Since $`\theta ^{}`$ is taken in $`\mathrm{\Theta }`$, one has that $`P_\theta _1(𝒳)`$, therefore $`\theta \mathrm{\Theta }`$ and thus $`\theta =\theta ^{}`$ by identifiability in $`\mathrm{\Theta }`$. Notice that $`\overline{\mathrm{\Theta }}`$ is the convex envelope of $`\mathrm{\Theta }\{0\}`$. The set $`\overline{\mathrm{\Theta }}`$ contains the null measure $`0`$, for which $`f_00`$ and thus $`m_0\mathrm{}`$. If $`\theta \overline{\mathrm{\Theta }}`$ with $`\theta 0`$, then $`f_\theta >0`$ on $`𝒳`$ since $`k>0`$, and thus $`m_\theta (x):=\mathrm{log}f_\theta (x)`$ is finite for any $`x𝒳`$. For any $`x𝒳`$, the map $`\theta \overline{\mathrm{\Theta }}m_\theta (x)`$ is continuous since $`k(x,)`$ is in $`𝒞_0(𝒵)`$.
For any $`\theta \overline{\mathrm{\Theta }}`$ with $`\theta 0`$, one can write $`\theta =\alpha \theta ^{}`$ with $`\theta ^{}\mathrm{\Theta }`$ and $`\alpha :=\theta (𝒵)[0,1]`$. One has then $`f_\theta =\alpha f_\theta ^{}`$ and thus $`m_\theta =\mathrm{log}\alpha +m_\theta ^{}`$. Therefore,
$$M_n(\theta )=\mathrm{log}\alpha +M_n(\theta ^{})M_n(\theta ^{}).$$
As a consequence, $`sup_{\theta \mathrm{\Theta }}M_n(\theta )=sup_{\theta \overline{\mathrm{\Theta }}}M_n(\theta )`$, and one may substitute $`\mathrm{\Theta }`$ by $`\overline{\mathrm{\Theta }}`$ in the definition (2). Now, let $`(\widehat{\theta }_n)_n`$ be a sequence in $`\mathrm{\Theta }`$ of asymptotic maximum likelihood estimators. Corollary 2.4 and Remark 2.6 for $`(P_\theta )_{\theta \overline{\mathrm{\Theta }}}`$ apply and give the $`P^{}`$-a.s. convergence for the vague topology of $`(\widehat{\theta }_n)_n`$ towards $`\theta ^{}`$. Since both the sequence and the limit are in $`\mathrm{\Theta }`$, the convergence holds for the Prohorov topology, and the desired result is established. ∎
###### Remark 3.2.
A mixture model can always be seen as a conditional model. The observed random variables $`X`$ with values in $`𝒳`$ is the first component of the couple $`(X,Z)`$ with values in $`𝒳\times 𝒵`$. The component $`Z`$ is not observed. However, the conditional law $`(X|Z=z)`$ is known, and has density $`k(,z)`$ with respect to $`Q`$ on $`𝒳`$. If $`\theta =(Z)`$, then $`(X)`$ has density $`f_\theta `$ with respect to $`Q`$ on $`𝒳`$.
## 4 Proof of main result
###### Lemma 4.1 (Reformulation).
The random sequence $`(\widehat{\theta }_n)_n`$ is a sequence of asymptotic M-estimators if and only if
$$\text{–a.s.},(\theta _n)_n\mathrm{\Theta }^{},\overline{\mathrm{lim}}_{n+\mathrm{}}\left(M_n(\theta _n)M_n(\widehat{\theta }_n)\right)0.$$
(5)
###### Proof.
The proof is done “$`\omega `$ by $`\omega `$”, and the a.s. sets in (2) and (5) are the same. Recall that $`(\widehat{\theta }_n)_n`$ is a sequence of asymptotic M-estimators if and only if (2) holds. Actually, the definition of the supremum gives $`sup_{\theta \mathrm{\Theta }}M_n(\theta )M_n(\widehat{\theta }_n)0`$. Therefore, (2) is equivalent to
$$\text{–a.s.},\overline{\mathrm{lim}}_{n+\mathrm{}}\left(\underset{\theta \mathrm{\Theta }}{sup}M_n(\theta )M_n(\widehat{\theta }_n)\right)0.$$
(6)
The Lemma is thus reduced to the equivalence between (6) and (5). We begin by the proof of the implication (6) $``$ (5). Let $`A`$ be some $``$–a.s. set such that (6) holds. We proceed by fixing $`\omega A`$. We hide the dependency on $`\omega `$ in the notation of $`M_n`$ and $`\widehat{\theta }_n`$ to lightweight the expressions. Let $`(\theta _n)_n`$ be a sequence in $`\mathrm{\Theta }`$. By definition of the supremum, we have $`M_n(\theta _n)sup_{\theta \mathrm{\Theta }}M_n(\theta )`$. Thus, we get
$$M_n(\theta _n)M_n(\widehat{\theta }_n)\underset{\theta \mathrm{\Theta }}{sup}M_n(\theta )M_n(\widehat{\theta }_n).$$
Taking the $`\overline{\mathrm{lim}}_{n+\mathrm{}}`$ of both sides and using (6) provides the expected result (5). It remains to establish the implication (5) $``$ (6). Let $`A`$ be some $``$–a.s. set such that (5) holds. Here again, we proceed by fixing $`\omega A`$, and we hide the dependency on $`\omega `$ in the notation of the random objects like $`M_n`$ and $`\widehat{\theta }_n`$. By definition of the supremum, there exists, for any $`n`$, an element $`\theta _n\mathrm{\Theta }`$ such that
$$\underset{\theta \mathrm{\Theta }}{sup}M_n(\theta )M_n(\theta _n)\frac{1}{n}0.$$
Notice that $`\theta _n`$ depends on $`\omega `$ since $`M_n`$ depends on $`\omega `$. This yields
$$\overline{\mathrm{lim}}_{n+\mathrm{}}\left(\underset{\theta \mathrm{\Theta }}{sup}M_n(\theta )M_n(\theta _n)\right)0.$$
(7)
Now we write the telescopic sum
$$\underset{\theta \mathrm{\Theta }}{sup}M_n(\theta )M_n(\widehat{\theta }_n)=\underset{\theta \mathrm{\Theta }}{sup}M_n(\theta )M_n(\theta _n)+M_n(\theta _n)M_n(\widehat{\theta }_n),$$
which gives
$$\begin{array}{c}\overline{\mathrm{lim}}_{n+\mathrm{}}\left(\underset{\theta \mathrm{\Theta }}{sup}M_n(\theta )M_n(\widehat{\theta }_n)\right)\hfill \\ \hfill \overline{\mathrm{lim}}_{n+\mathrm{}}\left(\underset{\theta \mathrm{\Theta }}{sup}M_n(\theta )M_n(\theta _n)\right)+\overline{\mathrm{lim}}_{n+\mathrm{}}\left(M_n(\theta _n)M_n(\widehat{\theta }_n)\right).\end{array}$$
The two terms of the right hand side are “$`0`$” by virtue of (7) and (5) respectively. This provides the desired result (6), as expected. ∎
###### Lemma 4.2 (Separation).
Assume that $``$–a.s., for any neighborhood $`U`$ of $`\theta ^{}`$, for any sequence $`(\theta _n)_n`$ in $`U^c`$, there exists a sequence $`(\theta _n^{})_n`$ in $`\mathrm{\Theta }`$ such that
$$\underset{¯}{\mathrm{lim}}_{n+\mathrm{}}\left(M_n(\theta _n^{})M_n(\theta _n)\right)>0.$$
(8)
Then, any asymptotic M-estimators sequence $`(\widehat{\theta }_n)_n`$ is strongly consistent.
###### Proof.
Suppose that (8) holds for some a.s. set $`A`$, and that $`(\widehat{\theta }_n)_n`$ is a sequence of asymptotic M-estimators which is not strongly consistent. Saying that $`(\widehat{\theta }_n)_n`$ is not strongly consistent means that for any $``$–a.s. set, there exists a neighborhood $`U`$ of $`\theta ^{}`$ and a subsequence $`(\widehat{\theta }_{n_k})_k`$ in $`U^c`$. In particular, on the a.s. set $`A`$, this gives a neighborhood $`U`$ of $`\theta ^{}`$ and a subsequence $`(\widehat{\theta }_{n_k})_k`$ in $`U^c`$. Now, by virtue of (8),
$$\text{–a.s},(\theta _{n_k}^{})_k\mathrm{\Theta }^{},\underset{¯}{\mathrm{lim}}_{k+\mathrm{}}\left(M_{n_k}(\theta _{n_k}^{})M_{n_k}(\widehat{\theta }_{n_k})\right)>0,$$
where the a.s. set is $`A`$. This contradicts (5) which holds $``$–a.s. too. ∎
###### Lemma 4.3 (The $`a^{}`$ map).
Assume that $`\mathrm{\Theta }`$ is compact and that there exists a map $`a^{}:\mathrm{\Theta }\mathrm{\Theta }`$ such that for any $`\theta \theta ^{}`$, there exists a neighborhood $`U_\theta `$ of $`\theta `$ such that
$$\text{–a.s.},\underset{¯}{\mathrm{lim}}_{n+\mathrm{}}\underset{U_\theta }{inf}\left(M_n(a^{})M_n\right)>0.$$
(9)
Then, any asymptotic M-estimators sequence $`(\widehat{\theta }_n)_n`$ is strongly consistent.
###### Proof.
Let us show that the assumptions of Lemma 4.2 are fulfilled. We will establish (8) for an a.s. set $`A`$ which does not depend on the neighborhood $`U`$ of $`\theta ^{}`$. Namely, let $`U`$ be an open neighborhood of $`\theta ^{}`$. For any $`\theta U^c`$, let $`U_\theta `$ and $`A_\theta `$ be the neighborhood of $`\theta `$ and the $``$–a.s. set for which (9) holds. Notice that $`A_\theta `$ depends on $`U_\theta `$. The set $`U^c_{\theta U^c}U_\theta `$ is compact as a closed subset of the compact set $`\mathrm{\Theta }`$. We can thus extract a finite sub-covering $`U^c_{i=1}^kU_{\theta _i}`$, and write
$`\underset{¯}{\mathrm{lim}}_n\underset{U^c}{inf}\left(M_n(a^{})M_n\right)`$ $`\underset{¯}{\mathrm{lim}}_n\underset{1ik}{\mathrm{min}}\underset{U_{\theta _i}}{inf}\left(M_n(a^{})M_n\right)`$
$`=\underset{1ik}{\mathrm{min}}\underset{¯}{\mathrm{lim}}_n\underset{U_{\theta _i}}{inf}\left(M_n(a^{})M_n\right).`$
By virtue of (9) we get from the above that
$$\text{–a.s.},\underset{¯}{\mathrm{lim}}_n\underset{U^c}{inf}\left(M_n(a^{})M_n\right)>0,$$
(10)
where the $``$–a.s. set is $`A_U:=_{i=1}^kA_{\theta _i}`$. Recall that $`U`$ was a freely chosen neighborhood of $`\theta ^{}`$. Consider now a countable base $`(U_k)_k`$ for $`\theta ^{}`$. Then (10) holds on the $``$–a.s. set $`A:=_{i=1}^{\mathrm{}}A_{U_k}`$, which does not depend on $`U`$. Notice at this step that
$$M_n(a^{}(\theta _n))M_n(\theta _n)\underset{U^c}{inf}\left(M_n(a^{})M_n\right)$$
as soon as $`\theta _nU^c`$ by definition of the infimum. This gives (8) from (10) on the $``$–a.s. set $`A`$ defined above, with $`(\theta _n^{})_n=(a^{}(\theta _n))_n`$. ∎
###### Proof of Theorem 2.1.
The desired result follows from Lemma 4.3. Namely, let us show that (9) is a consequence of (A1) and (A2). Let $`\theta \theta ^{}`$ and let $`a^{}`$ and $`V`$ as in (A2). Let $`V_k\{\theta \}`$ be a decreasing local base with $`V_0V`$. Let $`Z:=inf_V(m_a^{}m)`$ and $`Z_k:=inf_{V_k}(m_a^{}m)`$ and $`Z_{\mathrm{}}:=m_{a^{}(\theta )}m_\theta `$. By (A1) and the continuity of $`a^{}`$ and the separability of $`\mathrm{\Theta }`$, we get that $`Z_k:𝒳\overline{}`$ is measurable, and that
$$^{}\text{–a.s.,}ZZ_kZ_{\mathrm{}}.$$
Now, by (A2), we get that $`Z\mathrm{L}_{}^1(𝒳,P^{})`$ and $`Z_{\mathrm{}}\mathrm{L}_{}^1(𝒳,P^{})`$ and $`P^{}(Z_{\mathrm{}})>0`$. Observe that $`ZZ^{}\mathrm{L}^1(𝒳,P^{})`$. Thus, by the monotone convergence Theorem,
$$\underset{k}{lim}P^{}\left(Z_k\right)=P^{}(Z_{\mathrm{}})>0.$$
Therefore, $`P^{}(Z_k)>0`$ for some $`k`$ (actually for $`k`$ large enough). Let us denote $`U_\theta :=V_k`$. Now, by the law of large numbers
$$\text{–a.s.},\underset{n}{lim}_n\left(\underset{U_\theta }{inf}(m_a^{}m)\right)=P^{}\left(\underset{U_\theta }{inf}(m_a^{}m)\right)>0.$$
This gives finally (9) since for any $`n`$
$$\underset{U_\theta }{inf}(M_n(a^{})M_n)=\underset{U_\theta }{inf}_n(m_a^{}m)_n\left(\underset{U_\theta }{inf}(m_a^{}m)\right).$$
Acknowledgements. The article benefited from the comments and criticism of the Advisory Editor and two anonymous referees. The authors would like also to sincerely thank Professor Jon A. Wellner who has kindly answered to their questions during his visit in Toulouse.
Djalil Chafaï, corresponding author.
Address: UMR 181 INRA/ENVT Physiopathologie et Toxicologie Expérimentales,
École Nationale Vétérinaire de Toulouse,
23 Chemin des Capelles, F-31076, Toulouse Cedex 3, France.
E-mail: mailto:d.chafai(AT)envt.fr
Address: UMR 5583 CNRS/UPS Laboratoire de Statistique et Probabilités,
Institut de Mathématiques de Toulouse, Université Paul Sabatier,
118 route de Narbonne, F-31062, Toulouse, Cedex 4, France.
E-mail: mailto:chafai(AT)math.ups-tlse.fr
Web: http://www.lsp.ups-tlse.fr/Chafai/ |
warning/0507/math0507180.html | ar5iv | text | # Fast learning rates for plug-in classifiers under the margin condition
## 1 Introduction
Let $`(X,Y)`$ be a random couple taking values in $`𝒵^d\times \{0,1\}`$ with joint distribution $`P`$. We regard $`X^d`$ as a vector of features corresponding to an object and $`Y\{0,1\}`$ as a label indicating that the object belongs to one of the two classes. Consider the sample $`(X_1,Y_1),\mathrm{},(X_n,Y_n)`$, where $`(X_i,Y_i)`$ are independent copies of $`(X,Y)`$. We denote by $`P^n`$ the product probability measure according to which the sample is distributed, and by $`P_X`$ the marginal distribution of $`X`$.
The goal of a classification procedure is to predict the label $`Y`$ given the value of $`X`$, i.e., to provide a decision rule $`f:^d\{0,1\}`$ which belongs to the set $``$ of all Borel functions defined on $`^d`$ and taking values in $`\{0,1\}`$. The performance of a decision rule $`f`$ is measured by the misclassification error
$$R(f)P(Yf(X)).$$
The Bayes decision rule is a minimizer of the risk $`R(f)`$ over all the decision rules $`f`$, and one of such minimizers has the form $`f^{}(X)=1\mathrm{I}_{\{\eta (X)\frac{1}{2}\}}`$ where $`1\mathrm{I}_{\{\}}`$ denotes the indicator function and $`\eta (X)P(Y=1|X)`$ is the regression function of $`Y`$ on $`X`$ (here $`P(dY|X)`$ is a regular conditional probability which we will use in the following without further mention).
An empirical decision rule (a classifier) is a random mapping $`\widehat{f}_n:𝒵^n`$ measurable w.r.t. the sample. Its accuracy can be characterized by the excess risk
$$(\widehat{f}_n)=𝔼R(\widehat{f}_n)R(f^{})$$
where $`𝔼`$ is the sign of expectation. A key problem in classification is to construct classifiers with small excess risk for sufficiently large $`n`$ \[cf. Devroye, Györfi and Lugosi (1996), Vapnik (1998)\]. Optimal classifiers can be defined as those having the best possible rate of convergence of $`(\widehat{f}_n)`$ to 0, as $`n\mathrm{}`$. Of course, this rate, and thus the optimal classifier, depend on the assumptions on the joint distribution of $`(X,Y)`$. A standard way to define optimal classifiers is to introduce a class of joint distributions of $`(X,Y)`$ and to declare $`\widehat{f}_n`$ optimal if it achieves the best rate of convergence in a minimax sense on this class.
Two types of assumptions on the joint distribution of $`(X,Y)`$ are commonly used: complexity assumptions and margin assumptions.
Complexity assumptions are stated in two possible ways. First of them is to suppose that the regression function $`\eta `$ is smooth enough or, more generally, belongs to a class of functions $`\mathrm{\Sigma }`$ having a suitably bounded $`\epsilon `$-entropy. This is called a complexity assumption on the regression function (CAR). Most commonly it is of the following form.
Assumption (CAR). The regression function $`\eta `$ belongs to class $`\mathrm{\Sigma }`$ of functions on $`^d`$ such that
$$(\epsilon ,\mathrm{\Sigma },L_p)A_{}\epsilon ^\rho ,\epsilon >0,$$
with some constants $`\rho >0`$, $`A_{}>0`$. Here $`(\epsilon ,\mathrm{\Sigma },L_p)`$ denotes the $`\epsilon `$-entropy of the set $`\mathrm{\Sigma }`$ w.r.t. an $`L_p`$ norm with some $`1p\mathrm{}`$.
At this stage of discussion we do not identify precisely the value of $`p`$ for the $`L_p`$ norm in Assumption (CAR), nor the measure with respect to which this norm is defined. Examples will be given later. If $`\mathrm{\Sigma }`$ is a class of smooth functions with smoothness parameter $`\beta `$ on a compact in $`^d`$, for example, a Hölder class, as described below, a typical value of $`\rho `$ in Assumption (CAR) is $`\rho =d/\beta `$.
Assumption (CAR) is well adapted for the study of plug-in rules, i.e. of the classifiers having the form
$$\widehat{f}_n^{PI}(X)=1\mathrm{I}_{\{\widehat{\eta }_n(X)\frac{1}{2}\}}$$
(1.1)
where $`\widehat{\eta }_n`$ is a nonparametric estimator of the function $`\eta `$. Indeed, Assumption (CAR) typically reads as a smoothness assumption on $`\eta `$ implying that a good nonparametric estimator (kernel, local polynomial, orthogonal series or other) $`\widehat{\eta }_n`$ converges with some rate to the regression function $`\eta `$, as $`n\mathrm{}`$. In turn, closeness of $`\widehat{\eta }_n`$ to $`\eta `$ implies closeness of $`\widehat{f}_n`$ to $`f`$: for any plug-in classifier $`\widehat{f}_n^{PI}`$ we have
$$𝔼R(\widehat{f}_n^{PI})R(f^{})2𝔼|\widehat{\eta }_n(x)\eta (x)|P_X(dx)$$
(1.2)
(cf. Devroye, Györfi and Lugosi (1996), Theorem 2.2). For various types of estimators $`\widehat{\eta }_n`$ and under rather general assumptions it can be shown that, if (CAR) holds, the RHS of (1.2) is uniformly of the order $`n^{1/(2+\rho )}`$, and thus
$$\underset{P:\eta \mathrm{\Sigma }}{sup}(\widehat{f}_n^{PI})=O(n^{1/(2+\rho )}),n\mathrm{},$$
(1.3)
\[cf. Yang (1999)\]. In particular, if $`\rho =d/\beta `$ (which corresponds to a class of smooth functions with smoothness parameter $`\beta `$), we get
$$\underset{P:\eta \mathrm{\Sigma }}{sup}(\widehat{f}_n^{PI})=O(n^{\beta /(2\beta +d)}),n\mathrm{}.$$
(1.4)
Note that (1.4) can be easily deduced from (1.2) and standard results on the $`L_1`$ or $`L_2`$ convergence rates of usual nonparametric regression estimators on $`\beta `$-smoothness classes $`\mathrm{\Sigma }`$. The rates in (1.3), (1.4) are quite slow, always slower than $`n^{1/2}`$. In (1.4) they deteriorate dramatically as the dimension $`d`$ increases. Moreover, Yang (1999) shows that, under general assumptions, the bound (1.4) cannot be improved in a minimax sense. These results raised some pessimism about the plug-in rules.
The second way to describe complexity is to introduce a structure on the class of possible decision sets $`G^{}=\{x:f^{}(x)=1\}=\{x:\eta (x)1/2\}`$ rather than on that of regression functions $`\eta `$. A standard complexity assumption on the decision set (CAD) is the following.
Assumption (CAD). The decision set $`G^{}`$ belongs to a class $`𝒢`$ of subsets of $`^d`$ such that
$$(\epsilon ,𝒢,d_{\mathrm{}})A_{}\epsilon ^\rho ,\epsilon >0,$$
with some constants $`\rho >0`$, $`A_{}>0`$. Here $`(\epsilon ,𝒢,d_{\mathrm{}})`$ denotes the $`\epsilon `$-entropy of the class $`𝒢`$ w.r.t. the measure of symmetric difference pseudo-distance between sets defined by $`d_{\mathrm{}}(G,G^{})=P_X(G\mathrm{}G^{})`$ for two measurable subsets $`G`$ and $`G^{}`$ in $`^d`$.
The parameter $`\rho `$ in Assumption (CAD) typically characterizes the smoothness of the boundary of $`G^{}`$ \[cf. Tsybakov (2004a)\]. Note that, in general, there is no connection between Assumptions (CAR) and (CAD). Indeed, the fact that $`G^{}`$ has a smooth boundary does not imply that $`\eta `$ is smooth, and vice versa. The values of $`\rho `$ closer to 0 correspond to smoother boundaries (less complex sets $`G^{}`$). As a limit case when $`\rho 0`$ one can consider the Vapnik-Chervonenkis classes (VC-classes) for which the $`\epsilon `$-entropy is logarithmic in $`1/\epsilon `$.
Assumption (CAD) is suited for the study of empirical risk minimization (ERM) type classifiers introduced by Vapnik and Chervonenkis (1974), see also Devroye, Györfi and Lugosi (1996), Vapnik (1998). As shown in Tsybakov (2004a), for every $`0<\rho <1`$ there exist ERM classifiers $`\widehat{f}_n^{ERM}`$ such that, under Assumption (CAD),
$$\underset{P:G^{}𝒢}{sup}(\widehat{f}_n^{ERM})=O(n^{1/2}),n\mathrm{}.$$
(1.5)
The rate of convergence in (1.5) is better than that for plug-in rules, cf. (1.3) – (1.4), and it does not depend on $`\rho `$ (respectively, on the dimension $`d`$). Note that the comparison between (1.5) and (1.3) – (1.4) is not quite legitimate, because there is no inclusion between classes of joint distributions $`P`$ of $`(X,Y)`$ satisfying Assumptions (CAR) and (CAD). Nevertheless, such a comparison have been often interpreted as an argument in disfavor of the plug-in rules. Indeed, Yang’s lower bound shows that the $`n^{1/2}`$ rate cannot be attained under Assumption (CAD) suited for the plug-in rules. Recently, advantages of the ERM type classifiers, including penalized ERM methods, have been further confirmed by the fact that, under the margin (or low noise) assumption, they can attain fast rates of convergence, i.e. the rates that are faster than $`n^{1/2}`$ \[Mammen and Tsybakov (1999), Tsybakov (2004a), Massart and Nédélec (2003), Tsybakov and van de Geer (2005), Koltchinskii (2005), Audibert (2004)\].
The margin assumption (or low noise assumption) is stated as follows.
Assumption (MA). There exist constants $`C_0>0`$ and $`\alpha 0`$ such that
$$P_X\left(0<|\eta (X)1/2|t\right)C_0t^\alpha ,t>0.$$
(1.6)
The case $`\alpha =0`$ is trivial (no assumption) and is included for notational convenience. Assumption (MA) provides a useful characterization of the behavior of regression function $`\eta `$ in a vicinity of the level $`\eta =1/2`$ which turns out to be crucial for convergence of classifiers (for more discussion of the margin assumption see Tsybakov (2004a)). The main point is that, under (MA), fast classification rates up to $`n^1`$ are achievable. In particular, for every $`0<\rho <1`$ and $`\alpha >0`$ there exist ERM type classifiers $`\widehat{f}_n^{ERM}`$ such that
$$\underset{P:(CAD),(MA)}{sup}(\widehat{f}_n^{ERM})=O(n^{\frac{1+\alpha }{2+\alpha +\alpha \rho }}),n\mathrm{},$$
(1.7)
where $`sup_{P:(CAD),(MA)}`$ denotes the supremum over all joint distributions $`P`$ of $`(X,Y)`$ satisfying Assumptions (CAD) and (MA). The RHS of (1.7) can be arbitrarily close to $`O(n^1)`$ for large $`\alpha `$ and small $`\rho `$. Result (1.7) for direct ERM classifiers on $`\epsilon `$-nets is proved by Tsybakov (2004a), and for some other ERM type classifiers by Tsybakov and van de Geer (2005), Koltchinskii (2005) and Audibert (2004) (in some of these papers the rate of convergence (1.7) is obtained with an extra log-factor).
Comparison of (1.5) and (1.7) with (1.2) and (1.3) seems to confirm the conjecture that the plug-in classifiers are inferior to the ERM type ones. The main message of the present paper is to disprove this conjecture. We will show that there exist plug-in rules that converge with fast rates, and even with super-fast rates, i.e. faster than $`n^1`$ under the margin assumption (MA). The basic idea of the proof is to use exponential inequalities for the regression estimator $`\widehat{\eta }_n`$ (see Section 3 below) or the convergence results in the $`L_{\mathrm{}}`$ norm (see Section 5), rather than the usual $`L_1`$ or $`L_2`$ norm convergence of $`\widehat{\eta }_n`$, as previously described (cf. (1.2)). We do not know whether the super-fast rates are attainable for ERM rules or, more precisely, under Assumption (CAD) which serves for the study of the ERM type rules. It is important to note that our results on fast rates cover more general setting than just classification with plug-in rules. These are rather results about classification in the regression complexity context under the margin assumption. In particular, we establish minimax lower bounds valid for all classifiers, and we construct a “hybrid” plug-in/ ERM procedure (ERM based on a grid on a set regression functions $`\eta `$) that achieves optimality. Thus, the point is mainly not about the type of procedure (plug-in or ERM) but about the type of complexity assumption (on the regression function (CAR) or on the decision set (CAD)) that should be natural to impose. Assumption (CAR) on the regression function arises in a natural way in the analysis of several practical procedures of plug-in type, such as boosting and SVM \[cf. Blanchard, Lugosi and Vayatis (2003), Bartlett, Jordan and McAuliffe (2003), Scovel and Steinwart (2003), Blanchard, Bousquet and Massart (2004), Tarigan and van de Geer (2004)\]. These procedures are now intensively studied but, to our knowledge, only suboptimal rates of convergence have been proved in the regression complexity context under the margin assumption. The results in Section 4 point out this fact (see also Section 5), and establish the best achievable rates of classification that those procedures should expectedly attain.
## 2 Notation and definitions
In this section we introduce some notation, definitions and basic facts that will be used in the paper.
We denote by $`C,C_1,C_2,\mathrm{}`$ positive constants whose values may differ from line to line. The symbols $``$ and $`𝔼`$ stand for generic probability and expectation signs, and $`E_X`$ is the expectation w.r.t. the marginal distribution $`P_X`$. We denote by $`(x,r)`$ the closed Euclidean ball in $`^d`$ centered at $`x^d`$ and of radius $`r>0`$.
For any multi-index $`s=(s_1,\mathrm{},s_d)^d`$ and any $`x=(x_1,\mathrm{},x_d)^d`$, we define $`|s|=_{i=1}^ds_i`$, $`s!=s_1!\mathrm{}s_d!`$, $`x^s=x_1^{s_1}\mathrm{}x_d^{s_d}`$ and $`x(x_1^2+\mathrm{}+x_d^2)^{1/2}`$. Let $`D^s`$ denote the differential operator $`D^s\frac{^{s_1+\mathrm{}+s_d}}{x_1^{s_1}\mathrm{}x_d^{s_d}}.`$
Let $`\beta >0`$. Denote by $`\beta `$ the maximal integer that is strictly less than $`\beta `$. For any $`x^d`$ and any $`\beta `$ times continuously differentiable real valued function $`g`$ on $`^d`$, we denote by $`g_x`$ its Taylor polynomial of degree $`\beta `$ at point $`x`$:
$$g_x(x^{})\underset{|s|\beta }{}\frac{(x^{}x)^s}{s!}D^sg(x).$$
Let $`L>0`$. The $`(\beta ,L,^d)`$-Hölder class of functions, denoted $`\mathrm{\Sigma }(\beta ,L,^d)`$, is defined as the set of functions $`g:^d`$ that are $`\beta `$ times continuously differentiable and satisfy, for any $`x,x^{}^d`$ , the inequality
$$|g(x^{})g_x(x^{})|Lxx^{}^\beta .$$
Fix some constants $`c_0,r_0>0`$. We will say that a Lebesgue measurable set $`A^d`$ is $`(c_0,r_0)`$-regular if
$$\lambda \left[A(x,r)\right]c_0\lambda \left[(x,r)\right],0<rr_0,xA,$$
(2.1)
where $`\lambda [S]`$ stands for the Lebesgue measure of $`S^d`$. To illustrate this definition, consider the following example. Let $`d2`$. Then the set $`A=\{x=(x_1,\mathrm{},x_d)^d:_{j=1}^d|x_j|^q1\}`$ is $`(c_0,r_0)`$-regular with some $`c_0,r_0>0`$ for $`q1`$, and there are no $`c_0,r_0>0`$ such that $`A`$ is $`(c_0,r_0)`$-regular for $`0<q<1`$.
Introduce now two assumptions on the marginal distribution $`P_X`$ that will be used in the sequel.
###### Definition 2.1
Fix $`0<c_0,r_0,\mu _{\mathrm{max}}<\mathrm{}`$ and a compact $`𝒞^d`$. We say that the mild density assumption is satisfied if the marginal distribution $`P_X`$ is supported on a compact $`(c_0,r_0)`$-regular set $`A𝒞`$ and has a uniformly bounded density $`\mu `$ w.r.t. the Lebesgue measure: $`\mu (x)\mu _{\mathrm{max}},xA`$.
###### Definition 2.2
Fix some constants $`c_0,r_0>0`$ and $`0<\mu _{\mathrm{min}}<\mu _{\mathrm{max}}<\mathrm{}`$ and a compact $`𝒞^d`$. We say that the strong density assumption is satisfied if the marginal distribution $`P_X`$ is supported on a compact $`(c_0,r_0)`$-regular set $`A𝒞`$ and has a density $`\mu `$ w.r.t. the Lebesgue measure bounded away from zero and infinity on $`A`$:
$$\mu _{\mathrm{min}}\mu (x)\mu _{\mathrm{max}}\text{for}xA,\text{and}\mu (x)=0\text{otherwise}.$$
We finally recall some notions related to locally polynomial estimators.
###### Definition 2.3
For $`h>0`$, $`x^d`$, for an integer $`l0`$ and a function $`K:^d_+`$, denote by $`\widehat{\theta }_x`$ a polynomial on $`^d`$ of degree $`l`$ which minimizes
$$\underset{i=1}{\overset{n}{}}\left[Y_i\widehat{\theta }_x(X_ix)\right]^2K\left(\frac{X_ix}{h}\right).$$
(2.2)
The locally polynomial estimator $`\widehat{\eta }_n^{LP}(x)`$ of order $`l`$, or LP(l) estimator, of the value $`\eta (x)`$ of the regression function at point $`x`$ is defined by: $`\widehat{\eta }_n^{LP}(x)\widehat{\theta }_x(0)`$ if $`\widehat{\theta }_x`$ is the unique minimizer of (2.2) and $`\widehat{\eta }_n^{LP}(x)0`$ otherwise. The value $`h`$ is called the bandwidth and the function $`K`$ is called the kernel of the LP(l) estimator.
Let $`T_s`$ denote the coefficients of $`\widehat{\theta }_x`$ indexed by multi-index $`s^d`$: $`\widehat{\theta }_x(u)=_{|s|l}T_su^s.`$ Introduce the vectors $`T\left(T_s\right)_{|s|l}`$, $`V\left(V_s\right)_{|s|l}`$ where
$$\begin{array}{ccc}V_s_{i=1}^nY_i(X_ix)^sK\left(\frac{X_ix}{h}\right),\hfill & & \end{array}$$
(2.3)
$`U(u)\left(u^s\right)_{|s|l}`$ and the matrix $`Q\left(Q_{s_1,s_2}\right)_{|s_1|,|s_2|l}`$ where
$$\begin{array}{ccc}Q_{s_1,s_2}_{i=1}^n(X_ix)^{s_1+s_2}K\left(\frac{X_ix}{h}\right).\hfill & & \end{array}$$
(2.4)
The following result is straightforward (cf. Section 1.7 in Tsybakov (2004b) where the case $`d=1`$ is considered).
###### Proposition 2.1
If the matrix $`Q`$ is positive definite, there exists a unique polynomial on $`^d`$ of degree $`l`$ minimizing (2.2). Its vector of coefficients is given by $`T=Q^1V`$ and the corresponding LP(l) regression function estimator has the form
$$\widehat{\eta }_n^{LP}(x)=U^T(0)Q^1V=\underset{i=1}{\overset{n}{}}Y_iK\left(\frac{X_ix}{h}\right)U^T(0)Q^1U(X_ix).$$
## 3 Fast rates for plug-in rules: the strong density assumption
We first state a general result showing how the rates of convergence of plug-in classifiers can be deduced from exponential inequalities for the corresponding regression estimators.
In the sequel, for an estimator $`\widehat{\eta }_n`$ of $`\eta `$, we write
$$\left(|\widehat{\eta }_n(X)\eta (X)|\delta \right)P^n\left(|\widehat{\eta }_n(x)\eta (x)|\delta \right)P_X(dx),\delta >0,$$
i.e., we consider the probability taken with respect to the distribution of the sample $`(X_1,Y_1,\mathrm{}X_n,Y_n)`$ and the distribution of the input $`X`$.
###### Theorem 3.1
Let $`\widehat{\eta }_n`$ be an estimator of the regression function $`\eta `$ and $`𝒫`$ a set of probability distributions on $`𝒵`$ such that for some constants $`C_1>0,`$ $`C_2>0`$, for some positive sequence $`a_n`$, for $`n1`$ and any $`\delta >0`$, and for almost all $`x`$ w.r.t. $`P_X`$, we have
$$\underset{P𝒫}{sup}P^n\left(|\widehat{\eta }_n(x)\eta (x)|\delta \right)C_1\mathrm{exp}\left(C_2a_n\delta ^2\right).$$
(3.1)
Consider the plug-in classifier $`\widehat{f}_n=1\mathrm{I}_{\{\widehat{\eta }_n\frac{1}{2}\}}`$. If all the distributions $`P𝒫`$ satisfy the margin assumption (MA), we have
$$\underset{P𝒫}{sup}\left\{𝔼R(\widehat{f}_n)R(f^{})\right\}Ca_n^{\frac{1+\alpha }{2}}$$
for $`n1`$ with some constant $`C>0`$ depending only on $`\alpha `$, $`C_0`$, $`C_1`$ and $`C_2`$.
Proof. Consider the sets $`A_j^d,j=1,2,\mathrm{},`$ defined as
$$\begin{array}{ccc}A_0\hfill & \hfill & \{x^d:0<|\eta (x)\frac{1}{2}|\delta \},\hfill \\ A_j\hfill & \hfill & \{x^d:2^{j1}\delta <|\eta (x)\frac{1}{2}|2^j\delta \},\text{ for }j1.\hfill \end{array}$$
For any $`\delta >0`$, we may write
$$\begin{array}{ccc}𝔼R(\widehat{f}_n)R(f^{})\hfill & =\hfill & 𝔼\left(|2\eta (X)1|1\mathrm{I}_{\{\widehat{f}_n(X)f^{}(X)\}}\right)\hfill \\ & =\hfill & _{j=0}^{\mathrm{}}𝔼\left(|2\eta (X)1|1\mathrm{I}_{\{\widehat{f}_n(X)f^{}(X)\}}1\mathrm{I}_{\{XA_j\}}\right)\hfill \\ & \hfill & 2\delta P_X\left(0<|\eta (X)\frac{1}{2}|\delta \right)\hfill \\ & & +_{j1}𝔼\left(|2\eta (X)1|1\mathrm{I}_{\{\widehat{f}_n(X)f^{}(X)\}}1\mathrm{I}_{\{XA_j\}}\right).\hfill \end{array}$$
(3.2)
On the event $`\{\widehat{f}_nf^{}\}`$ we have $`|\eta \frac{1}{2}||\widehat{\eta }_n\eta |`$. So, for any $`j1`$, we get
$$\begin{array}{ccc}𝔼\left(|2\eta (X)1|1\mathrm{I}_{\{\widehat{f}_n(X)f^{}(X)\}}1\mathrm{I}_{\{XA_j\}}\right)\hfill & & \\ 2^{j+1}\delta 𝔼\left[1\mathrm{I}_{\{|\widehat{\eta }_n(X)\eta (X)|2^{j1}\delta \}}1\mathrm{I}_{\{0<|\eta (X)\frac{1}{2}|2^j\delta \}}\right]\hfill & & \\ 2^{j+1}\delta E_X\left[P^n\left(|\widehat{\eta }_n(X)\eta (X)|2^{j1}\delta \right)1\mathrm{I}_{\{0<|\eta (X)\frac{1}{2}|2^j\delta \}}\right]\hfill & & \\ C_12^{j+1}\delta \mathrm{exp}\left(C_2a_n(2^{j1}\delta )^2\right)P_X\left(0<|\eta (X)\frac{1}{2}|2^j\delta \right)\hfill & & \\ 2C_1C_02^{j(1+\alpha )}\delta ^{1+\alpha }\mathrm{exp}\left(C_2a_n(2^{j1}\delta )^2\right)\hfill & & \end{array}$$
where in the last inequality we used Assumption (MA). Now, from inequality (3.2), taking $`\delta =a_n^{1/2}`$ and using Assumption (MA) to bound the first term of the right hand side of (3.2), we get
$$\begin{array}{ccc}𝔼R(\widehat{f}_n)R(f^{})\hfill & \hfill & 2C_0a_n^{\frac{1+\alpha }{2}}+Ca_n^{\frac{1+\alpha }{2}}_{j2}2^{j(1+\alpha )}\mathrm{exp}\left(C_22^{2j2}\right)\hfill \\ & \hfill & Ca_n^{\frac{1+\alpha }{2}}.\hfill \end{array}$$
Inequality (3.1) is crucial to obtain the above result. This inequality holds true for various types of estimators and various sets of probability distributions $`𝒫`$. Here we focus on a standard case where $`\eta `$ belongs to the Hölder class $`\mathrm{\Sigma }(\beta ,L,^d)`$ and the marginal law of $`X`$ satisfies the strong density assumption. We are going to show that in this case there exist estimators satisfying inequality (3.1) with $`a_n=n^{\frac{2\beta }{2\beta +d}}`$. These can be, for example, locally polynomial estimators. Specifically, assume from now on that $`K`$ is a kernel satisfying
$`c>0:K(x)c1\mathrm{I}_{\{xc\}},x^d,`$ (3.3)
$`{\displaystyle _^d}K(u)𝑑u=1,`$ (3.4)
$`{\displaystyle _^d}\left(1+u^{4\beta }\right)K^2(u)𝑑u<\mathrm{},`$ (3.5)
$`\underset{u^d}{sup}\left(1+u^{2\beta }\right)K(u)<\mathrm{}.`$ (3.6)
Let $`h>0`$, and consider the matrix $`\overline{B}\left(\overline{B}_{s_1,s_2}\right)_{|s_1|,|s_2|\beta }`$ where $`\overline{B}_{s_1,s_2}=\frac{1}{nh^d}_{i=1}^n\left(\frac{X_ix}{h}\right)^{s_1+s_2}K\left(\frac{X_ix}{h}\right).`$ Define the regression function estimator $`\widehat{\eta }_n^{}`$ as follows. If the smallest eigenvalue of the matrix $`\overline{B}`$ is greater than $`(\mathrm{log}n)^1`$ we set $`\widehat{\eta }_n^{}(x)`$ equal to the projection of $`\widehat{\eta }_n^{LP}(x)`$ on the interval $`[0,1]`$, where $`\widehat{\eta }_n^{LP}(x)`$ is the LP($`\beta `$) estimator with a bandwidth $`h>0`$ and a kernel $`K`$ satisfying (3.3) – (3.6). If the smallest eigenvalue of $`\overline{B}`$ is less than $`(\mathrm{log}n)^1`$ we set $`\widehat{\eta }_n^{}(x)=0`$.
###### Theorem 3.2
Let $`𝒫`$ be a class of probability distributions $`P`$ on $`𝒵`$ such that the regression function $`\eta `$ belongs to the Hölder class $`\mathrm{\Sigma }(\beta ,L,^d)`$ and the marginal law of $`X`$ satisfies the strong density assumption. Then there exist constants $`C_1,C_2,C_3>0`$ such that for any $`0<hr_0/c`$, any $`C_3h^\beta <\delta `$ and any $`n1`$ the estimator $`\widehat{\eta }_n^{}`$ satisfies
$$\underset{P𝒫}{sup}P^n\left(\left|\widehat{\eta }_n^{}(x)\eta (x)\right|\delta \right)C_1\mathrm{exp}\left(C_2nh^d\delta ^2\right)$$
(3.7)
for almost all $`x`$ w.r.t. $`P_X`$. As a consequence, there exist $`C_1,C_2>0`$ such that for $`h=n^{\frac{1}{2\beta +d}}`$ and any $`\delta >0`$, $`n1`$ we have
$$\underset{P𝒫}{sup}P^n\left(\left|\widehat{\eta }_n^{}(x)\eta (x)\right|\delta \right)C_1\mathrm{exp}\left(C_2n^{\frac{2\beta }{2\beta +d}}\delta ^2\right)$$
(3.8)
for almost all $`x`$ w.r.t. $`P_X`$. The constants $`C_1,C_2,C_3`$ depend only on $`\beta `$, $`d`$, $`L`$, $`c_0`$, $`r_0`$, $`\mu _{\mathrm{min}},\mu _{\mathrm{max}}`$, and on the kernel $`K`$.
Proof. See Section 6.1.
###### Remark 3.1
We have chosen here the LP estimators of $`\eta `$ because for them the exponential inequality (3.1) holds without additional smoothness conditions on the marginal density of $`X`$. For other popular regression estimators, such as kernel or orthogonal series ones, similar inequality can be also proved if we assume that the marginal density of $`X`$ is as smooth as the regression function.
###### Definition 3.1
For a fixed parameter $`\alpha 0`$, fixed positive parameters $`c_0,r_0,C_0,\beta ,L,\mu _{\mathrm{max}}>\mu _{\mathrm{min}}>0`$ and a fixed compact $`𝒞^d`$, let $`𝒫_\mathrm{\Sigma }`$ denote the class of all probability distributions $`P`$ on $`𝒵`$ such that
1. the margin assumption (MA) is satisfied,
2. the regression function $`\eta `$ belongs to the Hölder class $`\mathrm{\Sigma }(\beta ,L,^d)`$,
3. the strong density assumption on $`P_X`$ is satisfied.
Theorem 3.1 and (3.8) immediately imply the next result.
###### Theorem 3.3
For any $`n1`$ the excess risk of the plug-in classifier $`\widehat{f}_n^{}=1\mathrm{I}_{\{\widehat{\eta }_n^{}\frac{1}{2}\}}`$ with bandwidth $`h=n^{\frac{1}{2\beta +d}}`$ satisfies
$$\underset{P𝒫_\mathrm{\Sigma }}{sup}\left\{𝔼R(\widehat{f}_n^{})R(f^{})\right\}Cn^{\frac{\beta (1+\alpha )}{2\beta +d}}$$
where the constant $`C>0`$ depends only on $`\alpha `$, $`C_0`$, $`C_1`$ and $`C_2`$.
For $`\alpha \beta >d/2`$ the convergence rate $`n^{\frac{\beta (1+\alpha )}{2\beta +d}}`$ obtained in Theorem 3.3 is a fast rate, i.e., it is faster than $`n^{1/2}`$. Furthermore, it is a super-fast rate (i.e., is faster than $`n^1`$) for $`(\alpha 1)\beta >d`$. We must note that if this condition is satisfied, the class $`𝒫_\mathrm{\Sigma }`$ is rather poor, and thus super-fast rates can occur only for very particular joint distributions of $`(X,Y)`$. Intuitively, this is clear. Indeed, to have a very smooth regression function $`\eta `$ (i.e., very large $`\beta `$) implies that when $`\eta `$ hits the level $`1/2`$, it cannot “take off” from this level too abruptly. As a consequence, when the density of the distribution $`P_X`$ is bounded away from $`0`$ at a vicinity of the hitting point, the margin assumption cannot be satisfied for large $`\alpha `$ since this assumption puts an upper bound on the “time spent” by the regression function near $`1/2`$. So, $`\alpha `$ and $`\beta `$ cannot be simultaneously very large. It can be shown that the cases of “too large” and “not too large” $`(\alpha ,\beta )`$ are essentially described by the condition $`(\alpha 1)\beta >d`$.
To be more precise, observe first that $`𝒫_\mathrm{\Sigma }`$ is not empty for $`(\alpha 1)\beta >d`$, so that the super-fast rates can effectively occur. Examples of laws $`P𝒫_\mathrm{\Sigma }`$ under this condition can be easily given, such as the one with $`P_X`$ equal to the uniform distribution on a ball centered at 0 in $`^d`$, and the regression function defined by $`\eta (x)=1/2Cx^2`$ with an appropriate $`C>0`$. Clearly, $`\eta `$ belongs to Hölder classes with arbitrarily large $`\beta `$ and Assumption (MA) is satisfied with $`\alpha =d/2`$. Thus, for $`d3`$ and $`\beta `$ large enough super-fast rates can occur. Note that in this example the decision set $`\{x:\eta (x)1/2\}`$ has the Lebesgue measure 0 in $`^d`$. It turns out that this condition is necessary to achieve classification with super-fast rates when the Hölder classes of regression functions are considered.
To explain this and to have further insight into the problem of super-fast rates, consider the following two questions:
* for which parameters $`\alpha ,`$ $`\beta `$ and $`d`$ is there a distribution $`P𝒫_\mathrm{\Sigma }`$ such that the regression function associated with $`P`$ hits<sup>1</sup><sup>1</sup>1 A function $`f:^d`$ is said to hit the level $`a`$ at $`x_0^d`$ if and only if $`f(x_0)=a`$ and for any $`r>0`$ there exists $`x(x_0,r)`$ such that $`f(x)a`$ . $`1/2`$ in the support of $`P_X`$?
* for which parameters $`\alpha ,`$ $`\beta `$ and $`d`$ is there a distribution $`P𝒫_\mathrm{\Sigma }`$ such that the regression function associated with $`P`$ crosses<sup>2</sup><sup>2</sup>2 A function $`f:^d`$ is said to cross the level $`a`$ at $`x_0^d`$ if and only if for any $`r>0`$, there exists $`x_{}`$ and $`x_+`$ in $`(x_0,r)`$ such that $`f(x_{})<a`$ and $`f(x_+)>a`$. $`1/2`$ in the interior of the support of $`P_X`$?
The following result gives a precise description of the constraints on $`(\alpha ,\beta )`$ leading to possibility or impossibility of the super-fast rates.
###### Proposition 3.4
* If $`\alpha (1\beta )>d`$, there is no distribution $`P𝒫_\mathrm{\Sigma }`$ such that the regression function $`\eta `$ associated with $`P`$ hits $`1/2`$ in the interior of the support of $`P_X`$.
* For any $`\alpha ,\beta >0`$ and integer $`d\alpha (1\beta )`$, any positive parameter $`L`$ and any compact $`𝒞^d`$ with non-empty interior, for appropriate positive parameters $`C_0,c_0,r_0,\mu _{\mathrm{max}}>\mu _{\mathrm{min}}>0`$, there are distributions $`P𝒫_\mathrm{\Sigma }`$ such that the regression function $`\eta `$ associated with $`P`$ hits $`1/2`$ in the boundary of the support of $`P_X`$.
* For any $`\alpha ,\beta >0`$, any integer $`d2\alpha `$, any positive parameter $`L`$ and any compact $`𝒞^d`$ with non-empty interior, for appropriate positive parameters $`C_0,c_0,r_0,\mu _{\mathrm{max}}>\mu _{\mathrm{min}}>0`$, there are distributions $`P𝒫_\mathrm{\Sigma }`$ such that the regression function $`\eta `$ associated with $`P`$ hits $`1/2`$ in the interior of the support of $`P_X`$.
* If $`\alpha (1\beta )>1`$ there is no distribution $`P𝒫_\mathrm{\Sigma }`$ such that the regression function $`\eta `$ associated with $`P`$ crosses $`1/2`$ in the interior of the support of $`P_X`$. Conversely, for any $`\alpha ,\beta >0`$ such that $`\alpha (1\beta )1`$, any integer $`d`$, any positive parameter $`L`$ and any compact $`𝒞^d`$ with non-empty interior, for appropriate positive parameters $`C_0,c_0,r_0,\mu _{\mathrm{max}}>\mu _{\mathrm{min}}>0`$, there are distributions $`P𝒫_\mathrm{\Sigma }`$ such that the regression function $`\eta `$ associated with $`P`$ crosses $`1/2`$ in the interior of the support of $`P_X`$.
Note that the condition $`\alpha (1\beta )>1`$ appearing in the last assertion is equivalent to $`\frac{\beta (1+\alpha )}{2\beta +d}>\frac{(2\beta )(\beta +1)}{2\beta +d}`$, which is necessary to have super-fast rates. As a consequence, in this context, super-fast rates cannot occur when the regression function crosses $`1/2`$ in the interior of the support. The third assertion of the proposition shows that super-fast rates can occur with regression functions hitting $`1/2`$ in the interior of the support of $`P_X`$ provided that the regression function is highly smooth and defined on a highly dimensional space and that a strong margin assumption holds (i.e. $`\alpha `$ large).
Proof. See Section 6.3.
The following lower bound shows optimality of the rate of convergence for the Hölder classes obtained in Theorem 3.3.
###### Theorem 3.5
Let $`d1`$ be an integer, and let $`L,\beta ,\alpha `$ be positive constants, such that $`\alpha \beta d`$. Then there exists a constant $`C>0`$ such that for any $`n1`$ and any classifier $`\widehat{f}_n:𝒵^n`$, we have
$$\begin{array}{ccc}\underset{P𝒫_\mathrm{\Sigma }}{sup}\left\{𝔼R(\widehat{f}_n)R(f^{})\right\}Cn^{\frac{\beta (1+\alpha )}{2\beta +d}}.\hfill & & \end{array}$$
Proof. See Section 6.2.
Note that the lower bound of Theorem 3.5 does not cover the case of super-fast rates ($`(\alpha 1)\beta >d`$).
Finally, we discuss the case where “$`\alpha =\mathrm{}`$”, which means that there exists $`t_0>0`$ such that
$$P_X\left(0<|\eta (X)1/2|t_0\right)=0.$$
(3.9)
This is a very favorable situation for classification. The rates of convergence of the ERM type classifiers under (3.9) are, of course, faster than under Assumption (MA) with $`\alpha <\mathrm{}`$ \[cf. Massart and Nédélec (2003)\], but they are not faster than $`n^1`$. Indeed, Massart and Nédélec (2003) provide a lower bound showing that, even if Assumption (CAD) is replaced by a very strong assumption that the true decision set belongs to a VC-class (note that both assumptions are naturally linked to the study the ERM type classifiers), the best achievable rate is of the order $`(\mathrm{log}n)/n`$. We show now that for the plug-in classifiers much faster rates can be attained. Specifically, if the regression function $`\eta `$ has some (arbitrarily low) Hölder smoothness $`\beta `$ the rate of convergence can be exponential in $`n`$. To show this, we first state a simple lemma which is valid for any plug-in classifier $`\widehat{f}_n`$.
###### Lemma 3.6
Let assumption (3.9) be satisfied, and let $`\widehat{\eta }_n`$ be an estimator of the regression function $`\eta `$. Then for the plug-in classifier $`\widehat{f}_n=1\mathrm{I}_{\{\widehat{\eta }_n\frac{1}{2}\}}`$ we have
$$𝔼R(\widehat{f}_n)R(f^{})\left(|\widehat{\eta }_n(X)\eta (X)|>t_0\right).$$
Proof. Following the argument similar to the proof of Theorem 3.1 and using condition (3.9) we get
$$\begin{array}{ccc}𝔼R(\widehat{f}_n)R(f^{})\hfill & \hfill & 2t_0P_X\left(0<|\eta (X)1/2|t_0\right)\hfill \\ & & +𝔼\left(|2\eta (X)1|1\mathrm{I}_{\{\widehat{f}_n(X)f^{}(X)\}}1\mathrm{I}_{\{|\eta (X)1/2|>t_0\}}\right)\hfill \\ & =\hfill & 𝔼\left(|2\eta (X)1|1\mathrm{I}_{\{\widehat{f}_n(X)f^{}(X)\}}1\mathrm{I}_{\{|\eta (X)1/2|>t_0\}}\right)\hfill \\ & \hfill & \left(|\widehat{\eta }_n(X)\eta (X)|>t_0\right).\hfill \end{array}$$
Lemma 3.6 and Theorem 3.2 immediately imply that, under assumption (3.9), the rate of convergence of the plug-in classifier $`\widehat{f}_n^{}=1\mathrm{I}_{\{\widehat{\eta }_n^{}\frac{1}{2}\}}`$ with a small enough fixed (independent of $`n`$) bandwidth $`h`$ is exponential. To state the result, we denote by $`𝒫_{\mathrm{\Sigma },\mathrm{}}`$ the class of probability distributions $`P`$ defined in the same way as $`𝒫_\mathrm{\Sigma }`$, with the only difference that in Definition 3.1 the margin assumption (MA) is replaced by condition (3.9).
###### Proposition 3.7
There exists a fixed (independent of $`n`$) $`h>0`$ such that for any $`n1`$ the excess risk of the plug-in classifier $`\widehat{f}_n^{}=1\mathrm{I}_{\{\widehat{\eta }_n^{}\frac{1}{2}\}}`$ with bandwidth $`h`$ satisfies
$$\underset{P𝒫_{\mathrm{\Sigma },\mathrm{}}}{sup}\left\{𝔼R(\widehat{f}_n^{})R(f^{})\right\}C_4\mathrm{exp}(C_5n)$$
where the constants $`C_4,C_5>0`$ depend only on $`t_0`$, $`\beta `$, $`d`$, $`L`$, $`c_0`$, $`r_0`$, $`\mu _{\mathrm{min}},\mu _{\mathrm{max}}`$, and on the kernel $`K`$.
Proof. Use Lemma 3.6, choose $`h>0`$ such that $`h<\mathrm{min}(r_0/c,(t_0/C_3)^{1/\beta })`$, and apply (3.7) with $`\delta =t_0`$.
Koltchinskii and Beznosova (2005) prove a result on exponential rates for the plug-in classifier with some penalized regression estimator in place of the locally polynomial one that we use here. Their result is stated under a less general condition, in the sense that they consider only the Lipschitz class of regression functions $`\eta `$, while in Proposition 3.7 the Hölder smoothness $`\beta `$ can be arbitrarily close to 0. Note also that we do not impose any complexity assumption on the decision set. However, the class of distributions $`𝒫_{\mathrm{\Sigma },\mathrm{}}`$ is quite restricted in a different sense. Indeed, for such distributions condition (3.9) should be compatible with the assumption that $`\eta `$ belongs to a Hölder class. A sufficient condition for that is the existence of a band or a “corridor” of zero $`P_X`$-measure separating the sets $`\{x:\eta (x)>1/2\}`$ and $`\{x:\eta (x)<1/2\}`$. We believe that this condition is close to the necessary one.
## 4 Optimal learning rates without the strong density assumption
In this section we show that if $`P_X`$ does not admit a density bounded away from zero on its support the rates of classification are slower than those obtained in Section 3. In particular, super-fast rates, i.e., the rates faster than $`n^1`$, cannot be achieved. Introduce the following class of probability distributions.
###### Definition 4.1
For a fixed parameter $`\alpha 0`$, fixed positive parameters $`c_0,r_0,C_0,\beta ,L,\mu _{\mathrm{max}}>0`$ and a fixed compact $`𝒞^d`$, let $`𝒫_\mathrm{\Sigma }^{}`$ denote the class of all probability distributions $`P`$ on $`𝒵`$ such that
1. the margin assumption (MA) is satisfied,
2. the regression function $`\eta `$ belongs to the Hölder class $`\mathrm{\Sigma }(\beta ,L,^d)`$,
3. the mild density assumption on $`P_X`$ is satisfied.
In this section we mainly assume that the distribution $`P`$ of $`(X,Y)`$ belongs to $`𝒫_\mathrm{\Sigma }^{}`$, but we also consider larger classes of distributions satisfying the margin assumption (MA) and the complexity assumption (CAR).
Clearly, $`𝒫_\mathrm{\Sigma }𝒫_\mathrm{\Sigma }^{}`$. The only difference between $`𝒫_\mathrm{\Sigma }^{}`$ and $`𝒫_\mathrm{\Sigma }`$ is that for $`𝒫_\mathrm{\Sigma }^{}`$ the marginal density of $`X`$ is not bounded away from zero. The optimal rates for $`𝒫_\mathrm{\Sigma }^{}`$ are slower than for $`𝒫_\mathrm{\Sigma }`$. Indeed, we have the following lower bound for the excess risk.
###### Theorem 4.1
Let $`d1`$ be an integer, and let $`L,\beta ,\alpha `$ be positive constants. Then there exists a constant $`C>0`$ such that for any $`n1`$ and any classifier $`\widehat{f}_n:𝒵^n`$ we have
$$\begin{array}{ccc}\underset{P𝒫_\mathrm{\Sigma }^{}}{sup}\left\{𝔼R(\widehat{f}_n)R(f^{})\right\}Cn^{\frac{(1+\alpha )\beta }{(2+\alpha )\beta +d}}.\hfill & & \end{array}$$
Proof. See Section 6.2.
In particular, when $`\alpha =d/\beta `$, we get slow convergence rate $`1/\sqrt{n}`$, instead of the fast rate $`n^{\frac{\beta +d}{2\beta +d}}`$ obtained in Theorem 3.3 under the strong density assumption. Nevertheless, the lower bound can still approach $`n^1`$, as the margin parameter $`\alpha `$ tends to $`\mathrm{}`$.
We now show that the rate of convergence given in Theorem 4.1 is optimal in the sense that there exist estimators that achieve this rate. This will be obtained as a consequence of a general upper bound for the excess risk of classifiers over a larger set $`𝒫`$ of distributions than $`𝒫_\mathrm{\Sigma }^{}`$.
Fix a Lebesgue measurable set $`𝒞^d`$ and a value $`1p\mathrm{}`$. Let $`\mathrm{\Sigma }`$ be a class of regression functions $`\eta `$ on $`^d`$ such that Assumption (CAR) is satisfied where the $`\epsilon `$-entropy is taken w.r.t. the $`L_p(𝒞,\lambda )`$ norm ($`\lambda `$ is the Lebesgue measure on $`^d`$). Then for every $`\epsilon >0`$ there exists an $`\epsilon `$-net $`𝒩_\epsilon `$ on $`\mathrm{\Sigma }`$ w.r.t. this norm such that
$$\mathrm{log}\left(\mathrm{card}𝒩_\epsilon \right)A^{}\epsilon ^\rho ,$$
where $`A^{}`$ is a constant. Consider the empirical risk
$$R_n(f)=\frac{1}{n}\underset{i=1}{\overset{n}{}}1\mathrm{I}_{\{f(X_i)Y_i\}},f,$$
and set
$$\epsilon _n=\epsilon _n(\alpha ,\rho ,p)\{\begin{array}{ccc}n^{\frac{1}{2+\alpha +\rho }}\hfill & \text{if}& \hfill p=\mathrm{},\\ n^{\frac{p+\alpha }{(2+\alpha )p+\rho (p+\alpha )}}\hfill & \text{if}& \hfill 1p<\mathrm{}.\end{array}$$
Define a sieve estimator $`\widehat{\eta }_n^S`$ of the regression function $`\eta `$ by the relation
$$\widehat{\eta }_n^S\mathrm{Argmin}_{\overline{\eta }𝒩_{\epsilon _n}}R_n(f_{\overline{\eta }})$$
(4.1)
where $`f_{\overline{\eta }}(x)=1\mathrm{I}_{\{\overline{\eta }(x)1/2\}}`$, and consider the classifier $`\widehat{f}_n^S=1\mathrm{I}_{\{\widehat{\eta }_n^S1/2\}}`$. Note that $`\widehat{f}_n^S`$ can be viewed as a “hybrid” plug-in/ ERM procedure: the ERM is performed on a set of plug-in rules corresponding to a grid on the class of regression functions $`\eta `$.
###### Theorem 4.2
Let $`𝒫`$ be a set of probability distributions $`P`$ on $`𝒵`$ such that
1. the margin assumption (MA) is satisfied,
2. the regression function $`\eta `$ belongs to a class $`\mathrm{\Sigma }`$ which satisfies the complexity assumption (CAR) with the $`\epsilon `$-entropy taken w.r.t. the $`L_p(𝒞,\lambda )`$ norm for some $`1p\mathrm{}`$,
3. for all $`P𝒫`$ the supports of marginal distributions $`P_X`$ are included in $`𝒞`$.
Consider the classifier $`\widehat{f}_n^S=1\mathrm{I}_{\{\widehat{\eta }_n^S1/2\}}`$. If $`p=\mathrm{}`$ for any $`n1`$ we have
$$\underset{P𝒫}{sup}\left\{𝔼R(\widehat{f}_n^S)R(f^{})\right\}Cn^{\frac{1+\alpha }{2+\alpha +\rho }}.$$
(4.2)
If $`1p<\mathrm{}`$ and, in addition, for all $`P𝒫`$ the marginal distributions $`P_X`$ are absolutely continuous w.r.t. the Lebesgue measure and their densities are uniformly bounded from above by some constant $`\mu _{\mathrm{max}}<\mathrm{}`$, then for any $`n1`$ we have
$$\underset{P𝒫}{sup}\left\{𝔼R(\widehat{f}_n^S)R(f^{})\right\}Cn^{\frac{(1+\alpha )p}{(2+\alpha )p+\rho (p+\alpha )}}.$$
(4.3)
Proof. See Section 6.4.
Theorem 4.2 allows one to get fast classification rates without any density assumption on $`P_X`$. Namely, define the following class of distributions $`P`$ of $`(X,Y)`$.
###### Definition 4.2
For fixed parameters $`\alpha 0`$, $`C_0>0,\beta >0,L>0`$, and for a fixed compact $`𝒞^d`$, let $`𝒫_\mathrm{\Sigma }^0`$ denote the class of all probability distributions $`P`$ on $`𝒵`$ such that
1. the margin assumption (MA) is satisfied,
2. the regression function $`\eta `$ belongs to the Hölder class $`\mathrm{\Sigma }(\beta ,L,^d)`$,
3. for all $`P𝒫`$ the supports of marginal distributions $`P_X`$ are included in $`𝒞`$.
If $`𝒞`$ is a compact the estimates of $`\epsilon `$-entropies of Hölder classes $`\mathrm{\Sigma }(\beta ,L,^d)`$ in the $`L_{\mathrm{}}(𝒞,\lambda )`$ norm can be obtained from Kolmogorov and Tikhomorov (1961), and they yield Assumption (CAR) with $`\rho =d/\beta `$. Therefore, from (4.2) we easily get the following upper bound.
###### Theorem 4.3
Let $`d1`$ be an integer, and let $`L,\beta ,\alpha `$ be positive constants. For any $`n1`$ the classifier $`\widehat{f}_n^S=1\mathrm{I}_{\{\widehat{\eta }_n^S1/2\}}`$ defined by (4.1) with $`p=\mathrm{}`$ satisfies
$$\underset{P𝒫_\mathrm{\Sigma }^0}{sup}\left\{𝔼R(\widehat{f}_n^S)R(f^{})\right\}Cn^{\frac{(1+\alpha )\beta }{(2+\alpha )\beta +d}}$$
with some constant $`C>0`$ depending only on $`\alpha `$, $`\beta `$, $`d`$, $`L`$ and $`C_0`$.
Since $`𝒫_\mathrm{\Sigma }^{}𝒫_\mathrm{\Sigma }^0`$, Theorems 3.5 and 4.3 show that $`n^{\frac{(1+\alpha )\beta }{(2+\alpha )\beta +d}}`$ is optimal rate of convergence of the excess risk on the class of distributions $`𝒫_\mathrm{\Sigma }^0`$.
## 5 Comparison lemmas
In this section we give some useful inequalities between the risks of plug-in classifiers and the $`L_p`$ risks of the corresponding regression estimators under the margin assumption (MA). These inequalities will be helpful in the proofs. They also illustrate a connection between the two complexity assumptions (CAR) and (CAD) defined in the Introduction and allow one to compare our study of plug-in estimators with that given by Yang (1999) who considered the case $`\alpha =0`$ (no margin assumption), as well as with the developments in Bartlett, Jordan and McAuliffe (2003) and Blanchard, Lugosi and Vayatis (2003).
Throughout this section $`\overline{\eta }`$ is a Borel function on $`^d`$ and
$$\overline{f}(x)=1\mathrm{I}_{\{\overline{\eta }(x)1/2\}}.$$
For $`1p\mathrm{}`$ we denote by $`_p`$ the $`L_p(^d,P_X)`$ norm. We first state some comparison inequalities for the $`L_{\mathrm{}}`$ norm.
###### Lemma 5.1
For any distribution $`P`$ of $`(X,Y)`$ satisfying Assumption $`(MA)`$ we have
$$R(\overline{f})R(f^{})2C_0\overline{\eta }\eta _{\mathrm{}}^{1+\alpha },$$
(5.1)
and
$$P_X\left(\overline{f}(X)f^{}(X),\eta (X)1/2\right)C_0\overline{\eta }\eta _{\mathrm{}}^\alpha .$$
(5.2)
Proof. To show (5.1) note that
$$\begin{array}{ccc}R(\overline{f})R(f^{})\hfill & =\hfill & 𝔼\left(|2\eta (X)1|1\mathrm{I}_{\{\overline{f}(X)f^{}(X)\}}\right)\hfill \\ & \hfill & 2𝔼\left(|\eta (X)\frac{1}{2}|1\mathrm{I}_{0<\{|\eta (X)\frac{1}{2}||\eta (X)\overline{\eta }(X)|\}}\right)\hfill \\ & \hfill & 2\eta \overline{\eta }_{\mathrm{}}P_X\left(0<|\eta (X)\frac{1}{2}|\eta \overline{\eta }_{\mathrm{}}\right)\hfill \\ & \hfill & 2C_0\eta \overline{\eta }_{\mathrm{}}^{1+\alpha }.\hfill \end{array}$$
Similarly,
$$\begin{array}{ccc}P_X\left(\overline{f}(X)f^{}(X),\eta (X)1/2\right)\hfill & \hfill & P_X\left(0<|\eta (X)\frac{1}{2}||\eta (X)\overline{\eta }(X)|\right)\hfill \\ & \hfill & P_X\left(0<|\eta (X)\frac{1}{2}|\eta \overline{\eta }_{\mathrm{}}\right)\hfill \\ & \hfill & C_0\eta \overline{\eta }_{\mathrm{}}^\alpha .\hfill \end{array}$$
###### Remark 5.1
Lemma 5.1 offers an easy way to obtain the result of Theorem 3.3 in a slightly less precise form, with an extra logarithmic factor in the rate. In fact, under the strong density assumption, there exist nonparametric estimators $`\widehat{\eta }_n`$ (for instance, suitably chosen locally polynomial estimators) such that
$$𝔼\left(\widehat{\eta }_n\eta _{\mathrm{}}^q\right)C\left(\frac{\mathrm{log}n}{n}\right)^{\frac{q\beta }{2\beta +d}},q>0,$$
uniformly over $`\eta \mathrm{\Sigma }(\beta ,L,^d)`$ \[see, e.g., Stone (1982)\]. Taking here $`q=1+\alpha `$ and applying the comparison inequality (5.1) we immediately get that the plug-in classifier $`\widehat{f}_n=1\mathrm{I}_{\{\widehat{\eta }_n1/2\}}`$ has the excess risk $`(\widehat{f}_n)`$ of the order $`\left(n/\mathrm{log}n\right)^{\beta (1+\alpha )/(2\beta +d)}`$.
Another immediate application of Lemma 5.1 is to get lower bounds on the risks of regression estimators in the $`L_{\mathrm{}}`$ norm from the corresponding lower bounds on the excess risks of classifiers (cf. Theorems 3.5 and 4.1). But here again we loose a logarithmic factor required for the best bounds.
We now consider the comparison inequalities for $`L_p`$ norms with $`1p<\mathrm{}`$.
###### Lemma 5.2
For any $`1p<\mathrm{}`$ and any distribution $`P`$ of $`(X,Y)`$ satisfying Assumption $`(MA)`$ with $`\alpha >0`$ we have
$$R(\overline{f})R(f^{})C_1(\alpha ,p)\eta \overline{\eta }_p^{\frac{p(1+\alpha )}{p+\alpha }},$$
(5.3)
and
$$P_X\left(\overline{f}(X)f^{}(X),\eta (X)1/2\right)C_2(\alpha ,p)\eta \overline{\eta }_p^{\frac{p}{p+\alpha }},$$
(5.4)
where $`C_1(\alpha ,p)=2(\alpha +p)p^1\left(\frac{p}{\alpha }\right)^{\frac{\alpha }{\alpha +p}}C_0^{\frac{p1}{\alpha +p}}`$, $`C_2(\alpha ,p)=(\alpha +p)p^1\left(\frac{p}{\alpha }\right)^{\frac{\alpha }{\alpha +p}}C_0^{\frac{p}{\alpha +p}}`$. In particular,
$$R(\overline{f})R(f^{})C_1(\alpha ,2)\left([\overline{\eta }(x)\eta (x)]^2P_X(dx)\right)^{\frac{1+\alpha }{2+\alpha }}.$$
(5.5)
Proof. For any $`t>0`$ we have
$`R(\overline{f})R(f^{})`$ (5.6)
$`=𝔼\left[|2\eta (X)1|1\mathrm{I}_{\{\overline{f}(X)f^{}(X)\}}\right]`$
$`=2𝔼\left[|\eta (X)1/2|1\mathrm{I}_{\{\overline{f}(X)f^{}(X)\}}1\mathrm{I}_{\{0<|\eta (X)1/2|t\}}\right]`$
$`+2𝔼\left[|\eta (X)1/2|1\mathrm{I}_{\{\overline{f}(X)f^{}(X)\}}1\mathrm{I}_{\{|\eta (X)1/2|>t\}}\right]`$
$`2𝔼\left[|\eta (X)\overline{\eta }(X)|1\mathrm{I}_{\{0<|\eta (X)1/2|t\}}\right]+2𝔼\left[|\eta (X)\overline{\eta }(X)|1\mathrm{I}_{\{|\eta (X)\overline{\eta }(X)|>t\}}\right]`$
$`2\eta \overline{\eta }_p\left[P_X(0<|\eta (X)1/2|t)\right]^{\frac{p1}{p}}+{\displaystyle \frac{2\eta \overline{\eta }_p^p}{t^{p1}}}`$
by Hölder and Markov inequalities. So, for any $`t>0`$, introducing $`E\eta \overline{\eta }_p`$ and using Assumption (MA) to bound the probability in (5.6) we obtain
$$R(\overline{f})R(f^{})2\left(C_0^{\frac{p1}{p}}t^{\frac{\alpha (p1)}{p}}E+\frac{E^p}{t^{p1}}\right).$$
Minimizing in $`t`$ the RHS of this inequality we get (5.3). Similarly,
$`P_X\left(\overline{f}(X)f^{}(X),\eta (X)1/2\right)`$ $``$ $`P_X\left(0<|\eta (X)1/2|t\right)+P_X(|\eta (X)\overline{\eta }(X)|>t)`$
$``$ $`C_0t^\alpha +{\displaystyle \frac{\eta \overline{\eta }_p^p}{t^p}},`$
and minimizing this bound in $`t`$ we obtain (5.4).
If the regression function $`\eta `$ belongs to the Hölder class $`\mathrm{\Sigma }(\beta ,L,^d)`$ there exist estimators $`\widehat{\eta }_n`$ such that, uniformly over the class,
$$𝔼\left\{\left[\widehat{\eta }_n(X)\eta (X)\right]^2\right\}Cn^{\frac{2\beta }{2\beta +d}}$$
(5.7)
for some constant $`C>0`$. This has been shown by Stone (1982) under the additional strong density assumption and by Yang (1999) with no assumption on $`P_X`$. Using (5.7) and (5.5) we get that the excess risk of the corresponding plug-in classifier $`\widehat{f}_n=1\mathrm{I}_{\{\widehat{\eta }_n1/2\}}`$ admits a bound of the order $`n^{\frac{2\beta }{2\beta +d}\frac{1+\alpha }{2+\alpha }}`$ which is suboptimal when $`\alpha 0`$ (cf. Theorems 4.2, 4.3). In other words, under the margin assumption, Lemma 5.2 is not the right tool to analyze the convergence rate of plug-in classifiers. On the contrary, when no margin assumption is imposed (i.e., $`\alpha =0`$ in our notation) inequality (1.2), which is a version of (5.5) for $`\alpha =0`$, is precise enough to give the optimal rate of classification \[Yang (1999)\].
Another way to obtain (5.5) is to use Bartlett, Jordan and McAuliffe (2003): it is enough to apply their Theorem 10 with (in their notation) $`\varphi (t)=(1t)^2,\psi (t)=t^2`$ and to note that for this choice of $`\varphi `$ we have $`R_\varphi (\overline{\eta })R_\varphi ^{}=\eta \overline{\eta }_2^2`$. Blanchard, Lugosi and Vayatis (2003) used the result of Bartlett, Jordan and McAuliffe (2003) to prove fast rates of the order $`n^{\frac{2(1+\alpha )}{3(2+\alpha )}}`$ for a boosting procedure over the class of regression functions $`\eta `$ of bounded variation in dimension $`d=1`$. Note that the same rates can be obtained for other plug-in classifiers using (5.5). Indeed, if $`\eta `$ is of bounded variation, there exist estimators of $`\eta `$ converging with the mean squared $`L_2`$ rate $`n^{2/3}`$\[cf. van de Geer (2000), Györfi et al. (2002)\], and thus application of (5.5) immediately yields the rate $`n^{\frac{2(1+\alpha )}{3(2+\alpha )}}`$ for the corresponding plug-in rule. However, Theorem 4.2 shows that this is not an optimal rate (here again we observe that inequality (5.5) fails to establish the optimal properties of plug-in classifiers). In fact, let $`d=1`$ and let the assumptions of Theorem 4.2 be satisfied, where instead of assumption (ii) we use its particular instance: $`\eta `$ belongs to a class of functions on $`[0,1]`$ whose total variation is bounded by a constant $`L<\mathrm{}`$. It follows from Birman and Solomjak (1967) that Assumption (CAR) for this class is satisfied with $`\rho =1`$ for any $`1p<\mathrm{}`$. Hence, we can apply (4.3) of Theorem 4.2 to find that
$$\underset{P𝒫}{sup}\left\{𝔼R(\widehat{f}_n^S)R(f^{})\right\}Cn^{\frac{(1+\alpha )p}{(2+\alpha )p+(p+\alpha )}}$$
(5.8)
for the corresponding class $`𝒫`$. If $`p>2`$ (recall that the value $`p[1,\mathrm{})`$ is chosen by the statistician), the rate in (5.8) is faster than $`n^{\frac{2(1+\alpha )}{3(2+\alpha )}}`$ obtained under the same conditions by Blanchard, Lugosi and Vayatis (2003).
## 6 Proofs
### 6.1 Proof of Theorem 3.2
Consider a distribution $`P`$ in $`𝒫_\mathrm{\Sigma }`$. Let $`A`$ be the support of $`P_X`$. Fix $`xA`$ and $`\delta >0`$. Consider the matrix $`B\left(B_{s_1,s_2}\right)_{|s_1|,|s_2|\beta }`$ with elements $`B_{s_1,s_2}_^du^{s_1+s_2}K(u)\mu (x+hu)𝑑u`$. The smallest eigenvalue $`\lambda _{\overline{B}}`$ of $`\overline{B}`$ satisfies
$$\begin{array}{ccc}\lambda _{\overline{B}}\hfill & =\hfill & \mathrm{min}_{W=1}W^T\overline{B}W\hfill \\ & \hfill & \mathrm{min}_{W=1}W^TBW+\mathrm{min}_{W=1}W^T(\overline{B}B)W\hfill \\ & \hfill & \mathrm{min}_{W=1}W^TBW_{|s_1|,|s_2|\beta }\left|\overline{B}_{s_1,s_2}B_{s_1,s_2}\right|.\hfill \end{array}$$
(6.1)
Let $`A_n\{u^d:uc;x+huA\}`$ where $`c`$ is the constant appearing in (3.3). Using (3.3), for any vector $`W`$ satisfying $`W=1`$, we obtain
$$\begin{array}{ccc}W^TBW\hfill & =\hfill & _^d\left(_{|s|\beta }W_su^s\right)^2K(u)\mu (x+hu)𝑑u\hfill \\ & \hfill & c\mu _{\mathrm{min}}_{A_n}\left(_{|s|\beta }W_su^s\right)^2𝑑u.\hfill \end{array}$$
By assumption of the theorem, $`chr_0`$. Since the support of the marginal distribution is $`(c_0,r_0)`$-regular we get
$$\begin{array}{ccc}\lambda [A_n]h^d\lambda \left[(x,ch)A\right]c_0h^d\lambda \left[(x,ch)\right]c_0v_dc^d,\hfill & & \end{array}$$
where $`v_d\lambda \left[(0,1)\right]`$ is the volume of the unit ball and $`c_0>0`$ is the constant introduced in the definition (2.1) of the $`(c_0,r_0)`$-regular set.
Let $`𝒜`$ denote the class of all compact subsets of $`(0,c)`$ having the Lebesgue measure $`c_0v_dc^d`$. Using the previous displays we obtain
$$\underset{W=1}{\mathrm{min}}W^TBWc\mu _{\mathrm{min}}\underset{W=1;S𝒜}{\mathrm{min}}_S\left(\underset{|s|\beta }{}W_su^s\right)^2𝑑u2\mu _0.$$
(6.2)
By the compactness argument, the minimum in (6.2) exists and is strictly positive.
For $`i=1,\mathrm{},n`$ and any multi-indices $`s_1,s_2`$ such that $`|s_1|,|s_2|\beta `$, define
$$\begin{array}{ccc}T_i\frac{1}{h^d}\left(\frac{X_ix}{h}\right)^{s_1+s_2}K\left(\frac{X_ix}{h}\right)_^du^{s_1+s_2}K(u)\mu (x+hu)𝑑u.\hfill & & \end{array}$$
We have $`𝔼T_i=0`$, $`|T_i|h^dsup_{u^d}\left(1+u^{2\beta }\right)K(u)\kappa _1h^d`$ and the following bound for the variance of $`T_i`$:
$$\begin{array}{ccc}𝕍\text{ar}T_i\hfill & \hfill & \frac{1}{h^{2d}}𝔼\left(\frac{X_ix}{h}\right)^{2s_1+2s_2}K^2\left(\frac{X_ix}{h}\right)\hfill \\ & =\hfill & \frac{1}{h^d}_^du^{2s_1+2s_2}K^2(u)\mu (x+hu)𝑑u\hfill \\ & \hfill & \frac{\mu _{\mathrm{max}}}{h^d}_^d\left(1+u^{4\beta }\right)K^2(u)𝑑u\frac{\kappa _2}{h^d}.\hfill \end{array}$$
From Bernstein’s inequality, we get
$$\begin{array}{ccc}P^n\left(|\overline{B}_{s_1,s_2}B_{s_1,s_2}|>ϵ\right)=P^n\left(\left|\frac{1}{n}_{i=1}^nT_i\right|>ϵ\right)2\mathrm{exp}\left\{\frac{nh^dϵ^2}{2\kappa _2+2\kappa _1ϵ/3}\right\}.\hfill & & \end{array}$$
This and (6.1) – (6.2) imply that
$$\begin{array}{ccc}P^n(\lambda _{\overline{B}}\mu _0)2M^2\mathrm{exp}\left(Cnh^d\right)\hfill & & \end{array}$$
(6.3)
where $`M^2`$ is the number of elements of the matrix $`\overline{B}`$. Assume in what follows that $`n`$ is large enough, so that $`\mu _0>(\mathrm{log}n)^1`$. Then for $`\lambda _{\overline{B}}>\mu _0`$ we have $`|\widehat{\eta }_n^{}(x)\eta (x)||\widehat{\eta }_n^{LP}(x)\eta (x)|`$. Therefore,
$$\begin{array}{ccc}P^n\left(\left|\widehat{\eta }_n^{}(x)\eta (x)\right|\delta \right)P^n\left(\lambda _{\overline{B}}\mu _0\right)+P^n\left(\left|\widehat{\eta }_n^{LP}(x)\eta (x)\right|\delta ,\lambda _{\overline{B}}>\mu _0\right).\hfill & & \end{array}$$
(6.4)
We now evaluate the second probability on the right hand side of (6.4). For $`\lambda _{\overline{B}}>\mu _0`$ we have $`\widehat{\eta }_n^{LP}(x)=U^T(0)Q^1V`$ (where $`V`$ is given by (2.3)). Introduce the matrix $`Z\left(Z_{i,s}\right)_{1in,|s|\beta }`$ with elements
$$\begin{array}{ccc}Z_{i,s}(X_ix)^s\sqrt{K\left(\frac{X_ix}{h}\right)}.\hfill & & \end{array}$$
The $`s`$-th column of $`Z`$ is denoted by $`Z_s`$, and we introduce $`Z^{(\eta )}_{|s|\beta }\frac{\eta ^{(s)}(x)}{s!}Z_s.`$ Since $`Q=Z^TZ`$, we get
$$\begin{array}{ccc}|s|\beta :U^T(0)Q^1Z^TZ_s=1\mathrm{I}_{\{s=(0,\mathrm{},0)\}},\hfill & & \end{array}$$
hence $`U^T(0)Q^1Z^TZ^{(\eta )}=\eta (x).`$ So we can write
$$\begin{array}{ccc}\widehat{\eta }_n^{LP}(x)\eta (x)=U^T(0)Q^1(VZ^TZ^{(\eta )})=U^T(0)\overline{B}^1𝐚\hfill & & \end{array}$$
where $`𝐚\frac{1}{nh^d}H(VZ^TZ^{(\eta )})^M`$ and $`H`$ is a diagonal matrix $`H\left(H_{s_1,s_2}\right)_{|s_1|,|s_2|\beta }`$ with $`H_{s_1,s_2}h^{s_1}1\mathrm{I}_{\{s_1=s_2\}}.`$ For $`\lambda _{\overline{B}}>\mu _0`$ we get
$$\begin{array}{ccc}\left|\widehat{\eta }_n^{LP}(x)\eta (x)\right|\overline{B}^1𝐚\lambda _{\overline{B}}^1𝐚\mu _0^1𝐚\mu _0^1M\mathrm{max}_s|a_s|,\hfill & & \end{array}$$
(6.5)
where $`a_s`$ are the components of the vector $`𝐚`$ given by
$$\begin{array}{ccc}a_s=\frac{1}{nh^d}_{i=1}^n\left[Y_i\eta _x(X_i)\right]\left(\frac{X_ix}{h}\right)^sK\left(\frac{X_ix}{h}\right).\hfill & & \end{array}$$
Define
$$\begin{array}{ccc}T_i^{(s,1)}\hfill & \hfill & \frac{1}{h^d}\left[Y_i\eta (X_i)\right]\left(\frac{X_ix}{h}\right)^sK\left(\frac{X_ix}{h}\right),\hfill \\ T_i^{(s,2)}\hfill & \hfill & \frac{1}{h^d}\left[\eta (X_i)\eta _x(X_i)\right]\left(\frac{X_ix}{h}\right)^sK\left(\frac{X_ix}{h}\right).\hfill \end{array}$$
We have
$$\begin{array}{ccc}|a_s|\left|\frac{1}{n}_{i=1}^nT_i^{(s,1)}\right|+\left|\frac{1}{n}_{i=1}^n\left[T_i^{(s,2)}𝔼T_i^{(s,2)}\right]\right|+\left|𝔼T_i^{(s,2)}\right|.\hfill & & \end{array}$$
(6.6)
Note that $`𝔼T_i^{(s,1)}=0`$, $`\left|T_i^{(s,1)}\right|\kappa _1h^d`$, and
$$\begin{array}{ccc}𝕍\text{ar}T_i^{(s,1)}\hfill & \hfill & 4^1h^du^{2s}K^2(u)\mu (x+hu)𝑑u(\kappa _2/4)h^d,\hfill \\ \left|T_i^{(s,2)}𝔼T_i^{(s,2)}\right|\hfill & \hfill & L\kappa _1h^{\beta d}+L\kappa _2h^\beta Ch^{\beta d},\hfill \\ 𝕍\text{ar}T_i^{(s,2)}\hfill & \hfill & h^dL^2h^{2\beta }u^{2s+2\beta }K^2(u)\mu (x+hu)𝑑uL^2\kappa _2h^{2\beta d}.\hfill \end{array}$$
From Bernstein’s inequality, for any $`ϵ_1,ϵ_2>0`$, we obtain
$$\begin{array}{ccc}P^n\left(\left|\frac{1}{n}_{i=1}^nT_i^{(s,1)}\right|ϵ_1\right)2\mathrm{exp}\left\{\frac{nh^dϵ_1^2}{\kappa _2/2+2\kappa _1ϵ_1/3}\right\}\hfill & & \end{array}$$
and
$$\begin{array}{ccc}P^n\left(\left|\frac{1}{n}_{i=1}^n\left[T_i^{(s,2)}𝔼T_i^{(s,2)}\right]\right|ϵ_2\right)2\mathrm{exp}\left\{\frac{nh^dϵ_2^2}{2L^2\kappa _2h^{2\beta }+2Ch^\beta ϵ_2/3}\right\}.\hfill & & \end{array}$$
Since also
$$\left|𝔼T_i^{(s,2)}\right|Lh^\beta u^{s+\beta }K^2(u)\mu (x+hu)𝑑uL\kappa _2h^\beta $$
we get, using (6.6), that if $`3\mu _0^1ML\kappa _2h^\beta \delta 1`$ the following inequality holds
$$\begin{array}{ccc}P^n\left(|a_s|\frac{\mu _0\delta }{M}\right)\hfill & \hfill & P^n\left(\left|\frac{1}{n}_{i=1}^nT_i^{(s,1)}\right|>\frac{\mu _0\delta }{3M}\right)+P^n\left(\left|\frac{1}{n}_{i=1}^n\left[T_i^{(s,2)}𝔼T_i^{(s,2)}\right]\right|>\frac{\mu _0\delta }{3M}\right)\hfill \\ & \hfill & 4\mathrm{exp}\left(Cnh^d\delta ^2\right).\hfill \end{array}$$
Combining this inequality with (6.3), (6.4) and (6.5), we obtain
$$\begin{array}{ccc}P^n\left(\left|\widehat{\eta }_n^{}(x)\eta (x)\right|\delta \right)C_1\mathrm{exp}\left(C_2nh^d\delta ^2\right)\hfill & & \end{array}$$
(6.7)
for $`3m^1ML\kappa _2h^\beta \delta `$ (for $`\delta >1`$ inequality (6.7) is obvious since $`\widehat{\eta }_n^{},\eta `$ take values in $`[0,1]`$). The constants $`C_1,C_2`$ in (6.7) do not depend on the distribution $`P_X`$, on its support $`A`$ and on the point $`xA`$, so that we get (3.7). Now, (3.7) implies (3.8) for $`Cn^{\frac{\beta }{2\beta +d}}\delta `$, and thus for all $`\delta >0`$ (with possibly modified constants $`C_1`$ and $`C_2`$).
### 6.2 Proof of Theorems 3.5 and 4.1
The proof of both theorems is based on Assouad’s lemma \[see, e.g., Korostelev and Tsybakov (1993), Chapter 2 or Tsybakov (2004b), Chapter 2\]. We apply it in a form adapted for the classification problem (Lemma 5.1 in Audibert (2004)).
For an integer $`q1`$ we consider the regular grid on $`^d`$ defined as
$$G_q\{(\frac{2k_1+1}{2q},\mathrm{},\frac{2k_d+1}{2q}):k_i\{0,\mathrm{},q1\},i=1,\mathrm{},d\}.$$
Let $`n_q(x)G_q`$ be the closest point to $`x^d`$ among points in $`G_q`$ (we assume uniqueness of $`n_q(x)`$: if there exist several points in $`G_q`$ closest to $`x`$ we define $`n_q(x)`$ as the one which is closest to 0). Consider the partition $`𝒳_1^{},\mathrm{},𝒳_{q^d}^{}`$ of $`[0,1]^d`$ canonically defined using the grid $`G_q`$ ($`x`$ and $`y`$ belong to the same subset if and only if $`n_q(x)=n_q(y)`$). Fix an integer $`mq^d`$. For any $`i\{1,\mathrm{},m\}`$, we define $`𝒳_i𝒳_i^{}`$ and $`𝒳_0^d_{i=1}^m𝒳_i`$, so that $`𝒳_0,\mathrm{},𝒳_m`$ form a partition of $`^d`$.
Let $`u:_+_+`$ be a nonincreasing infinitely differentiable function such that $`u=1`$ on $`[0,1/4]`$ and $`u=0`$ on $`[1/2,\mathrm{})`$. One can take, for example, $`u(x)=\left(_{1/4}^{1/2}u_1(t)𝑑t\right)^1_x^{\mathrm{}}u_1(t)𝑑t`$ where the infinitely differentiable function $`u_1`$ is defined as
$$\begin{array}{ccc}u_1(x)=\{\begin{array}{cc}\mathrm{exp}\left\{\frac{1}{(1/2x)(x1/4)}\right\}\hfill & \text{for }x(1/4,1/2),\hfill \\ 0\hfill & \text{otherwise.}\hfill \end{array}\hfill & & \end{array}$$
Let $`\varphi :^d_+`$ be the function defined as
$$\begin{array}{ccc}\varphi (x)C_\varphi u(x),\hfill & & \end{array}$$
where the positive constant $`C_\varphi `$ is taken small enough so ensure that $`|\varphi (x^{})\varphi _x(x^{})|Lx^{}x^\beta `$ for any $`x,x^{}^d`$. Thus, $`\varphi \mathrm{\Sigma }(\beta ,L,^d)`$.
Define the hypercube $`=\{_\stackrel{}{\sigma }:\stackrel{}{\sigma }=(\sigma _1,\mathrm{},\sigma _m)\{1,1\}^m\}`$ of probability distributions $`_\stackrel{}{\sigma }`$ of $`(X,Y)`$ on $`𝒵=^d\times \{0,1\}`$ as follows.
For any $`_\stackrel{}{\sigma }`$ the marginal distribution of $`X`$ does not depend on $`\stackrel{}{\sigma }`$, and has a density $`\mu `$ w.r.t. the Lebesgue measure on $`^d`$ defined in the following way. Fix $`0<wm^1`$ and a set $`A_0`$ of positive Lebesgue measure included in $`𝒳_0`$ (the particular choices of $`A_0`$ will be indicated later), and take: (i) $`\mu (x)=w/\lambda [(0,(4q)^1)]`$ if $`x`$ belongs to a ball $`(z,(4q)^1)`$ for some $`zG_d`$, (ii) $`\mu (x)=(1mw)/\lambda [A_0]`$ for $`xA_0`$, (iii) $`\mu (x)=0`$ for all other $`x`$.
Next, the distribution of $`Y`$ given $`X`$ for $`_\stackrel{}{\sigma }`$ is determined by the regression function $`\eta _\stackrel{}{\sigma }(x)=P(Y=1|X=x)`$ that we define as $`\eta _\stackrel{}{\sigma }(x)=\frac{1+\sigma _j\phi (x)}{2}`$ for any $`x𝒳_j`$, $`j=1,\mathrm{},m`$, and $`\eta _\stackrel{}{\sigma }1/2`$ on $`𝒳_0`$, where $`\phi (x)q^\beta \varphi \left(q[xn_q(x)]\right).`$ We will assume that $`C_\varphi 1`$ to ensure that $`\phi `$ and $`\eta _\stackrel{}{\sigma }`$ take values in $`[0,1]`$.
For any $`s^d`$ such that $`|s|\beta `$, the partial derivative $`D^s\phi `$ exists, and $`D^s\phi (x)=q^{|s|\beta }D^s\varphi \left(q[xn_q(x)]\right)`$. Therefore, for any $`i\{1,\mathrm{},m\}`$ and any $`x,x^{}𝒳_i`$, we have
$$|\phi (x^{})\phi _x(x^{})|Lxx^{}^\beta .$$
This implies that for any $`\stackrel{}{\sigma }\{1,1\}^m`$ the function $`\eta _\stackrel{}{\sigma }`$ belongs to the Hölder class $`\mathrm{\Sigma }(\beta ,L,^d)`$.
We now check the margin assumption. Set $`x_0=(\frac{1}{2q},\mathrm{},\frac{1}{2q})`$. For any $`\stackrel{}{\sigma }\{1,1\}^m`$ we have
$$\begin{array}{ccc}_\stackrel{}{\sigma }\left(0<\left|\eta _\stackrel{}{\sigma }(X)1/2\right|t\right)\hfill & =\hfill & m_\stackrel{}{\sigma }\left(0<\varphi [q(Xx_0)]2tq^\beta \right)\hfill \\ & =\hfill & m_{(x_0,(4q)^1)}1\mathrm{I}_{\{0<\varphi [q(xx_0)]2tq^\beta \}}\frac{w}{\lambda [(0,(4q)^1)]}𝑑x\hfill \\ & =\hfill & \frac{mw}{\lambda [(0,1/4)]}_{(0,1/4)}1\mathrm{I}_{\{\varphi (x)2tq^\beta \}}𝑑x\hfill \\ & =\hfill & mw1\mathrm{I}_{\{tC_\varphi /(2q^\beta )\}}.\hfill \end{array}$$
Therefore, the margin assumption (MA) is satisfied if $`mw=O(q^{\alpha \beta }).`$
According to Lemma 5.1 in Audibert (2004), for any classifier $`\widehat{f}_n`$ we have
$$\underset{P}{sup}\left\{𝔼R(\widehat{f}_n)R(f^{})\right\}mwb^{}(1b\sqrt{nw})/2$$
(6.8)
where
$$\begin{array}{ccc}b\hfill & \hfill & \left[1\left(_{𝒳_1}\sqrt{1\phi ^2(x)}\mu _1(x)𝑑x\right)^2\right]^{1/2}=C_\varphi q^\beta ,\hfill \\ b^{}\hfill & \hfill & _{𝒳_1}\phi (x)\mu _1(x)𝑑x=C_\varphi q^\beta \hfill \end{array}$$
with $`\mu _1(x)=\mu (x)/_{𝒳_1}\mu (z)𝑑z`$.
We now prove Theorem 3.5. Take $`q=\overline{C}n^{\frac{1}{2\beta +d}}`$, $`w=C^{}q^d`$ and $`m=C^{\prime \prime }q^{d\alpha \beta }`$ with some positive constants $`\overline{C},C^{},C^{\prime \prime }`$ to be chosen, and set $`A_0=[0,1]^d_{i=1}^m𝒳_i`$. The condition $`\alpha \beta d`$ ensures that the above choice of $`m`$ is not degenerate: we have $`m1`$ for $`C^{\prime \prime }`$ large enough. We now prove that $`𝒫_\mathrm{\Sigma }`$ under the appropriate choice of $`\overline{C},C^{},C^{\prime \prime }`$. In fact, select these constants so that the triplet $`(q,w,m)`$ meets the conditions $`mq^d`$, $`0<wm^1`$, $`mw=O(q^{\alpha \beta })`$. Then, in view of the argument preceding (6.8), for any $`\stackrel{}{\sigma }\{1,1\}^m`$ the regression function $`\eta _\stackrel{}{\sigma }`$ belongs to $`\mathrm{\Sigma }(\beta ,L,^d)`$ and Assumption (MA) is satisfied. We now check that $`P_X`$ obeys the strong density assumption. First, the density $`\mu (x)`$ equals to a positive constant for $`x`$ belonging to the union of balls $`_{i=1}^m(z_i,(4q)^1)`$ where $`z_i`$ is the center of $`𝒳_i`$, and $`\mu (x)=(1mw)/(1mq^d)=1+o(1)`$, as $`n\mathrm{}`$, for $`xA_0`$. Thus, $`\mu _{\mathrm{min}}\mu (x)\mu _{\mathrm{max}}`$ for some positive $`\mu _{\mathrm{min}}`$ and $`\mu _{\mathrm{max}}`$. (Note that this construction does not allow to choose any prescribed values of $`\mu _{\mathrm{min}}`$ and $`\mu _{\mathrm{max}}`$, because $`\mu (x)=1+o(1)`$. The problem can be fixed via a straightforward but cumbersome modification of the definition of $`A_0`$ that we skip here.) Second, the $`(c_0,r_0)`$-regularity of the support $`A`$ of $`P_X`$ with some $`c_0>0`$ and $`r_0>0`$ follows from the fact that, by construction, $`\lambda (A(x,r))=(1+o(1))\lambda ([0,1]^d(x,r))`$ for all $`xA`$ and $`r>0`$ (here again we skip the obvious generalization allowing to get any prescribed $`c_0>0`$). Thus, the strong density assumption is satisfied, and we conclude that $`𝒫_\mathrm{\Sigma }`$. Theorem 3.5 now follows from (6.8) if we choose $`C^{}`$ small enough.
Finally, we prove Theorem 4.1. Take $`q=Cn^{\frac{1}{(2+\alpha )\beta +d}}`$, $`w=C^{}q^{2\beta }/n`$ and $`m=q^d`$ for some constants $`C>0`$, $`C^{}>0`$, and choose $`A_0`$ as a Euclidean ball contained in $`𝒳_0`$. As in the proof of Theorem 3.5, under the appropriate choice of $`C`$ and $`C^{}`$, the regression function $`\eta _\stackrel{}{\sigma }`$ belongs to $`\mathrm{\Sigma }(\beta ,L,^d)`$ and the margin assumption (MA) is satisfied. Moreover, it is easy to see that the marginal distribution of $`X`$ obeys the mild density assumption (the $`(c_0,r_0)`$-regularity of the support of $`P_X`$ follows from considerations analogous to those in the proof of Theorem 3.5). Thus, $`𝒫_\mathrm{\Sigma }^{}`$. Choosing $`C^{}`$ small enough and using (6.8) we obtain Theorem 4.1.
### 6.3 Proof of Proposition 3.4
The following lemma describes how the smoothness constraint on the regression function $`\eta `$ at some point $`x^d`$ implies that $`\eta `$ “stays close” to $`\eta (x)`$ in the vicinity of $`x`$.
###### Lemma 6.1
For any distribution $`P𝒫_\mathrm{\Sigma }`$ with regression function $`\eta `$ and for any $`\kappa >0`$, there exist $`L^{}>0`$ and $`t_0>0`$ such that for any $`x`$ in the support of $`P_X`$ and $`0<tt_0`$, we have
$$P_X\left[\left|\eta (X)\eta (x)\right|t;X(x,\kappa t^{\frac{1}{1\beta }})\right]L^{}t^{\frac{d}{1\beta }}.$$
Proof of Lemma 6.1. Let $`A`$ denote the support of $`P_X`$. Let us first consider the case $`\beta 1`$. Then for any $`x,x^{}^d`$, we have $`\left|\eta (x^{})\eta (x)\right|Lx^{}x^\beta .`$ Let $`\kappa ^{}=\kappa L^{1/\beta }`$. For any $`0<tLr_0^\beta `$, we get
$$\begin{array}{ccc}P_X\left[\left|\eta (X)\eta (x)\right|t;X(x,\kappa t^{\frac{1}{1\beta }})\right]\hfill & & \\ =P_X\left[\left|\eta (X)\eta (x)\right|t;X(x,\kappa t^{\frac{1}{\beta }})A\right]\hfill & & \\ P_X\left[X(x,\kappa t^{\frac{1}{\beta }}\left(\frac{t}{L}\right)^{\frac{1}{\beta }})A\right]\hfill & & \\ \mu _{\mathrm{min}}\lambda \left[(x,\kappa ^{}t^{\frac{1}{\beta }})A\right]\hfill & & \\ c_0\mu _{\mathrm{min}}\lambda \left[(x,\kappa ^{}t^{\frac{1}{\beta }})\right]\hfill & & \\ c_0\mu _{\mathrm{min}}v_d(\kappa ^{})^dt^{\frac{d}{\beta }},\hfill & & \end{array}$$
which is the desired result with $`L^{}c_0\mu _{\mathrm{min}}v_d(\kappa ^{})^d`$ and $`t_0Lr_0^\beta `$.
For the case $`\beta >1`$, by assumption, $`\eta `$ is continuously differentiable. Let $`𝒞(A)`$ be the convex hull of the support $`A`$ of $`P_X`$. By compactness of $`𝒞(A)`$, there exists $`C>0`$ such that for any $`s^d`$ with $`|s|=1`$,
$$\underset{x𝒞(A)}{sup}\left|D^s\eta (x)\right|C.$$
So we have for any $`x,x^{}A`$,
$$|\eta (x)\eta (x^{})|Cxx^{}.$$
The rest of the proof is then similar to the one of the first case.
* We will now prove the first item of Proposition 3.4. Let $`P𝒫_\mathrm{\Sigma }`$ such that the regression function associated with $`P`$ hits $`1/2`$ at $`x_0\stackrel{}{A}`$, where $`\stackrel{}{A}`$ denotes the interior of the support of $`P_X`$. Let $`r>0`$ such that $`(x_0,r)A`$. Let $`x(x_0,r)`$ such that $`\eta (x)\frac{1}{2}`$. Let $`t_1=\left|\eta (x)1/2\right|.`$ For any $`0<tt_1`$, let $`x_t[x_0;x]`$ such that $`\left|\eta (x_t)1/2\right|=t/2.`$ We have $`x_tA`$ so that we can apply Lemma 6.1 (with $`\kappa =1`$ for instance) and obtain for any $`0<tt_1(4t_0)`$
$$P_X\left[0<\left|\eta (X)1/2\right|t\right]P_X\left[\left|\eta (X)\eta (x_t)\right|t/4\right]L^{}(t/4)^{\frac{d}{1\beta }}.$$
Now from the margin assumption, we get that for any small enough $`t>0`$ $`C_0t^\alpha L^{}(t/4)^{\frac{d}{1\beta }},`$ hence $`\alpha \frac{d}{1\beta }.`$
* For the second item of Proposition 3.4, to skip cumbersome details, we may assume that $`𝒞`$ contains the unit ball in $`^d`$. Consider the distribution such that
+ $`P_X`$ is the uniform measure on $`\{(x_1,\mathrm{},x_d)^d:|x_11/4|+|x_2|+\mathrm{}+|x_d|1/4\}`$
+ the regression function associated with $`P`$ is
$$\eta (x_1,\mathrm{},x_d)=\frac{1+C_\eta \text{sign}(x_1)|x_1|^{\beta 1}u(x_1)}{2},$$
where
$$u(t)=\{\begin{array}{ccc}\mathrm{exp}\left(\frac{1}{1t^2}\right)\hfill & \text{if }|t|<1\hfill & \\ 0\hfill & \text{otherwise},\hfill & \end{array}$$
and $`0<C_\eta 1`$ is small enough so that for any $`x,x^{}^d`$, $`\eta `$ satisfies
$$|\eta (x^{})\eta _x(x^{})|Lxx^{}^\beta .$$
For appropriate positive parameters $`c_0,r_0,\mu _{\mathrm{max}}>\mu _{\mathrm{min}}>0`$, the only non-trivial task in checking that $`P`$ belongs to $`𝒫_\mathrm{\Sigma }`$ is to check the margin assumption. For $`t`$ small enough, we have
$$P_X\left[\left|\eta (X)1/2\right|t\right]P_X\left[|X_1|^{\beta 1}Ct;|X_11/4|+|X_2|+\mathrm{}+|X_d|1/4\right]$$
for some $`C>0`$. Therefore, we have $`P_X\left[0<\left|\eta (X)1/2\right|t\right]Ct^{\frac{d}{\beta 1}}`$. So the margin assumption is satisfied for an appropriate $`C_0`$ whenever $`\alpha \frac{d}{\beta 1}.`$ Since $`\eta `$ hits $`1/2`$ at $`0_^d`$ which is in boundary of the support of $`P_X`$, we have proved the second assertion.
* For the third assertion of Proposition 3.4, to avoid cumbersome details again, we may assume that $`𝒞`$ contains the unit ball in $`^d`$. Consider the distribution such that
+ $`P_X`$ is the uniform measure on the unit ball,
+ the regression function associated with $`P`$ is
$$\eta (x)=\frac{1+C_\eta x^2u(x^2/2)}{2},$$
where $`0<C_\eta 1`$ is small enough so that for any $`x,x^{}^d`$, $`\eta `$ satisfies
$$|\eta (x^{})\eta _x(x^{})|Lxx^{}^\beta .$$
For appropriate positive parameters $`C_0,c_0,r_0,\mu _{\mathrm{max}}>\mu _{\mathrm{min}}>0`$, the distribution $`P`$ belongs to $`𝒫_\mathrm{\Sigma }`$ provided that $`\alpha d/2`$ (in order that the margin assumption holds). We have obtained the desired result since $`\eta `$ hits $`1/2`$ at $`0_^d`$ which is in the interior of the support of $`P_X`$.
* For the last item of Proposition 3.4, let $`P𝒫_\mathrm{\Sigma }`$ such that the regression function $`\eta `$ associated with $`P`$ crosses $`1/2`$ at $`x_0\stackrel{}{A}`$. For $`d=1`$, from the first item of the theorem, we necessarily have $`\alpha (\beta 1)1`$. Let us now consider the case: $`d>1`$.
Figure 1 will help to keep track of the following notation.
Let $`r_1>0`$ such that $`(x_0,3r_1)A`$. Introduce $`x_{}`$ and $`x_+`$ in $`(x_0,r_1)`$ such that $`\eta (x_{})<1/2`$ and $`\eta (x_+)>1/2.`$ Let $`t_1=\left(1/2\eta (x_{})\right)\left(\eta (x_+)1/2\right).`$ Define $`y=\frac{x_{}+x_+}{2}`$, $`e_d=\frac{x_+x_{}}{x_+x_{}}`$ and $`D=x_+x_{}`$. Let $`e_1,\mathrm{},e_{d1}`$ be unit vectors such that $`e_1,\mathrm{},e_d`$ is an orthonormal basis of $`^d`$. Let $`^{}(x,r)`$ (resp. $`𝒮^{}(x,r)`$) denote the ball (resp. the sphere) centered at $`x`$ and of radius $`r`$ wrt the norm $`x_{}=sup_{1id}|x,e_i|.`$
Since $`\eta `$ is continuous, there exists $`r_2>0`$ such that
$$\{\begin{array}{ccc}\eta (x)<1/2t_1/2\hfill & \text{for any }x^{}(x_{},r_2)\hfill & \\ \eta (x)>1/2+t_1/2\hfill & \text{for any }x^{}(x_+,r_2)\hfill & \end{array}$$
Let $`\zeta =\frac{1}{\beta 1}`$. For any $`k=(k_1,\mathrm{},k_{d1})^{d1}`$, introduce
$$y_k=y+t^\zeta \underset{i=1}{\overset{d1}{}}k_ie_i.$$
For any $`t`$ in $`]0;t_1[`$, consider the grid $`G=\{y_k;k^{d1},\underset{1id1}{\mathrm{max}}|k_i|\frac{D}{2\sqrt{d1}t^\zeta }\}.`$ For any $`y_k`$ in $`G`$, we have $`y_ky\sqrt{d1}\underset{1id1}{\mathrm{max}}|t^\zeta k_i|D/2r_1`$. Therefore, using that $`y(x_0,r_1)`$, the grid $`G`$ is included in $`(x_0,2r_1)`$. For any $`y_kG`$, let $`y_k^{}=[x_{};y_k]𝒮^{}(x_{},r_2)`$ and $`y_k^+=[x_+;y_k]𝒮^{}(x_+,r_2)`$. Since $`y_kyD/2`$, we have $`y_k^{}=x_{}+r_2e_d+\frac{2r_2}{D}t^\zeta _{i=1}^{d1}k_ie_i`$ and $`y_k^+=x_+r_2e_d+\frac{2r_2}{D}t^\zeta _{i=1}^{d1}k_ie_i`$.
For any $`y_k`$ in $`G`$, consider the continuous path formed by the segments $`[y_k^{};y_k]`$ and $`[y_k;y_k^+]`$. Since $`\eta `$ is continuous on this path, there exists $`w_k\gamma _k[y_k^{};y_k][y_k;y_k^+]`$ such that $`\eta (w_k)=1/2+t/2`$. Now let us show that when $`kk^{}`$, $`w_k`$ and $`w_k^{}`$ are at least $`\frac{\sqrt{2}r_2}{D}t^\zeta `$ away from each other. The distance between $`w_k`$ and $`w_k^{}`$ is not less than the distance between the paths $`\gamma _k`$ and $`\gamma _k^{}`$. Let $`U`$ denote the biggest integer smaller than or equal to $`\frac{D}{2\sqrt{d1}t^\zeta }`$. When $`y_ky_k^{}`$ in $`G`$, the distance between $`\gamma _k`$ and $`\gamma _k^{}`$ is minimum for $`k=K(U,\mathrm{},U)`$ and $`k^{}=K^{}(U1,U,\mathrm{},U)`$. This distance is equal to the distance between $`y_K^{}`$ and its orthogonal projection on $`[y_K^{}^{};y_K^{}]`$, which is the distance between $`y_K^{}`$ and the line $`(x_{};y_K^{})`$. Let $`K^{\prime \prime }=(0,U,\mathrm{},U)^{d1}`$. To compute this distance $`V`$, it suffices to look at the plane $`(x_{};y_{K^{\prime \prime }};y_K)`$ (see figure 2).
We obtain that the angle $`\theta `$ between $`y_K^{}x_{}`$ and $`y_{K^{\prime \prime }}x_+`$ is smaller than $`\pi /4`$. As a consequence, $`V=y_K^{}y_K^{}^{}\mathrm{cos}\theta \sqrt{2}r_2t^\zeta /D`$.
Finally, focusing on the behaviour of the regression function near the $`w_k`$’s, by using Lemma 6.1 with $`\kappa =\frac{4^\zeta r_2}{\sqrt{2}D}`$, we obtain that there exists $`L^{}>0`$ and $`t_0>0`$ such that for any $`0<t<4t_0t_1`$,
$$\begin{array}{ccc}C_0t^\alpha \hfill & \hfill & P_X\left[0<\left|\eta (X)\frac{1}{2}\right|t\right]\hfill \\ & \hfill & \underset{k^{d1}:\underset{1id1}{\mathrm{max}}|k_i|\frac{D}{2\sqrt{d1}t^\zeta }}{}P_X\left[\left|\eta (X)\eta (w_k)\right|t/4;X(w_k,\frac{r_2t^\zeta }{\sqrt{2}D})\right]\hfill \\ & \hfill & (2U+1)^{d1}L^{}(t/4)^{d\zeta }\hfill \\ & \hfill & \left(\frac{D}{2\sqrt{d1}t^\zeta }\right)^{d1}L^{}(t/4)^{d\zeta }\hfill \\ & \hfill & Ct^\zeta ,\hfill \end{array}$$
hence $`\alpha \zeta `$ (which is the desired result).
For the converse, the proof is similar to the ones of the second and third assertions of the proposition. Without loss of generality, we may assume that $`𝒮=\{(x_1,\mathrm{},x_d)^d:\underset{1id}{\mathrm{max}}|x_i|1/2\}`$ is a subset of $`𝒞`$. we consider the distribution $`P`$ such that
+ $`P_X`$ is the uniform measure on $`𝒮`$
+ the regression function associated with $`P`$ is
$$\eta (x_1,\mathrm{},x_d)=\frac{1+C_\eta \text{sign}(x_1)|x_1|^{\beta 1}u(x_1)}{2},$$
where $`0<C_\eta 1`$ is small enough so that for any $`x,x^{}^d`$, $`\eta `$ satisfies
$$|\eta (x^{})\eta _x(x^{})|Lxx^{}^\beta .$$
For small enough $`t>0`$, we have
$$P_X\left[\left|\eta (X)1/2\right|t\right]P_X\left[|X_1|^{\beta 1}Ct\right],$$
for some constant $`C>0`$, so that we have $`P_X\left[0<\left|\eta (X)\frac{1}{2}\right|t\right]2(Ct)^{\frac{1}{\beta 1}}.`$ As a consequence, for appropriate parameters $`C_0,c_0,r_0,\mu _{\mathrm{max}}>\mu _{\mathrm{min}}>0`$, the distribution $`P`$ belongs to $`𝒫_\mathrm{\Sigma }`$ whenever $`\alpha \frac{1}{\beta 1}.`$ Since $`\eta `$ crosses $`1/2`$ at $`0_^d`$ which is in the interior of the support of $`P_X`$, the converse holds.
### 6.4 Proof of Theorem 4.2
We prove the theorem for $`p<\mathrm{}`$. The proof for $`p=\mathrm{}`$ is analogous. For any decision rule $`f`$ we set $`d(f)R(f)R(f^{})`$ and
$$f^{}(x,f)\{\begin{array}{ccc}f^{}(x)\hfill & \text{if}& \hfill \eta (x)1/2,\\ f(x)\hfill & \text{if}& \hfill \eta (x)=1/2,\end{array}x^d.$$
###### Lemma 6.2
Under Assumption (MA) for any decision rule $`f`$ we have
$$P_X(f(X)f^{}(X,f))Cd(f)^{\alpha /(1+\alpha )}.$$
(6.9)
Proof. Note that $`f^{}(,f)`$ is a Bayes rule, and following the same lines as in Proposition 1 of Tsybakov (2004a) we get $`P_X(f(X)f^{}(X,f),\eta (X)1/2)Cd(f)^{\alpha /(1+\alpha )}`$. It remains to observe that $`P_X(f(X)f^{}(X,f),\eta (X)1/2)=P_X(f(X)f^{}(X,f))`$.
For a Borel function $`\overline{\eta }`$ on $`^d`$ define $`f_{\overline{\eta }}1\mathrm{I}_{\{\overline{\eta }1/2\}}`$, $`f_{\overline{\eta }}^{}()f^{}(,f_{\overline{\eta }})`$ and
$$Z_n(f_{\overline{\eta }})[R_n(f_{\overline{\eta }})R_n(f_{\overline{\eta }}^{})][R(f_{\overline{\eta }})R(f_{\overline{\eta }}^{})]=[R_n(f_{\overline{\eta }})R_n(f_{\overline{\eta }}^{})]d(f_{\overline{\eta }}).$$
Let $`\eta _n`$ be an element of $`𝒩_{\epsilon _n}`$ such that $`\eta _n\eta _{p,\lambda }\epsilon _n`$, where $`_{p,\lambda }`$ is the $`L_p(𝒞,\lambda )`$ norm. In view of the assumption on $`𝒫`$ we have $`\eta _n\eta _p\mu _{\mathrm{max}}^{1/p}\epsilon _n`$ where $`_p`$ is the $`L_p(^d,P_X)`$ norm. It follows from the comparison inequality (5.3) that $`d(f_{\eta _n})C\epsilon _n^{\frac{(1+\alpha )p}{p+\alpha }}\delta _n`$. Set
$$\mathrm{\Delta }_n=Cn^{\frac{(1+\alpha )p}{(2+\alpha )p+\rho (p+\alpha )}}$$
(i.e., $`\mathrm{\Delta }_n`$ is of the order of desired rate). Fix $`t>0`$ and introduce the set
$$𝒩_n^{}=\{\overline{\eta }𝒩_{\epsilon _n}:d(f_{\overline{\eta }})t\mathrm{\Delta }_n\}.$$
For any $`t>0`$ we have
$`(d(\widehat{f}_n^s)t\mathrm{\Delta }_n)`$ $``$ $`(\underset{\overline{\eta }𝒩_n^{}}{\mathrm{min}}[R_n(f_{\overline{\eta }})R_n(f_{\eta _n})]0)`$
$`=`$ $`(\underset{\overline{\eta }𝒩_n^{}}{\mathrm{min}}[Z_n(f_{\overline{\eta }})Z_n(f_{\eta _n})+d(f_{\overline{\eta }})d(f_{\eta _n})]0)`$
$``$ $`(\underset{\overline{\eta }𝒩_n^{}}{\mathrm{min}}[Z_n(f_{\overline{\eta }})Z_n(f_{\eta _n})+d(f_{\overline{\eta }})/2+t\mathrm{\Delta }_n/2d(f_{\eta _n})]0)`$
$``$ $`(\underset{\overline{\eta }𝒩_n^{}}{\mathrm{min}}[Z_n(f_{\overline{\eta }})+d(f_{\overline{\eta }})/2]0)`$
$`+(Z_n(f_{\eta _n})t\mathrm{\Delta }_n/2d(f_{\eta _n}))`$
$``$ $`(\underset{\overline{\eta }𝒩_n^{}}{\mathrm{min}}[Z_n(f_{\overline{\eta }})+d(f_{\overline{\eta }})/2]0)`$
$`+(Z_n(f_{\eta _n})t\mathrm{\Delta }_n/2\delta _n).`$
Since $`\mathrm{\Delta }_n`$ is of the same order as $`\delta _n`$, we can choose $`t`$ large enough to have $`t\mathrm{\Delta }_n/2\delta _nt\mathrm{\Delta }_n/4`$. Thus,
$`(d(\widehat{f}_n^s)t\mathrm{\Delta }_n)`$ $``$ $`\mathrm{card}𝒩_n^{}\underset{\overline{\eta }𝒩_n^{}}{\mathrm{max}}(Z_n(f_{\overline{\eta }})d(f_{\overline{\eta }})/2)`$
$`+(Z_n(f_{\eta _n})t\mathrm{\Delta }_n/4)`$
$``$ $`\mathrm{exp}(A^{}\epsilon _n^\rho )\underset{\overline{\eta }𝒩_n^{}}{\mathrm{max}}(Z_n(f_{\overline{\eta }})d(f_{\overline{\eta }})/2)`$
$`+(Z_n(f_{\eta _n})t\mathrm{\Delta }_n/4).`$
Note that for any decision rule $`f`$ the value $`Z_n(f)`$ is an average of $`n`$ i.i.d. bounded and centered random variables whose variance does not exceed $`P_X(f(X)f^{}(X,f))`$. Thus, using Bernstein’s inequality and (6.9) we obtain
$$(Z_n(f)a)\mathrm{exp}\left(\frac{Cna^2}{a+d(f)^{\alpha /(1+\alpha )}}\right),a>0.$$
Therefore, for $`\overline{\eta }𝒩_n^{}`$,
$`(Z_n(f_{\overline{\eta }})d(f_{\overline{\eta }})/2)`$ $``$ $`\mathrm{exp}(Cnd(f_{\overline{\eta }})^{(2+\alpha )/(1+\alpha )})`$
$``$ $`\mathrm{exp}(Cn(t\mathrm{\Delta }_n)^{(2+\alpha )/(1+\alpha )}).`$
Similarly, for $`t>C`$,
$`(Z_n(f_{\eta _n})t\mathrm{\Delta }_n/4)`$ $``$ $`\mathrm{exp}\left({\displaystyle \frac{Cn\mathrm{\Delta }_n^2}{\mathrm{\Delta }_n+d(f_{\eta _n})^{\alpha /(1+\alpha )}}}\right)`$
$``$ $`\mathrm{exp}\left({\displaystyle \frac{Cn\mathrm{\Delta }_n^2}{\mathrm{\Delta }_n+\delta _n^{\alpha /(1+\alpha )}}}\right)`$
$``$ $`\mathrm{exp}\left(Cn\mathrm{\Delta }_n^{(2+\alpha )/(1+\alpha )}\right).`$
The result of the theorem follows now from the above inequalities and the relation $`n\mathrm{\Delta }_n^{(2+\alpha )/(1+\alpha )}\epsilon _n^\rho `$.
<sup>1</sup>Center for Education and Research in Informatics
Ecole Nationale des Ponts et Chaussées
19, rue Alfred Nobel
Cité Descartes, Champs-sur-Marne
77455 Marne-La-Vall e, France
e-mail: audibert@certis.enpc.fr
<sup>2</sup>Laboratoire de Probabilités et Modèles
Aléatoires (UMR CNRS 7599),
Université Paris VI
4, pl.Jussieu, Boîte courrier 188,
75252 Paris, France
e-mail: tsybakov@ccr.jussieu.fr |
warning/0507/quant-ph0507035.html | ar5iv | text | # Bell states diagonal entanglement witnesses
## 1 Introduction
Entanglement is one of the most fascinating features of quantum mechanics. As Einstein, Podolsky and Rosen pointed out, the quantum states of two physically separated systems that interacted in the past can defy our intuitions about the outcome of local measurements. Moreover, it has recently been recognized that entanglement is a very important resource in quantum information processing. A bipartite mixed state is said to be separable (not entangled) if it can be written as a convex combination of pure product states.
A separability criterion is based on a simple property that can be shown to hold for every separable state. If some state does not satisfy this property, then it must be entangled. But the converse does not necessarily imply the state to be separable. One of the first and most widely used related criterion is the Positive Partial Transpose (PPT) criterion, introduced by Peres . Furthermore, the necessary and sufficient condition for separability in $`H_2H_2`$ and $`H_2H_3`$ was shown by Horodeckis , which was based on a previous work by Woronowicz . However, in higher dimensions, there are PPT states that are nonetheless entangled, as was first shown in , based on . These states are called bound entangled states because they have the peculiar property that no entanglement can be distilled from them by local operations .
Another approach to distinguish separable states from entangled states involves the so called entanglement witness (EW) . An EW for a given entangled state $`\rho `$ is an observable W whose expectation value is nonnegative on any separable state, but strictly negative on an entangled state $`\rho `$.
There is a correspondence relating entanglement witnesses to linear positive (but not completely positive) maps from the operators on $`H_A`$ to the operators on $`H_B`$ via Jamiolkowski isomorphism, or vice versa.
There has been much work on the separability problem, particularly from the Innsbruck-Hannover group, as reviewed in , that emphasizes convexity and proceeds by characterizing entanglement witnesses in terms of their extreme points, the so-called optimal entanglement witnesses, and PPT entangled states in terms of their extreme points, the edge PPT entangled states .
Having constructed the EW , one can decompose it into a sum of local measurements, then the expectation value can be measured with simple method. This decomposition has to be optimized in a certain way since we want to use the smallest number of measurements possible.
In this paper, we show that finding generic Bell states diagonal entanglement witnesses (BDEW) for $`d_1d_2\mathrm{}.d_n`$ systems reduces to a convex optimization problem. If the feasible region for this optimization problem constructs a polygon by itself, the corresponding boundary points of the convex hull will minimize exactly optimization problem. This problem is called linear programming , and the simplex method is the easiest way of solving it. If the feasible region is not a polygon, with the help of tangent planes in this region at points which are determined either analytically or numerically one can define new convex hull which is a polygon and has encircled the feasible region. The points on the boundary of the polygon can approximately determine the minimum value of optimization problem. Thus approximated value is obtained via linear programming. In general it is difficult to find this region and solve the optimization problem, thus it is difficult to find any generic multipartite EW. In the following sections we consider some simple but important examples which are solved exactly or approximately by linear programming method. Then we consider the multi-qubits and $`2N`$ with exactly minimum value by linear programming and $`33`$ systems with approximately minimum value by linear programming and then establish $`33`$ optimality condition together with non-decomposability properties for some particular choice of its parameters. Then we combine the optimal well known reduction map, and the optimal as well as the non-decomposable 3 $``$ 3 BDEW (i.e., the critical entanglement witnesses) to obtain further family of optimal and non-decomposable 3 $``$ 3 BDEW. Finally, using the critical entanglement witnesses some 3 $``$ 3 bound entangled states are detected and we consider the well known Choi map as a particular case of the positive map in connection with this witness via Jamiolkowski isomorphism which approximately is obtained via linear programming.
The paper is organized as follows:
In section 2 we give a brief review of entanglement witness. In section 3 we show that finding generic Bell states diagonal entanglement witnesses for $`d_1d_2\mathrm{}.d_n`$ systems reduces to a linear programming problem. In section 4, we consider BDEW for multi-qubit system. In section 5, we provide BDEW for $`2N`$. In section 6, we provide BDEW for $`33`$ systems. Section 7 is devoted to prove the n-d of critical EW and introduce a new family of optimal nd-EW via combining critical EW with the well known reduction maps. In section 8, using the critical EW, we will be able to detect a bound BD entangled state. In section 9, we consider the well known Choi map as a particular case of the positive map connect with this witness via Jamiolkowski isomorphism. Finally in section 10 using the optimal EW, we show that some separable Bell states diagonal lies at the boundary of separable region. The paper is ended with a brief conclusion together with three appendices devoted to the proof of A) the optimization of product distributions B)optimality of critical, reduction map C)simplex method for solving linear programming problem.
## 2 Entanglement witness
Here we mention briefly those concepts and definitions of EW that will be needed in the sequel, a more detailed treatment may be found for example in .
Let S be a convex compact set in a finite dimensional Banach space. Let $`\rho `$ be a point in the space with $`\rho \text{which is not in}S`$. Then there exists a hyperplane that separates $`\rho `$ from S.
A hermitian operator (an observable) W is called an entanglement witness (EW) iff
$$\rho \text{such that}Tr(\widehat{\rho }W)<0$$
(2-1)
$$\rho ^{}STr(\rho ^{}\widehat{W})0.$$
(2-2)
Definition 1: An EW is decomposable iff there exists operators P, Q such that
$$W=P+Q^{T_A}P,Q>0.$$
(2-3)
Decomposable EW can not detect PPT entangled states.
Definition 2: An EW is called non-decomposable entanglement witness (nd-EW) iff there exists at least one PPT entangled state which the witness detects.
Definition 3: The (decomposable) entanglement witness is tangent to S (P) iff there exists a $`\sigma S`$ ( $`\rho P`$) with $`Tr(W\sigma )=0(Tr(W\rho )=0)`$.
Using these definitions we can restate the consequences of the Hahn-Banach theorem in several ways:
Theorem:
1- $`\rho `$ is entangled iff there exists a witness W such that $`Tr(\rho W)<0`$.
2- $`\rho `$ is a PPT entangled state iff there exists an non-decomposable entanglement witness W such that $`Tr(\rho W)<0`$.
3- $`\sigma `$ is separable iff for all EW $`Tr(W\sigma )0`$.
From theoretical point of view this theorem is quite powerful. However, it is not useful to construct witnesses that detect a given state $`\rho `$.
We know that a strong relation was developed between entanglement witnesses and positive maps. Notice that an entanglement witness only gives one condition (namely $`Tr(W\rho )<0`$) while for the map $`(I_A\varphi )\rho `$ to be positive definite, there are many conditions that have to be satisfied. Thus the map is much stronger, while the witnesses are much weaker in detecting entanglement. It is shown that this concept is able to provide a more detailed classification of entangled states.
## 3 Bell states diagonal entanglement witnesses
As we know, one can expand any trace class observable in the Bell basis as
$$W=\underset{_{i_1i_2\mathrm{}i_n}}{}W_{_{i_1i_2\mathrm{}i_n}}|\psi _{_{i_1i_2\mathrm{}i_n}}\psi _{_{i_1i_2\mathrm{}i_n}}|$$
(3-4)
where $`|\psi _{_{i_1i_2\mathrm{}i_n}}`$($`0i_1d_1,0i_2d_2,\mathrm{},0i_nd_n,`$ and $`d_1d_2\mathrm{}d_n`$) stands for the orthonormal states for a $`d_1d_2\mathrm{}d_n`$ Bell state defined as
$$|\psi _{_{i_1i_2\mathrm{}i_n}}=(\mathrm{\Omega })^{i_1}(S)^{i_2}\mathrm{}(S)^{i_n}|\psi _{_{00\mathrm{}0}}$$
(3-5)
where $`\mathrm{\Omega }`$ and S are phase modules and shift operators for a $`d_1d_2\mathrm{}.d_n`$ defined as
$$\begin{array}{cc}\mathrm{\Omega }=\left(\begin{array}{ccccc}1& 0& 0& \mathrm{}& 0\\ 0& \omega & 0& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& \mathrm{}& \omega ^{d1}\end{array}\right),& S=\left(\begin{array}{ccccc}0& 1& 0& \mathrm{}& 0\\ 0& 0& 1& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 1& 0& 0& \mathrm{}& 0\end{array}\right),\end{array}$$
(3-6)
with $`\omega =exp(\frac{2\pi i}{d})`$ and
$$|\psi _{_{00\mathrm{}0}}=\frac{1}{\sqrt{d}}\underset{_{i=0}}{\overset{d_11}{}}|i_1|i_2\mathrm{}|i_n.$$
(3-7)
W is a trace one observable i.e., $`Tr(W)=1`$ and we have $`_{_{i_1i_2\mathrm{}i_n}}W_{_{i_1i_2\mathrm{}i_n}}=1`$.
Let us split the observable W into its positive and negative spectra as:
$$W=\underset{k=1}{\overset{n^+}{}}\lambda _k^+|\varphi _k^+\varphi _k^+|\underset{k=1}{\overset{n^{}}{}}\lambda _k^{}|\varphi _k^{}\varphi _k^{}|,$$
(3-8)
where $`\lambda _k^+(\lambda _k^{})`$ are the positive (negative) eigenvalues $`|\varphi _k^+(|\varphi _k^{})`$, and we have $`n^++n^{}=d^n`$. Denoting $`\lambda _k^{}=s>0`$ we can write (3-8) as:
$$W=(1+s)\rho ^+s\rho ^{},$$
(3-9)
where $`\rho ^\pm `$ are two normalized positive operators or density matrices defined as
$$\rho ^+=\frac{1}{1+s}\underset{k=1}{\overset{n^+}{}}(\lambda _k^+\varphi _k^+><\varphi _k^+),\rho ^{}=\frac{1}{s}\underset{k=1}{\overset{n^{}}{}}(\lambda _k^{}\varphi _k^{}><\varphi _k^{}).$$
(3-10)
Now, using the Lewenstein-Sanpera technique the identity operator $`\frac{I_{d_1d_2\mathrm{}d_n}}{d_1d_2\mathrm{}d_n}`$ can be written in terms of $`\rho ^{}`$ and the other positive states as
$$\frac{I_{d_1d_2\mathrm{}d_n}}{d_1d_2\mathrm{}d_n}=\lambda \rho ^{}+(1\lambda )\rho ^{^{}},\mathrm{\hspace{0.33em}\hspace{0.33em}0}<\lambda <1.$$
(3-11)
By using the above equation we can replace $`\rho ^{}`$ in Eq.(3-9) in terms of the identity operator. So, Eq.(3-9) is written as a sum of the identity and positive operators. Thus we have
$$W=𝐫\frac{I_{d_1d_2\mathrm{}d_n}}{d_1d_2\mathrm{}d_n}+(1𝐫)\rho ,$$
(3-12)
where
$$\rho =\frac{(1+s)\lambda }{\lambda +s}\rho ^++s(\frac{1\lambda }{s+\lambda })\rho ^{^{}},$$
(3-13)
and $`𝐫=\frac{s}{\lambda }<0`$.
In this paper we have considered only trace one observables which are diagonal in the Bell states. Hence we restrict ourselves to the Bell states diagonal $`\rho `$ defined as
$$\rho =\underset{_{i_1i_2\mathrm{}i_n}}{}q_{_{i_1i_2\mathrm{}i_n}}|\psi _{_{i_1i_2\mathrm{}i_n}}\psi _{_{i_1i_2\mathrm{}i_n}}|,q_{_{i_1i_2\mathrm{}i_n}}>0\text{and}\underset{_{i_1i_2\mathrm{}i_n}}{}q_{_{i_1i_2\mathrm{}i_n}}=1.$$
(3-14)
Finally, by substituting (3-14) in (3-12) the trace one Bell states diagonal W observables are
$$W=𝐫\frac{I_{d_1d_2\mathrm{}d_n}}{d_1d_2\mathrm{}d_n}+(1𝐫)\underset{_{i_1i_2\mathrm{}i_n}}{}q_{_{i_1i_2\mathrm{}i_n}}|\psi _{_{i_1i_2\mathrm{}i_n}}\psi _{_{i_1i_2\mathrm{}i_n}}|.$$
(3-15)
The observable given by (3-15) is not a positive operator and can not be an EW provided that its expectation value on any pure product state is positive. For a given product state $`|\gamma =|\alpha _1|\alpha _2\mathrm{}|\alpha _n`$ the non negativity of
$$Tr(W|\gamma \gamma |)0$$
(3-16)
implies that
$$\frac{d_1d_2\mathrm{}d_n_{_{i_1i_2\mathrm{}i_n}}q_{_{i_1i_2\mathrm{}i_n}}P_{_{i_1i_2\mathrm{}i_n}}}{1d_1d_2\mathrm{}d_n_{_{i_1i_2\mathrm{}i_n}}q_{_{i_1i_2\mathrm{}i_n}}P_{_{i_1i_2\mathrm{}i_n}}}𝐫0,$$
(3-17)
where $`P_{_{i_1i_2\mathrm{}i_n}}=<\gamma \psi _{_{i_1i_2\mathrm{}i_n}}>^2`$.
Denoting the summation in the numerator and the dominator in (3-17) by $`C(\gamma )=d_1d_2\mathrm{}d_n_{_{i_1i_2\mathrm{}i_n}}q_{_{i_1i_2\mathrm{}i_n}}P_{_{i_1i_2\mathrm{}i_n}}`$ we see that the least possible $`𝐫_0=\frac{C(\gamma )}{1C(\gamma )}`$ is the decreasing function of $`C(\gamma )`$ for $`C(\gamma )<1`$ (obviously for $`C(\gamma )>1`$ all r while being positive provide positive expectation value). Therefore, for given parameters $`q_{_{i_1i_2\mathrm{}i_n}}>0`$,with $`_{_{i_1i_2\mathrm{}i_n}}q_{_{i_1i_2\mathrm{}i_n}}=1`$, the least allowed value of the parameter $`𝐫`$, called the critical parameter (denoted by $`𝐫_c`$ ) is obtained from the product state $`\gamma `$ which minimizes $`C_\gamma =_{_{i_1i_2\mathrm{}i_n}}q_{_{i_1i_2\mathrm{}i_n}}P_{_{i_1i_2\mathrm{}i_n}}`$, with $`0P_{_{i_1i_2\mathrm{}i_n}}1`$ and the constraint $`_{_{i_1i_2\mathrm{}i_n}}P_{_{i_1i_2\mathrm{}i_n}}=1`$. As for the completeness of the Bell state $`_{_{i_1i_2\mathrm{}i_n}}|\psi _{_{i_1i_2\mathrm{}i_n}}\psi _{_{i_1i_2\mathrm{}i_n}}|=1`$, the determination of $`𝐫_c`$ reduces to the following optimization problem
$$\begin{array}{cc}\text{minimize}& C_\gamma =_{_{i_1i_2\mathrm{}i_n}}q_{_{i_1i_2\mathrm{}i_n}}P_{_{i_1i_2\mathrm{}i_n}}(\gamma )\\ & 0P_{_{i_1i_2\mathrm{}i_n}}(\gamma )\frac{1}{d_1}\\ & _{_{i_1i_2\mathrm{}i_n}}P_{_{i_1i_2\mathrm{}i_n}}(\gamma )=1.\end{array}$$
(3-18)
Always the distribution $`P_{_{i_1i_2\mathrm{}i_n}}`$ satisfies $`0P_{_{i_1i_2\mathrm{}i_n}}(\gamma )\frac{1}{d_1}`$ for all pure product states (the proof is given in the Appendix A). One can calculate the distributions $`P_{_{i_1i_2\mathrm{}i_n}}(\gamma )`$, consistent with the aforementioned optimization problem, from the information about the boundary of feasible region. To achieve the feasible region we obtain the extreme points corresponding to the product distributions $`P_{_{i_1i_2\mathrm{}i_n}}(\gamma )`$ for every given product states by applying the special conditions on $`q_{_{i_1i_2\mathrm{}i_n}}`$’s parameters. $`C_\gamma `$ themselves are functions of the product distributions, and they are in turn are functions of $`\gamma `$. They are not real variables of $`\gamma `$ but the product states will be multiplicative. If this feasible region constructs a polygon by itself, the corresponding boundary points of the convex hull will minimize exactly $`C_\gamma `$ in Eq. (3-18). This problem is called linear programming , and the simplex method is the easiest way of solving it. If the feasible region is not a polygon, with the help of tangent planes in this region at points which are determined either analytically or numerically one can define new convex hull which is a polygon and has encircled the feasible region. The points on the boundary of the polygon can approximately determine the minimum value $`C_\gamma `$ from Eq.(3-18), thus the problem is that of a linear programming again. In general it is difficult to find this region and solve the optimization problem, thus it is difficult to find any generic multipartite EW. In the following sections we consider some simple but important examples which are solved as linear programming problem.
## 4 Bell states diagonal entanglement witnesses for multi-qubit system
Here we provide a multi-qubit entanglement witness. From the previous section one can show that the Bell states diagonal observable W for multi qubit system is defined by
$$W=𝐫\frac{I_{2^n}}{2^n}+(1𝐫)\underset{i_1,\mathrm{},i_n=0}{\overset{1}{}}q_{i_1,i_2,\mathrm{},i_n}|\psi _{i_1,i_2,\mathrm{},i_n}\psi _{i_1,i_2,\mathrm{},i_n}|,$$
(4-19)
where $`|\psi _{i_1,i_2,\mathrm{},i_n}`$ is a Bell state:
$$|\psi _{i_1,i_2,\mathrm{},i_n}=(\sigma _z)^{i_1}(\sigma _x)^{i_2}\mathrm{}(\sigma _x)^{i_2}|\psi _{{}_{0}{}^{},_{0}^{},\mathrm{},_0},$$
(4-20)
with
$$|\psi _{{}_{0}{}^{},_{0}^{},\mathrm{},_0}=\frac{1}{\sqrt{2}}\underset{i=0}{\overset{1}{}}|i_1|i_2\mathrm{}|i_n,$$
(4-21)
and $`\sigma _z`$ and $`\sigma _x`$ are the Pauli operators. This observable is not a positive operator and can not be an EW provided that its expectation value on any product state $`|\gamma =|\alpha _1|\alpha _2\mathrm{}|\alpha _n`$ is positive.
We consider an easy case $`q_{_{00\mathrm{}00}}=0,q_{_{10\mathrm{}00}}=x`$ with all the other $`q`$’s being equal, i.e., $`q_{_{i_1,i_2,\mathrm{},i_n}}=\frac{1x}{2(2^{n1}1)}`$ except for $`i_1=i_2=\mathrm{}=i_n=0`$ and $`i_2=i_3=\mathrm{}=i_n=0,i_1=1`$. Then the observable W reduces to the following form
$$W=𝐫\frac{I_{2^n}}{2^n}+\frac{(1𝐫)}{2(2^{n1}1)}((1x)I_{2^n}(1x)|\psi _{{}_{0}{}^{},_{0}^{},\mathrm{},_0}\psi _{{}_{0}{}^{},_{0}^{},\mathrm{},_0}|+((2^n1)x1)|\psi _{{}_{1}{}^{},_{0}^{},\mathrm{},_0}\psi _{{}_{1}{}^{},_{0}^{},\mathrm{},_0}|).$$
(4-22)
We can calculate $`C_\gamma `$ from the non negativity of $`Tr(W|\gamma \gamma |)`$ for a given product state $`|\gamma `$
$$C_\gamma =\frac{1}{2(2^{n1}1)}((1x)(1x)P_{_{00\mathrm{}00}}+((2^n1)x1)P_{_{10\mathrm{}00}}).$$
(4-23)
According to the definition of product distributions, we have
$$\begin{array}{c}P_{_{00\mathrm{}0}}=\frac{1}{2}\alpha _1\alpha _2\mathrm{}\alpha _n+\beta _1\beta _2\mathrm{}\beta _n^2\\ P_{_{10\mathrm{}0}}=\frac{1}{2}\alpha _1\alpha _2\mathrm{}\alpha _n\beta _1\beta _2\mathrm{}\beta _n^2,\end{array}$$
(4-24)
where
$$|\alpha _i=\left(\begin{array}{c}\alpha _i\\ \beta _i\end{array}\right),i=1,2,\mathrm{},n.$$
(4-25)
For the most general given product states $`\gamma `$, we determine the extreme allowed values of these product distributions. Let
$$\begin{array}{c}|\alpha _1=|\alpha _2=\mathrm{}=|\alpha _n=\left(\begin{array}{c}1\\ 0\end{array}\right),\end{array}$$
(4-26)
then the product distributions are
$$\{\begin{array}{c}P_{_{00\mathrm{}00}}=\frac{1}{2}\\ P_{_{10\mathrm{}00}}=\frac{1}{2}\end{array}.$$
(4-27)
Another choice will make $`P_{_{10\mathrm{}0}}`$ minimal
$$|\alpha _1=|\alpha _2=\mathrm{}=|\alpha _n=\frac{1}{\sqrt{2}}\left(\begin{array}{c}1\\ 1\end{array}\right),$$
(4-28)
and the product distributions will become
$$\{\begin{array}{c}P_{_{10\mathrm{}0}}=\frac{1}{2^{n1}}\\ P_{_{10\mathrm{}00}}=0\end{array}.$$
(4-29)
Similarly we can also make $`P_{_{00\mathrm{}0}}`$ minimal
$$\begin{array}{cc}|\alpha _2=|\alpha _3=\mathrm{}=|\alpha _n=\frac{1}{\sqrt{2}}\left(\begin{array}{c}1\\ 1\end{array}\right),|\alpha _1=\frac{1}{\sqrt{2}}\left(\begin{array}{c}1\\ 1\end{array}\right)& \{\begin{array}{c}P_{_{00\mathrm{}00}}=0\\ P_{_{10\mathrm{}00}}=\frac{1}{2^{n1}}\end{array}\end{array}.$$
(4-30)
From the definition of the convex function we can show that the convex combination of these distributions provide a convex region called the feasible region, where all points in the interior of this region satisfy the positivity constraint of $`Tr(W|\gamma \gamma |)`$. Then we have
$$\{\begin{array}{cc}\text{maximize}& C_\gamma =\frac{1}{2(2^{n1}1)}((1x)+(1x)P_{_{00\mathrm{}00}}((2^n1)x1)P_{_{10\mathrm{}00}})\\ \text{subject to}& 2P_{_{00\mathrm{}00}}2P_{_{10\mathrm{}00}}(1\frac{1}{2^{n1}})\frac{1}{2^{n1}}\\ & 2P_{_{10\mathrm{}00}}2P_{_{00\mathrm{}00}}(1\frac{1}{2^{n1}})\frac{1}{2^{n1}}\\ & P_{_{00\mathrm{}00}}0,P_{_{10\mathrm{}00}}0.\end{array}$$
(4-31)
Now we must prove that the feasible region, constructed from the convex of these points, is a polygon. Let $`P_+=P_{_{00\mathrm{}0}}`$ and $`P_{}=P_{_{10\mathrm{}0}}`$. The equation of the line passing through $`(P_+=\frac{1}{2^{n1}},P_{}=0)`$ and $`(P_+=\frac{1}{2},P_{}=\frac{1}{2})`$ is
$$P_{}=(\frac{2^{n1}}{2^{n1}2})P_+\frac{1}{2^{n1}2}.$$
(4-32)
Let us further assume $`P_{}=\lambda P_+`$. By intersecting this equation with the one above we get
$$P_+=\frac{1}{2^{n1}\lambda (2^{n1}2)}.$$
(4-33)
Now if we assume $`\lambda =0`$ we arrive at the point $`(P_+=\frac{1}{2^{n1}},P_{}=0)`$, for $`\lambda =1`$ we conclude $`(P_+=\frac{1}{2},P_{}=\frac{1}{2})`$. One can rewrite Eq.(4-24) as
$$P_\pm =\frac{1}{2}(\alpha _1^2\alpha _2^2\mathrm{}\alpha _n^2+(1\alpha _1^2)(1\alpha _2^2)\mathrm{}(1\alpha _n^2)$$
$$\pm 2\alpha _1(1\alpha _1^2)\alpha _2(1\alpha _2^2)\mathrm{}\alpha _n(1\alpha _n^2)\mathrm{cos}(\varphi )).$$
(4-34)
Thus we write the Lagrangian as
$$=P_++\mu (P_{}\lambda P_+),$$
(4-35)
where $`\mu `$ is the Lagrange multiplier. With $`\alpha _i=cos\theta _i`$ we maximize $``$ with respect to $`\theta _i`$’s and $`\varphi `$
$$\{\begin{array}{c}\theta _1=\theta _2=\mathrm{}=\theta _n\alpha _1=\alpha _2=\mathrm{}=\alpha _n\\ \varphi =0\end{array},$$
(4-36)
such that
$$\mathrm{tan}^n\theta _i=\frac{1\sqrt{\lambda }}{1+\sqrt{\lambda }},$$
(4-37)
so that
$$P_+=(\frac{2}{1+\sqrt{\lambda }})\frac{1}{(1+(\frac{1\sqrt{\lambda }}{1+\sqrt{\lambda }})^{\frac{2}{n}})^n}.$$
(4-38)
As we see the equation for $`P_+`$ is less than the one in (4-33), moreover, this relation indicates the correctness of the result (4-26)-(4-30). Thus the convex hull is a polygon and the optimization problem will be converted into the linear programming one.
There is no simple analytical formula for solving a linear programming, however there are a variety of very effective methods, including the simplex method to solve them. So, minimization solutions of $`C_\gamma `$ is obtained by the simplex method and we have (see Appendix C):
I) For $`0x\frac{1}{2^{n1}+1}`$ the extreme points of the feasible region are $`P_{_{00\mathrm{}00}}=P_{_{10\mathrm{}00}}=\frac{1}{2}`$ and the minimum value of $`C_\gamma `$ is defined by $`(C_\gamma )_{min}=\frac{x}{2}`$. By substituting these values in (3-17) we have
$$\frac{2^{n1}x}{12^{n1}x}𝐫0𝐫_c=\frac{2^{n1}x}{12^{n1}x},$$
(4-39)
where $`𝐫_c`$ is called the critical r. By substituting $`𝐫_c`$ in (4-19) this observable has positive expectation value under any product state, thus it will be an EW called critical EW equal to
$$W_c(x)=\frac{1}{2(2^{n1}1)}(I_{2^n}\frac{1x}{12^{n1}x}|\psi _{{}_{0}{}^{},_{0}^{},\mathrm{},_0}\psi _{{}_{0}{}^{},_{0}^{},\mathrm{},_0}|+\frac{(2^n1)x1}{12^{n1}x}|\psi _{{}_{1}{}^{},_{0}^{},\mathrm{},_0}\psi _{{}_{1}{}^{},_{0}^{},\mathrm{},_0}|)),$$
(4-40)
which in the special case where $`x=\frac{1}{2^n1}`$ the $`W_c(x)`$ reduces to
$$W_{red}=\frac{1}{2(2^{n1}1)}(I_{2^n}2|\psi _{{}_{0}{}^{},_{0}^{},\mathrm{},_0}\psi _{{}_{0}{}^{},_{0}^{},\mathrm{},_0}|),$$
(4-41)
which is the well known reduction map.
II) For $`\frac{1}{2^{n1}+1}x1`$ the extreme points of the feasible region are $`P_{_{00\mathrm{}00}}=\frac{1}{2^{n1}}`$ and $`P_{_{10\mathrm{}00}}=0`$ respectively. Therefore, from the simplex method we get $`(C_\gamma )_{min}=\frac{1x}{2^n}`$, hence $`𝐫_c=\frac{1x}{x}`$ and the critical EW is calculated to be
$$W_c(x)=\frac{1}{(2^{n1}1)}(\frac{1x}{x}\frac{I_{2^n}}{2^n}\frac{1x}{2x}|\psi _{{}_{0}{}^{},_{0}^{},\mathrm{},_0}\psi _{{}_{0}{}^{},_{0}^{},\mathrm{},_0}|+\frac{(2^n1)x1}{2x}|\psi _{{}_{1}{}^{},_{0}^{},\mathrm{},_0}\psi _{{}_{1}{}^{},_{0}^{},\mathrm{},_0}|)).$$
(4-42)
Note that this choice of q is not the only way of defining a BDEW for multi-qubit system in the one parameter representation. Let us consider the alternative definition for the one parameter BDEW by studying the following example. Assume $`q_{_{00\mathrm{}01}}=x`$ and set all the other $`q`$’s to be equal. Thus we have
$$W=𝐫\frac{I_{2^n}}{2^n}+\frac{(1𝐫)}{2(2^{n1}1)}((1x)I_{2^n}(1x)|\psi _{{}_{0}{}^{},_{0}^{},\mathrm{},_0}\psi _{{}_{0}{}^{},_{0}^{},\mathrm{},_0}|+((2^n1)x1)|\psi _{{}_{0}{}^{},_{0}^{},\mathrm{},_1}\psi _{{}_{0}{}^{},_{0}^{},\mathrm{},_1}|).$$
(4-43)
Similarly we can find the extreme points of $`P_{_{00\mathrm{}00}}`$ and $`P_{_{00\mathrm{}01}}`$ as
$$\begin{array}{cc}|\alpha _1=|\alpha _2=\mathrm{}=|\alpha _n=\left(\begin{array}{c}1\\ 0\end{array}\right)& \{\begin{array}{c}P_{_{00\mathrm{}00}}=\frac{1}{2}\\ P_{_{00\mathrm{}01}}=0\end{array}\end{array}.$$
(4-44)
$$\begin{array}{cc}|\alpha _1=|\alpha _2=\mathrm{}=|\alpha _{n1}=\left(\begin{array}{c}1\\ 0\end{array}\right)\text{and}|\alpha _n=\left(\begin{array}{c}0\\ 1\end{array}\right)& \{\begin{array}{c}P_{_{00\mathrm{}00}}=0\\ P_{_{00\mathrm{}01}}=\frac{1}{2}\end{array}\end{array}.$$
(4-45)
Also we know that the convex combination of $`P_{_{00\mathrm{}00}}`$ and $`P_{_{00\mathrm{}01}}`$ provides a convex region or a feasible region. Then we have an optimization problem as follows:
$$\{\begin{array}{cc}\text{minimize}& C_\gamma =\frac{1}{2(2^{n1}1)}((1x)(1x)P_{_{00\mathrm{}00}}+((2^n1)x1)P_{_{00\mathrm{}01}})\\ \text{subject to}& \frac{1}{2}P_{_{00\mathrm{}00}}P_{_{00\mathrm{}01}}0\\ & P_{_{00\mathrm{}00}},P_{_{00\mathrm{}01}}0\end{array}.$$
(4-46)
Here the optimization is converted to the linear programming problem. To prove, we must show that the feasible region is a polygon. Let us suppose $`P_+=P_{_{00\mathrm{}00}}`$ and $`P_{}=P_{_{00\mathrm{}01}}`$ with
$$\{\begin{array}{c}P_+=\frac{1}{2}\alpha _1\alpha _2\mathrm{}\alpha _n+\beta _1\beta _2\mathrm{}\beta _n^2\\ P_{}=\frac{1}{2}\beta _1\alpha _2\mathrm{}\alpha _n+\alpha _1\beta _2\mathrm{}\beta _n^2,\end{array}$$
(4-47)
where as before we have $`P_{}=\lambda P_+`$ and $`\alpha _i=\mathrm{cos}\theta _i`$. By maximizing the Lagrangian we get $`\theta _2=\theta _3=\mathrm{}=\theta _n`$. Thus we have
$$P_++P_{}=\frac{1}{2}(\mathrm{cos}^{2n2}\theta _2+\mathrm{sin}^{2n2}\theta _2)\frac{1}{2}.$$
(4-48)
The line passing through the points $`(P_+=\frac{1}{2},P_{}=0)`$ and $`(P_+=0,P_{}=\frac{1}{2})`$ is $`P_{}=P_++\frac{1}{2}`$, which is always located above the curve obtained in the Eq.(4-48). Therefore, all the points are within the feasible region and this region constructs a polygon. Thus the above optimization problem reduces to a linear programming one. This minimization is exactly solved in the same way as mentioned above, and the critical EW is obtained as
$$𝐫_c=\frac{2^n(1x)}{2(2^{n1}(x+1)2)}$$
(4-49)
$$W_c=\frac{1}{(2^{n1}(x+1)2)}((1x)I_{2^n}2(1x)|\psi _{_{\mathrm{00..00}}}\psi _{_{\mathrm{00..00}}}|+2((2^n1)x1)|\psi _{_{\mathrm{00..01}}}\psi _{_{\mathrm{00..01}}}|).$$
## 5 Bell states diagonal entanglement witnesses for $`2N`$ system
Here, we will find a $`2N`$ entanglement witness. From the previous discussions we can define the Bell states diagonal observable W as
$$W=𝐫\frac{I_{2N}}{2N}+(1𝐫)\underset{i=0}{\overset{N1}{}}\underset{\alpha =0}{\overset{1}{}}q_{_{i\alpha }}|\psi _{_{i\alpha }}\psi _{_{i\alpha }}|,$$
(5-50)
where $`|\psi _{_{i\alpha }}=I_2(S)^i(\mathrm{\Omega })^\alpha |\psi _{_{00}}`$, with $`|\psi _{_{00}}=\frac{1}{\sqrt{2}}_{k=0}^1|k|k`$ and
$$\begin{array}{cc}\omega =\left(\begin{array}{ccccc}1& 0& 0& \mathrm{}& 0\\ 0& 1& 0& \mathrm{}& 0\\ 0& 0& 1& \mathrm{}& 0\\ \mathrm{}& & \mathrm{}& & \mathrm{}\\ 0& 0& 0& \mathrm{}& 1\end{array}\right)& S=\left(\begin{array}{ccccc}0& 1& 0& \mathrm{}& 0\\ 1& 0& 0& \mathrm{}& 0\\ 0& 0& 1& \mathrm{}& 0\\ \mathrm{}& & \mathrm{}& & \mathrm{}\\ 0& 0& 0& \mathrm{}& 1\end{array}\right).\end{array}$$
(5-51)
Similar to multi-qubit let $`q_{_{00}}=0`$ and $`q_{_{10}}=x`$ and let all the other q’s be equal to $`\frac{1x}{2N2}`$. Then by obtaining the expectation value of W on the product states and finding the product distributions we have
$$\begin{array}{cc}|\alpha _1=\left(\begin{array}{c}1\\ 0\end{array}\right),|\alpha _2=\left(\begin{array}{c}1\\ 0\\ \mathrm{}\\ 0\end{array}\right)& \{\begin{array}{c}P_{_{00}}=\frac{1}{2}\\ P_{_{10}}=\frac{1}{2}\end{array}\end{array}$$
(5-52)
$$\begin{array}{cc}|\alpha _1=\frac{1}{\sqrt{2}}\left(\begin{array}{c}1\\ 1\end{array}\right),|\alpha _2=\frac{1}{\sqrt{2}}\left(\begin{array}{c}1\\ 1\\ 0\\ \mathrm{}\\ 0\end{array}\right)& \{\begin{array}{c}P_{_{00}}=\frac{1}{2}\\ P_{_{10}}=0\end{array}\end{array}$$
(5-53)
$$\begin{array}{cc}|\alpha _1=\frac{1}{\sqrt{2}}\left(\begin{array}{c}1\\ 1\end{array}\right),|\alpha _2=\frac{1}{\sqrt{2}}\left(\begin{array}{c}1\\ 1\\ 0\\ \mathrm{}\\ 0\end{array}\right)& \{\begin{array}{c}P_{_{00}}=0\\ P_{_{10}}=\frac{1}{2}\end{array}\end{array}.$$
(5-54)
The feasible region is a rectangular and the optimization problem reduces to linear programming. Therefore, by using simplex method for $`0x\frac{1}{N+1}`$ we find the minimum value of $`C_\gamma =\frac{x}{2}`$ and critical $`𝐫`$ as $`𝐫_c=\frac{Nx}{1Nx}`$, and the critical EW is defined as
$$W_c=\frac{1}{2(N1)}(I_{2N}\frac{1x}{1Nx}|\psi _{_{00}}\psi _{_{00}}|+\frac{(2N1)x1}{1Nx}|\psi _{_{10}}\psi _{_{10}}|).$$
(5-55)
For the critical $`𝐫`$ we find $`𝐫_c=\frac{1x}{x}`$ in the region $`\frac{1}{N+1}x1`$ and the critical EW has the following form
$$W_c(x)=\frac{1}{2}(\frac{1x}{x}\frac{I_{2N}}{2N}\frac{1x}{2x}|\psi _{_{00}}\psi _{_{00}}|+\frac{(2N1)x1}{2x}|\psi _{_{10}}\psi _{_{10}}|).$$
(5-56)
In another one parameter EW example we assume that $`q_{_{01}}=x`$ and set all the other $`q`$’s to be equal so that we have
$$W=𝐫\frac{I_{2N}}{2N}+\frac{(1𝐫)}{2(N1)}((1x)I_{2N}(1x)|\psi _{_{00}}\psi _{_{00}}|+((2N1)x1)|\psi _{_{01}}\psi _{_{01}}|).$$
(5-57)
Similarly we can find the extreme points of $`P_{_{00}}`$ and $`P_{_{01}}`$ as
$$\begin{array}{cc}|\alpha _1=\left(\begin{array}{c}1\\ 0\end{array}\right),|\alpha _2=\left(\begin{array}{c}1\\ 0\\ \mathrm{}\\ 0\end{array}\right)& \{\begin{array}{c}P_{_{00}}=\frac{1}{2}\\ P_{_{01}}=0\end{array}\end{array},$$
(5-58)
$$\begin{array}{cc}|\alpha _1=\left(\begin{array}{c}1\\ 0\end{array}\right),|\alpha _2=\left(\begin{array}{c}0\\ 1\\ 0\\ \mathrm{}\\ 0\end{array}\right)& \{\begin{array}{c}P_{_{00}}=0\\ P_{_{01}}=\frac{1}{2}\end{array}\end{array}.$$
(5-59)
Then the critical EW is defined as
$$𝐫_c=\frac{2N(1x)}{2(N(x+1)2)}$$
(5-60)
$$W_c=\frac{1}{(N(x+1)2)}((1x)I_{2N}(1x)|\psi _{_{00}}\psi _{_{00}}|+((2N1)x1)|\psi _{_{01}}\psi _{_{01}}|).$$
## 6 Bell states diagonal entanglement witnesses for $`33`$ system
Here we provide a $`33`$ Bell diagonal entanglement witness. One can show that the Eq. (3-15) for a $`33`$ system reads as
$$W=𝐫\frac{I_9}{9}+(1𝐫)\underset{i_1,i_2=0}{\overset{2}{}}q_{i_1,i_2}|\psi _{i_1,i_2}\psi _{i_1,i_2}|.$$
(6-61)
It is difficult to prove whether or not the EW for a $`33`$ system is optimal. Also it is difficult to see for which value of the allowed $`𝐫`$, EW are (or are not) decomposable. Therefore, to investigate the optimality and non decomposability of these EW we restrict ourselves below to some particular choice of $`q_{ij}`$:
Because the distributions $`0P_{ij}\frac{1}{3}`$ and the minimum value of $`C_\gamma `$ are dependent on the coefficients $`q_{ij}`$, we consider a special case for the coefficients $`q_{ij}`$ defined by
$$q_{01}=q_{02}=q_{11}=q_{22}=q_{12}=q_{21}=\frac{1}{8},q_{10}=x\text{and}q_{20}=\frac{1}{4}x,\mathrm{\hspace{0.33em}\hspace{0.33em}0}x\frac{1}{4}.$$
(6-62)
By substituting these values in (3-15) we get
$$W(x)=𝐫\frac{I_9}{9}+(1𝐫)(\frac{I_9}{8}\frac{1}{8}|\psi _{_{00}}\psi _{_{00}}|\frac{8x1}{8}(|\psi _{_{10}}\psi _{_{10}}||\psi _{_{20}}\psi _{_{20}}|)).$$
(6-63)
By using (3-16) for non-negativity of the observable W we find the distributions $`P_{ij}`$ as a function of x. The minimum value of $`C_\gamma `$ is obtained from the boundary of the feasible region, i.e., we have
$$C_\gamma =\frac{1}{8}(1P_{_{00}}(8x1)(P_{_{10}}P_{_{20}})).$$
(6-64)
For given product states $`|\gamma =|\alpha _1|\alpha _2`$ one can obtain the extreme points of the product distributions as
$$\{\begin{array}{ccc}|\alpha _1=|\alpha _2=\frac{1}{\sqrt{3}}\left(\begin{array}{c}1\\ 1\\ 1\end{array}\right)& & (P_{00}=\frac{1}{3},P_{10}=0,P_{20}=0)\\ |\alpha _1=|\alpha _2=\frac{1}{\sqrt{3}}\left(\begin{array}{c}1\\ \omega \\ \overline{\omega }\end{array}\right)& & (P_{00}=0,P_{10}=\frac{1}{3},P_{20}=0)\\ |\alpha _1=|\alpha _2=\frac{1}{\sqrt{3}}\left(\begin{array}{c}1\\ \overline{\omega }\\ \omega \end{array}\right)& & (P_{00}=0,P_{10}=0,P_{20}=\frac{1}{3})\end{array}$$
$$\{\begin{array}{ccc}|\alpha _1=\frac{1}{\sqrt{2}}\left(\begin{array}{c}0\\ 1\\ 1\end{array}\right),|\alpha _2=\frac{1}{\sqrt{2}}\left(\begin{array}{c}0\\ 1\\ 1\end{array}\right)& & (P_{00}=0,P_{10}=\frac{1}{4},P_{20}=\frac{1}{4})\\ |\alpha _1=\frac{1}{\sqrt{2}}\left(\begin{array}{c}1\\ 0\\ \overline{\omega }\end{array}\right),|\alpha _2=\frac{1}{\sqrt{2}}\left(\begin{array}{c}1\\ 0\\ \overline{\omega }\end{array}\right)& & (P_{00}=\frac{1}{4},P_{10}=0,P_{20}=\frac{1}{4})\\ |\alpha _1=\frac{1}{\sqrt{2}}\left(\begin{array}{c}1\\ \omega \\ 0\end{array}\right),|\alpha _2=\frac{1}{\sqrt{2}}\left(\begin{array}{c}1\\ \omega \\ 0\end{array}\right)& & (P_{00}=\frac{1}{4},P_{10}=\frac{1}{4},P_{20}=0)\end{array}$$
(6-65)
and
$$\begin{array}{ccc}|\alpha _1=|\alpha _2=\left(\begin{array}{c}1\\ 0\\ 0\end{array}\right)& & (P_{00}=\frac{1}{3},P_{10}=\frac{1}{3},P_{20}=\frac{1}{3})\end{array}.$$
(6-66)
By convex combination of these points we obtain the possible region. Thus we have an optimization problem as
$$\{\begin{array}{cc}\text{minimize}& C_\gamma =\frac{1}{8}(1P_{_{00}}(8x1)(P_{_{10}}P_{_{20}}))\\ \text{subject to}& 13P_{_{00}}P_{_{10}}+P_{_{20}}0\\ & 1+P_{_{00}}3P_{_{10}}P_{_{20}}0\\ & 1P_{_{00}}+P_{_{10}}3P_{_{20}}0\\ & P_{_{00}},P_{_{10}},P_{_{20}}0.\end{array}$$
(6-67)
One can prove analytically that the region can be encircled with a polygon and the optimization problem is reduced to a linear programming. To prove, we begin from the definition of the product distributions
$$\begin{array}{c}P_{00}=\frac{1}{3}\alpha _1\beta _1+\alpha _2\beta _2+\alpha _3\beta _3^2\\ P_{10}=\frac{1}{3}\alpha _1\beta _1+\alpha _2\beta _2\omega +\alpha _3\beta _3\overline{\omega }^2\\ P_{20}=\frac{1}{3}\alpha _1\beta _1+\alpha _2\beta _2\overline{\omega }+\alpha _3\beta _3\omega ^2.\end{array}$$
(6-68)
Without loss of generality, one can assume that
$$\alpha _1\beta _1=x_1,\alpha _2\beta _2=x_2,\alpha _3\beta _3=x_3.$$
(6-69)
The Schwartz inequality yields
$$x_1+x_2+x_31.$$
(6-70)
Since we are looking the extreme points we will choose the maximum value in the inequality (6-70), that is $`x_1+x_2+x_3=1.`$ Thus the product distributions are written as
$$\begin{array}{c}P_{00}=\frac{1}{3}x_1+x_2e^{i\varphi _2}+x_3e^{i\varphi _3}^2\\ P_{10}=\frac{1}{3}x_1+x_2e^{i\varphi _2}\omega +x_3e^{i\varphi _3}\overline{\omega }^2\\ P_{20}=\frac{1}{3}x_1+x_2e^{i\varphi _2}\overline{\omega }+x_3e^{i\varphi _3}\omega ^2.\end{array}$$
(6-71)
Now, supposing that $`P_{00}`$ and $`P_{10}`$ are fixed values, we conclude
$$\varphi _2=\varphi _3,x_2=x_3.$$
(6-72)
Thus
$$\begin{array}{c}P_{00}=\frac{1}{3}1+2x_2(1\mathrm{cos}\varphi _2)^2\\ P_{10}=\frac{1}{3}1+2x_2(1\mathrm{cos}\varphi _2+\frac{2\pi }{3})^2\\ P_{20}=\frac{1}{3}1+2x_2(1\mathrm{cos}\varphi _2\frac{2\pi }{3})^2.\end{array}$$
(6-73)
We write down the equation for the planes passing through the obtained extreme points for $`P_{ij}`$ and maximize it with respect to the variables $`\varphi _2`$ and $`x_2`$. Where the obtained values for $`\varphi _2`$ and $`x_2`$ are indicative of the violation from the equations of planes, that is there are points which are located out of these planes, or in other word the planes have become convex. Thus, the equations of the planes (Fig 1) and their maximum violation(D) are obtained as follows
$$\{\begin{array}{ccc}1)3P_{10}P_{20}+P_{00}1=0& x_2=\frac{7}{61},\mathrm{cos}\varphi _2=\frac{1}{7}& D=\frac{2}{61}\\ 2)3P_{10}+P_{20}P_{00}1=0& x_2=\frac{7}{61},\mathrm{cos}\varphi _2=\frac{11}{14}& D=\frac{2}{61}\\ 3)3P_{20}+P_{10}P_{00}1=0& x_2=\frac{7}{61},\mathrm{cos}\varphi _2=\frac{11}{14}& D=\frac{2}{61}\\ 4)3P_{20}P_{10}+P_{00}1=0& x_2=\frac{7}{61},\mathrm{cos}\varphi _2=\frac{1}{7}& D=\frac{2}{61}\\ 5)3P_{00}+P_{20}P_{10}1=0& x_2=\frac{7}{61},\mathrm{cos}\varphi _2=\frac{13}{14}& D=\frac{2}{61}\\ 6)3P_{00}P_{20}+P_{10}1=0& x_2=\frac{7}{61},\mathrm{cos}\varphi _2=\frac{13}{14}& D=\frac{2}{61}.\end{array}$$
(6-74)
It is seen that the points thus obtained are located out of the considered plane. Thus, the equations of the planes passing through the new extreme points which are parallel to the above plane are obtained. For example, the equation of plane parallel to $`3p_{00}+P_{10}+p20=1`$ is $`3p_{00}+P_{10}+p20=1+\frac{2}{61}`$ which under permutation $`(P_{00},P_{10},P_{20})`$ will act similarly. Tack arbitrary any three of planes passing through the new extreme, $`(P_{00}=\frac{1}{3},P_{10}=\frac{1}{3},P_{20}=\frac{1}{3})`$, $`(P_{10}=0,P_{20}=0,P_{00}=0)`$ and $`P_{10}+P_{20}+P_{00}1=0`$ and intersecting with each other. Hence new extreme points will be produced. Thus we have encircled a polygon by its feasible region and the optimization problem will be reduced to that of a linear programming. The vertices of this polygon are the solutions of the problem provided that they obey $`(P_{00}\frac{1}{3},P_{10}\frac{1}{3},P_{20}\frac{1}{3})`$ and by substituting them into equation $`C_\gamma `$ one can determine its minimum value. Thus, for $`\frac{67}{756}x\frac{61}{378}`$, the extreme point is defined by $`(\frac{1}{3},\frac{1}{3},\frac{1}{3})`$ and finally $`(C_\gamma )_{min}=(\frac{1}{12})`$. Having found the critical $`𝐫`$ we substitute it in (3-15) and obtain a family of EW (called critical EW). Thus we have
$$p_c=3,W_c(x)=\frac{1}{2}(\frac{1}{3}I_9|\psi _{_{00}}\psi _{_{00}}|+(8x1)(|\psi _{_{10}}\psi _{_{10}}||\psi _{_{20}}\psi _{_{20}}|)),$$
(6-75)
where $`W_c(x)`$ reduces to the following well known reduced EW at $`x=\frac{1}{8}`$:
$$W_{red}=\frac{I_93|\psi _{00}\psi _{00}|}{6},$$
(6-76)
In the Appendix B it is shown that the above EW is optimal in contrast to the conclusion that it is a decomposable EW Ref..
In the Appendix B, we discuss the possible choice of x consistent with $`C_{mn}=\frac{1}{12}`$ and the optimality of the corresponding $`W_c(x)`$.
Also we prove in the following section that $`W_c`$ is nd-EW for all values of $`\frac{67}{756}x\frac{61}{378}`$, except for $`x=\frac{1}{8}`$. Besides taking a convex combination of $`W_c`$ and $`W_{red}`$ ,i.e.,
$$W_\mathrm{\Lambda }=\mathrm{\Lambda }W_c+(1\mathrm{\Lambda })W_{red},$$
(6-77)
we obtain a new EW which is optimal (see Appendix B) and is also an nd-EW for certain value of the parameter $`\mathrm{\Lambda }`$ as will be shown in section 7.
However we can consider other values for $`q_{ij}`$ in (3-15), e.g., $`q_{20}=q_{02}=q_{11}=q_{22}=q_{12}=q_{21}=\frac{1}{8},q_{10}=x`$ and $`q_{01}=\frac{1}{4}x,\mathrm{\hspace{0.33em}0}x\frac{1}{4}`$ then define the observable W by substituting the above condition in (3-15) as follows
$$W(x)=𝐫\frac{I_9}{9}+(1𝐫)(\frac{I_9}{8}\frac{1}{8}|\psi _{_{00}}\psi _{_{00}}|\frac{8x1}{8}(|\psi _{_{10}}\psi _{_{10}}||\psi _{_{01}}\psi _{_{01}}|)).$$
(6-78)
By using (3-16) for non-negativity of the observable W we find the distributions $`P_{ij}`$ as a function of x. The minimum value of $`C_\gamma `$ is obtained from the boundary of the feasible region, i.e., we have
$$C_\gamma =\frac{1}{8}(1P_{_{00}}(8x1)(P_{_{10}}P_{_{01}})).$$
(6-79)
For given product states $`|\gamma =|\alpha _1|\alpha _2`$ one can obtain the extreme points of the product distributions as
$$\{\begin{array}{ccc}|\alpha _1=|\alpha _2=\left(\begin{array}{c}1\\ 0\\ 0\end{array}\right)& & (P_{00}=\frac{1}{3},P_{10}=\frac{1}{3},P_{01}=0)\\ |\alpha _1=|\alpha _2=\frac{1}{\sqrt{3}}\left(\begin{array}{c}1\\ 1\\ 1\end{array}\right)& & (P_{00}=\frac{1}{3},P_{10}=0,P_{01}=\frac{1}{3})\\ |\alpha _1=\frac{1}{\sqrt{3}}\left(\begin{array}{c}\omega \\ \omega \\ \overline{\omega }\end{array}\right),|\alpha _2=\frac{1}{\sqrt{3}}\left(\begin{array}{c}\overline{\omega }\\ \omega \\ \overline{\omega }\end{array}\right)& & (P_{00}=0,P_{10}=\frac{1}{3},P_{01}=\frac{1}{3})\end{array}$$
$$\{\begin{array}{ccc}|\alpha _1=\frac{1}{\sqrt{3}}\left(\begin{array}{c}1\\ 1\\ \omega \end{array}\right),|\alpha _2=\frac{1}{\sqrt{3}}\left(\begin{array}{c}1\\ 1\\ \overline{\omega }\end{array}\right)& & (P_{00}=\frac{1}{3},P_{10}=0,P_{01}=0)\\ |\alpha _1=\frac{1}{\sqrt{3}}\left(\begin{array}{c}1\\ \omega \\ \overline{\omega }\end{array}\right),|\alpha _2=\frac{1}{\sqrt{3}}\left(\begin{array}{c}1\\ \omega \\ \overline{\omega }\end{array}\right)& & (P_{00}=0,P_{10}=\frac{1}{3},P_{01}=0)\\ |\alpha _1=\frac{1}{\sqrt{3}}\left(\begin{array}{c}1\\ \omega \\ 1\end{array}\right),|\alpha _2=\frac{1}{\sqrt{3}}\left(\begin{array}{c}1\\ 1\\ \overline{\omega }\end{array}\right)& & (P_{00}=0,P_{10}=0,P_{01}=\frac{1}{3})\end{array}.$$
(6-80)
Similar to $`(P_{00},P_{10},P_{20})`$ we obtain the extreme points. Then by convex combination of these points we obtain the feasible region (see Fig 2). The optimization problem is in the following form
$$\{\begin{array}{cc}\text{minimize}& C_\gamma =\frac{1}{8}(1P_{_{00}}(8x1)(P_{_{10}}P_{_{01}})\\ \text{subject to}& \frac{2}{3}P_{00}P_{10}P_{01}0\\ & P_{_{00}},P_{_{10}},P_{_{01}}0.\end{array}$$
(6-81)
We have been able to fined analytically the extreme points which at the same time don’t violate the plane $`(\frac{2}{3}P_{00}P_{10}P_{01}=0)`$. But we have failed to prove in general that no point lies out of the plane. Therefore, we have proved numerically that there is no violation from the plane. Thus, this feasible region is a convex hull or a polygon itself and reduces the optimization problem to a linear programming. So the vertices of the polygon are the solutions of the problem which by substituting them into equation $`C_\gamma `$ one can determine its minimum value as $`(C_\gamma )_{min}=(\frac{2(8x1)}{24})`$. We can find the critical r and by substituting the critical r in (3-15) we obtain a family of EW (called critical EW) resulting in
$$𝐫_c=\frac{3+24x}{1+24x}$$
$$W_c(x)=\frac{1}{3(1+24x))}((8x1))I_93|\psi _{_{00}}\psi _{_{00}}|3(8x1)(|\psi _{_{10}}\psi _{_{10}}||\psi _{_{01}}\psi _{_{01}}|)).$$
(6-82)
The obtained EW for two sets of triplet P, namely $`(P_{00},P_{20},P_{20})`$ and $`(P_{00},P_{20},P_{01})`$ can produce the most general form EW corresponding to combination of another triplet $`P_{ij}`$’s. Since under Fourier transform one can transform all the shifts in to the modulation. Moreover, the shift and modulation operators themselves can affect on the corresponding EW too, and so produce new combination of triplet.
## 7 Non-decomposible 3 $``$ 3 Bell states diagonal entanglement witnesses
By calculating the partial transpose of $`W_c(\frac{67}{756}x\frac{61}{378})`$ (for $`\{P_{00},P_{10},P_{20}`$ case) we prove that it is an nd-EW. The necessary and sufficient condition for non-decomposibility of $`W_c`$ reduces to the negativity of its partial transpose. Using the following relation
$$(\psi _{j^{}k^{}}><\psi _{jk})^{T_A}=\frac{1}{3}\underset{l,m}{}\omega ^{ml}\psi _{m+j^{},l+k^{}}><\psi _{m+j,3(lk)},$$
(7-83)
one can show that $`(W_c)^{T_A}`$ is a block diagonal, i.e., we have
$$(W_c)^{T_A}=\underset{j,k,k^{}}{}(O_j)_{kk^{}}\psi _{j^{}k^{}}><\psi _{jk},$$
with the matrices $`O_j`$ calculated as
$$O_j=\left(\begin{array}{ccc}0& 0& 0\\ 0& \frac{1}{6}& C_j\\ 0& \overline{C_j}& \frac{1}{6}\end{array}\right),$$
(7-84)
with
$$C_j=\frac{4}{3}(x\omega +(\frac{1}{4}x)\overline{\omega })\overline{\omega }^j,j=0,1,2.$$
(7-85)
Using the fact that $`C_2=C_1=C_0`$ one can show that the matrices $`O_j`$ have the same eigenvalues
$$\{\begin{array}{c}\lambda =0\\ \lambda _\pm =\lambda _\pm ^j=\frac{1}{6}\pm \frac{1}{6}\sqrt{4+48x(4x1)}.\end{array}$$
(7-86)
The above equation indicates that $`\lambda _{}`$ is negative except for the particular case in which $`x=\frac{1}{8}`$, i.e., $`W_{red}`$. Then different eigenvectors are so obtained
$$\{\begin{array}{ccc}\lambda =0& \varphi _j^0>=\psi _{j0}>,& \\ \lambda =\lambda _\pm & \varphi _j^\pm >=\frac{1}{\sqrt{\beta _\pm ^j^2+1}}(\beta _\pm ^j\psi _{j1}>+\psi _{j2}>),& \end{array},$$
(7-87)
where
$$\beta _\pm ^j=\frac{C_j\lambda _\pm }{\lambda _\pm ^2\lambda _\pm B}.$$
(7-88)
So we conclude that $`W_c^{T_A}`$ has three eigenvalues, namely $`\lambda _0,\lambda _\pm `$, each with degeneracy 3, and the following projection operators
$$\{\begin{array}{c}Q_+=_{j=0}^2\varphi _j^+><\varphi _j^+\\ Q_{}=_{j=0}^2\varphi _j^{}><\varphi _j^{}\\ Q_0=_{j=0}^2\varphi _j^0><\varphi _j^0.\end{array}$$
(7-89)
Here we have
$$W_c^{T_A}=\lambda _+Q_+\lambda _{}Q_{}.$$
(7-90)
The equation indicates that $`W_c^{T_A}`$ is not a positive definite operator except for the particular case $`W_{red}`$, hence it is non-decomposable entanglement witness.
We are interested in the n-d of EW given in (3-15) for the allowed values of p. Therefore, we write Eq.(3-15) as
$$W=\epsilon I_9/9+(1\epsilon )W_c,$$
(7-91)
with
$$\epsilon =\frac{𝐫+3}{4}.$$
(7-92)
Now, expanding $`I_9/9`$ in terms of the projection operator (7-89) as
$$I_9=Q_0^{T_A}+Q_{}^{T_A}+Q_+^{T_A},$$
(7-93)
the EW given by (3-15) can be written as
$$W=\epsilon /9Q_0^{T_A}+(\frac{\epsilon }{9}+(1\epsilon )\lambda _+)Q_+^{T_A}+(\frac{\epsilon }{9}(1\epsilon )\lambda _{})Q_{}^{T_A}.$$
(7-94)
The above form of EW indicates that its partial transpose $`W^{T_A}`$ is positive, i.e., it is decomposable EW if we have
$$W^{T_A}0(\frac{\epsilon }{9}(1\epsilon )\lambda _{})0𝐫\frac{3+9\lambda _{}}{1+9\lambda _{}},$$
(7-95)
for $`\frac{3+9\lambda _{}}{1+9\lambda _{}}𝐫3`$. It is not easy to tell where the EW is or is not decomposable. In the next section using some bound entangled state we will investigate their non-decomposability.
Now, in the remaining part of this section we try to obtain some nd-EW by taking the convex combination $`W_c(x)\text{for all }\frac{67}{756}x\frac{61}{378}`$ and $`W_{red}`$ (6-76) as
$$W_\mathrm{\Lambda }(x)=\mathrm{\Lambda }W_c(x)+(1\mathrm{\Lambda })W_{red},\mathrm{\Lambda }[0,1].$$
(7-96)
In order to test the positivity of $`W_\mathrm{\Lambda }^{T_A}(x)`$ we must first expand $`W_c`$ and $`W_{red}`$ in terms of the positive diagonal operators. Thus at first we write the projection operators defined in (7-89) in the following form
$$Q_\pm =\underset{k=0}{\overset{2}{}}|\chi _k^\pm \chi _k^\pm |,Q_0=\underset{k=0}{\overset{2}{}}|\psi _{ko}\psi _{k0}|$$
(7-97)
with
$$\chi _k^\pm >=(\psi _{k1}>\pm \omega ^k\psi _{k2}>).$$
(7-98)
Now writing $`I_9/9`$ in terms of the projection operator (7-97) and using the fact that
$$(\psi _{00}><\psi _{00})^{T_A}=\frac{1}{3}(\underset{k=0}{\overset{2}{}}\psi _{k0}><\psi _{k0}+\underset{k=0}{\overset{2}{}}\chi _k^+><\chi _k^+\underset{k=0}{\overset{2}{}}\chi _k^{}><\chi _k^{})$$
and
$$W_c^{T_A}(x)=\lambda _+\underset{k=0}{\overset{2}{}}\chi _k^+><\chi _k^+\lambda _{}\underset{k=0}{\overset{2}{}}\chi _k^{}><\chi _k^{},$$
(7-99)
we get for the partial transpose $`W_\mathrm{\Lambda }(x)`$in Eq.(7-96)
$$W_{\mathrm{\Lambda }}^{}{}_{}{}^{T_A}(x)=\mathrm{\Lambda }(\lambda _+)\underset{k=0}{\overset{2}{}}\chi _k^+><\chi _k^++(\mathrm{\Lambda }\lambda _{}+\frac{(1\mathrm{\Lambda })}{3})\chi _k^{}><\chi _k^{}.$$
(7-100)
This expression implies that $`W_{\mathrm{\Lambda }}^{}{}_{}{}^{T_A}(x)`$ is positive, since
$$\mathrm{\Lambda }\frac{1}{1+3\lambda _{}}.$$
(7-101)
Again, for $`\frac{1}{1+3\lambda _{}}\mathrm{\Lambda }1`$, it is not easy to talk about decomposable or non-decomposable $`W_\mathrm{\Lambda }(x)`$, and one needs to find some bound entangled states to show their non-decomposability, this will be done in the following section.
## 8 Detection of bound entangled state with Bell states diagonal entanglement witnesses
Now if we succeed to find any bound entangled state so that BDEW is able to detect this bound state corresponding to BDEW, from definition 2 in section 1 EW will be an nd-EW. Let a bound entangled Bell decomposable state be written as
$$\rho =\mu Q_0^{T_A}+\eta Q_+^{T_A}+\zeta Q_{}^{T_A},\rho ^{T_A}0\{\mu ,\eta ,\zeta \}0.$$
(8-102)
Optimal BDEW must detect this bound state, i.e.,
$$Tr[W_c\rho ]<0\eta \lambda _+<\zeta \lambda _{}.$$
(8-103)
On the other hand this bound state must be positive. For simplicity we use the operator $`W_c`$ and the identity operator $`I_9`$ in the bound state definition
$$Q_+^{T_A}=\frac{W_c+\lambda _{}(I_9Q_0^{T_A})}{\lambda _{}+\lambda _+},Q_{}^{T_A}=\frac{W_c+\lambda _+(I_9Q_0^{T_A})}{\lambda _{}+\lambda _+},$$
(8-104)
so that the bound state reduces to the following form
$$\rho =(\mu \frac{\eta \lambda _{}+\zeta \lambda _+}{\lambda _{}+\lambda _+})Q_0^{T_A}+(\frac{\eta \lambda _{}+\zeta \lambda _+}{\lambda _{}+\lambda _+})I_9+(\frac{\eta \zeta }{\lambda _{}+\lambda _+})W_c.$$
(8-105)
In this case $`Q_0=\psi _{00}><\psi _{00}+\psi _{10}><\psi _{10}+\psi _{20}><\psi _{20}`$ and by substituting this result in the Eq.(8-105) we get
$$\rho =(\mu \frac{\eta \zeta }{3(\lambda _{}+\lambda _+)})\psi _{00}><\psi _{00}+(\mu +(12x1)\frac{\eta \zeta }{3(\lambda _{}+\lambda _+)})\psi _{10}><\psi _{10}$$
$$+(\mu (12x2)\frac{\eta \zeta }{3(\lambda _{}+\lambda _+)})\psi _{20}><\psi _{20})+\frac{\eta (\lambda _{}\frac{1}{6})\zeta (\lambda _++\frac{1}{6})}{3(\lambda _{}+\lambda _+)})(\psi _{01}><\psi _{01}$$
$$+\psi _{02}><\psi _{02}+\psi _{11}><\psi _{11}+\psi _{22}><\psi _{22}+\psi _{12}><\psi _{12}+\psi _{21}><\psi _{21}).$$
(8-106)
The positivity of $`\rho `$ requires that all the Bell states diagonal operator coefficients to be positive, and that this condition be imposed on the coefficient $`\mu `$ only. So we get
$$\{\begin{array}{cc}x\frac{1}{8}& \mu \frac{(12x1)(\frac{1}{3}2\eta )}{(12x1)+3(\lambda _{}+\lambda _+)}\\ x\frac{1}{8}& \mu \frac{(212x)(\frac{1}{3}2\eta )}{(212x)+3(\lambda _{}+\lambda _+)},\end{array}$$
(8-107)
which, in this case means $`Q_0`$ is on the boundary. Now by using this bound entangled BD state we can find n-d condition for BDEW. We know EW will be an nd-EW if this EW is able to detect any bound state. Then by using the equations (7-94) and (8-102) we have
$$Tr(W\rho )=(\frac{\epsilon \mu }{3}+3(\frac{\epsilon }{9\lambda _+}+(1\epsilon ))\eta \lambda _++3(\frac{\epsilon }{9\lambda _{}}(1\epsilon ))\zeta \lambda _{}<0.$$
(8-108)
Now by substituting $`\epsilon `$ from Eq.(7-92) we obtain
$$𝐫<\frac{3+27(\zeta \lambda _{}\eta \lambda _+)}{1+27(\zeta \lambda _{}\eta \lambda _+)},$$
(8-109)
where the calculated r is greater than the represented r for EW in Eq.(7-95). Therefore, we can find one of the p’s corresponding to EW which is an nd-EW. Non-decomposable generalized EW for a general case is under investigation.
## 9 Choi map
Choi positive map $`\varphi (a,b,c):M^3M^3`$ is defined as
$$\varphi _{a,b,c}(\rho )=\left(\begin{array}{ccc}a\rho _{11}+b\rho _{22}+c\rho _{33}& 0& 0\\ 0& a\rho _{22}+b\rho _{33}+c\rho _{11}& 0\\ 0& 0& a\rho _{33}+b\rho _{11}+c\rho _{22}\end{array}\right)\rho ,$$
(9-110)
where $`\rho M^3`$. It was shown that $`\varphi (a,b,c)`$ is positive iff
$$a1,a+b+c3,\mathrm{\hspace{0.33em}\hspace{0.33em}1}a3.$$
(9-111)
Using Jamiolkowski isomorphism between the positive map and the operators we obtain the following $`33`$ EW corresponding to Choi map
$$W_{Choi}=\frac{1}{3(a+b+c1)}(a\underset{k=0}{\overset{2}{}}|\psi _{k0}\psi _{k0}|+b\underset{k=0}{\overset{2}{}}|\psi _{k2}\psi _{k2}|+c\underset{k=0}{\overset{2}{}}|\psi _{k1}\psi _{k1}|3|\psi _{00}\psi _{00}|).$$
(9-112)
Similar to BDEW we expand $`|\psi _{00}\psi _{00}|`$ using the identity operator and the other Bell diagonal states:
$$|\psi _{00}\psi _{00}|=I_9\underset{ij=0}{\overset{2}{}}|\psi _{ij}\psi _{ij}|.$$
(9-113)
Then we reduce EW to the following form
$$W_{Choi}=\frac{1}{3(a+b+c1)}((3a)I_9+3\underset{k=1}{\overset{2}{}}|\psi _{k0}\psi _{k0}|$$
$$+(b+3a)\underset{k=0}{\overset{2}{}}|\psi _{k2}\psi _{k2}|+(c+3a)\underset{k=0}{\overset{2}{}}|\psi _{k1}\psi _{k1}|).$$
(9-114)
Comparing with BDEW (3-15) we have
$$𝐫=\frac{3(3a)}{(a+b+c1)},$$
(9-115)
and the EW operator is defined as
$$W_{Choi}=𝐫I_9/9+(1𝐫)(\frac{1}{(82a+b+c)}\underset{k=1}{\overset{2}{}}|\psi _{k0}\psi _{k0}|+\frac{(b+3a)}{3(82a+b+c)}\underset{k=0}{\overset{2}{}}|\psi _{k2}\psi _{k2}|$$
$$+\frac{(c+3a)}{3(82a+b+c)}\underset{k=0}{\overset{2}{}}|\psi _{k1}\psi _{k1}|).$$
(9-116)
By comparing (9-114) with (3-15) we obtain the coefficients $`q_{ij}`$
$$q_{10}=q_{20}=\frac{1}{(82a+b+c)},q_{02}=q_{12}=q_{22}=\frac{b+3a}{3(82a+b+c)},$$
$$q_{01}=q_{11}=q_{21}=\frac{c+3a}{3(82a+b+c)}.$$
(9-117)
Note that if $`r`$ is negative, as introduced in EW above, this operator will be positive, but not a completely positive map. For $`𝐫0`$ we have $`1a3`$. By assuming $`abc`$, the minimum negative eigenvalue of choi EW (9-116) is given by
$$\frac{𝐫}{9}+(1𝐫)\frac{c+3a}{3(82a+b+c)}<0,$$
(9-118)
where upon substituting r from Eq.(9-117) we get $`1a2`$. This is equal to the introduced positivity condition of Choi map in .
By using (3-16) for non-negativity of the observable $`W_{choi}`$ we find the distributions $`P_{ij}`$ as a function of $`q_{ij}`$. The minimum value of $`C_\gamma `$ is obtained from the boundary of the feasible region, i.e., we have
$$(C_\gamma )=\frac{1}{(82a+b+c)}𝒫_1+\frac{(b+3a)}{3(82a+b+c)}𝒫_2+\frac{(c+3a)}{3(82a+b+c)}𝒫_3,$$
(9-119)
where $`𝒫_1=_{k=1}^2P_{_{k0}},𝒫_2=_{k=0}^2P_{_{k2}}`$ and $`𝒫_3=_{k=0}^2P_{_{k1}}`$. We can find the extreme value of $`(𝒫_1,𝒫_2,𝒫_3)`$ which is obtained under the product states $`|\gamma =|\alpha _1|\alpha _2`$ as
$$\{\begin{array}{c}𝒫_1=\alpha _1^2\beta _1^2+\alpha _2^2\beta _2^2+\alpha _3^2\beta _3^2\\ \frac{1}{3}\alpha _1\beta _1+\alpha _2\beta _2e^{i\varphi _2}+\alpha _3\beta _3e^{i\varphi _3}^2\\ 𝒫_2=\alpha _1^2\beta _1^2+\alpha _3^2\beta _1^2+\alpha _3^2\beta _2^2\\ 𝒫_3=\alpha _1^2\beta _2^2+\alpha _2^2\beta _3^2+\alpha _3^2\beta _1^2\end{array},$$
(9-120)
where $`|\alpha _1=\left(\begin{array}{c}\alpha _1\\ \alpha _2\\ \alpha _3\end{array}\right)`$ and $`|\alpha _1=\left(\begin{array}{c}\beta _1\\ \beta _2\\ \beta _3\end{array}\right)`$. One can obtain the extreme points of the $`(𝒫_1,𝒫_2,𝒫_3)`$ as
$$\{\begin{array}{ccc}|\alpha _1=|\alpha _2=\left(\begin{array}{c}1\\ \omega \\ \overline{\omega }\end{array}\right)& & (𝒫_1=\frac{1}{3},𝒫_2=\frac{1}{3},𝒫_3=\frac{1}{3})\\ |\alpha _1=\left(\begin{array}{c}1\\ 0\\ 0\end{array}\right),|\alpha _2=\left(\begin{array}{c}0\\ 0\\ 1\end{array}\right)& & (𝒫_1=0,𝒫_2=1,𝒫_3=0)\\ |\alpha _1=\left(\begin{array}{c}1\\ 0\\ 0\end{array}\right),|\alpha _2=\left(\begin{array}{c}0\\ \mathrm{1\; 0}\end{array}\right)& & (𝒫_1=0,𝒫_2=0,𝒫_3=1)\\ |\alpha _1=\left(\begin{array}{c}1\\ 0\\ 0\end{array}\right),|\alpha _2=\left(\begin{array}{c}1\\ 0\\ 0\end{array}\right)& & (𝒫_1=\frac{2}{3},𝒫_2=0,𝒫_3=0)\end{array}.$$
(9-121)
The convex combination of all extreme points provide a convex or a feasible region (Fig-3), then we have the following optimization problem
$$\{\begin{array}{cc}\text{minimize}& (C_\gamma )=(\frac{1}{(82a+b+c)}𝒫_1+\frac{(b+3a)}{3(82a+b+c)}𝒫_2+\frac{(c+3a)}{3(82a+b+c)}𝒫_3\\ \text{subject to}& 1\frac{3}{2}𝒫_1𝒫_2\frac{1}{2}𝒫_30\\ & 1\frac{3}{2}𝒫_1\frac{1}{2}𝒫_2𝒫_30\\ & 1𝒫_1𝒫_2𝒫_30\\ & 𝒫_1,𝒫_2,𝒫_30.\end{array}$$
(9-122)
Whether analytically we have been able to show that we will have violation only from the two planes
$$23𝒫_12𝒫_2𝒫_3=0$$
$$23𝒫_12𝒫_3𝒫_2=0.$$
Now let us assume that the maximum value of the violation from the planes is $`\mathrm{\Delta }<1`$. Thus, the equation of the plane passing through the new extreme points, parallel to the above plane, is obtained. Next we derive the intersection of the following adjacent planes
$$\{\begin{array}{cc}1)& 3𝒫_1+𝒫_2+2𝒫_3(2+\mathrm{\Delta })=0\\ 2)& 3𝒫_1+2𝒫_2+𝒫_3(2+\mathrm{\Delta })=0\\ 3)& 𝒫_1+𝒫_2+𝒫_31=0\\ 4)& 𝒫_1=0\\ 5)& 𝒫_2=0\\ 6)& 𝒫_3=0\\ 7)& 𝒫_1=\frac{2}{3}\\ 8)& 𝒫_2=1\\ 9)& 𝒫_3=1\end{array},$$
(9-123)
where new extreme points are obtained from intersecting the above planes. Next we calculate $`C_\gamma `$ for all the newly obtained extreme points and compare them with each other. Some easy calculations gives the minimum value of the parameter $`C_\gamma `$ which is independent from $`\mathrm{\Delta }`$
$$(C_\gamma )_{min}=\frac{6+2(ca)}{9(82a+b+c)},$$
(9-124)
then the critical value of the parameter r is obtained as
$$𝐫_c=\frac{6+2(ac)}{2+bc}.$$
(9-125)
For $`a=b=c=1`$ the parameter r reduces to $`r_c=3`$ corresponding to the well known reduction map. On the other hand, EW (9-116) must have positive trace under any product state $`|\gamma \gamma |`$. Thus the introduced r in EW must satisfy
$$rr_c\frac{3(3a)}{(a+b+c1)}\frac{6+2(ac)}{2+bc},$$
(9-126)
where the inequality is satisfied for all value of $`0a2`$ and $`abc`$.
## 10 Some separable states at the boundary of separable region
Here we introduce some set of separable states as
$$\rho _m=\underset{_k}{}|\psi _{km}\psi _{km}|=\underset{_l}{}|ll||l+ml+m|,$$
$$\rho _m^{}=\underset{_k}{}|\psi _{mk}\psi _{mk}|=\underset{_{l,l^{},k}}{}\omega ^{m(ll^{})}|ll^{}||l+kl^{}+k|,$$
$$\rho _n^{\prime \prime }=\underset{_k}{}|\psi _{nk,k}\psi _{nk,k}|=\underset{_{l,l^{},k}}{}\omega ^{nk(ll^{})}|ll^{}||l+kl^{}+k|,$$
(10-127)
where $`n=0,1,2,m=0,1,2`$. One can show that the convex sum of $`\rho _0,\rho _0^{\prime \prime }=\rho _0^{}`$, i.e, $`\rho _\mu ^S=\mu \rho _0+(1\mu )\rho _0^{},`$ is orthogonal to the optimal $`W_\mathrm{\Lambda }=\mathrm{\Lambda }W_c+(1\mathrm{\Lambda })W_{red},`$ i.e., we have $`Tr(W_\mathrm{\Lambda }\rho _S^\mu )=0`$.
Hence, $`\rho _\mu ^S`$ lie at the boundary of the separable region . On the other hand, one can show that by acting the local unitary operation $`U_{ij}`$ over $`W_\mathrm{\Lambda }`$ as $`(W_\mathrm{\Lambda })_{ij}=U_{ij}(W_\mathrm{\Lambda })U_{ij}^{}`$ he obtains a new set of optimal EW, $`(W_\mathrm{\Lambda })_{ij}`$, the application of which is not only to get a new set of bound entangled states by acting local unitary operation, but also to obtain some separable states $`(\rho _S^\mu )_{ij}=U_{ij}\rho _S^\mu U_{ij}^{}`$ as such which are the convex sum of separable states (10-127) at the boundary of separable states.
## 11 Conclusion
We have shown that finding generic Bell states diagonal entanglement witnesses (BDEW) for $`d_1d_2\mathrm{}.d_n`$ systems has reduced to a linear programming problem. Since solving linear programming for generic case is difficult we have considered the following special cases. Also we have considered BDEW for multi-qubit, $`2N`$ and $`33`$ systems and then have considered optimality condition for $`33`$ EW. Also, we have considered an n-d condition over $`33`$ BDEW and have obtained this condition for some special cases exactly. We have defined extensive group of nd-BDEW by combining critical EW and the reduction map (each with special coefficients). Then we have defined the Bell decomposable bound entangled state and have considered detection of this state with optimal BDEW and a general BDEW. Finally, we have considered Choi map as an example of BDEW. Optimality and non-decomposibility of EW for multi-qubit and $`2N`$ as well as EW for generic bipartite $`d_1d_2`$ systems and multipartite $`d_1d_2\mathrm{}d_n`$ are under investigation. As a physical implementation of EW we know that the optimization of decomposition of EW to find the smallest number of measurements possible for local measurement on a system can be used. Therefore to make use of this implementation of EW for the obtained EW’s is currently under investigation.
APPENDIX A
Minimization of the product distributions:
In Eq.(3-5) the Bell orthonormal states for a $`d_1d_2\mathrm{}d_n`$ ($`d_1d_2\mathrm{}d_n`$) have been introduced by applying local unitary operation on $`|\psi _{_{00}}`$. Let us further consider a pure product state $`|\gamma =|\alpha _1|\alpha _2\mathrm{}|\alpha _n`$. Then the product distributions can be written as
$$P_{_{i_1,i_2,\mathrm{},i_n}}(\gamma )=<\gamma \psi _{_{i_1,i_2,\mathrm{},i_n}}>^2.$$
(A-i)
It easily follows that
$$0P_{_{i_1,i_2,\mathrm{},i_n}}(\gamma )\frac{1}{d_1}.$$
(A-ii)
On the other hand, from the completeness of Bell states:
$$\underset{_{i_1,i_2,\mathrm{},i_n}}{}|\psi _{_{i_1,i_2,\mathrm{},i_n}}\psi _{_{i_1,i_2,\mathrm{},i_n}}|=I_{d_1}I_{d_2}\mathrm{}I_{d_n},$$
(A-iii)
we have $`_{_{i_1,i_2,\mathrm{},i_n}}P_{_{i_1,i_2,\mathrm{},i_n}}(\gamma )=1`$, which leads to
$$\underset{_{i_1,i_2,\mathrm{},i_n}}{}<\gamma \psi _{_{i_1,i_2,\mathrm{},i_n}}>^2=d_1.$$
(A-iv)
The above equation indicates that if we can show that for a particular choice of $`|\alpha _i`$’s, the $`d_1`$-number of $`<\gamma \psi _{_{i_1,i_2,\mathrm{},i_n}}>^2=P_{_{i_1,i_2,\mathrm{},i_n}}`$ can have their maximum value equal to $`\frac{1}{d_1}`$, then the remaining ones will be zero.
To minimize the summation $`C=_{ij}q_{ij}P_{ij}`$ for a $`33`$ system, assuming that $`q_{00}=0`$, let us first suppose that $`|\alpha =|\beta `$ so that $`P_{_{00}}=\frac{1}{3}`$. Then we find the set $`\alpha |U_{ij}|\beta ^2=1`$ for different possible choices of $`|\alpha `$ and $`U_{ij}`$:
$$|\alpha =\begin{array}{cccccccc}\left(\begin{array}{c}1\\ 1\\ 1\end{array}\right)& ,& \left(\begin{array}{c}1\\ \omega \\ \overline{\omega }\end{array}\right)& ,& \left(\begin{array}{c}1\\ \overline{\omega }\\ \omega \end{array}\right)& ,& |\psi _{01},|\psi _{02},& \text{min}(_{ij}q_{ij})=q_{_{01}}+q_{_{02}},\end{array}$$
$$|\alpha =\begin{array}{cccccccc}\left(\begin{array}{c}1\\ 0\\ 0\end{array}\right)& ,& \left(\begin{array}{c}0\\ 1\\ 0\end{array}\right)& ,& \left(\begin{array}{c}0\\ 0\\ 1\end{array}\right)& ,& |\psi _{10},|\psi _{20},& \text{min}(_{ij}q_{ij})=q_{_{10}}+q_{_{20}},\end{array}$$
$$|\alpha =\begin{array}{cccccccc}\left(\begin{array}{c}1\\ 1\\ \omega \end{array}\right)& ,& \left(\begin{array}{c}1\\ \omega \\ 1\end{array}\right)& ,& \left(\begin{array}{c}\omega \\ 1\\ 1\end{array}\right)& ,& |\psi _{11},|\psi _{22},& \text{min}(_{ij}q_{ij})=q_{_{11}}+q_{_{22}},\end{array}$$
$$|\alpha =\begin{array}{cccccccc}\left(\begin{array}{c}1\\ 1\\ \overline{\omega }\end{array}\right)& ,& \left(\begin{array}{c}1\\ \overline{\omega }\\ 1\end{array}\right)& ,& \left(\begin{array}{c}\overline{\omega }\\ 1\\ 1\end{array}\right)& ,& |\psi _{12},|\psi _{21},& \text{min}(_{ij}q_{ij})=q_{_{12}}+q_{_{21}}.\end{array}$$
The above relations imply that $`C_{mn}=\frac{1}{3}(q_1+q_2)`$, where $`q_1`$ and $`q_2`$ correspond to two of $`q_{ij}`$ appearing in the same row.
APPENDIX B
Critical entanglement witness is optimal:
According to the References , an EW will be optimal if for all positive operator P and $`\epsilon >0`$, the operator
$$W^{}=(1+\epsilon )W_c\epsilon P$$
(B-i)
is not an EW. In order to prove the critical EW given in (6-75) is optimal, we first show that
$$Tr(W_c|\alpha \alpha ||\alpha ^{}\alpha ^{}|)=0.$$
(B-ii)
It just suffices to check that for the product distribution $`P_{ij}=<\psi _{ij}|\alpha \alpha ||\alpha ^{}\alpha ^{}|\psi _{ij}>`$, we have $`P_{00}=\frac{1}{3},P_{01}=P_{02},P_{11}=P_{22},P_{12}=P_{21}`$.
Substituting $`P_{ij}`$ given above in (B-ii), it is easy to see that $`Tr(W_c|\alpha \alpha ||\alpha ^{}\alpha ^{}|)=0`$. Also it is straightforward to see that there exists no positive operator P with the constraint
$$Tr(P|\alpha \alpha ||\alpha ^{}\alpha ^{}|)=0,|\alpha $$
. Therefore, there exist no positive operator r to satisfy (B-i). Hence $`W_c`$, and in particular $`W_{red}`$, are optimal.
APPENDIX C
Simplex method for solving multi-qubit minimization problem
We know that simplex method is an elegant way for solving linear programming problems. As an example we obtain the $`P_{_{00\mathrm{}00}}`$ and $`P_{_{10\mathrm{}00}}`$ constraints in Eq.(4-31), thus we have two slack variables which are defined as
$$\omega _1=\frac{1}{2^{n1}}2P_{_{00\mathrm{}00}}+2P_{_{10\mathrm{}00}}(1\frac{1}{2^{n1}}),\omega _2=\frac{1}{2^{n1}}2P_{_{10\mathrm{}00}}+2P_{_{00\mathrm{}00}}(1\frac{1}{2^{n1}}).$$
(B-i)
We carry out this procedure to transform the inequality constraints (4-31) into equality
$$\{\begin{array}{cc}\text{maximize}& C_\gamma =\frac{1}{2(2^{n1}1)}((1x)+(1x)P_{_{00\mathrm{}00}}+((2^n1)x1)P_{_{10\mathrm{}00}})\\ \text{subject to}& \omega _1=\frac{1}{2^{n1}}2P_{_{00\mathrm{}00}}+2P_{_{10\mathrm{}00}}(1\frac{1}{2^{n1}})\\ & \omega _2=\frac{1}{2^{n1}}2P_{_{10\mathrm{}00}}+2P_{_{00\mathrm{}00}}(1\frac{1}{2^{n1}})\\ & P_{_{00\mathrm{}00}},P_{_{10\mathrm{}00}},\omega _1,\omega _20.\end{array}$$
(B-ii)
Now we rewrite the first equation in (B-ii) in terms of $`\omega _1`$ and $`\omega _2`$, making use of the slack variables:
$$C_\gamma =\frac{1}{2(2^{n1}1)}((1x)+\frac{(1x)a(2^n1)x+1}{2(a^21)}\omega _1+\frac{(1x)((2^n1)x1)a}{2(a^21)}\omega _2,$$
(B-iii)
where $`a=1\frac{1}{2^{n2}}`$. For $`0x\frac{1}{2^{n1}+1}`$ the coefficients $`\omega _1`$ and $`\omega _2`$ are both negative. Now from the simplex method we conclude $`\omega _1=\omega _2=0`$ , i.e., $`P_{_{00\mathrm{}00}}=P_{_{10\mathrm{}00}}=\frac{1}{2}`$. Thus the minimum value of $`C_\gamma =\frac{x}{2}`$. For $`\frac{1}{2^{n1}+1}x1`$, from (B-ii), we see that the coefficient $`P_{_{10\mathrm{}00}}`$ is negative, so that $`P_{_{10\mathrm{}00}}=0`$, hence $`P_{_{00\mathrm{}00}}=\frac{1}{2^{n1}}`$. Therefore, we find the minimum value of $`C_\gamma `$ as $`(C_\gamma )_{min}=\frac{1x}{2^n}`$.
Figure Captions
Figure-1: Feasible region for $`33`$ systems for particular choice of $`q_{_{00}}=0`$, $`q_{_{10}}=x`$, $`q_{_{20}}=\frac{1}{4}x`$ and others $`q`$’s being equal, i.e., when the linear programming variables are $`P_{_{00}}`$, $`P_{_{10}}`$ and $`P_{_{20}}`$ .
Figure-2: Feasible region for $`33`$ systems for particular choice of $`q_{_{00}}=0`$, $`q_{_{10}}=x`$, $`q_{_{01}}=\frac{1}{4}x`$ and others $`q`$’s being equal, i.e., when the linear programming variables are $`P_{_{00}}`$, $`P_{_{10}}`$ and $`P_{_{01}}`$ .
Figure-3: Feasible region for $`33`$ Choi map for particular choice of $`abc`$, i.e., when the linear programming variables are $`𝒫_1_{k=1}^2P_{_{k0}},𝒫_2_{k=0}^2P_{_{k2}}`$ and $`𝒫_1_{k=0}^2P_{_{k1}}`$. |
warning/0507/cond-mat0507142.html | ar5iv | text | # RANDOM WALKS ON GRAPHS: IDEAS, TECHNIQUES AND RESULTS
## 1 Introduction
A graph is the most general mathematical description of a set of elements connected pairwise by some kind of relation. Therefore, it is not surprising that graph theory has been successfully applied to a wide range of very different disciplines, from biology to social science, computing, psychology, economy, chemistry and physics.
In recent times, physicists have been mainly interested in graphs as models of complex systems, in condensed matter and in network theory. Indeed, these structures have proven to be very useful to describe inhomogeneous structures such as disordered materials, glasses, polymers, biomolecules as well as electric circuits, communication networks, statistical models of algorithms, and applications of statistical mechanics to different (non physical) systems.
The function of graphs in physics, however, is not purely descriptive. Geometry and topology have a deep influence on the physical properties of complex systems, where the presence of a large number of interacting degrees of freedom typically matters more than the interaction details. In fact, the specific interest of a physicist concerns the properties of a graph which most affect the dynamical and thermodynamical behaviour of the system it describes. On the other hand, the study of complex systems requires the introduction of statistical methods, to give an effective description of a number of quantities which, otherwise, would be too difficult to control.
Random walks are probably the simplest stochastic process affected by topology and, at the same time, the basic model of diffusion phenomena and non-deterministic motion. They have been extensively studied for decades on regular structures such as lattices, and most of the common wisdom concerning them relies on the results obtained in this particular geometry. The richer topology of a generic graph can have a dramatic effect on the properties of random walks, especially when considering infinite graphs, which are introduced to describe macroscopic systems in the thermodynamic limit. There, the asymptotic behaviour at long time typically exhibits universal features, only depending on large scale topology. On lattices, such features are known to be related to the Euclidean dimension only. On general graphs, universality allows to generalize the concept of dimension to inhomogeneous structures, providing a very powerful tool to investigate a large class of different physical models, apparently not connected to diffusion processes. On the other hand, a new and unexpected phenomenon arises in presence of strong inhomogeneity, namely the splitting between local and average properties. This provides a fundamental conceptual framework to investigate complex systems even from an experimental point of view.
Most results concerning random walks on graphs in physics have been obtained in the last two decades and are scattered in a large number of technical papers. This review is intended to provide the reader with a rigorous, self-contained and up to date account of the present knowledge about this subject. Particular attention has been paid to give a simple and general framework effectively resuming rather different results. As for specific calculations, we refer to bibliography, unless required by clarity reasons. The emphasis is always put on the physical meaning. The reader more interested in formal aspects can find a presentation focused on mathematics in another recent review .
The paper is organized as follows: In the first sections we give a brief mathematical description of graphs and random walks, introducing the language and the formalism we will use through the whole article. Then we present a simple treatment of the finite graphs case, before dealing with infinite graphs. The latter require the introduction of specific concepts, which are fully discussed in an introductory section. Then, the asymptotic behaviour of random walks on infinite graphs is studied and used to define the type problem and the spectral dimension. The difference between local and average properties is evidenced in the following sections. The concluding chapters are devoted to the analysis of a large class of specific graphs and to the relations of random walks with different physical problems.
## 2 Mathematical description of graphs
Let us begin by introducing the basic mathematical definitions and results concerning graphs .
A graph $`𝒢`$ is a countable set $`V`$ of vertices (or sites) $`(i)`$ connected pairwise by a set $`E`$ of unoriented links (or bonds) $`(i,j)=(j,i)`$. If the set $`V`$ is finite, $`𝒢`$ is called a finite graph and we will denote by $`N`$ the number of vertices of $`𝒢`$. Otherwise, when $`V`$ is infinite, $`𝒢`$ is called an infinite graph. A subgraph $`𝒮`$ of $`𝒢`$ is a graph whose set of vertices $`SV`$ and whose set of links $`E^{}E`$.
A path $`C_{ij}`$ in $`𝒢`$ connecting points $`i`$ and $`j`$ is a sequence of consecutive links $`\{(i,k)(k,h)\mathrm{}(n,m)(m,j)\}`$ and a graph is said to be connected, if for any two points $`i,jV`$ there is always a path joining them. In the following we will consider only connected graphs.
Every connected graph $`𝒢`$ is endowed with an intrinsic metric generated by the chemical distance $`r_{ij}`$ which is defined as the number of links in the shortest path connecting vertices $`i`$ and $`j`$.
A particular class of graphs, often occurring in physical applications, is characterized by the absence of closed paths containing an odd number of links. These graphs are called bipartite, since we can divide their sites into $`2`$ sets $`V_1`$ and $`V_2`$ such that the points of $`V_1`$ are connected by a link only to points $`V_2`$ and viceversa. Square and hypercubyc lattices are the most typical examples of bipartite graphs, as well as all trees (graphs without closed self-avoiding paths).
The graph topology can be algebraically represented introducing its adjacency matrix $`A_{ij}`$ given by:
$$A_{ij}=\{\begin{array}{cc}1& \mathrm{if}(i,j)E\hfill \\ 0& \mathrm{if}(i,j)E\hfill \end{array}$$
(1)
The Laplacian matrix $`\mathrm{\Delta }_{ij}`$ is defined by:
$$\mathrm{\Delta }_{ij}=z_i\delta _{ij}A_{ij}$$
(2)
where $`z_i=_jA_{ij}`$, the number of nearest neighbours of $`i`$, is called the coordination number of site $`i`$.
In order to describe disordered structures we introduce a generalization of the adjacency matrix given by the ferromagnetic coupling matrix $`J_{ij}`$, with $`J_{ij}0A_{ij}=1`$ and $`sup_{(i,j)}J_{ij}<\mathrm{}`$, $`inf_{(i,j)}J_{ij}>0`$. One can then define the generalized Laplacian:
$$L_{ij}=I_i\delta _{ij}J_{ij}$$
(3)
where $`I_i=_jJ_{ij}`$.
## 3 The random walk problem
Let us now introduce the so called simple random walk on a graph $`𝒢`$. Assuming the time $`(t)`$ to be discrete, we define at each time step $`t`$ the jumping probability $`p_{ij}`$ between nearest neighbour sites $`i`$ and $`j`$:
$$p_{ij}=\frac{A_{ij}}{z_i}=(Z^1A)_{ij}$$
(4)
where $`Z_{ij}=z_i\delta _{ij}`$.
This is the simplest case we can consider: the jumping probabilities are isotropic at each point and they do not depend on time; in addition the walker is forced to jump at every time step. As we will see later, the last condition, i.e. the impossibility of staying on site, although crucial for the short time behaviour, has no significant influence on the long time regime.
Usually, the random walk problem is considered to be completely solvable if, for any $`i,j𝒢`$ and $`tN`$, we are able to calculate the functions $`P_{ij}(t)`$, each representing the probability of being in site $`j`$ at time $`t`$ for a walker starting from site $`i`$ at time $`0`$. These probabilities are the elements of a matrix $`P=P_{ij}(t)`$ which is equal to the t-th power of the jumping probabilities matrix $`p=p_{ij}`$:
$$P_{ij}(t)=(p^t)_{ij}.$$
(5)
The relation (5) can be easily proven by induction on $`t`$. It also has an interesting physical interpretation as a sum over paths; developing the matrix products term by term we can write the whole expression as
$$P_{ij}(t)=(p^t)_{ij}=\underset{C_{ij}(t)}{}w(C_{ij}(t))$$
(6)
where the sum is over all $`t`$steps paths between $`i`$ and $`j`$. The weight $`w(C_{ij}`$ is the probability for the walker of going from $`i`$ to $`j`$ following exactly the path $`w(C_{ij})`$:
$$w(C_{ij}(t))=\underset{(k,l)C_{ij}(t)}{}p_{kl}$$
(7)
the product being over all the $`t`$ links belonging to the path.
The calculation of all $`P_{ij}(t)`$, which is straightforward as far as relatively small graphs are concerned, for large or infinite graphs becomes practically impossible and, above all, little significant. In fact, for large systems we are mainly interested in global and collective properties as it typically happens in statistical physics. Therefore, a small subset of all these quantities is usually chosen, together with some other related to them, which give an effective physical description of the random walker behaviour. The most relevant of them is from many points of view the probability $`P_{ii}(t)`$ of returning to the starting point after $`t`$ steps, also called the random walk autocorrelation function. As we will see, its asymptotic behaviour gives the most direct characterization of the large scale topology for infinite graphs. A related quantity is the average number $`P_{ii}`$ of returns to the starting point $`i`$, which can be generalized to the average number $`P_{ij}`$ of passages through $`j`$ starting from $`i`$:
$$P_{ij}\underset{t\mathrm{}}{lim}\underset{k=0}{\overset{t}{}}P_{ij}(k),$$
(8)
where the limit can be infinite.
The mean displacement $`r_i(t)`$ from the starting site $`i`$ after $`t`$ steps is deeply related to the diffusion properties and is defined as
$$r_i(t)\underset{j}{}r_{ij}P_{ij}(t)$$
(9)
Notice that, unlike the case of random walks in continuous Euclidean space, here we consider $`r`$ instead of $`r^2`$, the latter having no particular significance in absence of Euclidean metric.
The quantities introduced up to now are not ”sensible to the history”. Indeed, we can in principle determine all of them simply by considering the situation of the walker at time $`t`$ regardless of his previous behaviour. In order to keep track of what happened before the instant $`t`$, a different class of functions is introduced, starting with the first passage probability $`F_{ij}(t)`$. The latter denotes the conditional probability for a walker starting from $`i`$ of reaching for the first time the site $`ji`$ in $`t`$ steps. For $`i=j`$ the previous definition would not be interesting, being the walker in $`i`$ at $`t=0`$ by definition. Therefore, one defines $`F_{ii}(t)`$ to be the probability of returning to the starting point $`i`$ for the first time after $`t`$ steps and one sets $`F_{ii}(0)=0`$. In spite of the deeply different nature of $`P`$ and $`F`$, a fundamental relation can be established between them if all time steps form $`0`$ to $`t`$ are taken into account (in other words, we have to give up the time locality):
$$P_{ij}(t)=\underset{k=0}{\overset{t}{}}F_{ij}(k)P_{jj}(tk)+\delta _{ij}\delta _{t0}.$$
(10)
This can be easily obtained by considering that each walker which is in $`j`$ at time $`t`$ only has two possibilities: either it gets there for the first time, or it has reached $`j`$ for the first time at a previous time $`k`$ and then it has returned there after $`tk`$ steps. The first passage probability is in turn connected to other meaningful history dependent quantities. The probability $`F_{ij}`$ of ever reaching the site $`j`$ starting from $`i`$ (or of ever returning to i, if $`i=j`$) is given by
$$F_{ij}=\underset{t=0}{\overset{\mathrm{}}{}}F_{ij}(t)$$
(11)
By $`S_i(t)`$ we denote the average number of different sites visited after $`t`$ steps by a walker starting from $`i`$. Its relation to $`F_{ij}(t)`$ is
$$S_i(t)=1+\underset{k=1}{\overset{t}{}}\underset{j}{}F_{ij}(k)$$
(12)
Finally, the first passage time $`t_{ij}`$, i.e. the average time at which a walker starting from $`i`$, and passing at least once through $`j`$, reaches $`j`$ for the first time (or returns for the first time to $`i`$, if $`i=j`$) is
$$t_{ij}=\underset{t\mathrm{}}{lim}\frac{\underset{k=0}{\overset{t}{}}kF_{ij}(k)}{F_{ij}}$$
(13)
The simple random walk can be modified to give a richer behaviour and to describe more general physical problems. Indeed, one can introduce anisotropic jumping probabilities by substituting in (4) the adjacency matrix with a ferromagnetic coupling matrix:
$$p_{ij}=\frac{J_{ij}}{I_i}=(I^1J)_{ij}$$
(14)
where $`I_{ij}=I_i\delta _{ij}`$ and $`I_i=_kJ_{ik}`$. Depending on the specific properties of $`J_{ij}`$, this can produce only local effects or introduce a global bias which destroys the leading diffusive behaviour giving rise to transport phenomena.
Moreover, one can relax the constraint of jumping at each time step by introducing waiting and traps on the sites. The jumping probabilities are then modified to:
$$p_{i,j}=\frac{J_{i,j}+w_i\delta _{i,j}}{I_i+w_i+d_i}$$
(15)
where both $`w_i`$ and $`d_i`$ are real positive numbers. From (15), $`w_i/(I_i+w_i+d_i)`$ is the probability for the walker to stay on site $`i`$ instead of jumping away and $`d_i/(I_i+w_i+d_i)`$ the probability of disappearing (or dying, or being trapped forever) at site $`i`$. As we will see later, waiting only affects the short time behaviour, while traps can dramatically modify also the long time asymptotic properties.
## 4 The generating functions
Even if we consider only the few fundamental quantities devised at the end of the last section, their direct calculation can be in practice a hard or impossible task on general graphs. However, a powerful indirect mathematical technique exists allowing overcoming a series of typical difficulties: this is the discrete Laplace transform, which maps a time function into its generating function. The generating function $`\stackrel{~}{f}(\lambda )`$ of $`f(t)`$is defined by:
$$\stackrel{~}{f}(\lambda )=\underset{t=0}{\overset{\mathrm{}}{}}\lambda ^tf(t)$$
(16)
where $`\lambda `$ is a complex number. The inverse equation giving $`f(t)`$ from $`\stackrel{~}{f}(\lambda )`$ is
$$f(t)=\frac{^t\stackrel{~}{f}(\lambda )}{\lambda ^t}|_{\lambda =0}.$$
(17)
This equation is useful as fare as we are interested in small $`t`$ behaviour, but it becomes absolutely ineffectual in the study of asymptotic regimes for $`t\mathrm{}`$. In this case a very powerful tool is provided by the Tauberian theorems, relating the singularities of $`\stackrel{~}{f}(\lambda )`$ to the leading large $`t`$ behaviour of $`f(t)`$. We give here a rather general Tauberian theorem, which is particularly useful when dealing with random walks. The main assumption we make concerns the analytical form of the leading singularity: we only consider power laws and logarithmic behaviours, since all cases discussed in this paper as well as all physically meaningful cases belong to this class. Suppose that $`\stackrel{~}{f}(\lambda )`$ has its singularity nearest to $`\lambda =0`$ in $`\lambda =1`$. and that $`\stackrel{~}{f}(1ϵ)`$, for $`ϵ0^+`$ goes as
$$\stackrel{~}{f}(1ϵ)h(ϵ)+const\underset{i=0}{\overset{\mathrm{}}{}}\left({}_{}{}^{i}\mathrm{ln}(1/ϵ)\right)^{\alpha (i)}$$
(18)
where $`{}_{}{}^{i}\mathrm{ln}x\mathrm{ln}{}_{}{}^{i1}\mathrm{ln}x`$, with $`{}_{}{}^{0}\mathrm{ln}xx`$ e $`h(ϵ)`$ is finite for $`ϵ0^+`$.
Then, for $`t\mathrm{}`$
$$f(t)const^{}r^t\underset{i=0}{\overset{\mathrm{}}{}}{}_{}{}^{i}\mathrm{ln}_{}^{\beta (i)}(t)$$
(19)
where $`\beta (i)`$ are related to $`\alpha (i)`$ by
$$\beta (i)=\{\begin{array}{cc}\alpha (0)1\hfill & \text{for }i=0\hfill \\ \theta (im)(\alpha (i)+1)1\delta _{i,m}I(\stackrel{~}{d}/2)\hfill & \text{otherwise}\hfill \end{array}$$
(20)
where
$$m=\mathrm{min}\{i0|\beta (i)1\}$$
(21)
and
$$I\left(\stackrel{~}{d}/2\right)=\{\begin{array}{cc}1\hfill & \text{if }\stackrel{~}{d}/2\text{ is integer}\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}$$
(22)
The constant $`const^{}`$ is in general a function of $`const`$ and of all the exponents appearing in the previous formulas. We don’t give here its rather involved explicit expression, since it is not relevant to the purposes of this paper.
The generating functions are usually easier to calculate, since they allow exploiting some peculiar properties of random walks functions. Moreover, a series of relevant random walk parameters which are non-local in time, can be obtained directly from generating function, without calculating the corresponding time dependent quantities. A good example is given by $`P_{ij}`$, $`F_{ij}`$, and $`t_{ij}`$ which are related to $`\stackrel{~}{P_{ij}}(\lambda )`$ and $`\stackrel{~}{F_{ij}}(\lambda )`$ by
$$P_{ij}=\underset{\lambda 1^{}}{lim}\stackrel{~}{P}_{ij}(\lambda )$$
(23)
$$F_{ij}=\stackrel{~}{F}_{ij}(1)$$
(24)
$$t_{ij}=\underset{\lambda 1^{}}{lim}\frac{\mathrm{log}\stackrel{~}{F_{ij}}((\lambda )}{\lambda }$$
(25)
The basic property of random walks generating functions arise from the deconvolution of eq.(10), which after some straightforward steps becomes
$$\stackrel{~}{P}_{ij}(\lambda )\stackrel{~}{F}_{ij}(\lambda )\stackrel{~}{P}_{jj}(\lambda )+\delta _{ij}$$
(26)
In other words, the relation which was non-local in time becomes local in $`\lambda `$. As we will see in practical application, many iteration techniques for the analytical calculation of generating functions are based on this property.
## 5 Random walks on finite graphs
Finite graphs consist of a finite number of sites and links. In principle, every physical structure is composed of a finite number of elements, but it is well known that a series of behaviours occurring in macroscopic systems are better described in the thermodynamic limit. Indeed, the typical singularities and power laws characterizing phase transitions and asymptotic regimes, such as large scale, long times, low temperature and low frequency behaviours, can only be found on infinite graphs.
However, finite graphs are appropriate when dealing with mesoscopic structures and finite size effects. The random walk problem on finite graphs is simplified by the finiteness of the adjacency matrix. In fact, the analytical study is reducible to a spectral problem on a real finite-dimensional vector space, and numerical simulations are easily implemented by Monte Carlo techniques.
Let us first consider the case of random walk without traps, whose jumping probabilities are given by (15) with $`d_i=0`$ $`i`$. The matrix elements $`p_{ij}`$ satisfy the relations
$$p_{ij}0i,j$$
(27)
$$\underset{j=1}{\overset{N}{}}p_{ij}=1i$$
(28)
defining a stochastic matrix.
The stochastic matrices we are considering exhibits different properties according to the some general features of the graph and of the jumping probabilities . We distinguish two cases:
1. c1 If $`𝒢`$ is not bipartite, or if it has a staying probability on at least one site, then it has only one eigenvalue $`p_{max}`$ with maximal modulus and $`p_{max}=1`$. Moreover, the eigenvector corresponding to $`p_{max}`$ has the same entry on each site (usually one chooses $`v_{max}=(1,1,1,\mathrm{},1)`$ for simplicity)
2. c2 If $`𝒢`$ is bipartite without staying probabilities, then the spectrum of $`p`$ is symmetric with respect to the origin of the complex plane. Therefore, in addition to $`p_{max}`$ it has a second maximal modulus eigenvalue $`p_{min}=1`$. The eigenvector $`v_{max}`$ has the same properties of the previous case, while $`v_{min}`$ has all entries on $`V_1`$ equal to the same number $`v`$ and all entries on $`V_2`$ equal to $`v`$ (usually one chooses $`v=1`$)
In case c1, one can easily show that the random walk is ergodic, i.e. that it admits limit probabilities for $`t\mathrm{}`$:
$$P_{ij}^{\mathrm{}}=\underset{t\mathrm{}}{lim}P_{ij}(t)i,j$$
(29)
and that
$$P_{ij}^{\mathrm{}}=\frac{1}{N}i,j.$$
(30)
This means that, independently of the initial conditions, the asymptotic probabilities are the same over all the graph sites. Moreover, this uniform limit value is reached exponentially and the exponential decay of each matrix element is no slower than $`p_2^t`$, $`p_2`$ being the second greatest eigenvalue of $`p_{ij}`$.
In case c2, the random walk is not ergodic. In particular, we have
$$P_{ij}(t)=0$$
(31)
for all $`t`$ such that $`tr_{ij}`$ is odd. On the other hand, considering for each couple of sites $`i`$ and $`j`$ only the values $`t_{ij}^{}`$ of $`t`$ having the same parity as $`r_{ij}`$, one can show that
$$P_{ij}^{\mathrm{}}=\underset{t_{ij}^{}\mathrm{}}{lim}P_{ij}(t_{ij}^{})i,j$$
(32)
with
$$P_{ij}^{\mathrm{}}=\frac{2}{N}i,j$$
(33)
and the limit is reached exponentially as in case 1.
Similarly, one can easily prove that
$$\underset{t\mathrm{}}{lim}F_{ij}(t)=0i,j$$
(34)
and
$$\underset{t\mathrm{}}{lim}S_i(t)=Ni$$
(35)
the limit values being reached exponentially.
Moreover,
$$P_{ij}=\mathrm{}i,j$$
(36)
$$F_{ij}=1i,j$$
(37)
and
$$t_{ij}<\mathrm{}i,j$$
(38)
.
The introduction of at least one trap, setting $`d_i>0`$ for at least one site $`i`$ in (15), dramatically changes the random walk behaviour. The jumping probabilities matrix $`p`$ is no longer stochastic, since condition (28) is not satisfied. However, condition (27) (i.e. non-negativity) still holds, implying relevant properties. We can still distinguish between case $`1`$ and $`2`$, but the corresponding properties are modified as follows:
1. If $`𝒢`$ is not bipartite, or if it has a staying probability on at least one site, then it has only one eigenvalue $`p_{max}`$ with maximal modulus and $`p_{max}<1`$. Moreover, the entries $`v_{maxi}`$ of the eigenvector $`v_{max}`$, corresponding to $`p_{max}`$, have the same sign, and $`v_{max}`$ is the only eigenvector having such a property.
2. If $`𝒢`$ is bipartite without staying probabilities, then the spectrum of $`p`$ is symmetric with respect to the origin of the complex plane. Therefore, in addition to $`p_{max}<1`$ it has a second maximal modulus eigenvalue $`p_{min}=p_{max}`$. The eigenvector $`v_{max}`$ has the same properties of the previous case, while $`v_{min}`$ can be chosen in such a way that all its entries $`v_{mini}`$ on $`V_1`$ are equal to $`v_{maxi}`$ and all its entries $`v_{minj}`$ on $`V_2`$ equal to $`v_{maxj}`$.
In both cases the random walk is ergodic and the limit probabilities vanish:
$$P_{ij}^{\mathrm{}}=0i,j.$$
(39)
However, in case $`2`$ the time parity still has to be taken into account and (31) holds. Moreover, the asymptotic decay is exponential and no slower than $`p_{max}^t`$. Finally, as for the other random walk functions we get
$$\underset{t\mathrm{}}{lim}F_{ij}(t)=0i,j$$
(40)
$$\underset{t\mathrm{}}{lim}S_i(t)<Ni$$
(41)
the limit values being reached exponentially,
$$P_{ij}<\mathrm{}i,j$$
(42)
$$F_{ij}<1i,j$$
(43)
$$t_{ij}<\mathrm{}i,j$$
(44)
.
## 6 Infinite graphs
When dealing with macroscopic systems, composed of a very large number $`N`$ of sites, one usually takes the thermodynamic limit $`N\mathrm{}`$. This means that we have to consider infinite graphs, i.e. graphs composed by an infinite number of sites. This is particularly convenient for two main reasons:
* first of all, a single infinite structure effectively describes a very large (infinite, indeed) number of large structures having different sizes, but similar geometrical features
* the singularities in thermodynamic potentials typical of critical phenomena as well as a series of universal asymptotic behaviours only occur on infinite structures.
As for random walks on large real structures, the time dependence of physical quantities exhibits different features according to the time scale which is considered. For very long times, the walker can explore every site and its behaviour is described by the finite graphs laws introduced in the previous section. However, if the time is long enough to explore large portions of the system, but still too short to experience the finite size effects, many significant quantities are quite insensitive to local details and exhibit power law time dependence with universal exponents. Often, this is the most interesting regime in physical applications. On infinite graphs, this is the true asymptotic regime even for very large times; therefore, we can reproduce the universal behaviours of a huge variety of finite large structures simply by considering infinite graphs with similar topological features.
To deal with infinite graphs, some further mathematics has to be introduced. In particular, we need tools to ”explore” large scale topology. To this purpose, we define the Generalized Van Hove Spheres (GVHS): a GVHS $`𝒮_{o,r}`$ of centre $`o`$ and radius $`r`$ is the subgraph of $`𝒢`$, given by the set of vertices $`V_{o,r}=\{iV|r_{i,o}r\}`$ and by the set of links $`E_{o,r}=\{(i,j)E|iV_{o,r},jV_{o,r}\}`$.
Let us use $`|S|`$ to denote the number of elements of a set $`S`$. Then $`|V_{o,r}|`$, as a function of the distance $`r`$, describes the growth rate of the graph at the large scales . In particular, a graph is said to have a polynomial growth if $`oV_Xc,k`$, such that
$$|V_{o,r}|<cr^k.$$
(45)
For a graph satisfying (45), we define the upper growth exponent $`d_g^+`$ and the lower growth exponent $`d_g^{}`$ as
$$d_g^+=inf\{k||V_{o,r}|<c_1r^k,oV\}$$
(46)
and
$$d_g^{}=sup\{k||V_{o,r}|>c_2r^k,oV\}.$$
(47)
If $`d_g^+=d_g^{}`$, which usually happens on physically interesting structures, we call them the growth exponent $`d_g`$, or the connectivity dimension.
The connectivity dimension $`d_g`$ is known for a large class of graphs: on lattices, it coincides with the usual Euclidean dimension $`d`$, and for many fractals it has been exactly evaluated . In general, we can think of it as the analogous of the fractal dimension, when the chemical distance metric is considered instead of the usual Euclidean metric.
Infinite graphs are too general to describe systems of physical interest. Indeed, the discrete structures usually studied in physics are characterized by some important properties, often implicitly assumed in literature, which can be translated in mathematical requirements:
1. We consider only connected graphs, since any physical model on disconnected structures can be reduced to the separate study of the models defined on each connected component and hence to the case of connected graphs.
2. Since physical interactions are always bounded, the coordination numbers $`z_i`$, representing the number of neighbours interacting with the site $`i`$, have to be bounded; i.e. $`z_{max}|z_iz_{max}iV`$.
3. Real systems are always embedded in finite dimensional spaces. This constraint requires for the graph $`𝒢`$ the conditions:
1. $`𝒢`$ has a polynomial growth (Definition 45)
2. $$\underset{r\mathrm{}}{lim}\frac{|V_{o,r}|}{|V_{o,r}|}=0$$
(48)
where $`V_{o,r}`$ denotes the border of $`V_{o,r}`$, i.e. the set of points of $`V_{o,r}`$ not belonging to $`V_{o,r1}`$ (the existence itself of the limit is a physical requirement on $`𝒢`$). This condition is equivalent to require that boundary conditions are negligible in the thermodynamic limit.
Notice that some graphs studied in physical literature, such as the Bethe lattice, do not satisfy (a) and (b), while many random graphs do not fulfil B.
For a large class of physically interesting graphs we have considered so far, conditions (a) and (b) appear to be equivalent. However for the equivalence of the two conditions a rigorous result is still lacking. A graph satisfying A, B and C will be called physical graph. Conditions A and B represent strong constraints on $`𝒢`$ and, as we will see later, they have very important consequences.
## 7 Random walks on infinite graphs
Considering random walks on infinite structures, some further mathematical constraints are to be introduced to describe physical situations.
First of all, the problem of uniform boundedness comes into play. Indeed, in (14) the conditions
$$J_{min},J_{max}>0|J_{min}J_{i,j}J_{max}i,j$$
(49)
together with B are usually required to exclude the presence of a global bias, which would generate a non-diffusive behaviour.
Moreover, in (15), in presence of waiting and traps, analogous considerations lead to the following conditions
$$w_{min},w_{max}>0|\mathrm{either}w_i=0,\mathrm{or}w_{min}w_iw_{max}i$$
(50)
$$d_{min},d_{max}>0|\mathrm{either}d_i=0,\mathrm{or}d_{min}d_id_{max}i$$
(51)
In the case of finite graphs, the possibility of associating to any matrix an operator acting on a finite dimensional vector space allowed to obtain very general and rigorous results. In the infinite case, it is in general impossible to associate a linear operator acting on a Hilbert space to any matrix. However, when B holds, the jumping probabilities matrix is quite particular: indeed, it only has a limited number of non vanishing entries in each row and column. Due to this property, the elements of a matrix product are given by finite sums, as in the finite graphs case, instead of being sums of series. Therefore, the typical convergence problems of infinite dimension space do not arise, allowing for a simple and effective study of random walks properties.
Despite the increased mathematical complexity, many general results about infinite graphs have been rigorously proven. Some have correspondents in the finite graphs case, but most of them concern quantities and properties which cannot be even defined on finite structures. The rest of this section is devoted to the former: following the same format used in section 5 we resume the main differences with respect to the finite case. The new properties arising in the thermodynamic limit will be discussed in the following sections.
First of all, let us consider random walks without traps on infinite graphs satisfying A,B, (49) and (50). It can be shown that
$$P_{ij}^{\mathrm{}}=\underset{t\mathrm{}}{lim}P_{ij}(t)=0i,j$$
(52)
even for bipartite graphs. However, for bipartite graphs without staying probabilities, we still have
$$P_{ij}(t)=0$$
(53)
for all $`t`$ such that $`tr_{ij}`$ is odd. Therefore, to study the large times asymptotic behaviours, as in the finite case we usually consider, for any given couple of sites $`i`$ and $`j`$, only the values $`t_{ij}^{}`$ of $`t`$ having the same parity as $`r_{ij}`$. Unlike the finite case, the limit in (52) in general is not reached exponentially. Indeed, if C also holds, i.e. for physical graphs, the asymptotic behaviour is typically a power law, whose exponent only depends on topology, as we will discuss in details in the next sections. Notice that the widely studied case of Bethe lattices, not satisfying C, is still characterized by an exponential decay.
Similarly, one can easily prove that
$$\underset{t\mathrm{}}{lim}F_{ij}(t)=0i,j$$
(54)
and
$$\underset{t\mathrm{}}{lim}S_i(t)=\mathrm{}i$$
(55)
the asymptotic behaviour being always bounded from above by $`t`$.
As for the quantities concerning the number of visits and the first visit probabilities, the situation is far more complex. Indeed, dramatically different behaviours can occur, according to the graph topology. In particular
$$P_{ij}=\mathrm{}i,j\mathrm{or}<\mathrm{}i,j$$
(56)
$$F_{ij}=1i,j\mathrm{or}<1i,j$$
(57)
and
$$t_{ij}<\mathrm{}i,j\mathrm{or}=\mathrm{}i,j$$
(58)
.
The classification of infinite graphs according to these possible behaviours is the subject of the next section.
## 8 Recurrence and transience: the type problem
On finite graphs, in absence of traps, the probability of ever reaching (or returning to) a site, $`F_{ij}`$, is always 1. This means that the walker surely visits each site. This probability can be lowered only by adding traps, but in this case the total probability is not conserved, i.e. the walker asymptotically disappears. On infinite graphs, a third possibility arises, which is expressed in (56) and (57): the walker can escape forever from its starting point, or never reach a given site, even in absence of traps.
This phenomenon was first noticed by Polya in 1921 on lattices: he showed that, while in $`1`$ and $`2`$ dimensions $`F_{ij}=1`$, for $`d3`$ $`F_{ij}<1`$ . Since him, the former case has been called recurrent and the latter transient. Transience is an exclusive property of infinite graphs and it is fundamentally due to large scale topology. In other words, in the transient case, it happens that the number of paths leading the walker away from its starting point is large enough, with respect to the number of returning paths, to act as an asymptotic trap (still conserving the total probability).
As we will see in a while, transience and recurrence of random walks, when (49) and (50) are satisfied, only depend on the graph topology. Therefore they are intrinsic properties of a discrete structure and the classification of infinite graphs according to them is also known as the type problem.
Let us define the problem mathematically. First of all, a very general theorem on Markov chains states the following:
$$i,j𝒢|F_{ij}=1F_{hk}=1h,k𝒢$$
(59)
this means that recurrence is point independent, or, in other words, that if a walker surely reaches a point $`j`$ starting from a given point $`i`$, then it surely reaches any point $`k`$ starting from any point $`h`$. It is straightforward to see that an analogous result follows for the case $`F_{ij}<1`$. Therefore, recurrence and transience are global properties of a random walk.
Another important result relates $`F_{ij}`$ and $`P_{ij}`$. Indeed, from (23), (24), and (26), it follows that
$$F_{ij}=1P_{ij}=\mathrm{}$$
(60)
and
$$F_{ij}<1P_{ij}<\mathrm{}$$
(61)
i.e., a walk is recurrent (transient) if and only if any site is visited an infinite (finite) number of times. The latter can be taken as an alternative definition of recurrence and transience. However, as we will see, the situation is more complex when considering averages over all the sites. A consequent property concerns the way the walker explores the sites of $`𝒢`$. Indeed, it can be shown that
$$F_{ij}=1\underset{t\mathrm{}}{lim}\frac{S_i(t)}{t}=0$$
(62)
while
$$F_{ij}<10<\underset{t\mathrm{}}{lim}\frac{S_i(t)}{t}<1$$
(63)
In the first situation, where the number of distinct visited sites increases slower than the number of steps, is sometimes called compact exploration, since the subgraphs of the visited sites presents a negligible number of ”holes”.
Recurrent graphs exhibit a further relevant property: one can show that
$$\underset{\lambda 1^{}}{lim}\frac{\stackrel{~}{P_{ij}}(\lambda )}{\stackrel{~}{P_{hk}}(\lambda )}=\underset{t\mathrm{}}{lim}\frac{P_{ij}(t)}{P_{hk}(t)}=\frac{z_j}{z_k}i,j,h,k$$
(64)
The most important properties in the type problem concern its invariance with respect to a wide class of dynamical and topological transformations, establishing its independence of the graph details.
First of all, consider two different random walks (without traps) on the same graph $`𝒢`$, one (W) defined by the ferromagnetic coupling matrix $`J_{ij}`$ and by the waiting probabilities $`w_i`$, and one (W’) by $`J_{ij}^{}`$ and $`w_i^{}`$ . It can be shown that, if both satisfy (49) and (50), then W is recurrent if and only if W’ is. In other words, any local bounded rescaling of ferromagnetic couplings and waiting probabilities leaves the random walk type unchanged. Therefore, provided the previously mentioned boundedness conditions are satisfied, the walk type only depends on the graph topology.
Moreover, even the local topological details are irrelevant to determine the type of a graph. Indeed, it is possible to show that recurrence and transience are left invariant by adding and cutting of links satisfying the quasi-isometry conditions. More precisely, two graphs $`𝒢`$ and $`𝒢^{}`$ are called quasi-isometric if there are a mapping $`\phi :𝒢𝒢^{}`$ and constants $`A>0`$, $`B0`$ such that
$$A^1r_{ij}Br_{\phi i,\phi j}^{}Ar_{ij}+B$$
for all $`i,j𝒢`$, and
$$r_{i^{},\phi 𝒢}^{}B$$
for every $`i^{}𝒢^{}`$.
If $`B=0`$ then we say that $`𝒢`$ and $`𝒢^{}`$ are metrically equivalent. Quasi-isometries can be defined between arbitrary metric spaces and represent the most general local topology deformations. Typical examples of them are given by the decimation transformations used on fractals and in real-space renormalization. In some sense, we can consider quasi-isometries as their extension to general networks.
All the results presented so far refer to random walks without traps, i.e. to jumping probabilities given by (15) with $`d_i=0`$. The introduction of at least one trap, setting $`d_i>0`$ for at least one site $`i`$, has a very general and simple influence on the random walk behaviour: indeed transience is left unchanged, while recurrent random walks always become transient.
## 9 The local spectral dimension
As well as it happens for recurrence and transience properties, large scale topology affects the long time dependence of random walks quantities on infinite graphs. Indeed, it has been known for many years that, on regular (translation invariant) lattices, the exponents of the asymptotic power laws of random walks only depend on the lattice (Euclidean) dimension $`d`$. For example,
$$P_{ii}(t)t^{d/2}\mathrm{for}t\mathrm{},i$$
(65)
and
$$S_i(t)t^{min(1,d/2)}\mathrm{for}t\mathrm{},i,\mathrm{for}d2$$
(66)
(while, for $`d=2`$, $`S_i(t)t/\mathrm{ln}t)`$. As we mentioned before, these laws typically present power behaviour even on general physical graphs, and the exponents of such powers can be used to define a generalized dimension.
Let us consider a random walk without waitings and traps satisfying (49), and suppose that, for a given $`i𝒢`$
$$P_{ii}(t)t^{\stackrel{~}{d}/2}\mathrm{for}t\mathrm{}$$
(67)
then it can be shown that
$$P_{hk}(t)t^{\stackrel{~}{d}/2}\mathrm{for}t\mathrm{},h,k$$
(68)
(for bipartite graphs, the usual assumptions on the parity $`t`$ are understood).
This means that the exponent of the power law is site independent and, therefore, it is a parameter characterizing the whole random walk. Since $`\stackrel{~}{d}=d`$ on regular lattice, we can consider it as a dimension associated to the random walk on $`𝒢`$. More precisely, we shall call local spectral dimension the limit
$$\stackrel{~}{d}=2\underset{t\mathrm{}}{lim}\frac{\mathrm{ln}P_{ii}(t)}{\mathrm{ln}t}$$
(69)
when it exists.
Notice that the existence of this limit for a given $`i`$ implies it exists and has the same value for any $`j𝒢`$. Moreover the definition given in (69) is more general than (67), since it includes the case of possible multiplicative correction to the asymptotic behaviour, provided they are slower than any power law (e.g. logarithmic corrections).
¿From an historical point of view, the term ”spectral dimension” was first introduced by Alexander and Orbach in 1982 , studying the anomalous vibrational dynamics on fractals. In the same work, they suggested that even the random walks should be ruled by the same parameter and wrote eq. (67). Then the definition was generalized to general networks by Hattori, Hattori and Watanabe . Later, it has been shown that the anomalous dimension involved in vibrational dynamics is the average spectral dimension (), we shall discuss in further sections, which coincides with $`\stackrel{~}{d}`$ only for particular graphs, such as exactly decimable fractals.
As for the existence of the limit (69), a general theorem is still lacking, but it can be easily proven that the asymptotic decay of $`P_{ii}(t)`$ is always bounded from above and from below by power laws. In any case, on all known cases of random walks on physical graphs, the local spectral dimension has been shown to exist. Notice that for the Bethe lattice, which does not fulfil the polynomial growth condition, the limit (69) is infinite. ¿From now on, we shall consider random walks on graphs where $`\stackrel{~}{d}`$ is defined. Then, one can easily derive the following results:
* Random walks are recurrent if $`\stackrel{~}{d}<2`$ and transient if $`\stackrel{~}{d}>2`$. For $`\stackrel{~}{d}=2`$, if (67) holds, random walks are recurrent. However, subleading corrections to the power law can change the type to transient.
* When (67) holds,
$$S_i(t)t^{min(1,\stackrel{~}{d}/2)}\mathrm{for}t\mathrm{},i,\mathrm{for}\stackrel{~}{d}2$$
(70)
otherwise, in general,
$$\underset{t\mathrm{}}{lim}\frac{\mathrm{ln}S_i(t)}{\mathrm{ln}t}=min(1,\stackrel{~}{d}/2)i,\mathrm{for}\stackrel{~}{d}2$$
(71)
* When (67) holds,
$$F_{ij}(t)t^{min(\stackrel{~}{d}/22,\stackrel{~}{d}/2)}\mathrm{for}t\mathrm{},i,j\mathrm{for}\stackrel{~}{d}2$$
(72)
otherwise, in general,
$$\underset{t\mathrm{}}{lim}\frac{\mathrm{ln}F_{ij}(t)}{\mathrm{ln}t}=min(\stackrel{~}{d}/22,\stackrel{~}{d}/2)i,j\mathrm{for}\stackrel{~}{d}2$$
(73)
The case $`\stackrel{~}{d}=2`$ is rather particular. Indeed $`\stackrel{~}{d}=2`$ is a critical dimension for random walks, discriminating recurrence from transience. The asymptotic behaviours of $`S_i(t)`$ and $`F_{ij}(t)`$ have a different dependence on $`\stackrel{~}{d}`$ for $`\stackrel{~}{d}<2`$ and $`\stackrel{~}{d}>2`$. In particular, the probability of first visit has the same time decay of $`P_{ij}`$ for $`\stackrel{~}{d}>2`$ while it decays faster for $`\stackrel{~}{d}<2`$. When $`\stackrel{~}{d}=2`$, the behaviours of $`S_i(t)`$ and $`F_{ij}(t)`$ are strongly affected by subleading corrections.
As for the type problem, local spectral dimension presents interesting invariance properties. First of all, it can be shown that waitings satisfying (50) do not affect its value . Moreover, for $`\stackrel{~}{d}<2`$, on a given $`𝒢`$ it is the same for all ferromagnetic couplings satisfying (49) . Unfortunately, an analogous result has not been proven for $`\stackrel{~}{d}>2`$. However, we will see in later sections that the average spectral dimension has also this universality property.
The introduction of a finite number of traps do not affect $`\stackrel{~}{d}`$ if $`\stackrel{~}{d}>2`$. If $`\stackrel{~}{d}<2`$ a finite number of traps (even only one) changes $`\stackrel{~}{d}`$ to $`\stackrel{~}{d}+1`$. If the traps are infinite the behaviour is more complex and depends on their distribution.
## 10 Averages on infinite graphs
Usually, infinite graphs describing real systems are inhomogeneous, i.e., in mathematical terms, they are not invariant with respect to a transitive symmetry group. In simpler words, this means that the topology is seen in a different way from every site. The main effect of inhomogeneity is that the numerical values of physical quantities are site dependent. Therefore, one is typically interested in taking averages over all sites. This requires the introduction of suitable mathematical tools.
First of all, the average in the thermodynamic limit $`\overline{\varphi }`$ of a function $`\varphi _i`$ defined on each site $`i`$ of the infinite graph $`𝒢`$ is defined by:
$$\overline{\varphi }\underset{r\mathrm{}}{lim}\frac{{\displaystyle \underset{iS_{o,r}}{}}\varphi _i}{N_{o,r}}.$$
(74)
The measure $`|S|`$ of a subset $`S`$ of $`V`$ is the average value $`\overline{\chi (S)}`$ of its characteristic function $`\chi _i(S)`$ defined by $`\chi _i(S)=1`$ if $`iS`$ and $`\chi _i(S)=0`$ if $`iS`$. The measure of a subset of links $`E^{}E`$ is given by:
$$|E^{}|\underset{r\mathrm{}}{lim}\frac{E_r^{}}{N_{o,r}}.$$
(75)
where $`E_r^{}`$ is the number of links of $`E^{}`$ contained in the sphere $`S_{o,r}`$. The normalized trace $`\overline{\mathrm{Tr}}B`$ of a matrix $`B_{ij}`$ is:
$$\overline{\mathrm{Tr}}B\overline{b}$$
(76)
where $`b_iB_{ii}`$. If condition C holds, then we can prove that the averages of a bounded from below function $`\varphi _i`$ are independent from the centre $`o`$ of the spheres sequence, using the fact that $`\chi _i(S)`$ is bounded and that measures of subsets are always well defined.
Now, due to this site independence, we have a good definition of averages which we will use in dealing with properties of random walks on infinite graphs. As we shall see in the next section, on inhomogeneous networks the averages of site dependent functions can have a very different behaviour from their local counterparts, giving rise to rather unexpected phenomena.
## 11 The type problem on the average
In the last few years it has become clear that bulk properties are affected by the average values of random walks return probabilities over all starting sites: this is the case for spontaneous breaking of continuous symmetries , critical exponents of the spherical model , harmonic vibrational spectra . Therefore the classification of discrete structure in terms of recurrence on the average and transience on the average appears to be the most suitable. Unfortunately, while for regular lattices the two classifications are equivalent, on more general networks they can be different and one has to study a Type-Problem on the Average .
This is defined using the return probabilities on the average $`\overline{P}`$ and $`\overline{F}`$, which are given by:
$$\overline{P}=\underset{\lambda 1^{}}{lim}\overline{\stackrel{~}{P}(\lambda )}\underset{\lambda 1^{}}{lim}\overline{\mathrm{Tr}}\stackrel{~}{P}(\lambda )$$
(77)
$$\overline{F}=\underset{\lambda 1^{}}{lim}\overline{\stackrel{~}{F}(\lambda )}\underset{\lambda 1^{}}{lim}\overline{\mathrm{Tr}}\stackrel{~}{F}(\lambda )$$
(78)
A graph $`𝒢`$ is called recurrent on the average (ROA) if $`\overline{F}=1`$, while it is transient on the average (TOA) when $`\overline{F}<1`$.
Recurrence and transience on the average are in general independent of the corresponding local properties. The first example of this phenomenon occurring on inhomogeneous structures was found in a class of infinite trees called NTD (see sect. 13.3) which are locally transient but recurrent on the average .
Moreover, while for local probabilities (26) gives:
$$\stackrel{~}{P}_{ii}(\lambda )\stackrel{~}{F}_{ii}(\lambda )\stackrel{~}{P}_{ii}(\lambda )+1$$
(79)
an analogous relation for (78) and (77) does not hold since averaging (79) over all sites $`i`$ would involves the average of a product, which due to correlations is in general different from the product of the average. Therefore the double implication $`\stackrel{~}{F}_i(1)=1lim_{\lambda 1}\stackrel{~}{P}_i(\lambda )=\mathrm{}`$ is not true. Indeed there are graphs for which $`\overline{F}<1`$ but $`\overline{P}=\mathrm{}`$ (an example is shown in Fig.1) and the study of the relation between $`\overline{P}`$ and $`\overline{F}`$ is a non trivial problem.
A detailed study of this relation shows that a complete picture of the behaviour of random walks on graphs can be given by dividing transient on the average graphs into two further classes, which are called pure and mixed transient on the average (TOA).
First, considering a ROA graph, it can be proven that if $`\overline{F}=1`$ then $`\overline{P}=\mathrm{}`$. The proof can be easily generalized to graphs in which there is a positive measure subset $`S`$ such that: $`lim_{\lambda 1}\overline{\chi (S)\stackrel{~}{F}(\lambda )}=|S|`$. Indeed in an analogous way it can be proven that:
$$\overline{P}\underset{\lambda 1}{lim}\overline{\chi (S^{})\stackrel{~}{P}(\lambda )}=\mathrm{}S^{}S,|S^{}|>0$$
(80)
We call mixed transient on the average a TOA graph having a positive measure subset $`S`$ such that:
$$\underset{\lambda 1}{lim}\overline{\chi (S)\stackrel{~}{F}(\lambda )}=|S|.$$
(81)
while a graph is called pure TOA, if:
$$\underset{\lambda 1}{lim}\overline{\chi (S)\stackrel{~}{F}(\lambda )}<|S|SV,|S|>0$$
(82)
Examples of pure TOA graphs are all the $`d`$dimensional cubic lattices with $`d>2`$, while the ”haired cube” of Fig.1 is a typical mixed TOA graph. Notice that a relevant theorem establishes that for mixed TOA graphs we have $`\overline{P}=\mathrm{}`$, while for pure TOA graphs $`\overline{P}`$ is finite. A further important property, characterizing mixed TOA graphs, allows simplifying the study of statistical models on these very inhomogeneous structures. It can be shown that, in this case, the graph $`𝒢`$ can be always decomposed in a pure TOA subgraph $`𝒮`$ and a ROA subgraph $`\overline{𝒮}`$ with independent jumping probabilities by cutting a zero measure set of links $`𝒮\{(i,j)E|i𝒮j\overline{𝒮}\}`$. The separability property implies that the two subgraphs are statistically independent and that their thermodynamic properties can be studied separately. Indeed, in the thermodynamic limit, the partition functions referring to the two subgraphs factorize .
To conclude this section, we note that the same invariance properties of the local type problem under addition of waiting probabilities, coupling rescaling and quasi-isometries still hold for the type problem on the average. This means, as for the local case, that recurrence and transience on the average are intrinsic properties of a graph and not only of a specific random walk defined on it. On the other hand, the introduction of a finite number of traps does not change the type on the average. Notice also that a slightly different definition of the type problem on the average can be found in mathematical literature ; it is more convenient for the formal development of the theory, but it is not directly related to statistical models on graphs.
## 12 The average spectral dimension
The asymptotic time dependence of the return probability on the average can be used to define a new intrinsic dimension which turns out to be very strictly related to the physical behaviour of statistical models on graphs , as we will briefly discuss in the last section.
Indeed, even if the asymptotic time decay of $`P_{ii}(t)`$ is always the same for all sites $`i`$, when the graph topology is strongly inhomogeneous it happens that its average over all the sites decays according to a different law. The average spectral dimension $`\overline{d}`$ is defined for physical graphs, like the local one in (67) and (69), by
$$\overline{P}(t)t^{\overline{d}/2}\mathrm{for}t\mathrm{}$$
(83)
when the asymptotic behaviour is a power law without subleading corrections, or, more generally, by
$$\overline{d}=2\underset{t\mathrm{}}{lim}\frac{\mathrm{ln}\overline{P}(t)}{\mathrm{ln}t}.$$
(84)
Notice, however, that, differently from the local case, no physical graphs are known, up to now, where the long time decay is not given by (83). Considerations analogous to those presented for the local case hold here, concerning the existence of the limit (84). Obviously, in all cases where the local type is different from the average type, also the local spectral dimension differs from the average spectral dimension. A typical example, and, historically, the first one, is given again by NTD (see sect. 13.3) for a detailed account). However, the relations between $`\overline{d}`$ and the type problem on the average are not the same as in the local case. Indeed, while if $`\overline{d}>2`$ the walk is always pure TOA, random walks with $`\overline{d}<2`$ can be either pure ROA or mixed TOA.
The most relevant property of $`\overline{d}`$ is without any doubt its strong invariance with respect to a very large class of dynamical and topological transformations, making it a unique universal parameter associated to a graph $`𝒢`$ .
These transformations can be divided into three main classes:
1. Dynamical transformations leaving the graph topology unchanged. These consist in addition of waitings and of a finite number of traps, as well as in bounded local rescaling of ferromagnetic couplings.
2. Topological transformations modifying the number of links but leaving the sites unchanged. These include ”addition transformations” and ”cutting transformations”. The additions transforms consist in adding links joining sites up to an arbitrary but finite chemical distance from any site, while the cutting transform are defined to be their inverse. The most general transformations consist in a combination of addition and cutting. Notice that even an infinite number of links can be modified with respect to the original graph.
3. Topological rescaling, i.e. topological transformations modifying both links and sites. The most general topological rescaling can be realized through two independent steps. The first one is the partition and consists in dividing the graph $`𝒢`$ in an infinite family of connected subgraphs $`𝒢_\alpha `$, with uniformly bounded number of points. The second one is the substitution and consists in generating a new graph $`𝒢^{}`$ by replacing some or all $`𝒢_\alpha `$ by a different (connected) graph $`𝒮_\alpha `$, whose number of points ranges from 1 to a fixed $`N_{max}`$, and by adding links connecting different $`𝒮_\alpha `$ in such a way that two generic $`𝒮_\alpha `$ and $`𝒮_\beta `$ are connected by some links if and only if $`𝒢_\alpha `$ and $`𝒢_\beta `$ were. The simplest topological rescaling occurs when every $`𝒮_\alpha `$ is composed by just one point. In this case the resulting graph $`𝒢_𝓂`$ is called the minimal structure of the partition $`\{𝒢_\alpha \}`$.
These three very general classes of geometrical transformations (together with even more general ones violating conditions B and C and therefore not discussed here can be applied in all possible sequences to a graph, leading to an overall transformation on coupling strength, number of links and degrees of freedom which does not change its spectral dimension $`\overline{d}`$. We will call such a transformation an $`\mathrm{𝑖𝑠𝑜𝑠𝑝𝑒𝑐𝑡𝑟𝑎𝑙𝑖𝑡𝑦}`$.
Notice that isospectralities include quasi-isometries as a particular case. Indeed, isospectralities include most part of currently used transformations.
As an example, the usual decimation procedure on fractals is a topological rescaling. In particular, for all exactly decimable fractals (such as e.g. Sierpinski gaskets and T-fractals, as discussed in the next sections), the minimal structure of the graph coincides with the graph itself. Again, an isospectrality relates the usual two dimensional square lattice, the hexagonal lattice and the triangular lattice, which therefore all have dimension $`2`$. In other words, isospectralities are the theoretical formalization of the intuitive idea of invariance with respect to bounded scale perturbations and disorder and the isospectrality classes, defined as the classes of graphs related by such transformations, are the practical realization of the apparently abstract concept of non integer dimension.
Now, since most dynamical and thermodynamical properties of generic discrete structures depend only on $`\overline{d}`$, isospectralities provide a very powerful tool to reduce a very complicated geometrical structure to the simplest one having the same $`\overline{d}`$. The latter turns out to be much simpler to study and still presents the same universal properties.
Moreover, not only an isospectrality can be used to reduce and simplify structures and problems. It can also be applied, with the opposite aim, to build complicated structures with controlled dynamical and thermodynamical properties, starting from simple deterministic geometrical models. This is the point of view of spectral dimension engineering, providing a very interesting field of possibilities to polymer physicists and material scientists dealing with non-crystalline materials. In Fig.2 we give explicit examples of isospectral structures obtained applying isospectral transformations (without long range couplings) to the T-fractal and to the square lattice.
On macroscopically inhomogeneous graphs, it can happen that the average value of $`P_{ii}(t)`$ on infinite subgraphs of $`𝒢`$ with positive measure decays with a power law different from (83) . In such cases, it is interesting to look for the maximal (positive-measure) subgraphs having no (positive measure) parts with different power law decay. These are called spectral classes and each is characterized by its own spectral dimension. A theorem rather relevant in physical applications establishes that spectral classes can be separated from each other by cutting a zero-measure set of links, implying the same statistical independence property we discussed for mixed TOA graphs .
## 13 A survey of analytical results on specific networks
Apart from the well known case of regular lattices, where it is completely solved , the random walk problem has been studied analytically only on some specific classes of infinite graphs. In these cases, one usually focuses on the asymptotic properties of random walks autocorrelation functions and on the calculation of the local and average spectral dimension. As we discussed in the previous sections, these are the most important quantities in statistical physics and thermodynamics. On lattices, the random walk problem is solved by using the translation invariance of the structure, and this allows applying powerful mathematical tools, such as the Fourier transform. On general graphs these methods do not apply. Therefore due to the lack of translation invariance, one has to introduce new and alternative techniques, which can be grouped in three main classes: renormalization techniques, combinatorial techniques and mixed techniques. In the next subsections we will review recent and significant results obtained with these techniques.
### 13.1 Renormalization techniques
Renormalization techniques have been successfully applied on deterministic fractals networks, where one can take advantage of the decimation transformations which connect two consecutive generations. In particular, a well studied class of fractals is that of exactly decimable fractals. On these structures, exact renormalization group calculations based on a real space decimation procedure allow obtaining all the relevant random walks quantities.
Let us consider a random walk without traps and sources defined by the jumping probabilities (4) and let us write the master equation for the probability $`P_{0i}`$ of being at site $`i`$ after $`t`$ steps for a random walker starting from an origin site $`0`$ at time $`0`$:
$$P_{0i}(t+1)P_{0i}(t)=\underset{j}{}A_{0j}(\frac{P_{0j}(t)}{z_j}\frac{P_{0j}(t)}{z_i})+\delta _{i0}\delta _{t0}$$
(85)
Equation (85) can be written in terms of the generating function $`\stackrel{~}{P}_{ij}(\lambda )`$ when $`\lambda 1^{}`$ by setting $`\lambda =1ϵ`$, writing:
$$\stackrel{~}{P}_{ij}(ϵ)=\underset{t=0}{\overset{\mathrm{}}{}}(1+ϵ)^tP_{ij}(t)$$
(86)
and taking $`ϵ=0`$:
$$ϵ\stackrel{~}{P}_{0i}(ϵ)=\underset{j}{}A_{0j}(\frac{\stackrel{~}{P}_{0j}(ϵ)}{z_j}\frac{\stackrel{~}{P}_{0j}(ϵ)}{z_i})+\delta _{0i}.$$
(87)
Notice that the system (87) is inhomogeneous and corresponds to a Cauchy problem, which has only one solution. The behaviour of such solution for $`ϵ0`$ is what we need to obtain the local spectral dimension $`\stackrel{~}{d}`$, as defined in (67), through the Tauberian Theorems. On the other hand, to calculate the average spectral dimension $`\overline{d}`$ we will need to average over all starting points the solution of equation (87), strongly modifying its asymptotic behaviour on inhomogeneous graphs, as we will see in the following.
Exactly decimable fractals are a restricted class of self similar structures (i.e. not all self similar structures are exactly decimable) which are geometrically invariant under site decimation. This invariance is explicitly applied in analytical calculations for random walks. A geometrical structure is decimation invariant if it is possible to eliminate a subset of points (and all the bonds connecting these points) obtaining a network with the same geometry of the starting one. ¿From a mathematical point of view this corresponds to the possibility of eliminating by substitution a set of equations from system (85) or (87) obtaining a system which is similar to the initial one after a suitable redefinition of the coupling parameters. Examples of exactly decimable fractals are the Sierpinski Gasket (Fig.3), , the $`T`$fractal, shown in Fig.4 , the branched Koch curves, in Fig.5 . In general, all deterministic finitely-ramified fractals are exactly decimable. Notice that exact decimation is a particular case of isospectrality, as we discussed in previous sections.
Let us consider now the general procedure to decimate the set of equation (85). After eliminating a set of points and substituting the corresponding equation, one finds:
$$ϵϵ^{}(ϵ)a^2ϵ$$
(88)
The presence of the term $`\delta _{i0}`$ in (87) requires a redefinition of the quantities $`\stackrel{~}{P}_{ij}(ϵ)`$ to assure that, even after the decimation, the initial condition will correspond to the probability of being in a fixed site equal to 1. One introduces a new parameter $`c`$ and writes the transformation law for $`\stackrel{~}{P}_{ij}(ϵ)`$ as:
$$\stackrel{~}{P}_{ij}(ϵ)\stackrel{~}{P}_{ij}^{}(ϵ^{})\frac{1}{c}\stackrel{~}{P}_{ij}(ϵ)$$
(89)
From the rescaling of $`ϵ`$ and $`\stackrel{~}{P}_{ij}(ϵ)`$, the local spectral dimension $`\stackrel{~}{d}`$ is obtained by using a suitable expression for $`\stackrel{~}{P}_{00}(ϵ)`$:
$$\stackrel{~}{P}_{00}(ϵ)ϵ^{\stackrel{~}{d}/21}$$
(90)
which holds only for $`\stackrel{~}{d}<2`$. This is always the case for exactly decimable fractals. Using expression (90) one easily finds:
$$\stackrel{~}{d}=2\frac{\mathrm{log}a^2/c}{\mathrm{log}a^2}$$
(91)
As for the average spectral dimension $`\overline{d}`$ , by using the relation between equation (87) and the equation for harmonic oscillations to be discussed later , one has that
$$\overline{d}=\frac{\mathrm{log}r}{\mathrm{log}a}$$
(92)
where $`r`$ is the decimation ratio used in the renormalization procedure. Therefore $`\stackrel{~}{d}=\overline{d}`$ if
$$r=a^2/c$$
(93)
This can be shown to be the case for all exactly decimable fractals, using results obtained for the Gaussian model.
Equation (92) allows calculating the spectral dimension on all exactly decimable fractals, once the decimation procedure is identified, recovering known results.
One of the most studied fractal is without any doubt the Sierpinski gasket and its generalizations. For the simplest case one has $`r=3`$ and $`a=\sqrt{5}`$, leading to:
$$\stackrel{~}{d}=2\frac{\mathrm{log}3}{\mathrm{log}5}$$
(94)
For d-dimensional generalized Sierpinski gaskets, which are built from a $`d`$ dimensional hypertetrahedron of side length $`b`$ filled with $`b`$ layers of smaller hypertetrahedra of unit site length, Hilfer and Blumen have shown that for $`b=2`$
$$\stackrel{~}{d}=2\frac{\mathrm{log}d+1}{\mathrm{log}d+3}$$
(95)
and for $`b=3`$
$$\stackrel{~}{d}=2\frac{\mathrm{log}((d+1)(d+2)/2)}{\mathrm{log}((d+2)(2d^2+9d+19)/(4d+6))}$$
(96)
Due to the self-similarity of the structure, the return probabilities on the Sierpinski gasket show a remarkable effect, which has been pointed out in . Indeed, the coefficients have an oscillatory behaviour, which is given by:
$$P_{00}(t)=t^{\stackrel{~}{d}/2}F(\frac{\mathrm{log}t}{\mathrm{log}5})$$
(97)
where $`F`$ is a periodic $`C^{\mathrm{}}`$-function of period 1 whose Fourier series is given by
$$F(x)=\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}\mathrm{\Gamma }(1\frac{log3}{log5}+\frac{2\pi ki}{log5})^1exp(2\pi kix)$$
(98)
Interestingly, it can be shown that the oscillation of the coefficients disappears in the probability of return on the average.
The renormalization techniques can be applied to all exactly decimable fractals. For example, for the $`T`$fractal , which is a particular case of hierarchical combs , one has $`r=3`$ and $`a=\sqrt{6}`$.
### 13.2 Combinatorial techniques
Renormalization procedures cannot be applied on non self-similar graphs. Therefore one has to develop alternative techniques to study the random walk problem. This is the case of bundled structures , a large class of very interesting graphs used in condensed matter as realistic models for the geometry and dynamics of polymers and other inhomogeneous systems. Given two graphs $``$ and $``$ , not necessarily different, and a site $`F`$ of $``$ , we call bundled graph with base $``$ and fibre $``$ the graph built by joining to each site of $``$ a copy of $``$ in such a way that $`F`$ is the only site $``$ and $``$ have in common (Fig.6). Examples of bundled structures are comb polymers (Fig.7), brush polymers, shown in Fig.8, and many kinds of branched aggregates (Fig.9). For these graphs a purely combinatorial technique allows calculating of the asymptotic properties of the random walk autocorrelation functions.
Let us consider a walker starting from a point belonging to the base and let us restrict ourselves to base graphs with constant coordination number $`z_{}`$. By decomposing the motion of the walker on the fibre and on the base, one can obtain:
$$P_0(t)=\underset{t_{_{}}=0}{\overset{\mathrm{}}{}}\underset{t_1=0}{\overset{\mathrm{}}{}}\mathrm{}\underset{t_{_{}}+1=0}{\overset{\mathrm{}}{}}P_{_{}}(t_{_{}})\left(\frac{z_{_{}}}{z_{_{}}+z__F}\right)^t_{_{}}P_{_{}}^{}(t_1)\mathrm{}P_{_{}}^{}(t_{t_{_{}}+1})\delta _{t,t_{_{}}+_{i=1}^{t_{_{}}+1}t_i}$$
(99)
where $`P_{_{}}^{}`$ refers to a random walk on $``$ with a trap in the starting point of $``$. In terms of the generating functions equation, (99) becomes:
$$\stackrel{~}{P}_0(\lambda )=\underset{t_{_{}}=0}{\overset{\mathrm{}}{}}P_{_{}}(t_{_{}})\left(\frac{\lambda z_{_{}}}{z_{_{}}+z__F}\right)^t_{_{}}\left(\stackrel{~}{P}_{_{}}^{}(\lambda )\right)^{t_{_{}}+1}=\stackrel{~}{P}_{_{}}^{}(\lambda )\stackrel{~}{P}_{_{}}(\lambda ^{})$$
(100)
with
$$\lambda ^{}\frac{\lambda z_{_{}}}{z_{_{}}+z__F}\stackrel{~}{P}_{_{}}^{}(\lambda )$$
(101)
and
$$\stackrel{~}{P}_{_{}}^{}(\lambda )=\left(1\frac{z__F}{z_{_{}}+z__F}\left(1\left(\stackrel{~}{P}_{_{}}(\lambda )\right)^1\right)\right)^1$$
(102)
with $`\stackrel{~}{P}_{_{}}(\lambda )`$ being the generating function of the probability of returning to the starting point $`F`$ on $``$ without the trap. From this relations one obtains the values for the local spectral dimension on general bundled graphs:
$$\stackrel{~}{d}=\{\begin{array}{cc}\stackrel{~}{d}_{_{}}\hfill & \text{if }\stackrel{~}{d}_{_{}}2\hfill \\ & \\ 4\stackrel{~}{d}_{_{}}\hfill & \text{if }\stackrel{~}{d}_{_{}}2\text{ e }\stackrel{~}{d}_{_{}}4\hfill \\ \stackrel{~}{d}_{_{}}+\stackrel{~}{d}_{_{}}\frac{\stackrel{~}{d}_{_{}}\stackrel{~}{d}_{_{}}}{2}\hfill & \text{if }\stackrel{~}{d}_{_{}}2\text{ e }\stackrel{~}{d}_{_{}}4\hfill \end{array}$$
(103)
where $`\stackrel{~}{d}_{_{}}`$ and $`\stackrel{~}{d}_{_{}}`$ are the local spectral dimension of the base and of the fibre. If the coordination number of the base is not constant, it can be shown that this amounts to introduce waiting probabilities on the points connecting the fibre and the base, which, as shown in the previous section, does not change the value of the spectral dimension.
As for the average spectral dimension, it is easy to show that if the fibre is an infinite graph, the average spectral dimension of the whole graph is the spectral dimension of the fibre. On the other hand, if the fibre is a finite graph, the average spectral dimension coincides with that of the base.
¿From equations (100), (102) one also obtains the asymptotic laws for the probability of returning to the starting point, which on these structures can contain logarithmic corrections. Indeed, writing
$$P_0(t)\underset{i=0}{\overset{\mathrm{}}{}}{}_{}{}^{i}\mathrm{ln}_{}^{\beta (i)}(t)$$
(104)
and setting
$$m=\mathrm{min}\{i0|\beta (i)1\}$$
(105)
and
$$I\left(\stackrel{~}{d}/2\right)=\{\begin{array}{cc}1\hfill & \text{if }\stackrel{~}{d}/2\text{ is an integer}\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}$$
(106)
one has:
a) if $`\stackrel{~}{d}_{}<4`$ and $`\stackrel{~}{d}_{}<2`$
$$\beta (i)=\{\begin{array}{cc}1\hfill & \text{for }0<i<m\hfill \\ & \\ \left(1\frac{\stackrel{~}{d}_{_{}}}{2}\right)\left[\beta _{_{}}(m_{_{}})+I\left(\frac{\stackrel{~}{d}_{_{}}}{2}\right)\right]I\left(\frac{\stackrel{~}{d}_{_{}}}{2}\right)\hfill & \text{for }i=mm_{_{}}\hfill \\ & \\ \left(1\frac{\stackrel{~}{d}_{_{}}}{2}\right)\beta _{_{}}(i)+\theta \left(im_{_{}}m_{_{}}\right)\beta _{_{}}\left(im_{_{}}\right)+\hfill & \\ & \\ +\delta _{_{im_{_{}},m_{_{}}}}I\left(\frac{\stackrel{~}{d}}{2}\right)\delta _{_{i,m}}I\left(\frac{\stackrel{~}{d}}{2}\right)\hfill & \text{otherwise}\hfill \end{array}$$
(107)
where $`m_{}`$ and $`m_{}`$ refers to the base and to the fibre respectively while $`m`$ refers to the whole graph and is determined by:
$$m=m_{_{}}+\delta _{\stackrel{~}{d}_{_{}},2}m_{_{}}$$
(108)
b) if $`\stackrel{~}{d}_{}>4`$ and $`\stackrel{~}{d}_{}<2`$
$$\beta (i)=\{\begin{array}{cc}\beta _{_{}}(i)2\delta _{i,m_{_{}}}I(\stackrel{~}{d}_{_{}}/2)\hfill & \text{for }im_{_{}}\hfill \\ \beta _{_{}}(i)\hfill & \text{for }0<i<m_{_{}}\hfill \end{array}$$
(109)
c) if $`\stackrel{~}{d}_{}=4`$ and $`\stackrel{~}{d}_{}<2`$ and $`d_{}>1`$
$`\beta (i)`$ has to be determined as in a).
d) if $`\stackrel{~}{d}_{}=4`$ and $`\stackrel{~}{d}_{}<2`$ e $`m_{}<1`$
$`\beta (i)`$ has to be determined as in b).
e) if $`\stackrel{~}{d}_{}>2`$
$$\beta (i)=\beta _{_{}}(i)i$$
(110)
The case $`\stackrel{~}{d}_{}=2`$ has to be treated separately, as the case $`\stackrel{~}{d}_{}<2`$ if the fibre is a recurrent graph or as the case $`\stackrel{~}{d}_{}>2`$ if it is transient.
Another interesting way of combining together two graphs to obtain a more complex structure is the Cartesian product. The Cartesian product of two graphs $`X,Y`$ has vertex set $`X\times Y`$, and two pairs $`xy`$, $`x^{}y^{}`$ are adjacent if $`xx^{}`$ and $`y=y^{}`$, or $`x=x^{}`$ and $`yy^{}`$. An example of an interesting Cartesian product is that of the Toblerone graph , shown in Fig.10, which is obtained from the product of a line with a Sierpinski gasket. Using combinatorial techniques analogous to those presented for bundled graphs, it can be show that the local and the average spectral dimension on the whole graph are the sum of the corresponding dimensions of the two initial graphs .
### 13.3 Mixed techniques
The random walk problem on some very interesting cases of graphs cannot be studied simply by one of the above cited techniques and it requires instead a ”mixed” use of the two, which gives rise to very interesting phenomena. Indeed, the first example of a difference between the local and the average spectral dimension, the ”dynamical dimension splitting”, was observed on the quasi self-similar graphs $`NT_D`$, where the asymptotic properties of the random walk were found by a mixed technique .
The fractal trees known as $`NT_D`$ can be recursively defined as follows: an origin point $`O`$ (Fig.11) is connected to a point $`1`$ by a link, of unitary length; from $`1`$, the tree splits in $`k`$ branches of length $`2`$ (i.e. consisting of two consecutive links); the ends of these branches split in $`k`$ branches of length $`4`$ and so on; each endpoint of a branch of length $`2^n`$ splits in $`k`$ branches of length $`2^{n+1}`$.
As one can easily verify, $`NT_D`$ are not exactly decimable and therefore the simple decimation techniques cite above cannot be applied. Indeed, after a simple decimation starting from the origin $`O`$, one obtains $`k`$ copies of the original structure joined together in a point instead of the same $`NT_D`$. However, $`NT_D`$ are invariant under a more complex transformation $`T=DC`$, consisting of the product of a cutting transform $`C`$ and a decimation $`D`$, that can be described as follows. Let us cut the log of the tree in point $`1`$ and separate the $`k`$ branches (cutting transform). Now, each branch can be obtained from the initial $`NT_D`$ by a dilatation with a factor 2. Eliminating all branches but one and decimating it (decimation transform), one obtains the original $`NT_D`$.
The $`T`$ transform can now be used to solve the random walks problem. Let us sketch the main points of the calculation. The cutting transform gives a relation between random walks on the whole tree and random walks on one of its branches; more precisely one relates $`\stackrel{~}{P}_O^{tree}(\lambda )`$, the generating function of the probability of returning to point $`O`$ after a random walk on the $`NT_D`$ tree, and $`\stackrel{~}{P}_1^{branch}(\lambda )`$, the generating function of the probability of returning to the starting point $`1`$ after a random walk on one of the branches. This relation is given by :
$$\stackrel{~}{P}_O^{tree}(\lambda )=\frac{\stackrel{~}{P}_1^{branch}(\lambda )+k}{2\lambda \stackrel{~}{P}_1^{branch}(\lambda )+k}$$
(111)
Now, the decimation transformation is performed using a time-rescaling technique. Indeed, the motion of the random walker on the branch considered only after an even number of steps can be exactly mapped in the motion of a random walker on the tree after the introduction of a staying probability $`p_{ii}=1/2`$ in every site $`i`$. This equivalence can be translated in terms of generating functions through the substitutions:
$$\stackrel{~}{P_O}(\lambda )\frac{\lambda }{2\lambda }\stackrel{~}{P_O}\left(\frac{2}{2\lambda }\right)$$
(112)
$$\lambda \lambda ^2$$
(113)
Equations (112) and (113) can be used to rewrite (111) as:
$$\stackrel{~}{P_O}^{tree}(\lambda )=\frac{\frac{2}{2\lambda ^2}\stackrel{~}{P_O}^{tree}\left(\frac{\lambda ^2}{2\lambda ^2}\right)+k}{(1\lambda ^2)\frac{2}{2\lambda ^2}\stackrel{~}{P_O}^{tree}\left(\frac{\lambda ^2}{2\lambda ^2}\right)+k}$$
(114)
Choosing a suitable power law expression for the singularity of $`P_{OO}^{tree}(\lambda )`$ for $`\lambda 1^{}`$ we obtain:
$$\stackrel{~}{d}=1+\frac{\mathrm{log}k}{\mathrm{log}2}$$
(115)
To obtain the average spectral dimension, one has to calculate the normalized trace of the return probability $`\stackrel{~}{P_O}^{tree}(\lambda )`$. It can be shown that $`\overline{d}=1`$ and this can be intuitively understood noting that the topology of $`NT_D`$ is dominated by linear chains which become longer and longer in the outer branches . Therefore, while $`NT_D`$ are locally transient if the ramification $`k`$ is greater that $`2`$, they are always recurrent on the average. This result has been generalized. Indeed, recently it has been shown that all physical trees, satisfying conditions A, B and C are recurrent on the average .
The cutting decimation transform can be applied to a large class of non-exactly decimable fractals which correspond to more general cases of the $`NT_D`$. This are built with the same recurrence procedure as the $`NT_D`$ and we shall call them $`2^mNT_D`$, $`nNT_D`$ and $`ppolygonNT_D`$, depending on the growth rules for the branches .
The first generalization is that of $`2^mNT_D`$. The $`2^mNT_D`$ are infinite fractal trees that can be recursively built using the same recipe as for $`NT_d`$ but, from point $`1`$, the log splits in $`k`$ branches of length $`2^m`$ (i.e. made of $`2^m`$ consecutive links) which, in turn, split in $`k`$ branches of length $`2^{2m}`$ and so on in such a way that each branch of length $`2^{nm}`$ splits in $`k`$ branches of length $`2^{(n+1)m}`$. The case $`m=1`$ corresponds to the usual $`NT_D`$ previously studied. If $`m>1`$ the time rescaling procedure which led to (112) and (113) must be iterated $`m`$ times obtaining:
$$\stackrel{~}{P_O}^{tree}(\lambda )=\frac{\left(\underset{i=1}{\overset{m}{}}\frac{2}{2\lambda _i^2}\right)\stackrel{~}{P_O}^{tree}(\lambda _{i+1})+k}{(1\lambda ^2)\left(_{i=1}^m\frac{2}{2\lambda _i^2}\right)\stackrel{~}{P_O}^{tree}(\lambda _{i+1})+k}$$
(116)
with
$$\lambda _i=\{\begin{array}{cc}\lambda \hfill & i=1\hfill \\ & \\ \frac{\lambda _{i1}^2}{2\lambda _{i1}^2}\hfill & i>1\hfill \end{array}$$
$`i`$ being the iteration step. This gives, with the same steps as for $`m=1`$:
$$\stackrel{~}{d}_{2^m}=1+\frac{\mathrm{ln}k}{\mathrm{ln}2^m}$$
(117)
which represent the generalization of the result obtained for $`m=1`$.
The previous results can be extended to $`nNT_D`$, where now $`n`$ is an integer and not necessarily a power of 2, and to $`ppolygonNT_D`$, where the branches of $`NT_D`$ are replaced by $`p`$-vertices regular polygons (Fig.12).
Let us consider $`nNT_D`$ first. While relation (111) for the cutting transform still holds, the exact time-rescaling procedure can not be applied to the branch of generic length $`n`$. However even in this case it is possible to obtain an asymptotic recursion relation applying the Renormalization Group techniques usually implemented on exactly decimable fractals. Although this procedure cannot give an exact equation for $`\stackrel{~}{P_O}^{tree}(\lambda )`$ as in the previous case, nevertheless it can be used to obtain the exact value of $`\stackrel{~}{d}`$ via an asymptotic estimation.
Indeed, in this case the branch of the $`nNT_D`$ can be considered as a tree with a dilatation factor equal to $`n`$. The log of this tree can be reduced to a unitary length log after the suppression of the $`n2`$ sites between the edges and introducing a new link connecting the edges. The same operation can be repeated for branches of every length suppressing the inner $`n2`$ consecutive sites in every sequence of $`n`$ sites and introducing a new link between the surviving points. The final structure is equal to the original tree and the generating function $`\stackrel{~}{P_1}^{branch}(\lambda )`$ becomes $`\stackrel{~}{P_1}^{{}_{}{}^{}branch}(\lambda ^{})`$ where:
$$\lambda ^{}=n^2\lambda $$
(118)
$$\stackrel{~}{P_1}^{{}_{}{}^{}branch}(\lambda ^{})=\frac{1}{n}\stackrel{~}{P_1}^{branch}(\lambda )$$
(119)
Now $`\stackrel{~}{P_1}^{{}_{}{}^{}branch}(\lambda ^{})`$ coincides with $`\stackrel{~}{P_O}^{tree}(\lambda ^{})`$ since our branch has been transformed into a tree and (111) can be rewritten as:
$$\stackrel{~}{P_O}^{tree}(\lambda )=\frac{n\stackrel{~}{P_O}^{tree}(n^2\lambda )+k}{2\lambda n\stackrel{~}{P_O}^{tree}(n^2\lambda )+k}$$
(120)
Using the procedure described in the previous section for $`2^mNT_D`$, from (120) it follows that for an $`nNT_D`$ the spectral dimension is given by:
$$\stackrel{~}{d}_n=1+\frac{\mathrm{ln}k}{\mathrm{ln}n}$$
(121)
An analogous technique can be used for $`ppolygonNT_D`$ (Fig.12). The log polygon has now $`p`$ faces of unitary length; from each of $`p1`$ of its vertices $`k`$ polygons depart, whose faces have length $`n`$ and so on. These structures, though similar to $`NT_D`$ are no longer loopless structures nor necessarily bipartite graphs (e.g. the $`3polygon`$ tree). The Cutting-Decimation transform can be applied to $`ppolygonNT_D`$ as in the case of $`NT_D`$ with the same substitutions (118) and (119). Indeed, even if (111) does not hold in this case, a new relation between the generating functions of the tree and that of one of its branches can be obtained using bundled structures theory discussed above . Let us consider a $`ppolygonNT_D`$ and suppose to attach $`k`$ branches also in the free vertex of the log (the root of the tree): we obtain a bundled structure having the log polygon as base and the graph made of $`k`$ branches as fibre. Since for a $`p`$-polygon:
$$\stackrel{~}{P_O}(\lambda )\frac{1}{p(1\lambda )}$$
(122)
as $`\lambda 1`$, we obtain for our bundled structure:
$$\stackrel{~}{P_O}^{b.s.}(\lambda )=\frac{1}{1\frac{k}{k+1}\stackrel{~}{F_1}^{branch}(\lambda )}\frac{1}{p}\left(1\frac{\lambda }{k+1}\frac{1}{1\frac{k}{k+1}\stackrel{~}{F_1}^{branch}(\lambda )}\right)^1$$
(123)
where $`\stackrel{~}{P_O}^{b.s.}(\lambda )`$ is the generating function of the probability of returning to point $`O`$ (one of the vertices of the log polygon) after a random walk on the bundled structure and $`\stackrel{~}{F_1}^{branch}(\lambda )`$ is the generating function of the probability of returning for the first time to the point of connection with the base after a random walk on the fibre. Now,
$$\stackrel{~}{F_O}^{b.s.}(\lambda )=\frac{k}{k+1}\stackrel{~}{F_1}^{branch}(\lambda )+\frac{1}{k+1}\stackrel{~}{F_O}^{tree}(\lambda )$$
(124)
where $`F_{O}^{}{}_{}{}^{tree}(\lambda )`$ refers to the $`ppolygonNT_D`$. From (123), and (124) and using the usual relation between $`\stackrel{~}{F_1}^{branch}(\lambda )`$ and $`\stackrel{~}{P_1}^{branch}(\lambda )`$, a relation between $`\stackrel{~}{P_O}^{tree}(\lambda )`$ and $`\stackrel{~}{P_1}^{branch}(\lambda )`$ follows, which represents the cutting transformation. It is now possible to perform the Cutting-Decimation transform to $`ppolygonNT_D`$ and get:
$$\stackrel{~}{d}_p=1+\frac{\mathrm{ln}k(p1)}{\mathrm{ln}n}$$
(125)
In the same way we can calculate the spectral dimension of an $`NT_D`$ built with $`d`$dimensional simplexes instead of $`ppolygons`$. A $`d`$-dimensional simplex is a complete graph of $`d+1`$ points i.e. a graph where each point is nearest neighbour of all other points. The $`2`$dimensional case is the triangle, the $`3`$dimensional one is the tetrahedron and so on. Since for $`d`$simplex $`\stackrel{~}{P_O}(\lambda )1/(d+1)(1\lambda )`$ the spectral dimension is:
$$\stackrel{~}{d}_d=1+\frac{\mathrm{ln}kd}{\mathrm{ln}n}$$
(126)
## 14 Relation with other physical problems
As we have shown in previous chapters, the random walk problem is strictly related to graph topology. Indeed, the main physical quantities are simple functions of the adjacency matrix $`A`$, which algebraically describes the graph structure. Now, the Hamiltonians of a series of fundamental statistical models are linear in $`A`$, therefore even their behaviour is deeply influenced by topology and it can be expressed in terms of random walks functions. Due to this reason, the main concepts and parameters characterizing random walks, such as recurrence and transience, as well as the spectral dimension, determine also the properties of these models, which have very different physical origins. This provides a very powerful tool to investigate and classify geometrically disordered and inhomogeneous systems, where the usual techniques and ideas developed for lattices do not apply.
### 14.1 The oscillating network
Probably, the physical model whose connection with random walks has been most extensively explored is the so-called oscillating network.
The harmonic oscillations of a generic network of masses $`m`$ linked by springs of elastic constant $`K`$ can be studied by writing the equations of motion of the displacements $`x_i`$ of each mass from its equilibrium position:
$$m\frac{d^2}{dt}x_i=K\underset{j}{}A_{ij}(x_ix_j)=K\underset{j}{}\mathrm{\Delta }_{ij}x_j$$
(127)
which after Fourier transforming with respect to the time reads:
$$\frac{\omega ^2}{\omega _0^2}\stackrel{~}{x}_i=\underset{j}{}\mathrm{\Delta }_{ij}\stackrel{~}{x}_j$$
(128)
where $`\omega _0^2K/m`$. In other words, the determination of the normal modes and of the normal frequencies of the oscillating network reduces to the diagonalization of the Laplacian operator $`\delta `$.
Noticing that $`\mathrm{\Delta }=Z(\mathrm{𝟏}P)`$, where 1 the identity matrix and $`P`$ is given by (4), it is not difficult to establish mathematical correspondences with random walks. In particular, using the universality properties discussed in the previous sections, one can show a fundamental result concerning the density $`\rho (\omega )`$ of normal modes at low frequencies:
$$\rho (\omega )\omega ^{\overline{d}1}\mathrm{for}\omega 0$$
(129)
This basic connection between random walks and harmonic oscillations was first introduced by Alexander and Orbach in 1982 for the case of fractals. Notice that at that time the splitting between local and average spectral dimension on inhomogeneous structures was not yet known and the exponent describing the scaling of the density of states at low frequencies was simply called spectral dimension, since it was related to the vibrational spectrum. Due to the already mentioned universality properties, the above result hold for the very general case where oscillating masses and elastic constants may have different values on different sites and links, provided they are bounded by positive numbers. More precisely, considering the equations of motions
$$m_i\frac{d^2}{dt}x_i=K\underset{j}{}J_{ij}(x_ix_j)=K\underset{j}{}L_{ij}x_j$$
(130)
for the same graph of (127), if (49) holds together with
$$m_{min},m_{max}>0|m_{min}m_im_{max}i$$
(131)
then the asymptotic behaviour of the density of vibrational states is still given by (49).
¿From all the above properties, it follows that also the spectrum of the Laplacian operator $`L`$ depends on $`\overline{d}`$: indeed it can be shown that the spectral density $`\rho (l)`$ of $`L`$ at low eigenvalues behaves as $`\rho (l)l^{\overline{d}/21}`$.
The average spectral dimension is crucial in determining the behaviour of the oscillating network in equilibrium with a thermal bath at temperature $`T`$. Considering the Hamiltonian of the system given by (129)
$$H=\underset{i}{}\frac{p_i^2}{2m_i}+\frac{1}{2}m\omega _0^2\underset{ij}{}J_{ij}x_ix_j$$
(132)
and calculating the thermodynamic averages with respect to the Gibbs weight $`\mathrm{exp}(H/kT)`$, where $`k`$ denotes the Boltzmann constant, one can show that, for positive $`T`$,
$$\overline{<x^2>}=\mathrm{}\mathrm{for}\overline{d}2$$
(133)
while
$$\overline{<x^2>}<\mathrm{}\mathrm{for}\overline{d}>2.$$
(134)
This is the generalization to graphs of the fundamental Peierls result about the thermodynamic instability of oscillating crystals in low dimensions. In other words, for an infinite oscillating network with $`\overline{d}2`$ in equilibrium with a thermal bath, the mean square displacement of masses from their equilibrium positions would diverge.
### 14.2 The Gaussian model
The Gaussian model is the simplest statistical model used to study magnetic systems on lattices. Even if it is not realistic, its properties are fundamental to understand more complex and phenomenologically significant models. In field theory it is also known as ”free scalar field”. The Gaussian model on $`𝒢`$ is defined by the Hamiltonian:
$$H=\frac{1}{2}\underset{ij}{}\varphi _i(JL_{ij}+m_i^2\delta _{ij})\varphi _jh\underset{i}{}\varphi _i$$
(135)
where $`\varphi _i`$ is a real field, $`J>0`$ a ferromagnetic coupling, $`h`$ an external magnetic field and $`m_i^2=\alpha _im^2`$, with $`1/K<\alpha _i<K`$ for some positive $`K`$. Its specific free energy $`f_G`$ is given by
$$f_G(J,m_i^2,h)=\underset{N\mathrm{}}{lim}\frac{1}{N}F=\underset{N\mathrm{}}{lim}\frac{1}{N}\mathrm{log}Z$$
(136)
where $`Z`$ is the partition function calculated according to the Boltzmann weight $`\mathrm{exp}(H)`$. The spectral dimension is related to the singular part of $`f_G`$ for $`h=0`$ and $`m^20`$ by:
$$Sing(f)m^{\overline{d}}.$$
(137)
The covariance of this Gaussian process reads
$$\varphi _i\varphi _jC_{ij}(m^2)=(\mathrm{\Delta }+m^2\eta )_{ij}^1$$
(138)
and hence it satisfies by definition the Schwinger–Dyson (SD) equation
$$(J_i+m^2\eta _i)C_{ij}(m^2)\underset{kG}{}J_{ik}C_{kj}(m^2)=\delta _{ij}$$
(139)
Setting
$$C_{ij}=\frac{(1W)_{ij}^1}{J_i+m^2\eta _i},W_{ij}=\frac{J_{ij}}{J_j+m^2\eta _j}$$
(140)
one obtains the standard connection with the random walk (RW) over $`𝒢`$ :
$$(1W)_{ij}^1=\underset{t=0}{\overset{\mathrm{}}{}}(W^t)_{ij}=\underset{\gamma :ij}{}W[\gamma ]$$
(141)
where the last sum runs over all paths from $`j`$ to $`i`$, each weighted by the product along the path of the one–step probabilities in $`W`$:
$$\gamma =(i,k_{t1},\mathrm{},k_2,k_1,j)W[\gamma ]=W_{ik_{t1}}W_{k_{t1}k_{t2}},\mathrm{},W_{k_2k_1}W_{k_1j}$$
(142)
Notice that, as long as $`m>0`$, we have $`_i(W^t)_{ij}<1`$ for any $`t`$, namely the walker has a non-zero death probability. This implies that $`C_{ij}`$ is a smooth functions of $`m^2`$ for $`mϵ>0`$. In the limit $`m0`$ the walker never dies and the sum over paths in eq. (141) is dominated by the infinitely long paths which sample the large scale structure of the entire graph (“large scale” refers here to the metric induced by the chemical distance alone). This typically reflects itself into a singularity of $`C_{ij}`$ at $`m=0`$ whose nature does not depend on the detailed form of $`J_{ij}`$ or $`\eta _i`$, as long these stay uniformly positive and bounded.
Of particular importance is the leading singular infrared behaviour, as $`m^20`$, of the average $`[C(m^2)]_G`$ of $`C_{ii}(m^2)`$, which is a positive definite quantity, over all points $`i`$ of the graph $`𝒢`$, which we may write in general as
$$\mathrm{Sing}[C(m^2)]_Gc(m^2)^{\overline{d}/21}$$
(143)
### 14.3 Spherical model and $`O(n)`$ models
The spherical model is again a magnetic model with no direct connection to phenomenology. Nevertheless, is a little more complex than the Gaussian one and, most important, it exhibits phase transitions at finite temperature for $`\overline{d}>2`$. Moreover, its critical exponents can be exactly determined and they turn out to be simple functions of $`\overline{d}`$, pointing out the crucial role of the average spectral dimension in phase transitions and critical phenomena. The spherical model can be defined on a generic graph through the Hamiltonian (135) with the generalized spherical constraint $`_iz_i\varphi _i^2=N`$. We assume the coordination numbers to be bound: $`1z_iz_{\mathrm{max}}`$. Its free energy and correlation functions can be expressed in terms of the Gaussian ones. Then the critical behaviour is obtained from the infrared singularities of the latter, i.e. in terms of the long time behaviour of random walks. The results concerning the critical exponents are summarized in the following table, where $`T_c=0`$ for $`\overline{d}2`$:
The so called $`O(n)`$ models are defined, for positive integer $`n`$, by the Boltzmann weight $`\mathrm{exp}(\beta H_n)`$, where
$$H_n[𝐒]=\frac{1}{2}\underset{<ij>}{}J_{ij}(𝐒_i𝐒_j)^2$$
(144)
the sum extends to all links of a certain graph $`𝒢`$, $`J_{ij}>0`$ are ferromagnetic interactions, which may vary from link to link, and $`𝐒_i`$ is an $`n`$dimensional vector of fixed length normalized by $`𝐒_i𝐒_i=n`$. They represent more realistic magnetic models, but their exact solution is in general impossible. However, a series of complex but powerful inequalities, relating their correlation functions to the random walks generating functions, allow proving some very general results shedding light on the complicated phenomena concerning phase transitions on graphs. In particular it has been proven that:
* they cannot have phase transitions at $`T>0`$ if $`𝒢`$ is recurrent on the average;
* they exhibit phase transitions at $`T>0`$ if $`𝒢`$ is transient on the average;
* for $`n\mathrm{}`$ their critical exponents tend to the spherical ones. |
warning/0507/cs0507048.html | ar5iv | text | # Redundancy in Logic III: Non-Mononotonic Reasoning
## 1 Introduction
In this paper, we study the problem of whether a circumscriptive \[McC80\] or default \[Rei80\] theory is redundant, that is, it contains unnecessary parts. Formally, a theory is redundant if it is equivalent to one of its proper subsets; a part is redundant in a theory if the theory is not semantically changed by the removal of the part. The redundancy of propositional theories in CNF, 2CNF, and Horn form has already been analyzed in other papers \[Lib05b, Lib05c\] where motivations are also given. Other problems related to redundancy have been considered by various authors \[Gin88, SS97, MS72, Mai80, ADS86, HK93, HW97, LM00, Uma98, GF93, BZ05, PW88, FKS02, Bru03\].
Circumscription and default logic are two forms of non-monotonic reasoning, as opposite to classical logic, which is monotonic. A logic is monotonic if the consequences of a set of formulae monotonically non-decrease with the set. In other words, all formulae that are entailed by a set are also entailed by every superset of it. Circumscription and default logic do not have this property, and are therefore non-monotonic.
The difference between monotonic and non-monotonic logic is important in the study of redundancy. In propositional logic, as in all forms of monotonic logic, if a set does not entail a formula, the same is true for all of its subsets. As a result, if $`\mathrm{\Pi }`$ is a set of clauses, $`\gamma `$ is one of its clauses, and $`\mathrm{\Pi }\backslash \{\gamma \}`$ is not equivalent to $`\mathrm{\Pi }`$, no subset of $`\mathrm{\Pi }\backslash \{\gamma \}`$ is equivalent to $`\mathrm{\Pi }`$. In the other way around, if all clauses of a set of clauses $`\mathrm{\Pi }`$ are irredundant, then $`\mathrm{\Pi }`$ is irredundant. We call this property local redundancy. The converse of this property is obviously true: if a clause is redundant in a formula, the formula is redundant because it is equivalent to the subset composed of all its clauses but the redundant one.
Local redundancy holds for all monotonic logic. In nonmonotonic logics, removing a clause from a formula might result in a decrease of the set of consequences, which can however grow to the original one when another clause is further removed. However, some nonmonotonic logics have the local redundancy property. We prove that local redundancy holds for circumscriptive entailment and for the redundancy of the background theory in default logic when all defaults are categorical (prerequisite-free) and normal. In the general case, default logic does not have the local redundancy property.
Since redundancy is defined in terms of equivalence (namely, equivalence of a formula with a proper subset of it), it is affected by the kind of equivalence used. In particular, equivalence can be defined in two ways for default logic: equality of extensions and equality of consequences. This lead to two different definitions of redundancy in default logic.
Regarding the complexity results, we show that checking whether a clause is redundant in a formula according to circumscriptive inference is $`\mathrm{\Pi }_2^p`$-complete. For default logic, we mainly considered redundancy in the background theory according to Reiter semantic using both kinds of equivalence, but we also considered justified default logic \[Luk88\], constrained default logic \[Sch92, DSJ94\], and rational default logic \[MT95\]. The results are as follows: the redundancy of a clause in the background theory is $`\mathrm{\Pi }_2^p`$-complete and $`\mathrm{\Pi }_3^p`$-complete for equivalence based on extensions and consequences, respectively. The problems of redundancy of the background theory are $`\mathrm{\Sigma }_3^p`$-complete and $`\mathrm{\Sigma }_4^p`$-complete, respectively. The proofs of the latter two results are of some interest, as they are done by first showing that the problems are $`\mathrm{\Pi }_2^p`$-complete and $`\mathrm{\Pi }_3^p`$-complete, respectively, and then showing that such complexity results can be raised of one level in the polynomial hierarchy. This technique allows for a proof of hardness for a class such as $`\mathrm{\Sigma }_4^p`$ without involving complicated QBFs such as $`WXYZ.F`$.
We also considered the redundancy of defaults in a default theory. We show that these problems are at least as hard as the corresponding problems for the redundancy of the background theory for Reiter and justified default logics.
## 2 Preliminaries
If $`\mathrm{\Pi }`$ and $`\mathrm{\Gamma }`$ are sets, $`\mathrm{\Pi }\backslash \mathrm{\Gamma }`$ denotes the set of elements that are in $`\mathrm{\Pi }`$ but not in $`\mathrm{\Gamma }`$. This operator is often called set subtraction, because the elements of $`\mathrm{\Gamma }`$ are “subtracted” from $`\mathrm{\Pi }`$. An alternative definition of this operator is: $`\mathrm{\Pi }\backslash \mathrm{\Gamma }=\mathrm{\Pi }\overline{\mathrm{\Gamma }}`$, where $`\overline{\mathrm{\Gamma }}`$ is the complement of $`\mathrm{\Gamma }`$.
All formulae considered in this paper are propositional and finite Boolean formulae over a finite alphabet. We typically use formulae in CNF, that is, sets of clauses. We simply refer to sets of clauses as formulae. We assume that no clause is tautological (e.g., $`x\neg x`$): formulae containing tautological clauses can be simplified in linear time. By $`Var(\mathrm{\Pi })`$ we mean the set of variables mentioned in the formula $`\mathrm{\Pi }`$.
In some places, we use the notation $`\neg \gamma `$, where $`\gamma `$ is a clause, to denote the formula $`\{\neg l|l\gamma \}`$. Note that $`\gamma `$ is a clause, while both $`\{\gamma \}`$ and $`\neg \gamma `$ are formulae (sets of clauses). A clause is positive if and only if it contains only positive literals.
A propositional model is an assignment from a set of propositional variables to the set $`\{\mathrm{𝗍𝗋𝗎𝖾},\mathrm{𝖿𝖺𝗅𝗌𝖾}\}`$. We denote a model by the set of variables it assigns to $`\mathrm{𝗍𝗋𝗎𝖾}`$. We use the notation $`Mod(\mathrm{\Pi })`$ to denote the set of models of a formula $`\mathrm{\Pi }`$. We sometimes use models as formulae, e.g., $`\mathrm{\Pi }\omega `$ where $`\mathrm{\Pi }`$ is a formula and $`\omega `$ is a model. In the context where a formula is expected, a model $`\omega `$ represents the formula $`\{x|x\omega \}\{\neg x|x\omega \}`$. If $`\mathrm{\Pi }`$ is a formula and $`\omega _X`$ is a model over the variables $`X`$, we denote by $`\mathrm{\Pi }|_{\omega _X}`$ the formula obtained by replacing each variable of $`X`$ with its value assigned by $`\omega _X`$ in $`\mathrm{\Pi }`$.
A quasi-order is a reflexive and transitive relation (formally, a quasi-order is a pair composed of a set and a reflexive and transitive relation on this set, but the set will be implicit in this paper). The set containment relation $``$ among models is a quasi-order. According to our definition, a model is a set of positive literals; as a result $`MM^{}`$ holds if and only if $`M`$ assigns to false all variables that $`M^{}`$ assign to false.
A clause is (classically) redundant in a CNF formula $`\mathrm{\Pi }`$ if $`\mathrm{\Pi }\backslash \{\gamma \}\gamma `$. A CNF formula is (classically) redundant if it is equivalent to one of its proper subsets. Propositional logic has the local redundancy property: a formula is redundant if and only if it contains a redundant clause. The local redundancy property is defined as follows.
###### Definition 1 (Local redundancy)
A logic has the local redundancy property if, in this logic, a theory is redundant only if it contains a redundant clause.
Propositional logic has the local redundancy property. This is however not true for all logics.
## 3 Circumscription
Circumscriptive inference is based on the minimal models of a theory, i.e., the models that assign the maximum quantity of literals to false. Formally, we define the set of minimal models as follows.
###### Definition 2
The set of minimal models of a propositional formula $`\mathrm{\Pi }`$, denoted by $`\mathrm{CIRC}(\mathrm{\Pi })`$, is defined as follows.
$$\mathrm{CIRC}(\mathrm{\Pi })=\underset{}{\mathrm{min}}(Mod(\mathrm{\Pi }))$$
We define $`\mathrm{CIRC}(\mathrm{\Pi })`$ to be a set of models instead of a formula, although the latter is more common in the literature. Circumscriptive entailment is defined like classical entailment but only minimal models are taken into account.
###### Definition 3
The circumscriptive inference $`_M`$ is defined by: $`\mathrm{\Pi }_M\mathrm{\Gamma }`$ if and only if $`\mathrm{\Gamma }`$ is satisfied by all minimal models of $`\mathrm{\Pi }`$:
$$\mathrm{\Pi }_M\mathrm{\Gamma }\text{ if and only if }\mathrm{CIRC}(\mathrm{\Pi })Mod(\mathrm{\Gamma })$$
Equivalence in propositional logic can be defined in two equivalent ways: either by equality of the models or by equality of the sets of entailed formulae. These two definitions of equivalence coincide for circumscriptive inference as well. We define $`_M`$ as follows: $`\mathrm{\Pi }_M\mathrm{\Gamma }`$ if and only if $`\mathrm{CIRC}(\mathrm{\Pi })=\mathrm{CIRC}(\mathrm{\Gamma })`$. Redundancy of a clause is defined as follows.
###### Definition 4
A clause $`\gamma \mathrm{\Pi }`$ is $`\mathrm{CIRC}`$-redundant in the CNF formula $`\mathrm{\Pi }`$ if and only if $`\mathrm{\Pi }\backslash \{\gamma \}_M\mathrm{\Pi }`$.
A formula is redundant if some of its clauses can be removed without changing its semantics.
###### Definition 5
A formula is $`\mathrm{CIRC}`$-redundant if it is $`_M`$-equivalent to one of its proper subsets.
A formula is therefore redundant if some clauses can be removed from it while preserving equivalence. In the next section we show that a formula is $`\mathrm{CIRC}`$-redundant if and only if it contains a $`\mathrm{CIRC}`$-redundant clause, that is, circumscription has the local redundancy property.
### 3.1 Clause-Redundancy vs. Formula-Redundancy
Propositional logic has the local redundancy property. Showing why is interesting for comparison with logics not allowing the same proof to be used. If $`\mathrm{\Pi }`$ does not contain a redundant clause, then $`\mathrm{\Pi }\backslash \{\gamma \}\overline{)}\mathrm{\Pi }`$ for any clause $`\gamma \mathrm{\Pi }`$. Therefore, $`Mod(\mathrm{\Pi })Mod(\mathrm{\Pi }\backslash \{\gamma \})`$. Since $`\mathrm{\Pi }\backslash \{\gamma \}`$ is a subset of $`\mathrm{\Pi }`$, we have $`Mod(\mathrm{\Pi })Mod(\mathrm{\Pi }\backslash \{\gamma \})`$ in general and $`Mod(\mathrm{\Pi })Mod(\mathrm{\Pi }\backslash \{\gamma \})`$ in this case. If $`\mathrm{\Pi }^{}\mathrm{\Pi }`$ then $`\mathrm{\Pi }^{}\mathrm{\Pi }\backslash \{\gamma \}`$ for a clause $`\gamma `$. Therefore, $`Mod(\mathrm{\Pi })Mod(\mathrm{\Pi }\backslash \{\gamma \})Mod(\mathrm{\Pi }^{})`$, which proves that $`\mathrm{\Pi }`$ and $`\mathrm{\Pi }^{}`$ are not equivalent.
This proof does not work for circumscription because the set of minimal models of a formula can grow or shrink in response to a clause deletion. In principle, $`\mathrm{\Pi }`$ and $`\mathrm{\Pi }\backslash \{\gamma \}`$ might have different sets of minimal models and yet $`\mathrm{\Pi }`$ and $`\mathrm{\Pi }^{}\mathrm{\Pi }\backslash \{\gamma \}`$ have the same minimal models. We show that this is not possible. The proof is based on the following simple result about quasi-orders (reflexive and transitive relations.)
###### Lemma 1
If $``$ is a quasi-order (a reflexive and transitive relation) and $`A`$ and $`B`$ are two finite sets such that $`AB`$ and $`\mathrm{min}_{}(A)\mathrm{min}_{}(B)`$, then $`\mathrm{min}_{}(B)\backslash A`$ is not empty.
Proof. Since $`\mathrm{min}_{}(A)\mathrm{min}_{}(B)`$, then either $`\mathrm{min}_{}(A)\backslash \mathrm{min}_{}(B)`$ or $`\mathrm{min}_{}(B)\backslash \mathrm{min}_{}(A)`$ is not empty. We consider these two cases separately.
Let $`x\mathrm{min}_{}(B)\backslash \mathrm{min}_{}(A)`$. We prove that $`xA`$. Since $`x`$ is minimal in $`B`$, there is no element of $`yB`$ such that $`y<x`$. Since $`AB`$, the same holds for every element of $`A`$ in particular. As a result, if $`xA`$ then $`x\mathrm{min}_{}(A)`$, contradicting the assumption.
Let us instead assume that $`\mathrm{min}_{}(A)\backslash \mathrm{min}_{}(B)`$ is not empty. Let $`x\mathrm{min}_{}(A)\backslash \mathrm{min}_{}(B)`$. Since $`xA`$, it holds $`xB`$. Since $`xB`$, $`x\mathrm{min}_{}(B)`$, and $`B`$ is a finite set, there exists $`y\mathrm{min}_{}(B)`$ such that $`y<x`$. Since $`x`$ is minimal in $`A`$, we have that $`yA`$
The order $``$ on propositional models is a quasi-order. As a result, if $`A`$ and $`B`$ are two sets of models such that $`AB`$ and the set of minimal elements of $`A`$ and $`B`$ are different, then $`B`$ has a minimal element that is not in $`A`$. When applied to circumscription, this result tells that a formula can be non-equivalent to a stronger one only because of a minimal model that is not a model of the stronger formula. In the other way around, if a formula is weakened, the set of minimal models either remains the same or acquires a new element.
###### Theorem 1
If $`Mod(\mathrm{\Pi })Mod(\mathrm{\Pi }^{})Mod(\mathrm{\Pi }^{\prime \prime })`$ and $`\mathrm{CIRC}(\mathrm{\Pi })\mathrm{CIRC}(\mathrm{\Pi }^{})`$ then $`\mathrm{CIRC}(\mathrm{\Pi })\mathrm{CIRC}(\mathrm{\Pi }^{\prime \prime })`$.
Proof. Let us assume that $`\mathrm{CIRC}(\mathrm{\Pi })=\mathrm{CIRC}(\mathrm{\Pi }^{\prime \prime })`$. Since $`\mathrm{CIRC}(\mathrm{\Pi })\mathrm{CIRC}(\mathrm{\Pi }^{})`$, we have $`\mathrm{CIRC}(\mathrm{\Pi }^{})\mathrm{CIRC}(\mathrm{\Pi }^{\prime \prime })`$. Since $``$ on propositional models is a quasi-order, Lemma 1 applies: there exists $`M`$ such that $`M\mathrm{CIRC}(\mathrm{\Pi }^{\prime \prime })`$ and $`MMod(\mathrm{\Pi }^{})`$. Since $`Mod(\mathrm{\Pi })Mod(\mathrm{\Pi }^{})`$, the latter implies $`MMod(\mathrm{\Pi })`$. As a result, $`M\mathrm{CIRC}(\mathrm{\Pi })`$. Since $`M\mathrm{CIRC}(\mathrm{\Pi }^{\prime \prime })`$, we have $`\mathrm{CIRC}(\mathrm{\Pi })\mathrm{CIRC}(\mathrm{\Pi }^{\prime \prime })`$
The redundancy of a formula and the presence of a redundant clause in the formula are related by the following theorem, which is an application of the above to the case in which $`\mathrm{\Pi }^{}=\mathrm{\Pi }\backslash \{\gamma \}`$ and $`\mathrm{\Pi }^{\prime \prime }`$ is a subset of $`\mathrm{\Pi }^{}`$.
###### Theorem 2
A CNF formula $`\mathrm{\Pi }`$ is CIRC-redundant if and only if it contains a CIRC-redundant clause.
Proof. The “if” direction is obvious: if $`\gamma `$ is redundant in $`\mathrm{\Pi }`$, then $`\mathrm{\Pi }\backslash \{\gamma \}_M\gamma `$, and $`\mathrm{\Pi }\backslash \{\gamma \}`$ is therefore a strict subset of $`\mathrm{\Pi }`$ that is equivalent to it. The “only if” direction is a consequence of the above theorem. Assume that $`\mathrm{CIRC}(\mathrm{\Pi }\backslash \{\gamma \})\mathrm{CIRC}(\mathrm{\Pi })`$ holds for every $`\gamma \mathrm{\Pi }`$. Let us consider $`\mathrm{\Pi }^{\prime \prime }\mathrm{\Pi }`$: we prove that $`\mathrm{CIRC}(\mathrm{\Pi }^{\prime \prime })\mathrm{CIRC}(\mathrm{\Pi })`$. Since $`\mathrm{\Pi }^{\prime \prime }\mathrm{\Pi }`$, there exists $`\gamma \mathrm{\Pi }^{\prime \prime }\backslash \mathrm{\Pi }`$. Consider one such clause $`\gamma `$. Since $`\mathrm{\Pi }\mathrm{\Pi }\backslash \{\gamma \}\mathrm{\Pi }^{\prime \prime }`$, we have that $`Mod(\mathrm{\Pi }^{\prime \prime })Mod(\mathrm{\Pi }\backslash \{\gamma \})Mod(\mathrm{\Pi })`$. We are thus in the conditions to apply the above theorem: since $`\mathrm{CIRC}(\mathrm{\Pi }\backslash \{\gamma \})\mathrm{CIRC}(\mathrm{\Pi })`$, we have that $`\mathrm{CIRC}(\mathrm{\Pi }^{\prime \prime })\mathrm{CIRC}(\mathrm{\Pi })`$. Therefore, $`\mathrm{\Pi }`$ is not equivalent to any of its proper subsets.
This theorem shows that circumscription, although nonmonotonic, has the local redundancy property.
### 3.2 Redundant Clauses
The following lemma characterizes the clauses that are redundant in a formula.
###### Lemma 2
The following three conditions are equivalent:
1. the clause $`\gamma \mathrm{\Pi }`$ is $`\mathrm{CIRC}`$-redundant in $`\mathrm{\Pi }`$;
2. for each $`MMod(\mathrm{\Pi }\backslash \{\gamma \}\neg \gamma )`$ there exists $`M^{}Mod(\mathrm{\Pi })`$ such that $`M^{}M`$;
3. for each $`MMod(\mathrm{\Pi }\backslash \{\gamma \}\neg \gamma )`$ there exists $`M^{}Mod(\mathrm{\Pi }\backslash \{\gamma \})`$ such that $`M^{}M`$.
Proof. The models of $`\mathrm{\Pi }\backslash \{\gamma \}`$ that are not models of $`\mathrm{\Pi }`$ are exactly the models of $`\mathrm{\Pi }\backslash \{\gamma \}\neg \gamma `$. The two formulae $`\mathrm{\Pi }`$ and $`\mathrm{\Pi }\backslash \{\gamma \}`$ are $`_M`$-equivalent if none of these models (if any) is minimal, that is, all these models contain other models of $`\mathrm{\Pi }`$. In other words, $`\gamma `$ is redundant if and only if every model of $`\mathrm{\Pi }\backslash \{\gamma \}\neg \gamma `$ contains a model of $`\mathrm{\Pi }`$.
The fact that we can check $`M^{}Mod(\mathrm{\Pi }\backslash \{\gamma \})`$ instead of $`M^{}Mod(\mathrm{\Pi })`$ follows from the fact that $`Mod(\mathrm{\Pi }\backslash \{\gamma \})`$ is composed of all models of $`\mathrm{\Pi }`$ and all models of $`\mathrm{\Pi }\backslash \{\gamma \}\neg \gamma `$. Consider a model $`M`$ that is a minimal model of $`\mathrm{\Pi }\backslash \{\gamma \}\neg \gamma `$. The condition $`M^{}M`$ implies that $`M^{}`$ is not a model of $`\mathrm{\Pi }\backslash \{\gamma \}\neg \gamma `$, and is therefore a model of $`\mathrm{\Pi }`$. By transitivity, the condition that there exists $`M^{}Mod(\mathrm{\Pi })`$ such that $`M^{}M`$ holds for all models of $`\mathrm{\Pi }\backslash \{\gamma \}\neg \gamma `$
Computationally, checking the second or third condition of this lemma can be done by checking whether for all $`M\mathrm{}`$ there exists $`M^{}\mathrm{}`$ such that a simple condition is met. As a result, the problem is in $`\mathrm{\Pi }_2^p`$. For positive clauses, checking $`\mathrm{CIRC}`$-redundancy is easier, as it amounts to checking classical redundancy.
###### Lemma 3
A positive clause is $`\mathrm{CIRC}`$-redundant in $`\mathrm{\Pi }`$ if and only if it is classically redundant in $`\mathrm{\Pi }`$.
Proof. If a clause is redundant in $`\mathrm{\Pi }`$ it is also $`\mathrm{CIRC}`$-redundant in $`\mathrm{\Pi }`$. Let us now prove the converse: assume that $`\gamma `$ is a positive clause that is $`\mathrm{CIRC}`$-redundant in $`\mathrm{\Pi }`$. By Lemma 2, every model of $`\mathrm{\Pi }\backslash \{\gamma \}\neg \gamma `$ contains a model of $`\mathrm{\Pi }`$, which is the same as $`\mathrm{\Pi }\backslash \{\gamma \}\{\gamma \}`$. Let $`\gamma =x_{i_1}\mathrm{}x_{i_k}`$. Since $`\gamma `$ is $`\mathrm{CIRC}`$ redundant, each model of $`\mathrm{\Pi }\backslash \{\gamma \}\{\neg x_{i_i},\mathrm{},\neg x_{i_k}\}`$ contains at least a model of $`\mathrm{\Pi }\backslash \{\gamma \}\{x_{i_i}\mathrm{}x_{i_k}\}`$. All models of the latter formula contain at least a variable among $`x_{i_i},\mathrm{},x_{i_k}`$ while no models of the former contain any of them. Therefore, no model of the first formula contains a model of the second. Therefore, the condition can be true only if $`\mathrm{\Pi }\backslash \{\gamma \}\{\neg x_{i_1}\mathrm{}\neg x_{i_k}\}`$ has no models, that is, $`\mathrm{\Pi }\backslash \{\gamma \}\gamma `$: the clause $`\gamma `$ is classically redundant in $`\mathrm{\Pi }`$
Intuitively, positive clauses only exclude models with all their literals assigned to false. Therefore, whenever a positive clause is irredundant w.r.t. $``$, it is because such models were not otherwise excluded; therefore, it is also irredundant w.r.t. minimal models.
According to this argument, it may look like all negative clauses are redundant because they exclude models with positive literals, and these models are not minimal. This is however not the case: a model with some positive literals might be minimal because no other model of the formula has less positive literal. Consider, for example, the following formula:
$$\mathrm{\Pi }=\{\neg x_1\neg x_2,x_1x_3,x_2x_3\}$$
The clause $`\neg x_1\neg x_2`$, although negative, is irredundant. Indeed, $`\mathrm{\Pi }\backslash \{\neg x_1\neg x_2\}=\{x_1x_3,x_2x_3\}`$, and this formula has $`\{x_1,x_2\}`$ and $`\{x_3\}`$ as its minimal models. The first one is not a model of $`\mathrm{\Pi }`$ because of the clause $`\neg x_1\neg x_2`$. Therefore, $`\neg x_1\neg x_2`$ is CIRC-irredundant in $`\mathrm{\Pi }`$.
Intuitively, a negative clause excludes the possibility of setting all variables to true, while minimal inference only tries to set variables to false. Therefore, removing the clause may generate a model that have its variables set to true ($`\{x_1,x_2\}`$ in the example), but is minimal because of the values of the other variables ($`x_3`$ in the example).
Lemma 3 can be extended to clauses containing negative literals via the addition of new clauses and new variables. To this aim, the following property of quasi-orders is needed.
###### Lemma 4
If $``$ is a quasi-order, $`X\mathrm{min}_{}(A)`$, $`XB`$, and $`BA`$, then $`X\mathrm{min}_{}(B)`$.
Proof. Since $`X\mathrm{min}_{}(A)`$, there is no element $`YA`$ such that $`Y<X`$. Since $`BA`$, there is no element of $`B`$ with the same property. Since $`X`$ is an element of $`B`$ such that $`Y<X`$ does not hold for any $`YB`$, it holds $`X\mathrm{min}_{}(B)`$ by definition.
Applied to formulae: if $`M`$ is a minimal model of $`\mathrm{\Pi }`$ and satisfies $`\mathrm{\Pi }^{}`$, then $`M`$ is a minimal model of $`\mathrm{\Pi }\mathrm{\Pi }^{}`$.
###### Lemma 5
A clause $`\gamma `$ is classically redundant in $`\mathrm{\Pi }`$ if and only if it is $`\mathrm{CIRC}`$-redundant in $`\mathrm{\Pi }\{xx^{}|\neg x\gamma \}`$.
Proof. If $`\gamma `$ is redundant in $`\mathrm{\Pi }`$ then $`\mathrm{\Pi }\backslash \{\gamma \}\gamma `$ and therefore $`\mathrm{\Pi }\{xx^{}|\neg x\gamma \}\backslash \{\gamma \}\gamma `$. Since $`\gamma `$ is redundant in $`\mathrm{\Pi }\{xx^{}|\neg x\gamma \}`$, it is also $`\mathrm{CIRC}`$-redundant.
Let us now assume that $`\gamma `$ is irredundant in $`\mathrm{\Pi }`$, that is, $`\mathrm{\Pi }\backslash \{\gamma \}\neg \gamma `$ has some models. Let $`M`$ be one such model. Since this model satisfies $`\neg \gamma `$, it assigns false to any variable $`x`$ such that $`x\gamma `$ and true to any variable $`x`$ such that $`\neg x\gamma `$. Extending $`M`$ to assign false to all variables $`x^{}`$, this model also satisfies $`\mathrm{\Pi }\{xx^{}|\neg x\gamma \}\backslash \{\gamma \}\neg \gamma `$.
We show that $`M`$ cannot contain a model of $`\mathrm{\Pi }\{xx^{}|\neg x\gamma \}`$. This model assigns false to all $`x\gamma `$ and also false to all $`x^{}`$ such that $`\neg x\gamma `$. On the other hand, $`\{\gamma \}\{xx^{}|\neg x\gamma \}`$ entails the clause $`\{x|x\gamma \}\{x^{}|\neg x\gamma \}`$; this can be proved for example by iteratively resolving upon all literals $`x`$ such that $`\neg x\gamma `$. As a result, no model of $`\mathrm{\Pi }\{xx^{}|\neg x\gamma \}`$ has a model that assign false to all $`x\gamma `$ and all $`x^{}`$ such that $`\neg x\gamma `$. Since this is instead done by $`M`$, it follows that no model of $`\mathrm{\Pi }\{xx^{}|\neg x\gamma \}`$ is contained in $`M`$
Note that the clauses $`xx^{}`$ are not necessarily $`\mathrm{CIRC}`$-irredundant in the considered formula. On the other hand, Lemma 3 can be applied to them: they are $`\mathrm{CIRC}`$-redundant if and only if they are classically redundant.
### 3.3 Complexity of Clause Redundancy
Let us now turn to the hardness of the problem of checking the redundancy of a clause in a formula. We first show a reduction that proves the hardness of the problem of redundancy of a clause and then show how this result can be used to prove that the problem of redundancy of a formula has the same complexity.
###### Theorem 3
Checking the CIRC-redundancy of a clause in a formula is $`\mathrm{\Pi }_2^p`$-complete.
Proof. Lemma 2 proves that the redundancy of a clause in a formula can be checked by solving a $``$QBF (for all $`M`$… there exists $`M^{}`$…), and is therefore in $`\mathrm{\Pi }_2^p`$.
Let us now show hardness. We show that the QBF formula $`XY.\mathrm{\Gamma }`$, where $`\mathrm{\Gamma }=\{\delta _1,\mathrm{},\delta _m\}`$ and $`n=|Y|`$, is valid if and only if $`\gamma `$ is $`\mathrm{CIRC}`$-redundant in $`\mathrm{\Pi }`$, where:
$`\mathrm{\Pi }`$ $`=`$ $`\{x_ip_i\}\{\neg ay_i\}\{a\delta _i|\delta _i\mathrm{\Gamma }\}\{\gamma \}`$
$`\gamma `$ $`=`$ $`\neg a\neg y_1\mathrm{}\neg y_n`$
The clause $`\gamma `$ is CIRC-redundant in $`\mathrm{\Pi }`$ if and only if all minimal models of $`\mathrm{\Pi }\backslash \{\gamma \}\neg \gamma `$ contain some models of $`\mathrm{\Pi }`$. The following equivalences holds:
$`\mathrm{\Pi }`$ $``$ $`\{x_ip_i\}\{\neg a\}\mathrm{\Gamma }`$
$`\mathrm{\Pi }\backslash \{\gamma \}\neg \gamma `$ $``$ $`\{x_ip_i\}\{a,y_1,\mathrm{},y_n\}`$
The first equivalence holds because $`\{\neg ay_i\}\{\neg a\neg y_1\mathrm{}\neg y_n\}`$ is equivalent to $`\neg a`$, as can be checked by resolving upon each $`y_i`$ in turn. The second equivalence holds because $`\neg \gamma =\{a,y_1,\mathrm{},y_n\}`$ and this set implies all clauses $`\neg ay_i`$ and $`a\delta _i`$.
The formula $`\mathrm{\Pi }\backslash \{\gamma \}\neg \gamma \{x_ip_i\}\{a,y_1,\mathrm{},y_n\}`$ has a minimal model for each truth evaluation $`\omega _X`$ over the variables $`x_i`$:
$$I_{\omega _X}=\omega _X\{p_i|x_i\omega _X\}\{a\}\{y_i|1in\}$$
We show that the model $`I_{\omega _X}`$ contains a model of $`\mathrm{\Pi }`$ if and only if $`\mathrm{\Gamma }|_{\omega _X}`$ is satisfiable. By Lemma 2, the redundancy of $`\gamma `$ corresponds to this condition being true for all possible models of $`\mathrm{\Pi }^{}\neg \gamma `$. This would therefore prove that the QBF is valid if and only if $`\gamma `$ is redundant in $`\mathrm{\Pi }`$.
Since $`\mathrm{\Pi }\{x_ip_i\}\{\neg a\}\mathrm{\Gamma }`$, if $`\mathrm{\Gamma }`$ has a model with a given value of $`\omega _X`$ then $`\mathrm{\Pi }`$ has a model that is strictly contained in $`I_{\omega _X}`$: add to the satisfying assignment of $`\mathrm{\Gamma }`$ the setting of every $`p_i`$ to the opposite of $`x_i`$ and $`a`$ to false.
On the converse, if $`\mathrm{\Pi }`$ contains a model that is strictly contained in $`I_{\omega _X}`$, this model must have exactly the same value of $`XP`$ because $`\mathrm{\Pi }`$ contains $`x_ip_i`$ and either $`x_i`$ or $`y_i`$ is false in $`I_{\omega _X}`$. On the other hand, this model of $`\mathrm{\Pi }`$ must also set $`a`$ to false and satisfy $`\mathrm{\Gamma }`$, thus showing that there exists an assignment extending $`\omega _X`$ and satisfying $`\mathrm{\Gamma }`$.
### 3.4 Complexity of Formula Redundancy
In order to characterize the complexity of the problem of checking the CIRC-redundancy of a formula, we use the fact that a formula is CIRC-redundant if and only if it contains a CIRC-redundant clause by Theorem 2. In particular, Lemma 3 shows that the problem of checking the $`\mathrm{CIRC}`$-redundancy of a clause $`\gamma `$ in $`\mathrm{\Pi }`$ is $`\mathrm{\Pi }_2^p`$-hard. In order for this result to be used as a proof of hardness for the problem of CIRC-redundancy of formulae, we need to modify the formula $`\mathrm{\Pi }`$ in such a way all its clauses but $`\gamma `$ are made CIRC-irredundant. This is the corresponding of Lemma 4 of the paper of redundancy of propositional CNF formulae \[Lib05b\], which has been useful because it allows to “localize” problems about redundancy.
###### Lemma 6
For every consistent formula $`\mathrm{\Pi }`$ and $`\mathrm{\Pi }^{}\mathrm{\Pi }`$, the only $`\mathrm{CIRC}`$-redundant clauses of $`I(\mathrm{\Pi },\mathrm{\Pi }^{})`$ are the clauses $`\neg s\neg t\gamma _i`$ such that $`\gamma _i\mathrm{\Pi }^{}`$ and $`\gamma _i`$ is $`\mathrm{CIRC}`$-redundant in $`\mathrm{\Pi }`$.
$`I(\mathrm{\Pi },\mathrm{\Pi }^{})`$ $`=`$ $`\{st\}\{sa,tb\}`$
$`\{\neg stc_id_i\}\{\neg s\neg c_i\}`$
$`\{\neg tc_i\gamma _i|\gamma _i\mathrm{\Pi }\backslash \mathrm{\Pi }^{}\}\{s\neg txx^{}|xVar(\mathrm{\Pi })\}`$
$`\{\neg s\neg t\gamma _i|\gamma _i\mathrm{\Pi }^{}\}`$
Proof. There are four possible assignment to the variables $`s`$ and $`t`$. Since the models of $`I(\mathrm{\Pi },\mathrm{\Pi }^{})`$ can be partitioned into the models of $`I(\mathrm{\Pi },\mathrm{\Pi }^{})\{\neg s,\neg t\}`$, $`I(\mathrm{\Pi },\mathrm{\Pi }^{})\{s,\neg t\}`$, $`I(\mathrm{\Pi },\mathrm{\Pi }^{})\{\neg s,t\}`$, and $`I(\mathrm{\Pi },\mathrm{\Pi }^{})\{s,t\}`$, the minimal models of $`I(\mathrm{\Pi },\mathrm{\Pi }^{})`$ are necessarily some of the minimal models of these formulae.
In the table below we show what remains of $`I(\mathrm{\Pi },\mathrm{\Pi }^{})\backslash \{st\}`$ in each of the four possible assignment to $`s`$ and $`t`$ after removing entailed clauses and false literals. We also show the minimal models of the resulting formulae.
| assignment | subformula | minimal models |
| --- | --- | --- |
| $`\{\neg s,\neg t\}`$ | $`\{a,b\}`$ | $`\{a,b\}`$ |
| $`\{s,\neg t\}`$ | $`\{b\}\{c_id_i\}\{\neg c_i\}`$ | $`\{s,b\}\{d_i\}`$ |
| $`\{\neg s,t\}`$ | $`\{a\}\{c_i\gamma _i|\gamma _i\mathrm{\Pi }\backslash \mathrm{\Pi }^{}\}`$ | $`\{t,a\}+\text{some subsets of }(CXX^{})`$ |
| | $`\{xx^{}|xVar(\mathrm{\Pi })\}`$ | |
| $`\{s,t\}`$ | $`\{\neg c_i\}\{c_i\gamma _i|\gamma _i\mathrm{\Pi }\backslash \mathrm{\Pi }^{}\}\mathrm{\Pi }^{}`$ | $`\{s,t\}+\text{a minimal model of }\mathrm{\Pi }`$ |
The four subformulae are all satisfiable. Moreover, no minimal model of one is contained in the minimal models of the other ones because of either the values of $`\{s,t\}`$ and $`\{a,b\}`$. As a result, the minimal models of $`I(\mathrm{\Pi },\mathrm{\Pi }^{})\backslash \{st\}`$ are exactly the minimal models of the four subformulae. The clause $`st`$ is irredundant because its addition deletes the minimal model $`\{a,b\}`$. The minimal models of $`I(\mathrm{\Pi },\mathrm{\Pi }^{})`$ are therefore exactly the minimal models of the remaining three subformulae.
We show that the remaining clauses but the ones derived from $`\mathrm{\Pi }^{}`$ are irredundant. This is shown by removing a clause from the set and showing that some of the minimal models of a subformula can be removed some elements. Since the minimal models of these three subformulae are exactly the minimal models of $`\mathrm{\Pi }`$, this is a proof that the clause is irredundant.
1. The clauses $`sa`$ and $`tb`$ are irredundant because their removal would allow $`a`$ and $`b`$ to be set to false in the minimal models of the third and second subformula, respectively.
2. The clauses $`\neg stc_id_i`$ and $`\neg s\neg c_i`$ are irredundant because their removal would allow $`d_i`$ to be set to false in the minimal model of the second subformula.
3. The clauses $`\neg tc_i\gamma _i`$ and $`s\neg txx^{}`$ require a longer analysis. In the third assignment, $`I(\mathrm{\Pi },\mathrm{\Pi }^{})`$ becomes:
$$C=\{a\}\{c_i\gamma _i\}\{xx^{}|xVar(\mathrm{\Pi })\}$$
The clauses $`xx^{}`$ are positive. By Lemma 3, they are $`\mathrm{CIRC}`$-redundant if and only if they are redundant. In turn, they are not redundant because $`\{a\}\{c_i\}\{y|yx\}`$ is a model of all clauses but $`xx^{}`$.
Since $`c_i`$ occurs positive in $`c_i\gamma _i`$, Lemma 5 ensures that this clause is $`\mathrm{CIRC}`$-redundant in $`C`$ if and only if it is redundant in $`C\backslash \{xx^{}|\neg x\gamma _i\}`$. This is false because the removal of $`c_i\gamma _i`$ creates the following new model:
$$M=\{c_j|ji\}\{x^{}\}\{x|\neg x\gamma _i\}$$
This model $`M`$ satisfies $`C\backslash \{xx^{}|\neg x\gamma _i\}\backslash \{c_i\gamma _i\}`$: all clauses $`c_j\gamma _j`$ are satisfied because $`c_jM`$ and all clauses $`xx^{}`$ are satisfied because $`x^{}M`$. On the other hand, $`M`$ does not satisfy $`c_i\gamma _i`$ because it assigns all its literals to false.
The only clauses that can therefore be redundant are those corresponding to the clauses of $`\mathrm{\Pi }^{}`$. In particular, these clauses only occur in the fourth subformula, which is equivalent to $`\{c\}\{\neg c_i\}\mathrm{\Pi }`$. A clause $`\neg s\neg t\gamma _i`$ with $`\gamma _i\mathrm{\Pi }^{}`$ is therefore $`\mathrm{CIRC}`$-redundant in $`I(\mathrm{\Pi },\mathrm{\Pi }^{})`$ if and only if $`\gamma _i`$ is $`\mathrm{CIRC}`$-redundant in $`\mathrm{\Pi }`$
More precisely, this theorem shows a way to make the clauses of $`\mathrm{\Pi }^{}`$ necessary, that is, contained in all equivalent subsets of $`\mathrm{\Pi }`$. The theorem allows to characterize the complexity of formula CIRC-redundancy.
###### Theorem 4
The problem of $`\mathrm{CIRC}`$-redundancy is $`\mathrm{\Pi }_2^p`$-complete.
Proof. By Theorem 2, $`\mathrm{\Pi }`$ is redundant if and only if it contains a redundant clause. Therefore, we have to solve a linear number of problems in $`\mathrm{\Pi }_2^p`$. Since these problems can be solved in parallel, the whole problem is in $`\mathrm{\Pi }_2^p`$.
Hardness is proved by reduction from the problem of $`\mathrm{CIRC}`$-redundancy of a single clause. By Lemma 6, a clause $`\gamma `$ is CIRC-redundant in $`\mathrm{\Pi }`$ if and only if $`\neg s\neg t\gamma `$ is CIRC-redundant in $`I(\mathrm{\Pi },\{\gamma \})`$ and all other clauses of $`I(\mathrm{\Pi },\{\gamma \})`$ are irredundant.
## 4 Default Logic
A default theory is a pair $`D,W`$, where $`W`$ is formula and $`D`$ is a set of default rules, each rule being in the form:
$$\frac{\alpha :\beta }{\gamma }$$
The formulae $`\alpha `$, $`\beta `$, and $`\gamma `$ are called the precondition, the justification, and the consequence of the default, respectively. In this paper, we assume that $`W`$ is a CNF finite formula (a finite set of clauses) and that the set of variables and defaults are finite. We also assume that each default has a single justification, rather than a set of justifications. Given a default $`d=\frac{\alpha :\beta }{\gamma }`$, its parts are denoted by $`\mathrm{prec}(d)=\alpha `$, $`\mathrm{just}(d)=\beta `$, and $`\mathrm{cons}(d)=\gamma `$.
We use the operational semantics of default logics \[AS94, Ant99, FM92, FM94\], which is based on sequences of defaults with no duplicates. If $`\mathrm{\Pi }`$ is such a sequence, we denote by $`\mathrm{\Pi }[d]`$ the sequence of defaults preceeding $`d`$ in $`\mathrm{\Pi }`$, and by $`\mathrm{\Pi }[d]`$ the sequence obtained by adding $`d`$ at the end of $`\mathrm{\Pi }`$. We extend the notation from defaults to sequences, so that $`\mathrm{prec}(\mathrm{\Pi })`$ is the conjunction of all preconditions of the defaults in $`\mathrm{\Pi }`$, $`\mathrm{just}(\mathrm{\Pi })`$ is the conjunction of all justifications, and $`\mathrm{cons}(\mathrm{\Pi })`$ is the conjunction of all consequences.
Implication is denoted by $``$, $``$ indicates (combined) consistency, and $``$ indicates inconsistency. For example, $`AB`$ means that $`AB`$ is consistent, while $`AB`$ means that $`AB`$ is inconsistent.
Default logic can be defined in terms of the selected processes, that are the sequences of defaults that are considered applicable by the semantics \[Ant99\]. A sequence of defaults $`\mathrm{\Pi }`$ is a process if $`W\mathrm{cons}(\mathrm{\Pi }[d])\mathrm{prec}(d)`$ holds for any $`d\mathrm{\Pi }`$. A default $`d`$ is locally applicable in a sequence $`\mathrm{\Pi }`$ if $`\mathrm{cons}(\mathrm{\Pi })W\mathrm{prec}(d)`$ and $`\mathrm{cons}(\mathrm{\Pi })W\mathrm{just}(d)`$. Global applicability also requires $`\mathrm{cons}(\mathrm{\Pi })W\mathrm{just}(\mathrm{\Pi }[d])`$. Each semantics defines the sequences of defaults that are applied in a particular theory. Formally, the definitions are as follows:
a process $`\mathrm{\Pi }`$ is selected if $`\mathrm{cons}(\mathrm{\Pi })W\mathrm{just}(d)`$ for each $`d\mathrm{\Pi }`$ and no default $`d^{}\mathrm{\Pi }`$ is locally applicable in $`\mathrm{\Pi }`$;
a process is selected if it is a maximal process such that $`\mathrm{cons}(\mathrm{\Pi })W\mathrm{just}(d)`$ for each $`d\mathrm{\Pi }`$;
a process is selected if it is a maximal process such that $`\mathrm{cons}(\mathrm{\Pi })W\mathrm{just}(\mathrm{\Pi })`$;
a process is selected if $`\mathrm{cons}(\mathrm{\Pi })W\mathrm{prec}(d)`$ and no default $`d^{}\mathrm{\Pi }`$ is globally applicable in $`\mathrm{\Pi }`$.
The conditions on selected processes can be all broken in two parts: success (the consistency condition) and closure (the non-extendibility of the process). For example, for constrained default logic the condition of success is $`\mathrm{cons}(\mathrm{\Pi })W\mathrm{just}(d)`$ and the condition of closure is that $`\mathrm{\Pi }[d]`$ is not successful for any $`d^{}\mathrm{\Pi }`$.
Remarkably, the conditions above only mention the background theory $`W`$ in conjunction with $`\mathrm{cons}(\mathrm{\Pi })`$, that is, $`W`$ only occurs in subformulae of the form $`W\mathrm{cons}(\mathrm{\Pi })`$. The only conditions for which this is not true is that of $`\mathrm{\Pi }`$ being a process.
If $`\mathrm{\Pi }`$ is a selected process of $`D,W`$, the formula $`Cn(W\mathrm{cons}(\mathrm{\Pi }))`$ is an extension of $`D,W`$. We denote by $`\mathrm{Ext}(D,W)`$ or $`\mathrm{Ext}_D(W)`$ the set of all formulae that are equivalent to an extension of $`D,W`$. Including formulae that are equivalent to the extensions in this set allows to write $`E\mathrm{Ext}_D(W)`$ to denote the equivalence of $`E`$ with an extension of $`D,W`$.
A default theory $`D,W`$ entails a formula $`W^{}`$ if and only if $`EW^{}`$ for every $`E\mathrm{Ext}_D(W)`$. This condition is equivalent to $`\mathrm{Ext}_D(W)W^{}`$; as a result, the set of all consequences of a default theory is equivalent to $`\mathrm{Ext}_D(W)`$. The condition that $`D,W`$ entails $`W^{}`$ is denoted by $`D,WW^{}`$ or $`W_DW^{}`$. The latter notation emphasizes that every fixed set of defaults $`D`$ induces a nonmonotonic inference operator $`_D`$.
Some semantics of default logic do not assign any extension to some theories. In this paper, we try to derive the hardness results using only theories having extensions.
### 4.1 Equivalence in Default Logics
The monotonic inference operator $`_D`$ induced by a set of default $`D`$ is a consequence relation. Therefore, the definitions of redundancy of a clause and of a formula for $``$ and $`_M`$ can be given for $`_D`$ as well: a clause $`\gamma `$ of a formula $`W`$ is redundant w.r.t. default $`D`$ if and only if $`W`$ and $`W\backslash \{\gamma \}`$ are equivalent; a formula $`W`$ is redundant if there exists $`W^{}W`$ that is equivalent to it.
Both definitions are based on equivalence of two formulae, and in particular the equivalence of a formula with one of its proper subsets. In this section, we show that three different form of equivalence can be defined; we compare them in general and in the particular case of equivalence of a formula with one of its proper subsets. The first form of equivalence is based on entailment.
###### Definition 6 (Entailment and Mutual Equivalence)
For a given set of defaults $`D`$, formula $`W`$ entails $`W^{}`$, denoted by $`W_DW^{}`$, if $`\mathrm{Ext}_D(W)W^{}`$. These two formulae are mutual equivalent, denoted by $`W_D^mW^{}`$, if $`W_DW^{}`$ and $`W^{}_DW`$.
In classical logic, this definition of equivalence is the same as $`W`$ and $`W^{}`$ having the same set of consequences and the same set of models. In default logic, this is not the case. We define the equivalence based on the set of consequences as follows.
###### Definition 7 (Consequence-Entailment and Consequence-Equivalence)
For a given set of defaults $`D`$, formula $`W`$ consequence-entails $`W^{}`$, denoted $`W_D^cW^{}`$, if $`\mathrm{Ext}_D(W)\mathrm{Ext}_D(W^{})`$. These two two formulae are consequence-equivalent, denoted $`W_D^cW^{}`$, if $`\mathrm{Ext}_D(W)\mathrm{Ext}_D(W^{})`$.
Note that the comparison of the two formulae is based on a fixed set of defaults $`D`$. A more stringent condition of equivalence of two defaults theories is that of having the same extensions.
###### Definition 8 (Faithful Entailment and Faithful Equivalence)
For a given set of defaults $`D`$, a formula $`W`$ faithfully entails $`W^{}`$, denoted $`W_D^eW^{}`$, if $`\mathrm{Ext}_D(W)\mathrm{Ext}_D(W^{})`$. These two formulae are faithfully equivalent, denoted $`W_D^eW^{}`$, if $`\mathrm{Ext}_D(W)=\mathrm{Ext}_D(W^{})`$.
For all three definition, equivalence is the same as each formula implying the other one. We are especially interested into $`_D^c`$ and $`_D^e`$, that is, equality of consequences and equality of extensions. Mutual equivalence has been defined for technical reasons. Redundancy in default logic is defined as follows.
###### Definition 9 (Redundancy of a Clause)
For a given set of defaults $`D`$, a clause $`\gamma `$ is redundant in a formula $`W`$ according to equivalence $`_D^x`$ if $`W_D^xW\backslash \{\gamma \}`$.
###### Definition 10 (Redundancy of a Formula)
For a given set of defaults $`D`$, a formula $`W`$ is redundant according to equivalence $`_D^x`$ if there exists $`W^{}W`$ such that $`W_D^xW^{}`$.
In both cases, we are comparing for equivalence a formula and one of its subsets. In the following section we study the equivalence of $`W^{}`$ and $`W`$ when $`W^{}W`$.
#### 4.1.1 Correspondence, General
We now compare the three forms of equivalence defined above. The following chain of implications is easy to prove:
$$W^{}_D^eWW^{}_D^cWW^{}_DW$$
The latter implication is proved by the following lemma.
###### Lemma 7
If $`W^{}_D^cW`$ then $`W^{}_DW`$.
Proof. By assumption, $`\mathrm{Ext}_D(W^{})\mathrm{Ext}_D(W)`$. Since every extension of $`D,W`$ entails $`W`$, we have $`\mathrm{Ext}_D(W)W`$. As a result, $`\mathrm{Ext}_D(W^{})W`$, which is by definition $`W^{}_DW`$
Redundancy is defined in terms of equivalence of two formulae, one contained in the other. As a result, it makes sense to study the conditions of equivalence in the particular case in which $`W^{}W`$. We prove that the above chain of implication can be wrapped around in this case, thus proving that the three conditions are equivalent.
###### Lemma 8
If $`W^{}W`$ and $`W^{}_DW`$, then $`W^{}_D^eW`$.
Proof. Let $`\mathrm{\Pi }`$ be a selected process of $`D,W^{}`$. We prove that it is also a selected process of $`D,W`$. Since $`W^{}_DW`$, the formula $`W`$ is entailed by every extension in $`\mathrm{Ext}_D(W^{})`$. In particular, $`W^{}\mathrm{cons}(\mathrm{\Pi })W`$. Therefore, $`W^{}\mathrm{cons}(\mathrm{\Pi })W\mathrm{cons}(\mathrm{\Pi })`$. As a result, all conditions (such as success and closure) where $`W^{}`$ only occurs in the subformula $`W^{}\mathrm{cons}(\mathrm{\Pi })`$ are not changed by replacing $`W^{}`$ with $`W`$. This is in particular true for all considered conditions of success and closure.
The only condition that mentions the background theory not in conjunction with $`\mathrm{cons}(\mathrm{\Pi })`$ is the condition of a sequence being a process: $`\mathrm{\Pi }`$ is a process of $`D,W^{}`$ if and only if $`W^{}\mathrm{cons}(\mathrm{\Pi }[d])\mathrm{prec}(d)`$ for any $`d\mathrm{\Pi }`$. The same condition is however true for $`W`$ because $`W^{}W`$ implies $`WW^{}`$
The following is a consequence of the above.
###### Corollary 1
If $`W^{}W`$, then:
$$W^{}_DWW^{}_D^cWW^{}_D^eW$$
#### 4.1.2 Non-Correspondence, General
The last corollary proves that the three definitions of entailment from $`W^{}`$ to $`W`$ are equivalent if $`W^{}W`$. The same does not hold for equivalence, and therefore does not hold for entailment from $`W`$ to $`W^{}`$.
###### Theorem 5
There exists $`D`$, $`W`$, and $`W^{}W`$ such that $`W^{}_D^mW`$ and $`W^{}_D^cW`$
Proof. Since $`W^{}_DW`$, every extension of $`D,W^{}`$ implies $`W`$. A wrong proof of $`W^{}_D^eW`$ could then be based on the fact that, once $`W`$ is derived from $`D,W^{}`$ applying some defaults, we can proceed by applying the defaults of an arbitrary process of $`D,W`$.
This figure shows why a process $`\mathrm{\Pi }^{}`$ of $`D,W^{}`$ and a process $`\mathrm{\Pi }`$ of $`D,W`$ cannot always be concatenated: while $`\mathrm{\Pi }^{}`$ allows the derivation of $`W`$, this process might also derive another formula $`W^{\prime \prime }`$ that makes the process $`\mathrm{\Pi }`$ inapplicable. An example in which this situation arises is the following one:
$`D`$ $`=`$ $`\{d_1,d_2\}`$
where:
$`d_1={\displaystyle \frac{:a\neg b}{a\neg b}}`$
$`d_2={\displaystyle \frac{a:b}{b}}`$
$`W`$ $`=`$ $`\{a\}`$
$`W^{}`$ $`=`$ $`\mathrm{}`$
The only process of $`D,W^{}`$ is $`[d_1]`$, which generates the extension $`Cn(a\neg b)`$. This extension entails $`W`$, but it also entails $`\neg b`$. The theory $`D,W`$ has also the process $`[d_2]`$, generating the extension $`Cn(ab)`$. These two processes cannot however be concatenated, as the consequence $`\neg b`$ of $`d_1`$ is inconsistent with the justification of $`d_2`$.
Since $`W^{}W`$, we have that $`W_DW^{}`$. In this example, we also have $`W^{}_DW`$ because the single extension of $`W^{}`$ entails $`W=\{a\}`$. However, $`W^{}`$ and $`W`$ have different set of extensions; in particular, $`\mathrm{Ext}_D(W^{})a\neg b`$ and $`\mathrm{Ext}_D(W)a`$
A similar result can be proved about $`_D^c`$ and $`_D^e`$.
###### Theorem 6
There exists $`D`$, $`W`$, and $`W^{}W`$ such that $`W^{}_D^cW`$ but $`W^{}_D^eW`$ in Reiter and justified default logic.
Proof. Rather than the counterexample itself, it is interesting to show how it has been derived. The idea is the same as that of Theorem 5: a process of $`D,W^{}`$ that cannot be concatenated with a process of $`D,W`$. The proof for this case, however, is complicated by the fact that we assume $`W^{}_D^cW`$, that is, $`\mathrm{Ext}_D(W^{})`$ and $`\mathrm{Ext}_D(W)`$ are equivalent. The theories used in Theorem 5 do not work, as any other pair of theories having one extension each: in this cases, indeed, $`\mathrm{Ext}_D(W)`$ is equivalent to $`\mathrm{Ext}_D(W)`$. In order for the counterexample to work, $`D,W^{}`$ must have multiple processes, each entailing $`W`$ and some other formula.
In order for the counterexample to work, some defaults of $`\mathrm{\Pi }`$ cannot be applied after $`\mathrm{\Pi }_1`$ or $`\mathrm{\Pi }_2`$ because their justifications are inconsistent with $`W_1`$ or $`W_2`$. In order for $`W_D^cW^{}`$ to hold, however, $`W\mathrm{cons}(\mathrm{\Pi })`$ must be equivalent to $`W(W_1W_2)`$. As a result, every model of $`W\mathrm{cons}(\mathrm{\Pi })`$ is a model of $`WW_1`$ or $`WW_2`$.
The precondition of the first default of $`\mathrm{\Pi }`$ is entailed by $`WW_1`$ and $`WW_2`$. Since the justifications of the defaults in $`\mathrm{\Pi }`$ are consistent with the consequences of $`\mathrm{\Pi }`$, there is a model that is both a model of $`\mathrm{cons}(\mathrm{\Pi })W`$ and a model of $`\mathrm{just}(d)`$ for any $`d\mathrm{\Pi }`$. But this is also a model of $`WW_1`$ or $`WW_2`$. As a result, the default $`d`$ is applicable in $`\mathrm{\Pi }_i`$.
This arguments cannot be extended further, however. Indeed, $`\mathrm{\Pi }`$ may be composed of two defaults, one applicable to $`WW_1`$ and one applicable to $`WW_2`$. This is possible in a selected process because Reiter and justified semantics does not enforce joint consistency of justifications.
A minimal counterexample requires two defaults that can be applied in $`W^{}`$ leading to two disjoint extensions $`WW_1`$ and $`WW_2`$, and two other defaults that can be applied in sequence from $`W`$, but not from $`WW_1`$ or $`WW_2`$.
The background theory $`W`$ of this counterexample is composed of the four possible models $`A`$, $`B`$, $`C`$, and $`D`$. We define $`D`$ so that $`W`$ has a process that generates an extension $`E`$ having $`A`$ and $`B`$ as its models. In order for the consequences to be the same of those of $`W^{}`$, both $`A`$ and $`B`$ have to be part of some extensions of $`W^{}`$. We define the defaults so that $`W^{}`$ has two processes generating $`A`$ and $`B`$, respectively. Namely, the first process generates $`A`$, but its justifications are satisfiable because the extension contains $`C`$; the other one contains $`B`$, but consistency with justifications are ensured by the model $`D`$. This trick is necessary to avoid these processes to be extended with the defaults that generate $`W^{}`$.
In order to make the discussion more intuitive, we identify models with terms, and define formulae and defaults based on terms. We then convert terms into real formulae. The defaults that are applicable from $`W^{}`$ are defined as follows.
$`d_1`$ $`=`$ $`{\displaystyle \frac{:C}{AC}}`$
$`d_2`$ $`=`$ $`{\displaystyle \frac{:D}{BD}}`$
Both $`d_1`$ and $`d_2`$ can be applied from $`W^{}`$, leading to a consequence that is inconsistent with the justification of the other default. Moreover, each extension contains a model of $`\{A,B\}`$, as required. For example, $`d_1`$ produces $`A`$. However, the consistency of the extension with the justification is ensured by the other model $`C`$. This is necessary to avoid these defaults to be applicable in the new extension $`E`$ of $`W`$. Let us now define the two defaults that are applicable from $`W`$ only and generate this extension $`E`$.
$`d_3`$ $`=`$ $`{\displaystyle \frac{W:A}{ABC}}`$
$`d_4`$ $`=`$ $`{\displaystyle \frac{ABC:B}{AB}}`$
The processes of $`W^{}`$ are $`[d_1d_3]`$ and $`[d_2]`$. There is no way to avoid the first default $`d_3`$ of the new extension to be part of some process of $`W^{}`$ as well. However, as in this case, it may have no effects. Indeed, the extensions are $`AC`$ and $`BD`$. Let us now consider which defaults can be applied to $`W`$. Both processes of $`W^{}`$ are still processes of $`W`$. However, $`d_3`$ is applicable to $`W`$, leading to a state in which both $`d_1`$ and $`d_2`$ are applicable. However, both processes only contain models that are among the previous ones.
The above terms can be translated into the following formuale:
$`A`$ $`=`$ $`abc`$
$`B`$ $`=`$ $`ab\neg c`$
$`C`$ $`=`$ $`a\neg bc`$
$`D`$ $`=`$ $`a\neg b\neg c`$
The defaults will be then defined as follows.
$`d_1`$ $`=`$ $`{\displaystyle \frac{:\neg bc}{ac}}`$
$`d_2`$ $`=`$ $`{\displaystyle \frac{:\neg b\neg c}{a\neg c}}`$
$`d_3`$ $`=`$ $`{\displaystyle \frac{a:bc}{bc}}`$
$`d_4`$ $`=`$ $`{\displaystyle \frac{a(bc):b\neg c}{b}}`$
In $`W^{}`$, only $`d_1`$ and $`d_2`$ are applicable. The first one leads to $`ac`$, which is consistent with the justification of $`d_3`$. The first selected process of $`W^{}`$ is therefore $`[d_1d_3]`$, leading to the extension $`ac`$.
The second process from $`W^{}`$ starts with $`d_2`$, which generates $`a\neg c`$, which is not consistent with $`d_1`$ and $`d_3`$, and does not imply the precondition of $`d_4`$. As a result, $`[d_2]`$ is the second selected process of $`W^{}`$, leading to the extension $`a\neg c`$. We therefore have $`\mathrm{Ext}_D(W^{})(ac)(a\neg c)a`$.
Let us now consider the extensions from $`W`$. All selected processes of $`W^{}`$ are also selected processes of $`W`$. However, we can now apply $`d_3`$, as $`a`$ is true in the background theory. We therefore obtain $`bc`$. This conclusion is inconsistent with the justification of $`d_2`$, but $`d_1`$ and $`d_4`$ can be applied. The first one leads to the extension $`ac`$, which is also an extension of $`D,W^{}`$. On the other hand, $`[d_3d_4]`$ leads to $`ab`$, which is a new extension. Nevertheless, $`\mathrm{Ext}_D(W)(ac)(a\neg c)(ab)a`$: the theory $`D,W`$ has some extensions that $`D,W^{}`$ does not have, but the skeptical consequences are the same.
In the proof, we used two defaults that are applicable in $`W`$ but not in the processes of $`D,W^{}`$. These two defaults cannot have mutually consistent justifications; otherwise, they would be both applicable in some process of $`D,W^{}`$ thanks to the fact that any extension of $`D,W`$ contains only models of some extensions of $`D,W^{}`$. This proof does not work for constrained and rational default logic; however, the same claim can be proved in a different way.
###### Theorem 7
There exists $`D`$, $`W`$, and $`W^{}W`$ such that $`W^{}_D^cW`$ but $`W^{}_D^eW`$, for constrained and rational default logic.
Proof. The idea is as follows: the constrained extensions of a default theory are each characterized by a model that is consistent with all justifications and consequences of the defaults used to generate the extensions \[Lib05d\]. Therefore, we might have a situation like the one depicted below:
The arrows indicate the model associated to each extension: $`C`$ is the model associated with $`E_1`$, etc. Note that $`E_1E_2E_3E_1E_2`$. As a result, a default theory having only the extensions $`E_1`$ and $`E_2`$ is equivalent to a theory having all three extensions, but yet these two theories have the same consequences.
We define $`D,W`$ in such a way it has all three extensions, but $`D,\mathrm{}`$ does not have the extension $`E_3`$ because the model $`B`$ is excluded by a default that generates $`W`$. The above condition can be realized using two variables $`a`$ and $`b`$ to distinguish the four models $`A`$, $`B`$, $`C`$, and $`D`$, and a variable $`x`$ to distinguish $`W`$ from $`W^{}`$.
$`D`$ $`=`$ $`\{d_1,d_2,d_3,d_4\}`$
$`W`$ $`=`$ $`\{x\}`$
$`W^{}`$ $`=`$ $`\mathrm{}`$
the defaults are:
generates $`W`$ from $`W^{}`$ $`d_1={\displaystyle \frac{:xa}{x}}`$
generates the extension $`E_1=Cn(xb)`$ $`d_2={\displaystyle \frac{x:ab}{b}}`$
generates the extension $`E_2=Cn(x\neg b)`$ $`d_3={\displaystyle \frac{x:a\neg b}{\neg b}}`$
generates the extension $`E_3=Cn(x)`$ $`d_4={\displaystyle \frac{x:\neg a\neg b}{x}}`$
The justification of $`d_2`$, $`d_3`$, and $`d_4`$ are mutually inconsistent. In $`D,W`$, the three extensions are generated by the processes $`[d_2,d_1]`$, $`[d_3,d_1]`$, and $`[d_4]`$. The presence of $`d_1`$ in these processes do not change their generated extensions, as $`\mathrm{cons}(d_1)=x`$, which is already in $`W`$. We have $`E_1E_2E_3x`$.
Let us now consider $`D,W^{}`$. The only default that is applicable in $`W^{}=\mathrm{}`$ is $`d_1`$, which generates $`x`$ but also have $`a`$ as a justification. As a result, the defaults $`d_2`$ and $`d_3`$ are still applicable, but $`d_4`$ is not. As a result, the only extensions of $`D,W^{}`$ are $`E_1`$ and $`E_2`$. We therefore have $`W^{}_D^eW`$. On the other hand, $`E_1E_2x`$, which is equivalent to $`E_1E_2E_3`$. As a result, $`W^{}_D^cW`$
#### 4.1.3 Correspondence, Particular Cases
While $`_D^c`$ and $`_D^e`$ are not the same in general, they coincide when all defaults are normal and one formula is contained in the other one.
###### Theorem 8
If $`W^{}W`$ and $`D`$ is a set of normal defaults, then $`W^{}_D^cW`$ implies $`W^{}_D^eW`$ in constrained default logic.
Proof. Given the previous result, we only have to prove that $`W_D^cW^{}`$ implies that $`\mathrm{Ext}_D(W)\mathrm{Ext}_D(W^{})`$, that is, $`D,W`$ does not have any extension that is not an extension of $`D,W^{}`$.
To the contrary, assume that such extension exists. Let $`\mathrm{\Pi }`$ be the process that generates the extension of $`D,W`$ that is not an extension of $`D,W^{}`$. By definition of process, $`\mathrm{cons}(\mathrm{\Pi })W\mathrm{just}(\mathrm{\Pi })`$ is consistent. Therefore, it has a model $`M`$. Since this model satisfies both $`W`$ and $`\mathrm{cons}(\mathrm{\Pi })`$, it is a model of the extension generated by $`\mathrm{\Pi }`$.
Since the conclusions of the two theories are the same, every model of the extension generated by $`\mathrm{\Pi }`$ is a model of some extensions of $`D,W^{}`$. Let $`\mathrm{\Pi }^{}`$ be the a process of $`D,W^{}`$ that generates an extension that contains the model $`M`$. We prove that all defaults of $`\mathrm{\Pi }`$ are in $`\mathrm{\Pi }^{}`$.
Since $`M`$ is a model of the extension generated by $`\mathrm{\Pi }^{}`$, it is a model of $`\mathrm{cons}(\mathrm{\Pi }^{})W^{}`$. Therefore, it is a model of $`\mathrm{cons}(\mathrm{\Pi }^{})`$, and a model of $`\mathrm{just}(\mathrm{\Pi }^{})`$ because defaults are normal. We have already proved that $`M`$ is a model of $`\mathrm{cons}(\mathrm{\Pi })`$ and $`\mathrm{just}(\mathrm{\Pi })`$ and of $`W`$. As a result, the set $`\mathrm{cons}(\mathrm{\Pi })\mathrm{cons}(\mathrm{\Pi }^{})Wjust(\mathrm{\Pi })\mathrm{just}(\mathrm{\Pi }^{})`$ is consistent. As a result, we can add all defaults of $`\mathrm{\Pi }`$ to $`\mathrm{\Pi }^{}`$ without contradicting the justifications.
As a result, the defaults of $`\mathrm{\Pi }`$ are not in $`\mathrm{\Pi }^{}`$ only if their preconditions are not entailed from the consequences of $`\mathrm{\Pi }`$. This is impossible: since $`\mathrm{\Pi }`$ is a process of $`D,W`$, we have $`W\mathrm{prec}(d)`$, where $`d`$ is the first default of $`\mathrm{\Pi }`$. As a result, $`d`$ must be part of $`\mathrm{\Pi }^{}`$, otherwise $`\mathrm{\Pi }^{}`$ would not be a maximal process. The consequences of $`d`$ are therefore part of $`\mathrm{cons}(\mathrm{\Pi }^{})W^{}`$. Repeating the argument with the second default of $`\mathrm{\Pi }`$ we get the same result. We can therefore conclude that all defaults of $`\mathrm{\Pi }`$ are in $`\mathrm{\Pi }^{}`$
Since Reiter, justified, constrained, and rational default logics coincide on normal default theories, the equality of the definitions of equivalence holds when defaults are normal.
###### Theorem 9
If $`D`$ is a set of normal defaults and $`W^{}W`$, then $`W^{}_D^cW`$ if and only if $`W^{}_D^eW`$.
When all defaults are categorical (prerequisite-free), the following lemma allows proving that the three considered forms of equivalence coincide.
###### Lemma 9
If $`D`$ is a set of categorical defaults, $`W^{}W`$, and $`W^{}_DW`$, then $`W_D^eW^{}`$ in constrained default logic.
Proof. Let $`\mathrm{\Pi }`$ be a selected process of $`D,W`$. We prove that $`\mathrm{\Pi }`$ is a selected process of $`D,W^{}`$ generating the same extension.
Since $`\mathrm{\Pi }`$ is a selected process of $`D,W`$, it holds that $`W\mathrm{cons}(\mathrm{\Pi })\mathrm{just}(\mathrm{\Pi })`$ is consistent. Since $`W^{}W`$, it also holds that $`W^{}\mathrm{cons}(\mathrm{\Pi })\mathrm{just}(\mathrm{\Pi })`$ is consistent. Since no default has preconditions, $`\mathrm{\Pi }`$ is a successful process of $`D,W^{}`$. Since constrained default logic is a failsafe semantics \[Lib05e\], there exists $`\mathrm{\Pi }^{}`$ such that $`\mathrm{\Pi }\mathrm{\Pi }^{}`$ is a selected process of $`D,W`$.
Since every extension of $`W^{}`$ entails $`W`$, this is in particular true for the extension generated by $`\mathrm{\Pi }\mathrm{\Pi }^{}`$. In other words, $`W^{}\mathrm{cons}(\mathrm{\Pi }\mathrm{\Pi }^{})W`$. As a result, $`W^{}\mathrm{cons}(\mathrm{\Pi }\mathrm{\Pi }^{})W\mathrm{cons}(\mathrm{\Pi }\mathrm{\Pi }^{})`$. Since $`\mathrm{\Pi }\mathrm{\Pi }^{}`$ is a process of $`D,W^{}`$, we have that $`W^{}\mathrm{cons}(\mathrm{\Pi }\mathrm{\Pi }^{})\mathrm{just}(\mathrm{\Pi }\mathrm{\Pi }^{})`$ is consistent, which is therefore equivalent to the consistency of $`W\mathrm{cons}(\mathrm{\Pi }\mathrm{\Pi }^{})\mathrm{just}(\mathrm{\Pi }\mathrm{\Pi }^{})`$. Therefore, $`\mathrm{\Pi }\mathrm{\Pi }^{}`$ is a successful process of $`D,W`$. Since $`\mathrm{\Pi }`$ is by assumption a maximal successful process of $`D,W`$, it must be $`\mathrm{\Pi }^{}=[]`$, that is, $`\mathrm{\Pi }\mathrm{\Pi }^{}=\mathrm{\Pi }`$. We have already proved that $`W^{}\mathrm{cons}(\mathrm{\Pi }\mathrm{\Pi }^{})W\mathrm{cons}(\mathrm{\Pi }\mathrm{\Pi }^{})`$, that is, $`\mathrm{\Pi }`$ generates the same extension in $`W`$ and in $`W^{}`$
Since constrained and Reiter default logics coincide on normal default theories, we have the following consequence.
###### Corollary 2
If $`D`$ is a set of normal and categorical defaults and $`W^{}W`$, the conditions $`W^{}_D^mW`$, $`W^{}_D^cW`$, and $`W^{}_D^eW`$ are equivalent.
### 4.2 Redundancy of Clauses vs. Theories
The redundancy of a clause $`\gamma `$ in a formula $`W`$ is defined as the equivalence of $`W`$ and $`W\backslash \{\gamma \}`$. The redundancy of a formula $`W`$ is defined as its equivalence to one of its proper subsets. A formula containing a redundant clause is redundant, but the converse is not always true: a formula might contain no redundant clause but yet it is equivalent to one of its proper subsets.
In this section, we compare the redundancy of a set of clauses with the redundancy of a single clause in default logic. In propositional logics, these two concepts are the same: $`\mathrm{\Pi }`$ is equivalent to one of its proper subsets if and only if it contains a redundant clause. In default logic, it may be that $`\gamma _1`$ and $`\gamma _2`$ are both irredundant in $`\{\gamma _1,\gamma _2\}`$ while $`\{\gamma _1,\gamma _2\}`$ is redundant, as shown by the theory $`D,W`$ defined below.
$`W`$ $`=`$ $`\{a,b\}`$
$`D`$ $`=`$ $`\{d_1,d_2,d_3\}`$
where
$`d_1={\displaystyle \frac{a:\neg b}{\neg b}}`$
$`d_2={\displaystyle \frac{b:\neg a}{\neg a}}`$
$`d_3={\displaystyle \frac{:ab}{ab}}`$
The theory $`D,W`$ has the single extension $`Cn(\{a,b\})`$. Indeed, $`d_1`$ and $`d_2`$ are not applicable because their justifications are inconsistent with $`W`$. The third default is applicable, but its consequence is $`ab`$, which is already in the theory.
The theory $`D,W\backslash \{b\}`$ still has the extension $`\{a,b\}`$, which results from the application of $`d_3`$, which then blocks the application of $`d_1`$ and $`d_2`$. However, it also has a new extension: since $`d_1`$ is applicable, it generates $`\neg b`$, which blocks the application of $`d_3`$. This produces the extension $`\{a,\neg b\}`$. In the same way, $`D,W\backslash \{a\}`$ has the two extensions $`\{a,b\}`$ and $`\{\neg a,b\}`$
The theory $`D,W\backslash \{a,b\}`$ has again a single extension: $`d_3`$ is the only applicable default, leading to the addition of $`ab`$. Neither $`d_1`$ nor $`d_2`$ are applicable. Therefore, $`\{a,b\}`$ is the only extension of this theory.
The set of extensions of the theory is changed by removing any single clause, but is not changed by the removal of both clauses. In other words, both $`a`$ and $`b`$ are irredundant in $`\{a,b\}`$, but $`\{a,b\}`$ is redundant. Since $`D`$ is a set of normal default, this counterexample holds even for normal default theories.
The two theories obtained by removing a single clause of $`W`$ differ from $`W`$ because of a new extension. This can be proved to be always the case if the removal of both clauses leads to the original set of extensions. This is proved by first showing a sort of “continuity” of extensions.
###### Lemma 10
If $`E`$ is an extension of $`D,W^{}`$ and of $`D,W`$ with $`W^{}W`$, then every selected process of $`D,W^{}`$ generating $`E`$ is a selected process of $`D,W`$ generating $`E`$.
Proof. Let $`\mathrm{\Pi }`$ be a selected process of $`D,W^{}`$ that generates $`E`$. Since $`W^{}W`$ we have $`WW^{}`$; therefore, $`\mathrm{\Pi }`$ is a process of $`D,W`$. Remains to prove that it is also selected. However, all conditions for a process to be selected in $`D,W^{}`$ contains $`W^{}`$ only in the subformula $`W^{}\mathrm{cons}(\mathrm{\Pi })`$. Since $`E=W^{}\mathrm{cons}(\mathrm{\Pi })`$, and $`E`$ is an extension of $`D,W`$, we have that $`EW`$. As a result, $`W^{}\mathrm{cons}(\mathrm{\Pi })W\mathrm{cons}(\mathrm{\Pi })`$. Therefore, every condition for $`\mathrm{\Pi }`$ in $`D,W^{}`$ is equivalent to the same condition for $`D,W`$
The following lemma relates the selected processes of three formulae.
###### Lemma 11
If $`\mathrm{\Pi }`$ is a selected process of both $`D,W^{}`$ and $`D,W`$ and generates the same extension in both theories, it is also a selected process of every $`D,W^{\prime \prime }`$ with $`W^{}W^{\prime \prime }W`$ and generates the same extension in $`D,W^{\prime \prime }`$.
Proof. If $`W^{}W`$ does not hold, the claim is trivially true because there is no $`W^{\prime \prime }`$ such that $`W^{}W^{\prime \prime }W`$.
Since $`\mathrm{\Pi }`$ is a process of $`D,W^{}`$, it is also a process of $`D,W^{\prime \prime }`$ because $`W^{\prime \prime }W^{}`$. Since $`\mathrm{\Pi }`$ generates the same extensions in $`W^{}`$ and $`W`$, we have that $`W^{}\mathrm{cons}(\mathrm{\Pi })W\mathrm{cons}(\mathrm{\Pi })`$. Since $`WW^{\prime \prime }W^{}`$, we also have $`W^{}\mathrm{cons}(\mathrm{\Pi })W^{\prime \prime }\mathrm{cons}(\mathrm{\Pi })`$. Therefore, every condition that is true for $`W^{}\mathrm{cons}(\mathrm{\Pi })`$ is also true for $`W^{\prime \prime }\mathrm{cons}(\mathrm{\Pi })`$
The following lemma proves that extensions of both a theory and a subset of it are also extensions of any theory “between them”.
###### Lemma 12
If $`E`$ is an extension of both $`D,W^{}`$ and $`D,W`$, with $`W^{}W`$, then it is an extension of any $`D,W^{\prime \prime }`$ with $`W^{}W^{\prime \prime }W`$
Proof. By Lemma 10, every selected process $`\mathrm{\Pi }`$ of $`D,W^{}`$ that generates $`E`$ is also a selected process of $`D,W`$ and generates the same extension $`E`$ in this theory. As a result, Lemma 11 applies, and $`\mathrm{\Pi }`$ is a selected process of $`D,W^{\prime \prime }`$ and generates the same extension.
This theorem shows that extensions have a form of “partial monotonicity”: an extension of both a subset and a superset of a formula is also an extension of the formula. This is important to our aims, as it shows that the equivalence $`W^{}_D^eW`$ implies that all default theories $`D,W^{\prime \prime }`$ with $`W^{}W^{\prime \prime }W`$ have the same extensions of $`D,W`$. Therefore, $`D,W^{\prime \prime }`$ can differ from $`D,W`$ only because of new extensions.
###### Corollary 3
If $`W^{}_D^eW`$, $`W^{}W^{\prime \prime }W`$, and $`W^{\prime \prime }_D^eW`$, then $`\mathrm{Ext}_D(W)\mathrm{Ext}_D(W^{\prime \prime })`$.
The existence of extensions of $`W^{\prime \prime }`$ that are not extensions of $`W`$ does not imply that $`W^{\prime \prime }`$ theory is not consequence-equivalent to $`W`$ and $`W^{\prime \prime }`$. On the other hand, $`W^{\prime \prime }_D^cW`$ implies $`W^{\prime \prime }_D^eW`$, which leads to the following consequence.
###### Corollary 4
If $`W^{}_D^eW`$, $`W^{}W^{\prime \prime }W`$, and $`W^{\prime \prime }_D^cW`$, then $`\mathrm{Ext}_D(W)\mathrm{Ext}_D(W^{\prime \prime })`$.
While it is not true that the irredundancy of two clauses proves the irredundancy of the set composed of them, it is however true that this can only happen because of new extensions that are created by removing a single clause. For some special cases of default logics, such creation is not possible.
###### Theorem 10
If $`D`$ is a set of normal and categorical defaults, then $`W^{}_D^eW`$ implies that $`W_D^eW^{\prime \prime }`$ for any $`W^{\prime \prime }`$ such that $`W^{}W^{\prime \prime }W`$.
Proof. By Lemma 12, all extensions of $`D,W`$ are also extensions of $`D,W^{\prime \prime }`$. We therefore only have to prove the converse. Let $`\mathrm{\Pi }`$ be a process of $`D,W^{\prime \prime }`$. Since the theory has no preconditions, all defaults $`d\mathrm{\Pi }`$ satisfy $`W^{}\mathrm{cons}(\mathrm{\Pi }[d])d`$. In other words, $`\mathrm{\Pi }`$ is a process of $`D,W^{}`$. Since $`W^{\prime \prime }\mathrm{cons}(\mathrm{\Pi })\mathrm{just}(d)`$ is consistent for every $`d\mathrm{\Pi }`$, and $`W^{}`$ is logically weaker than $`W^{\prime \prime }`$, the same condition is true for $`W^{}`$. Since normal default logic is fail-safe \[Lib05e\], there exists $`\mathrm{\Pi }^{}`$ such that $`\mathrm{\Pi }\mathrm{\Pi }^{}`$ is a selected process of $`D,W^{}`$. By Lemma 10 and Lemma 12, $`W\mathrm{cons}(\mathrm{\Pi }\mathrm{\Pi }^{})`$ is an extension of $`D,W^{\prime \prime }`$.
Let us assume that $`W\mathrm{cons}(\mathrm{\Pi }\mathrm{\Pi }^{})`$ and $`W^{\prime \prime }\mathrm{cons}(\mathrm{\Pi })`$ are not equivalent. This is only possible if $`W\mathrm{cons}(\mathrm{\Pi }\mathrm{\Pi }^{})W^{\prime \prime }\mathrm{cons}(\mathrm{\Pi })`$ but not vice versa. This is however impossible, because in normal default logic all extensions are mutually inconsistent \[Rei80\]
Since the two definitions of equivalence are the same on normal default theories, as proved in Theorem 9, this result extends to the definition of redundancy based on consequences.
###### Corollary 5
If $`D`$ is a set of normal and categorical defaults, then $`W_D^cW^{\prime \prime }`$ implies that $`W_D^cW^{}`$ for any $`W^{}`$ such that $`W^{\prime \prime }W^{}W`$.
This result does not hold for normal default theories with preconditions, as the counterexample at the beginning of the section is only composed of normal defaults with preconditions.
###### Corollary 6
If $`D`$ is a set of normal and categorical defaults, then a formula is redundant if and only if it contains a redundant clause.
In other words, default logic restricted to the case of normal and categorical defaults has the local redundancy property.
### 4.3 Making Clauses Irredundant
Modifying a theory in order to make some parts irredundant proved useful for classical and circumscriptive logics. We show a similar result for default logic.
###### Definition 11
The $`M`$-irredundant version of a default theory $`D,W`$, where $`MW=\{\gamma _1,\mathrm{},\gamma _m\}`$, is the following theory, where $`\{c_1,\mathrm{},c_m\}`$ are new variables.
$`I(D,W,M)=D^{},W^{}`$
where:
$`W^{}`$ $`=`$ $`\{c_i\gamma _i|\gamma _iW\}`$
$`D^{}`$ $`=`$ $`D_1D_2D_3`$
$`D_1`$ $`=`$ $`\left\{{\displaystyle \frac{:\neg c_1\mathrm{}\neg c_m}{\neg c_1\mathrm{}\neg c_m}}\right\}`$
$`D_2`$ $`=`$ $`\left\{{\displaystyle \frac{c_i\gamma _i:c_i\{\neg c_i|1jk,ii\}}{c_i\{\neg c_i|1jk,ii\}}}\right|\gamma _iM\}`$
$`D_3`$ $`=`$ $`\left\{{\displaystyle \frac{\neg c_1\mathrm{}\neg c_m\alpha :\beta }{\gamma }}\right|{\displaystyle \frac{\alpha :\beta }{\gamma }}D\}`$
The clauses of $`M`$ are made irredundant by this transformation, while the redundancy of the other clauses does not change.
###### Theorem 11
If $`MW`$, $`D^{},W^{}=I(D,W,M)`$, and $`W^{\prime \prime }W`$, then $`W^{}_D^{}^eW^{\prime \prime }`$ holds if and only if $`\{\gamma _i|c_i\gamma _iW^{\prime \prime }\}_D^eW`$ and $`W^{\prime \prime }`$ contains all clauses $`c_i\gamma _i`$ such that $`\gamma _iM`$. The same holds for consequence-equivalence.
Proof. The default of $`D_1`$ can be applied provided that $`W`$ is consistent. Its application makes all defaults of $`D_2`$ inapplicable and makes the background theory and the defaults of $`D_3`$ equivalent to $`W`$ and to $`D`$, respectively. As a result, $`I(D,W,M)`$ has an extension $`Cn(E\neg c_1\mathrm{}\neg c_m)`$ for any extension $`E`$ of $`D,W`$. As a result, a subset of $`W^{}`$ has this extension if and only if the corresponding subset of $`W`$ has the same extension.
The $`i`$-th default of $`D_2`$ is applicable to $`W`$ because $`c_i\gamma _i`$ is in the background theory. Since all clauses of $`W`$ contain the literals $`c_j`$ only positively, these literals cannot be removed by resolution. As a result, every non-tautological consequence of $`W\backslash \{c_i\gamma _i\}`$ is disjoined with at least a variable $`c_j`$ with $`ji`$. As a result, no subset of $`W^{}`$ allows for the application of this default unless it contains the clause $`c_i\gamma _i`$.
The application of this default makes all other defaults of $`D^{}`$ inapplicable. The generated extension is moreover inconsistent with all other extensions of the theory. As a result, any subset of $`W^{}`$ not containing $`c_i\gamma _i`$ necessarily has a different set of extensions and consequences than $`W^{}`$
### 4.4 Complexity of Clause Redundancy
In this section, we analyze the complexity of checking the redundancy of a clause in a formula. Formally, this is the problem of whether $`W\backslash \{\gamma \}`$ is equivalent to $`W`$ according to $`_D^c`$ or $`_D^e`$. By Corollary 1, these two forms of equivalence are related, as $`W^{}_D^cW`$ is equivalent to $`W^{}_D^eW`$ and also to $`W^{}_DW`$, if $`W^{}W`$. As a result, checking whether $`W^{}_DW`$ allows for telling whether the “first part of equivalence” between $`W^{}`$ and $`W`$ holds, for both kinds of equivalence. In other words, in order to check whether $`W^{}`$ and $`W`$ are equivalent with $`W^{}W`$, we can first check whether $`W^{}_DW`$; if this condition is true, we then proceed checking whether $`W_D^cW^{}`$ or $`W_D^eW^{}`$ depending on which equivalence is considered.
Lemma 8 tells that $`W^{}_DW`$ implies that all processes of $`D,W^{}`$ are also processes of $`D,W`$. This condition does not imply equivalence because $`D,W`$ may contain some other processes, as in the default theory $`D,W`$ below.
$`W`$ $`=`$ $`\{a\}`$
$`D`$ $`=`$ $`\{d_1,d_2\}`$
where:
$`d_1={\displaystyle \frac{a:b}{b}}`$
$`d_2={\displaystyle \frac{:a\neg b}{a\neg b}}`$
The theory $`D,W`$ has two extensions: applying either $`d_1`$ or $`d_2`$, the other is not applicable. The resulting extensions are $`Cn(ab)`$ and $`Cn(a\neg b)`$. Let $`W^{}=\mathrm{}`$. The only default that is applicable in $`W^{}`$ is $`d_2`$, leading to the only extension $`Cn(a\neg b)`$. This extension implies $`W`$. As a result, we have that $`W^{}_DW`$ but $`W^{}`$ and $`W`$ do not have the same extensions and the same consequences. In particular, $`W`$ has some extensions that $`W^{}`$ does not have. This is always the case if $`W^{}_DW`$ but $`W`$ and $`W^{}`$ are not equivalent.
In order to check equivalence of $`W^{}`$ and $`W`$ with $`W^{}W`$, two conditions have to be checked:
1. $`W^{}_DW`$; and
2. $`W_D^cW^{}`$ or $`W_D^eW^{}`$.
An upper bound on the complexity of checking the redundancy of a clause is given by the following theorem.
###### Theorem 12
Checking whether $`W^{}_D^eW`$ in Reiter and justified default logic is in $`\mathrm{\Pi }_2^p`$ if $`W^{}W`$.
Proof. Checking whether $`W^{}_DW`$ is in $`\mathrm{\Pi }_2^p`$. The other condition to be checked is $`W_D^eW^{}`$. The converse of this condition is that there exists a formula $`EW\mathrm{cons}(D)`$ such that $`E`$ is an extension of $`D,W`$ but is not an extension of $`D,W^{}`$. Since checking whether a formula is a Reiter or justified default logic is in $`\mathrm{\Delta }_2^p[\mathrm{log}n]`$ \[Ros99\], the whole problem is in $`\mathrm{\Sigma }_2^p`$. Its converse is therefore in $`\mathrm{\Pi }_2^p`$. The problem of redundancy of a clause can be solved by solving two problems in $`\mathrm{\Pi }_2^p`$ in parallel.
The hardness of the problem for the same class is proved by the following theorem.
###### Theorem 13
Checking whether $`W^{}_D^eW`$ is $`\mathrm{\Pi }_2^p`$-hard even if $`W=W^{}\{a\}`$ and all defaults are categorical and normal.
Proof. The claim could be proved from the fact that entailment in default logic is $`\mathrm{\Pi }_2^p`$-hard even if the formula to entail is a single positive literal, and all defaults are categorical and normal \[Got92, Sti92\]. If all defaults are categorical and normal, Corollary 2 proves that $`W^{}_D^mW`$ is equivalent to the two other forms of equivalence.
We however use a new reduction from $``$QBF because this is required by the proof of $`\mathrm{\Sigma }_3^p`$-hardness of formula redundancy. The formula $`XY.F`$ is valid if and only if $`a`$ is redundant in the theory below:
$$\{\frac{:x_i}{x_i},\frac{:\neg x_i}{\neg x_i}\}\left\{\frac{:Fa}{a}\right\},\{a\}$$
This theory has an extension for every possible truth evaluation over the variables $`X`$. For each such extension, the last default can be applied only if $`F`$ is consistent with the given evaluation of $`X`$. As a result, if $`F`$ is consistent with every truth evaluation over the variables $`X`$, then $`a`$ can be removed from the background theory without changing the consequences of these extensions. Otherwise, the removal of $`a`$ would cause some of these extensions not to entail $`a`$ any more.
We now consider the problem of redundancy of clauses when consequence-equivalence is used. The difference between the two kinds of equivalence is that two sets of extensions may be different but yet their disjunctions are the same. The necessity of calculating the disjunction of all extensions intuitively explains why checking redundancy for consequence-equivalence is harder than for faithful equivalence.
###### Theorem 14
The problem of checking whether $`W^{}_D^cW`$ is in $`\mathrm{\Pi }_3^p`$ if $`W^{}W`$.
Proof. $`W^{}`$ and $`W`$ are consequence-equivalent if $`W^{}_DW`$ and $`W_D^cW^{}`$. The first problem is in $`\mathrm{\Pi }_2^p`$. We prove that the converse of the second condition is in $`\mathrm{\Sigma }_3^p`$. By definition, $`W\vDash ̸_D^cW^{}`$ holds if and only if $`\mathrm{Ext}_D(W)\vDash ̸\mathrm{Ext}_D(W^{})`$. In terms of models, we have $`\{Mod(E)|E\mathrm{Ext}_D(W)\}\{Mod(E)|E\mathrm{Ext}_D(W^{})\}`$, that is, there exists $`M`$ and $`E`$ such that $`MMod(E)`$, $`E\mathrm{Ext}_D(W)`$, but $`M`$ is not a model of any extension of $`W^{}`$. The whole condition can therefore be expressed by the following formula.
$$ME.MMod(E)E\mathrm{Ext}_D(W)(E^{}.E^{}\mathrm{Ext}_D(W^{})MMod(E^{}))$$
Since $`E^{}\mathrm{Ext}_D(W^{})`$ is in $`\mathrm{\Delta }_2^p[\mathrm{log}n]`$ for Reiter \[Ros99\] and justified default logic and in $`\mathrm{\Pi }_2^p`$ for constrained and rational \[Lib05a\], the problem of checking $`W_D^cW^{}`$ is in $`\mathrm{\Pi }_3^p`$. Therefore, the problem of consequence-equivalence is in $`\mathrm{\Pi }_3^p`$ as well for all four considered semantics.
We show that the problem is hard for the same class.
###### Theorem 15
The problem of checking whether $`W^{}_D^cW`$ for Reiter and justified default logics is $`\mathrm{\Pi }_3^p`$-hard even if $`W=W^{}\{a\}`$.
Proof. Since checking whether $`W^{}_D^eW`$ is in $`\mathrm{\Pi }_2^p`$, a proof of $`\mathrm{\Pi }_3^p`$-hardness necessarily requires the use of theories having different extensions but might have the same consequences.
We prove that the problem of non-equivalence of default theories is $`\mathrm{\Sigma }_3^p`$-hard by reduction from QBF. We reduce a formula $`XYZ.F`$ into the problem of checking whether $`W_D^cW^{}`$, where $`W^{}=\mathrm{}`$, $`W=\{a\}`$, and $`D=D_1D_2D_3D_4D_5D_6`$. We show each $`D_i`$ at time. First, we generate a complete evaluation over the variables $`X`$ using the following defaults.
$`D_1`$ $`=`$ $`\{{\displaystyle \frac{:x_i}{x_ih_i}},{\displaystyle \frac{:\neg x_i}{\neg x_ih_i}}\}`$
Since these defaults have no preconditions, they can be applied regardless of whether $`W`$ or $`W^{}`$ is the background theory. They generate a process for any truth evaluation $`\omega _X`$ over the variables in $`X`$. The variables $`h_i`$ are all true only when all variables $`x_i`$ have been set to a value.
The processes of $`W`$ and $`W^{}`$ are so far the same. Once all $`h_i`$ are true, we can apply the defaults of $`D_2=\{d_1,d_2,d_3,d_4\}`$, which are the ones used in Theorem 6 to show two theories that have the same consequences but different extensions:
$`D_2`$ $`=`$ $`\{d_1,d_2,d_3,d_4\}`$
$`d_1`$ $`=`$ $`{\displaystyle \frac{h_1\mathrm{}h_n:\neg bc}{ac}}`$
$`d_2`$ $`=`$ $`{\displaystyle \frac{h_1\mathrm{}h_n:\neg b\neg c}{a\neg c}}`$
$`d_3`$ $`=`$ $`{\displaystyle \frac{h_1\mathrm{}h_na:bc}{bc}}`$
$`d_4`$ $`=`$ $`{\displaystyle \frac{h_1\mathrm{}h_na(bc):b\neg c}{b}}`$
Since these default have $`h_1\mathrm{}h_n`$ as a precondition, they can only be applied once a truth assignment over $`X`$ has been generated by the previous defaults. They act as in the proof of Theorem 6. Only $`[d_1d_3]`$ and $`[d_2]`$ are processes of $`W^{}`$; their consequences are $`ac`$ and $`a\neg c`$. The theory $`W`$ has the same processes, but also $`[d_3d_1]`$ and $`[d_3d_4]`$, which generate the extensions $`ac`$ and $`ab`$, respectively. While the first is also an extension of $`W^{}`$, the second is not. The disjunction of all extensions is equivalent to $`a`$ for both $`W`$ and $`W^{}`$.
The idea is as follows: from $`ab\omega _X`$, which is obtained from $`W`$ but not from $`W^{}`$, we always generate the extension $`ab\omega _Xdϵ_Y`$, where $`ϵ_Y`$ is the assignment of $`\mathrm{𝖿𝖺𝗅𝗌𝖾}`$ to all variables $`Y`$; from the two other points $`a\neg c\omega _X`$ and $`ac\omega _X`$ we instead generate an arbitrary assignment $`\omega _Y`$, which then has $`ab\omega _Xdϵ_Y`$ as a model only if $`F`$ is satisfiable.
This way, if there exists a value $`\omega _X`$ such that for all $`\omega _Y`$ the formula $`F`$ is satisfiable, then there is no extension of $`W^{}`$ having the model $`abd\omega _Xϵ_Y`$. Vice versa, if there exists even a single $`\omega _Y`$ such that $`F`$ is unsatisfiable, an extension $`ac\omega _X\mathrm{}`$ for $`W^{}`$ will be generated, and this extension has the model $`ab\omega _Xdϵ_Y`$.
The required defaults are the following ones. First, we generate the considered model from the process that has generated $`ab\omega _X`$:
$$D_3=\left\{\frac{b:}{d\neg y_1\mathrm{}\neg y_n}\right\}$$
From $`a\neg c\omega _X`$ and $`ac\omega _X`$ we generate an arbitrary truth evaluation over $`Y`$. Since the model $`abd\omega _Xϵ_Y`$ assigns false to all variables $`y_i`$, we cannot simply add $`y_i`$ as a conclusion. A similar effect can be achieved by the following defaults.
$`D_4`$ $`=`$ $`\left\{{\displaystyle \frac{\neg c:\neg dy_i}{d(y_il_i)}},{\displaystyle \frac{\neg c:\neg d\neg y_i}{d(\neg y_il_i)}}\right|1in\}`$
$`D_5`$ $`=`$ $`\left\{{\displaystyle \frac{c:\neg dy_i}{d(y_il_i)}},{\displaystyle \frac{c:\neg d\neg y_i}{d(\neg y_il_i)}}\right|1in\}`$
The two defaults associated with $`y_i`$ and $`\neg y_i`$ cannot be applied both at the same time, as the consequence of one contains the negation of the justification of the other one. Since the following defaults can only be applied when $`d(l_1\mathrm{}l_n)`$ has been derived, the current extensions before their application are $`a\neg c\omega _X(d(\omega _YL))`$ and $`ac\omega _X(d(\omega _YL))`$, where $`\omega _Y`$ is an arbitrary truth assignment over $`Y`$.
These extensions have all models of $`ab\omega _Xdϵ_Y`$. The following default removes these models from the extensions if and only if $`F`$ is satisfiable for these given assignments over $`X`$ and $`Y`$.
$$D_6=\left\{\frac{d(l_1\mathrm{}l_n):\neg dF}{\neg d}\right\}$$
This default is not applicable from $`ab\omega _Xdϵ_Y`$ because its justification contains $`\neg d`$. It is applicable from the other processes but only after the $`i`$-th default of $`D_4`$ or $`D_5`$ has been applied for each $`i`$ and only if the consequences of the applied defaults of $`D_4`$ or $`D_5`$ are consistent with $`\neg dF`$. In other words, $`(d(\omega _YL))\neg dF`$ must be consistent, which is equivalent to the consistency of $`\omega _YF`$ because $`d`$ and $`L`$ are not mentioned in $`\omega _Y`$ and $`F`$.
We can therefore conclude that:
1. for each truth assignment $`\omega _X`$, three “partial extensions” are generated from $`W`$: $`a\neg c\omega _X`$, $`ac\omega _X`$, and $`ab\omega _X`$; the first two ones are also generated by $`W^{}`$;
2. from $`ab\omega _X`$, the extension $`ab\omega _Xdϵ_Y`$ is generated; if the models of this extension are not models of an extension of $`W^{}`$, equivalence between $`W`$ and $`W^{}`$ does not hold;
3. from $`a[\neg ]c\omega _X`$ we generate $`a[\neg ]c\omega _X(d(\omega _YL))`$ for each truth evaluation $`\omega _Y`$ on the variables $`Y`$; to this formula, $`\neg d`$ is added if and only if $`F`$ is consistent with $`\omega _X`$ and $`\omega _Y`$.
As a result, the models of $`ab\omega _Xdϵ_Y`$ are not models of an extension of $`W^{}`$ if and only if $`F\omega _X`$ is satisfiable for every truth evaluation of $`Y`$. Since non-equivalence has to be checked for every $`\omega _X`$, we have that non-equivalence holds if and only if $`XYZ.F`$
A similar proof holds for constrained or rational default logics by replacing the default theory of Theorem 6 with that of Theorem 7. The proof can also slightly simplified in this case, as the defaults of $`D_4`$ and $`D_5`$ can be modified with justifications $`y_i`$ or $`\neg y_i`$ and consequence $`dl_i`$.
Since we have proved that the problem of clause redundancy w.r.t. consequence-equivalence is both in $`\mathrm{\Pi }_3^p`$ and hard for the same class, we have the following theorem.
###### Theorem 16
The problem of checking whether $`W^{}_D^cW`$ is $`\mathrm{\Pi }_3^p`$-complete if $`W^{}W`$; hardness holds even if $`W=W^{}\{a\}`$.
### 4.5 Complexity of Formula Redundancy
The next problem to analyze is whether a formula (a set of clauses) is redundant, for a fixed set of defaults. The complexity of formula redundancy w.r.t. faithful and consequence-equivalence is in $`\mathrm{\Sigma }_3^p`$ and $`\mathrm{\Sigma }_4^p`$, respectively.
###### Theorem 17
The problem of formula redundancy for faithful and consequence-equivalence is in $`\mathrm{\Sigma }_3^p`$ and $`\mathrm{\Sigma }_4^p`$, respectively.
Proof. Both problems can be expressed as the existence of a subset $`W^{}W`$ such that $`W^{}`$ is equivalent to $`W`$. Since equivalence is in $`\mathrm{\Pi }_2^p`$ and $`\mathrm{\Pi }_3^p`$, respectively, for faithful and consequence-equivalence, the claim follows.
Regarding hardness, we first show a theorem characterizing the complexity of the problem for the case of faithful equivalence. We then show a more general technique allowing an hardness result to be raised one level in the polynomial hierarchy.
###### Theorem 18
The problem of formula redundancy based on faithful equivalence is $`\mathrm{\Sigma }_3^p`$-hard.
Proof. We reduce the problem of validity of $`XYZ.F`$ to the problem of redundancy of a formula. Let $`n=|X|`$. The default theory corresponding to the formula $`XYZ.F`$ is the theory $`D,W`$ defined as follows.
$`W`$ $`=`$ $`\{s_i|1in\}\{r_i|1in\}`$
$`D`$ $`=`$ $`D_1D_2D_3D_4D_5D_6`$
$`D_1`$ $`=`$ $`\left\{{\displaystyle \frac{s_ir_i:\neg s_j}{a}},{\displaystyle \frac{s_ir_i:\neg r_j}{a}}\right|{\displaystyle \genfrac{}{}{0pt}{}{1in}{1jn}}\}`$
$`D_2`$ $`=`$ $`\left\{{\displaystyle \frac{:\neg s_i\neg r_i}{a}}\right|{\displaystyle \genfrac{}{}{0pt}{}{1in}{1jn}}\}`$
$`D_3`$ $`=`$ $`\left\{{\displaystyle \frac{:y_i}{y_ih_i}},{\displaystyle \frac{:\neg y_i}{\neg y_ih_i}}\right|1in\}`$
$`D_4`$ $`=`$ $`\left\{{\displaystyle \frac{:x_i}{p_ix_i}},{\displaystyle \frac{:\neg x_i}{p_i\neg x_i}}\right|1in\}`$
$`D_5`$ $`=`$ $`\left\{{\displaystyle \frac{x_ir_i:}{W}},{\displaystyle \frac{\neg x_is_i:}{W}}\right|1in\}`$
$`D_6`$ $`=`$ $`\left\{{\displaystyle \frac{p_1\mathrm{}p_nh_1\mathrm{}h_n:F}{W}}\right\}`$
The defaults of $`D_1`$ and $`D_2`$ cannot be applied from $`W`$. The defaults of $`D_3`$ and $`D_4`$ generates an extension for every possible truth evaluation over $`XY`$; this extension also contains all variables $`h_i`$ and $`p_i`$. Whether or not the last default is applicable, its consequence is equivalent to the background theory.
Let $`W^{}W`$. If there is an index $`i`$ such that both $`s_i`$ and $`r_i`$ are in $`W^{}`$, one of the defaults of $`D_1`$ is applicable, generating $`a`$. Therefore, $`W^{}`$ is not equivalent to $`W`$. If there exists an index $`i`$ such that neither $`s_i`$ nor $`r_i`$ is in $`W^{}`$, the $`i`$-th default of $`D_2`$ is applicable, still generating $`a`$.
In order to check for redundancy, we therefore only have to consider subsets $`W^{}W`$ for which either $`s_iW`$ or $`r_iW`$ but not both. Let $`\omega _X`$ be the assignment on the variables $`X`$ such that $`x_i`$ is assigned to true if $`s_iW^{}`$ and to false if $`r_iW^{}`$. The defaults of $`D_3`$ and $`D_4`$ generate an arbitrary truth evaluation of the variables $`XY`$. If the assignment on $`X`$ is not equal to $`\omega _X`$, the formula $`W`$ is generated, thus leading to an extension that is also an extension of $`W`$. As a result, all extensions of $`W^{}`$ that do not match the value $`\omega _X`$ are also extensions of $`W`$. If the same holds also for the extensions for which the values of $`X`$ match $`\omega _X`$, then $`W^{}`$ is equivalent to $`W`$.
For a given $`W^{}`$ we consider the extensions consistent with $`\omega _X`$. There is exactly one such extension for each possible truth evaluation over $`Y`$. If the default of $`D_6`$ can be applied, it generates $`W`$, thus making $`W^{}`$ equivalent to $`W`$. In turn, the default of $`D_6`$ can be applied for all truth evaluation over $`Y`$ if and only if for all such truth evaluation, $`F`$ is satisfiable. As a result, $`W^{}`$ is equivalent to $`W`$ if and only if, for all possible truth evaluations over $`Y`$, the formula $`F`$ is satisfiable. Since there exists a relevant $`W^{}`$ for each truth evaluation over $`X`$, the formula $`W`$ is redundant if and only if there exists a truth evaluation over $`X`$ such that, for all possible truth evaluations over $`Y`$, the formula $`F`$ is satisfiable.
In order to prove the $`\mathrm{\Sigma }_4^p`$-hardness of the problem of formula redundancy under consequence-equivalence, we should provide a reduction from $``$QBF validity into this problem. A simpler proof can however be given, based on the following consideration: checking clause redundancy has been proved $`\mathrm{\Pi }_2^p`$-hard or $`\mathrm{\Pi }_3^p`$-hard using reductions from QBFs that results in default theories having $`W=\{a\}`$ as the background theory. As a result, these reductions also prove that formula redundancy is $`\mathrm{\Pi }_2^p`$-hard or $`\mathrm{\Pi }_3^p`$-hard. In other words, we can reduce the validity of a $``$QBF or a $``$QBF into the problems of formula redundancy. What we show is that, if such reductions satisfy some assumptions, we can obtain new reductions from QBFs having an additional existential quantifier in the front. The assumptions are that the default theory resulting from the reduction is such that:
1. the background theory that results from the reduction is classically irredundant;
2. the matrix of the QBF is only used in the justification of a single default.
The reductions used for proving the hardness of clause redundancy satisfies both assumptions. In particular, $`XYZ.F`$ is valid if and only if the background theory of the following theory is consequence-redundant, where $`D`$, $`\alpha `$, $`\beta `$, $`\gamma `$, do not depend on $`F`$ but only on the quantifiers of the QBF and $`W`$ is classically irredundant.
$$D\left\{\frac{\alpha :\beta F}{\gamma }\right\},W$$
The fact that the matrix of the QBF is copied “verbatim” in the default theory is exploited as follows: if $`\omega _w`$ is a truth evaluation over the variable $`w`$, then $`XYZ.F|_{\omega _w}`$ is valid if and only if the background theory of $`D\{\frac{\alpha :\beta F\omega _w}{\gamma }\},W`$ is redundant. This default theory can be modified in such a way the subsets of the background theory are in correspondence with the truth evaluations over $`\omega _w`$. This way, the resulting theory is redundant if and only if $`wXYZ.F`$. The resulting default theory still satisfies the two assumptions above on the background theory and on the use of the matrix of the QBF; therefore, this procedure can be iterated to obtain a reduction from $``$QBF validity into the problem of formula redundancy under consequence-equivalence. A similar technique can be used for faithful equivalence.
The details of this technique are in the following three lemmas. The first one shows that a literal can be moved from the justification of a default to the background theory and vice versa, under certain conditions.
###### Lemma 13
If the variable of the literal $`l`$ is not mentioned in $`W`$, $`D`$, $`\mathrm{prec}(d)`$, and $`\mathrm{cons}(d)`$, the processes of the following two theories are the same modulo the replacement of $`d`$ with $`d^{}`$ and vice versa.
$`D\{d\},W\{l\}`$
$`D\{d^{}\},W`$
where
$`d^{}={\displaystyle \frac{\mathrm{prec}(d):\mathrm{just}(d)l}{\mathrm{cons}(d)}}`$
Proof. The literal $`l`$ and its negation only occur in the background theory $`W\{l\}`$ and in the justification of $`d`$ and $`d^{}`$. The conditions on a process of the first theory being selected either involve $`(W\{l\})\mathrm{just}(d)`$ or $`W\{l\}`$ with other formulae not containing $`l`$. As a result, moving $`l`$ from the background theory to the justification of $`d`$ or vice versa does not change these conditions.
Note that the processes are the same, but the extensions are different in that $`l`$ is in all extensions of the first theory but not in the extensions of the second.
The second lemma is an obvious consequence of the above: under the same conditions, moving a literal from the justification of a default to the background theory or vice versa does not change the redundancy of a theory.
###### Lemma 14
If $`W^{}W`$, it holds $`W^{}_D^{}^eW`$ if and only if $`W^{}\{l\}_{D^{\prime \prime }}^eW\{l\}`$, where $`l`$ is a literal that is not mentioned in $`W`$, $`D`$, $`\mathrm{prec}(d)`$, and $`\mathrm{cons}(d)`$, where $`D^{}`$ and $`D^{\prime \prime }`$ are as follows.
$`D^{}`$ $`=`$ $`D\left\{{\displaystyle \frac{\mathrm{prec}(d):\mathrm{just}(d)l}{\mathrm{cons}(d)}}\right\}`$
$`D^{\prime \prime }`$ $`=`$ $`D\{d\}`$
Proof. Obvious consequence of the lemma above: $`D^{},W^{}`$ and $`D^{},W`$ have the same processes of $`D^{\prime \prime },W^{}\{l\}`$ and $`D^{\prime \prime },W\{l\}`$, respectively.
A consequence of this lemma is that $`W`$ is redundant according to $`D^{}`$ if and only if $`W\{l\}`$ is redundant according to $`D^{\prime \prime }`$. Indeed, $`l`$ is not mentioned in the consequences of the defaults; therefore, a subset of $`W\{l\}`$ can only be equivalent to $`W\{l\}`$ if it contains $`l`$. The lemma is formulated in the more complicated way because it is necessary for proving the following lemma. The same property can be proved using consequence-equivalence because moving $`l`$ from the justification of the default to the background theory has the only effect of adding $`l`$ to all extensions.
###### Lemma 15
If $`W`$ is classically irredundant, then there exists $`W^{}W`$ such that $`W^{}\{w\}_D^eW\{w\}`$ or $`W^{}\{\neg w\}_D^eW\{\neg w\}`$ if and only if the following theory is redundant:
$`D_w\left\{{\displaystyle \frac{p\alpha :\beta }{\gamma }}\right|{\displaystyle \frac{\alpha :\beta }{\gamma }}D\},W\{w^+,w^{}\}`$
where:
$`D_w=\{{\displaystyle \frac{w^+w^{}:\neg W}{\neg p}},{\displaystyle \frac{:\neg w^+\neg w^{}}{\neg p}},{\displaystyle \frac{w^+:wp}{wp}},{\displaystyle \frac{w^{}:\neg wp}{\neg wp}}\}`$
and $`w^+`$, $`w^{}`$, and $`p`$ are new variables.
Proof. Since $`w^+`$ and $`w^{}`$ are new variables not contained in $`W`$ and $`W`$ is classically irredundant, $`W\{w^+,w^{}\}`$ is classically irredundant as well.
We now consider the processes that can be generated from $`W\{w^+,w^{}\}`$ and from its subsets. From $`W\{w^+,w^{}\}`$ we can apply only one of the last two defaults of $`D_w`$, generating either $`wp`$ or $`\neg wp`$. From this point on, we have exactly the same processes of $`D,W\{w\}`$ and $`D,W\{\neg w\}`$, the generated extensions only differing because of the addition of $`p`$ and $`w^+`$ or $`w^{}`$.
The proper subsets of $`W\{w^+,w^{}\}`$ are $`W^{}\{w^+,w^{}\}`$ where $`W^{}W`$, $`W^{}\{w^+\}`$, $`W^{}\{w^{}\}`$, and $`W^{}`$, where $`W^{}W`$. The fourth subset $`W^{}`$ is not equivalent to $`W`$ because the second default of $`D_w`$ allows the derivation of $`\neg p`$, which is not derivable from $`W`$. If $`W^{}W`$, since $`W`$ is (classically) irredundant, $`W^{}\{w^+,w^{}\}`$ allows for the application of the first default of $`D_w`$, deriving $`\neg p`$; therefore, this subset is not equivalent to the background theory.
The only two other subsets to consider are $`W^{}\{w^+\}`$ and $`W^{}\{w^{}\}`$. In the first subset, only $`wp`$ can be generated. In the second subset, only $`\neg wp`$ can be generated. From this point on, we have exactly the same processes of $`W^{}\{w\}`$ and $`W^{}\{\neg w\}`$ according to $`D`$. The generated extensions are the same but for the addition of $`p`$
These three lemmas together proves that a reduction from QBF to formula redundancy can be “raised” by the addition of an existential quantifier in the front of the QBF.
###### Lemma 16
If there exists a polynomial reduction from formulae $`Q.E`$, where $`Q`$ is a sequence of quantifier of a given class, to the problem of formula redundancy of a default theory in the following form, then there exists a polynomial reduction from formulae of the form $`wQ.F`$ to the formula redundancy of a theory in the following form.
$$D\left\{\frac{\alpha :\beta F}{\gamma }\right\},W$$
The formulae in $`D`$, $`\{\alpha ,\beta ,\gamma \}`$, and $`W`$ do not depend on $`F`$. The background theory $`W`$ is classically irredundant.
Proof. Let $`Q`$ be a sequence of quantifiers so that the validity of the formula $`Q.E`$ can be reduced to formula redundancy of a default theory of the above form. We show a reduction from the validity of $`wQ.F`$ to formula redundancy of a default theory of the same form.
By definition, both $`Q.F|_{w=\mathrm{𝗍𝗋𝗎𝖾}}`$ and $`Q.F|_{w=\mathrm{𝖿𝖺𝗅𝗌𝖾}}`$ can be reduced to the problem of formula redundancy. These two formulae only differ on their matrixes, which are $`F|_{w=\mathrm{𝗍𝗋𝗎𝖾}}`$ and $`F|_{w=\mathrm{𝖿𝖺𝗅𝗌𝖾}}`$. Therefore, the resulting default theories are:
$`D\left\{{\displaystyle \frac{\alpha :\beta F|_{w=\mathrm{𝗍𝗋𝗎𝖾}}}{\gamma }}\right\},W`$
$`D\left\{{\displaystyle \frac{\alpha :\beta F|_{w=\mathrm{𝖿𝖺𝗅𝗌𝖾}}}{\gamma }}\right\},W`$
Since $`w`$ does not occur anywhere else in the theory, we can replace $`F|_{w=\mathrm{𝗍𝗋𝗎𝖾}}`$ and $`F|_{w=\mathrm{𝖿𝖺𝗅𝗌𝖾}}`$ with $`Fw`$ and $`F\neg w`$, respectively. Indeed, justifications are only checked for consistency, and for any formula $`R`$ not containing $`w`$, the consistency of $`RF|_{w=\mathrm{𝗍𝗋𝗎𝖾}}`$ is the same as the consistency of $`R(Fw)`$, and the consistency of $`RF|_{w=\mathrm{𝖿𝖺𝗅𝗌𝖾}}`$ is the same as the consistency of $`R(F\neg w)`$.
$`D\left\{{\displaystyle \frac{\alpha :\beta Fw}{\gamma }}\right\},W`$
$`D\left\{{\displaystyle \frac{\alpha :\beta F\neg w}{\gamma }}\right\},W`$
By Lemma 14, formula redundancy of these two theories corresponds to formula redundancy of the same theories with $`w`$ or $`\neg w`$ moved to the background theory. More precisely, the redundancy of the first theory correspond to the existence of a subset $`W^{}W`$ such that $`W^{}\{w\}`$ is equivalent to $`W\{w\}`$ according to the defaults $`D\left\{\frac{\alpha :\beta F}{\gamma }\right\}`$. The same holds for the second theory.
$`D\left\{{\displaystyle \frac{\alpha :\beta F}{\gamma }}\right\},W\{w\}`$
$`D\left\{{\displaystyle \frac{\alpha :\beta F}{\gamma }}\right\},W\{\neg w\}`$
By Lemma 15, since $`W`$ is classically irredundant, we have that either the first or the second of the two theories are redundant if and only if the following theory is redundant:
$$D_w\left\{\frac{p\alpha :\beta }{\gamma }\right|\frac{\alpha :\beta }{\gamma }D\},W\{w^+,w^{}\}$$
where $`D_w`$ is defined in the statement of Lemma 15. As a result, this formula is redundant if and only if either $`Q.F|_{w=\mathrm{𝗍𝗋𝗎𝖾}}`$ is valid or $`Q.F|_{w=\mathrm{𝖿𝖺𝗅𝗌𝖾}}`$ is valid, that is, $`wQ.F`$ is valid.
In order to complete the lemma, we have to show that the background theory of the above theory is classically irredundant, and the theory is in the form specified by the statement of the theorem. Since $`W`$ is classically irredundant by assumption and $`w^+`$ and $`w^{}`$ are new variables, $`W\{w^+,w^{}\}`$ is classically irredundant. In the above theory, the matrix $`F`$ of the QBF is only mentioned in the justification of the default $`\frac{p\alpha :\beta F}{\gamma }`$. Therefore, the theory that results from the reduction is in the form specified by the theorem.
The above lemmas are also valid for consequence-equivalence. In both cases, we have that the hardness of formula redundancy is one level higher in the polynomial hierarchy than clause redundancy.
###### Theorem 19
Formula redundancy is $`\mathrm{\Sigma }_3^p`$-hard for faithful equivalence and $`\mathrm{\Sigma }_4^p`$-hard for consequence-equivalence.
Proof. The reduction shown after Theorem 13 and the reduction used in Theorem 15 are reductions from $``$QBF and $``$QBF, respectively, into the problem of formula redundancy. These reductions produce a default theory in which the background theory contains a single non-tautological clause, and is therefore irredundant, and the matrix of the QBF only occurs in the justification of a single default. These are the conditions of Lemma 16. As a result, one can reduce an $``$QBF or an $``$QBF to the problem of formula redundancy by iteratively applying the modification of Lemma 16 for all variables of the first existential quantifier.
### 4.6 Redundancy of Defaults
The redundancy of a default is defined in the same way as redundancy of clauses.
###### Definition 12 (Redundancy of a default)
A default $`d`$ is redundant in $`D,W`$ if and only if $`D\backslash \{d\},W`$ is equivalent to $`D,W`$.
This definition depends on the kind of equivalence used. Therefore, a default can be redundant w.r.t. faithful or consequence-equivalence. The redundancy of defaults is defined as follows.
###### Definition 13 (Redundancy of a theory)
A default theory $`D,W`$ is default redundant if and only if there exists $`D^{}D`$ such that $`D^{},W`$ is equivalent to $`D,W`$.
#### 4.6.1 Making Defaults Irredundant
The following lemma is the version of Theorem 11 to the case of default redundancy rather than clause redundancy. It proves that some defaults can be made irredundant while not changing the redundancy status of the other ones.
###### Lemma 17
For every default theory $`D,W`$, set of defaults $`D_ID`$, and $`D_1`$, $`D_2`$, $`D_3`$ defined as follows:
$`D_1`$ $`=`$ $`\{d+,d\}`$
where:
$`d+={\displaystyle \frac{:pq}{pq}}`$
$`d={\displaystyle \frac{:\neg pq}{\neg pq}}`$
$`D_2`$ $`=`$ $`\left\{{\displaystyle \frac{q(\neg p\alpha ):\neg p\beta }{(pv_i)(\neg p\gamma )}}\right|{\displaystyle \frac{\alpha :\beta }{\gamma }}D_I\}`$
$`D_3`$ $`=`$ $`\left\{{\displaystyle \frac{qp\alpha :\beta }{\gamma }}\right|{\displaystyle \frac{\alpha :\beta }{\gamma }}D\backslash D_I\}`$
if $`D,W`$ has extensions and $`W`$ is consistent, it holds that:
1. the processes of $`D_1D_2D_3,W`$ are (modulo the transformation of the defaults) the same of $`D,W`$ with $`d+`$ added to the front and a number of processes composed of $`d`$ and a sequence containing all defaults of $`D_2`$;
2. the extension of $`D_1D_2D_3,W`$ are the same of $`D,W`$ with $`\{p,q\}`$ added plus the single extension $`\{\neg p,q\}\{v_i\}`$;
3. a subset of $`D_1D_2D_3`$ is equivalent to it if and only if it contains $`D_1D_2`$ and the set of original defaults corresponding to those of $`D_2D_3`$ is equivalent to $`D`$.
Proof. Since all defaults of $`D_2D_3`$ have $`q`$ as a precondition, they are not applicable from $`W`$. The only defaults that are applicable to $`W`$ are therefore $`d+`$ and $`d`$, which are mutually exclusive.
Let us consider the processes with $`d`$ in first position. Since $`d`$ generates $`\neg p`$, the defaults of $`D_3`$ are not applicable. We prove that $`[d]\mathrm{\Pi }_2`$ is a successful process, where $`\mathrm{\Pi }_2`$ is an arbitrary sequence containing all defaults of $`D_2`$. The preconditions of all defaults of $`D_2`$ are entailed by $`q\neg p`$. The union of the justifications and consequences of all defaults of this process is $`\{\neg p,q\}\{\neg p\beta ,pv_i,\neg p\gamma \}`$, which is equivalent to $`\{\neg p,q\}\{v_i\}`$. This set is consistent with the background theory, which does not contain the variables $`p`$, $`q`$, and $`v_i`$.
If a subset of $`D_1D_2D_3`$ does not contain $`d`$, the literal $`\neg q`$ cannot derived because no other default has $`\neg q`$ as a conclusion. If a subset of $`D_1D_2D_3`$ does not contain a default of $`D_2`$, the corresponding variable $`v_i`$ is not in this extension. As a result, every subset of $`D_1D_2D_3`$ that is equivalent to it contains $`\{d\}D_2`$.
Let us now consider the processes with $`d+`$ in first position. Such a process cannot contain $`d`$. Since $`p`$ and $`q`$ are generated, the defaults of $`D_2D_3`$ can be simplified to $`\frac{\alpha :\beta }{\gamma }`$ by removing all clauses containing $`p`$ or $`q`$ and all literals $`\neg p`$ and $`\neg q`$ from the clauses containing them. As a result, the processes having $`d+`$ in first position correspond to the processes of the original theory.
Provided that the original theory has extensions, every subset of $`D_1D_2D_3`$ not containing $`d+`$ lacks these extensions. The defaults of $`D_3`$ are redundant if and only if they are redundant in the original theory. More precisely, a subset $`D^{}D_1D_2D_3`$ is equivalent to $`D_1D_2D_3`$ if and only if $`D^{}`$ contains $`D_1D_2`$, and the set of original defaults $`D^{\prime \prime }`$ corresponding to the defaults of $`D^{}(D_2D_3)`$ is equivalent to $`D`$
#### 4.6.2 Redundancy of Defaults vs. Sets of Defaults
While a formula is classically redundant if and only if it contains a redundant clause, the same does not happen for default redundancy. The following theorem indeed proves that Reiter and rational default logic do not have the local redundancy property w.r.t. redundancy of defaults.
###### Theorem 20
There exists a set of defaults $`D`$ such that, according to Reiter and rational default logic:
1. for any $`dD`$, the theory $`D\backslash \{d\},\mathrm{}`$ has extensions and $`D\backslash \{d\},W_D^cD,\mathrm{}`$;
2. there exists $`D^{}D`$ such that $`D^{},\mathrm{}_D^eD,\mathrm{}`$.
Proof. We use a pair of defaults that lead to failure is they are together in the same process. Removing one of them from the default theory leads to a new extension, while removing both of them lead to the original set of extensions. The following defaults are a realization of this idea.
$`D`$ $`=`$ $`\{d_1,d_2,d_3\}`$
where:
$`d_1={\displaystyle \frac{:b}{bc}}`$
$`d_2={\displaystyle \frac{:b}{b\neg c}}`$
$`d_3={\displaystyle \frac{:\neg b}{\neg b}}`$
The extensions of some $`D^{},\mathrm{}`$, with $`D^{}D`$, are as follows:
we can either apply both $`d_1`$ and $`d_2`$ (leading to a failure) or $`d_3`$ alone; the only extension of this theory is therefore $`\neg b`$;
both $`d_1`$ and $`d_3`$ can be applied, but not both; that results in two processes having conclusions $`\mathrm{Ext}(D,\mathrm{})=(b\neg c)\neg b`$;
same as above: $`\mathrm{Ext}(D,\mathrm{})=(bc)\neg b`$;
the only selected process is $`[d_3]`$, which leads to $`\mathrm{Ext}(D^{},W)=\neg b`$.
As a result, $`D\backslash \{d\},\mathrm{}`$ has extensions for every $`dD`$. The default $`d_3`$ is not irredundant, but can be made so by the transformation of Theorem 17, which preserves processes almost exactly; an alternative is to replace $`d_3`$ with $`\frac{:\neg b}{bd}`$ and $`\frac{:\neg b}{be}`$. The resulting set of default has no redundant default, but has an equivalent subset.
The same result holds for constrained default logic.
###### Theorem 21
There exists a set of defaults $`D`$ such that, according to constrained default logic:
1. for any $`dD`$, it holds $`D\backslash \{d\},W_D^cD,\mathrm{}`$;
2. there exists $`D^{}D`$ such that $`D^{},\mathrm{}_D^eD,\mathrm{}`$.
Proof. The defaults are the following ones:
$`D`$ $`=`$ $`\{d_1,d_2,d_3\}`$
where:
$`d_1={\displaystyle \frac{:x}{a}}`$
$`d_2={\displaystyle \frac{:x}{b}}`$
$`d_3={\displaystyle \frac{:\neg x\neg y}{ab}}`$
The theory $`D,\mathrm{}`$ has two selected processes (modulo permutation of defaults): $`[d_1,d_2]`$ and $`[d_3]`$, both generating the extension $`ab`$. Removing either $`d_1`$ or $`d_2`$ causes the first process to become $`[d_1]`$ or $`[d_2]`$, thus creating a new extension that is either $`a`$ or $`b`$. On the other hand, removing both $`d_1`$ and $`d_2`$ makes the only remaining process to be $`[d_3]`$, which generates the only extension $`ab`$ of the original theory. The default $`d_3`$ is not redundant, but can be made so by applying the transformation of Theorem 17
Justified default logic has the local redundancy property w.r.t. default redundancy. This is a combination of two factors: first, justified default logic is failsafe \[Lib05e\] (every successful process can be made selected by adding some defaults); second, every extension is generated by an unique set of defaults. The proofs requires two lemmas. The first one is about extendibility of processes when new defaults are added to a theory.
###### Lemma 18
In justified default logic, if $`\mathrm{\Pi }`$ is a selected process of $`D^{},W`$ and $`D^{}D`$, then there exists a sequence $`\mathrm{\Pi }^{}`$ of defaults of $`D\backslash D^{}`$ such that $`\mathrm{\Pi }\mathrm{\Pi }^{}`$ is a selected process of $`D,W`$.
Proof. Let $`\mathrm{\Pi }`$ be a selected process of $`D^{},W`$. By definition, it holds $`W\mathrm{cons}(\mathrm{\Pi }[d])\mathrm{prec}(d)`$ and $`W\mathrm{cons}(\mathrm{\Pi })\mathrm{just}(d)`$ for every $`d\mathrm{\Pi }`$. As a result, $`\mathrm{\Pi }`$ is a also a successful process of $`D,W`$. Therefore, there exists $`\mathrm{\Pi }^{}`$ such that $`\mathrm{\Pi }\mathrm{\Pi }^{}`$ is a selected process of $`D,W`$ because justified default logic is failsafe \[Lib05e\]. If $`\mathrm{\Pi }^{}`$ contains defaults of $`D^{}`$, then $`\mathrm{\Pi }`$ would not be a closed process of $`D^{},W`$
In order for proving the second lemma, we need an intermediate result, which is already well known.
###### Lemma 19
In justified default logic, the selected processes of $`D,W`$ generating the extension $`E`$ are composed of exactly the defaults of the following set:
$$GEN(E,D)=\{dD|E\mathrm{prec}(d)\text{ and }E\mathrm{just}(d)\mathrm{cons}(d)\}$$
Proof. Assume that $`\mathrm{\Pi }`$ is a selected process generating $`E`$ that does not contain a default $`dGEN(E,D)`$. Since $`E\mathrm{prec}(d)`$, $`E\mathrm{just}(d)\mathrm{cons}(d)`$, and $`E=W\mathrm{cons}(\mathrm{\Pi })`$, we have that $`W\mathrm{cons}(\mathrm{\Pi })\mathrm{prec}(d)`$ and $`W\mathrm{cons}(\mathrm{\Pi })\mathrm{just}(d)\mathrm{cons}(d)`$. As a result, $`\mathrm{\Pi }[d]`$ is a successful process, contradicting the assumption.
Let $`\mathrm{\Pi }`$ be a selected process containing a default $`d`$ not in $`GEN(E,D)`$. By definition, either $`E\vDash ̸\mathrm{prec}(d)`$ or $`E\mathrm{just}(d)\mathrm{cons}(d)`$. The first condition implies that $`W\mathrm{cons}(\mathrm{\Pi }[d])\vDash ̸d`$ whichever the position of $`d`$ in $`\mathrm{\Pi }`$ is. The second condition implies $`W\mathrm{cons}(\mathrm{\Pi })\mathrm{just}(d)\mathrm{cons}(d)`$: the process $`\mathrm{\Pi }`$ is not successful contrary to the assumption.
The next lemma relates the processes of two theories when they are assumed to have the same extension. In this lemma and in the following theorem, when a process is used in a place where a set of defaults is expected, it means the set of defaults of the process. For example, if $`\mathrm{\Pi }`$ is a sequence of defaults and $`D^{}`$ a set of defaults, $`\mathrm{\Pi }D^{}`$ is the set of defaults that are both in $`\mathrm{\Pi }`$ and in $`D`$.
###### Lemma 20
In justified default logic, if $`D^{}D`$, $`D^{},W_D^eD,W`$, and $`\mathrm{\Pi }`$ is a selected process $`D,W`$, then there exists a selected process of $`D^{},W`$ made exactly of the defaults of $`\mathrm{\Pi }D^{}`$ and generating the same extension generated by $`\mathrm{\Pi }`$.
Proof. Let $`E=W\mathrm{cons}(\mathrm{\Pi })`$ be the extension that is generated by $`\mathrm{\Pi }`$. By the lemma above, it is generated by the defaults in $`GEN(E,D)`$. Since $`E`$ is also an extension of $`D^{},W`$, it is generated by a process $`\mathrm{\Pi }^{}`$ made exactly of the defaults of $`GEN(E,D^{})=GEN(E,D)E^{}=\mathrm{\Pi }D^{}`$
The main theorem relating the extensions of theories differing for the set of defaults in justified default logic is the following one.
###### Theorem 22
If $`D^{}D^{\prime \prime }D`$ and $`D^{},W^eD,W`$ then $`D^{},W^eD^{\prime \prime },W`$ for justified default logic.
Proof. We first show that every extension of $`D^{\prime \prime },W`$ is also an extension of $`D,W`$, and then show the converse.
Let $`E`$ be an extension of $`D^{\prime \prime },W`$. Let $`\mathrm{\Pi }`$ is one its generating processes. By Lemma 18, there exists a sequence $`\mathrm{\Pi }^{}`$ of defaults of $`D\backslash D^{\prime \prime }`$ such that $`\mathrm{\Pi }\mathrm{\Pi }^{}`$ is a selected process of $`D,W`$. Let $`E^{}`$ be its generated extension. Since $`E`$ is generated by $`\mathrm{\Pi }`$ and $`E^{}`$ is generated by $`\mathrm{\Pi }\mathrm{\Pi }^{}`$, we have $`E^{}E`$. We prove that $`EE^{}`$, which implies $`EE^{}`$.
By Lemma 20, since $`\mathrm{\Pi }`$ is a selected process of $`D,W`$ and this theory is faithfully equivalent to $`D^{},W`$, there exists a selected process $`\mathrm{\Pi }^{\prime \prime }`$ of $`D^{},W`$ made of the defaults of $`(\mathrm{\Pi }\mathrm{\Pi }^{})D^{}`$ and generating the extension $`E^{}`$. Since $`\mathrm{\Pi }^{}`$ is made of defaults of $`D\backslash D^{\prime \prime }`$ and $`D^{}D^{\prime \prime }`$, we have that $`(\mathrm{\Pi }\mathrm{\Pi }^{})D^{}=\mathrm{\Pi }D^{}`$. As a result, $`\mathrm{\Pi }^{\prime \prime }`$ is only made of defaults in $`\mathrm{\Pi }D^{}`$. Since $`\mathrm{\Pi }^{\prime \prime }`$ generates $`E^{}`$ and $`\mathrm{\Pi }`$ generates $`E`$, we have $`EE^{}`$. We can therefore conclude that $`EE^{}`$.
Let us now prove the converse: we assume that $`E`$ is an extension of $`D,W`$ and prove that it is also an extension of $`D^{\prime \prime },W`$. Let $`\mathrm{\Pi }`$ be the process of $`D,W`$ that generates $`E`$. By definition, the following two properties are true:
1. $`W\mathrm{cons}(\mathrm{\Pi })\mathrm{prec}(d)`$ for every $`d\mathrm{\Pi }`$;
2. $`W\mathrm{cons}(\mathrm{\Pi })\mathrm{just}(d)\mathrm{cons}(d)`$ for every $`d\mathrm{\Pi }`$.
By Lemma 20, the theory $`D^{},W`$ has a selected process $`\mathrm{\Pi }^{}`$ that is composed exactly of the defaults of $`\mathrm{\Pi }D^{}`$ and that generates the same extension $`E`$. Since $`W\mathrm{cons}(\mathrm{\Pi }^{})W\mathrm{cons}(\mathrm{\Pi })`$, the two properties are equivalent to the following two ones:
1. $`W\mathrm{cons}(\mathrm{\Pi }^{})\mathrm{prec}(d)`$ for every $`d\mathrm{\Pi }`$;
2. $`W\mathrm{cons}(\mathrm{\Pi }^{})\mathrm{just}(d)\mathrm{cons}(d)`$ for every $`d\mathrm{\Pi }`$.
The first property implies that every default $`d\mathrm{\Pi }(D^{\prime \prime }\backslash D^{})`$ is applicable to $`\mathrm{\Pi }^{}`$: this is because the precondition of $`d`$ is entailed by $`W\mathrm{cons}(\mathrm{\Pi }^{})`$ and the process $`\mathrm{\Pi }^{}[d]`$ is successful because so is $`\mathrm{\Pi }`$, which contains all default of $`\mathrm{\Pi }^{}[d]`$. The second property implies that no default of $`D^{\prime \prime }\backslash \mathrm{\Pi }`$ is applicable to $`\mathrm{\Pi }^{}`$. As a result, $`\mathrm{\Pi }^{}`$ and the sequence composed of all defaults of $`\mathrm{\Pi }(D^{\prime \prime }\backslash D^{})`$ in any order form a selected process of $`D^{\prime \prime }`$. The extension generated by this process is equivalent to $`E`$ because this process is composed of a superset of the defaults of $`\mathrm{\Pi }^{}`$ and a subset of the defaults of $`\mathrm{\Pi }`$, and these two processes both generate $`E`$
We therefore have as a corollary that justified default logic has the local redundancy property when redundancy of defaults is considered.
###### Corollary 7
Justified default logic has the local redundancy property w.r.t. redundancy of defaults.
#### 4.6.3 Redundancy of Clauses and of Defaults
For Reiter and rational default logic, an upper bound on complexity can be given by showing a reduction from the complexity of clause or formula redundancy to the corresponding problems for defaults. This is possible thanks to the following lemma.
###### Lemma 21
$`D,W\{\gamma \}`$ has the same Reiter and rational extensions of $`D\{d_\gamma \},W`$, where $`d_\gamma =\frac{:}{\gamma }`$.
Proof. Since $`d_\gamma `$ has no precondition and a tautological justification, it is always applicable. Therefore, every process of $`D\{d_\gamma \},W`$ contains this default, and therefore generates $`\gamma `$
This lemma can be iterated for all clauses of $`W`$, leading to the following result.
###### Theorem 23
For Reiter and rational default logic, the problems of checking the redundancy of a default or of a default theory are at least as hard as the corresponding problems for clause redundancy.
Proof. The clause $`\gamma `$ is redundant in $`D,W`$ if and only if $`d_\gamma `$ is redundant in $`D\{d_\gamma \},W\backslash \{\gamma \}`$. Indeed, $`D,W`$ has the same extensions of $`D\{d_\gamma \},W\backslash \{\gamma \}`$, and removing $`\gamma `$ from the first theory or removing $`d_\gamma `$ from the second theory lead both to $`D,W\backslash \{\gamma \}`$.
The problems of formula redundancy can be reduced to default redundancy by first applying Lemma 21 to all clauses of $`W`$, and then making all original defaults irredundant using the transformation of Lemma 17
The complexity of redundancy for defaults can be therefore characterized as follows.
###### Corollary 8
For Reiter and justified default logic, the problem of redundancy of a default is $`\mathrm{\Pi }_2^p`$-hard and $`\mathrm{\Pi }_3^p`$-hard for faithful and consequence-equivalence, respectively; the problem of redundancy of a default theory is $`\mathrm{\Sigma }_3^p`$-hard and $`\mathrm{\Sigma }_4^p`$-hard for faithful and consequence-equivalence, respectively.
Equivalence of extensions can be proved to be in $`\mathrm{\Pi }_2^p`$ even if the defaults or the background theories are not even related.
###### Theorem 24
Checking whether $`D,W_D^eD^{},W^{}`$ is in $`\mathrm{\Pi }_2^p`$ for Reiter and justified default logic.
Proof. The contrary of the statement amounts to checking whether any of the two theories have an extension that the other one does not have. The number of possible extensions, however, is limited by the fact that any extension is generated by the set of consequences of some defaults.
Checking whether $`D,W`$ has an extension that $`D^{},W^{}`$ has not can be done as follows: guess a subset $`D^{\prime \prime }D`$, and let $`E=\mathrm{cons}(D^{\prime \prime })`$; check whether $`E`$ is an extension of $`D,W`$ but is not an extension of $`D^{},W^{}`$.
Checking whether a formula $`E`$ is an extension of a default theory can be done with a logarithmic number of satisfiability tests \[Ros99, Lib05a\]. As a result, the problem can also be expressed as a QBF formula $``$QBF. In order to check whether there exists $`D^{\prime \prime }`$ such that $`E=\mathrm{cons}(D^{\prime \prime })`$ is in this condition, we only have to add an existential quantifier to the front of this formula. The problem is therefore in $`\mathrm{\Pi }_2^p`$
The problem of checking the default redundancy of a theory is obviously in $`\mathrm{\Sigma }_3^p`$, as it can be solved by guessing a subsets of defaults and then checking equivalence.
###### Corollary 9
The problem of checking the redundancy of a default or the default redundancy of a theory are $`\mathrm{\Pi }_2^p`$-complete and $`\mathrm{\Sigma }_3^p`$-complete, respectively, for Reiter and justified default logic for faithful equivalence.
Consequence equivalence can also be proved to have the same complexity as for the case studied for clauses.
###### Theorem 25
Checking the consequence-equivalence for Reiter and justified default logic is in $`\mathrm{\Pi }_3^p`$.
Proof. The converse of the problem can be expressed as: there exists a model $`M`$ that is a model of an extension of the first theory but not of the second, or vice versa. This corresponds to two quantifications over extensions and a check for whether a formula is an extension. The latter is in $`\mathrm{\Delta }_2^p[\mathrm{log}n]`$ for the two considered semantics \[Ros99, Lib05a\]. Therefore, the whole problem is in $`\mathrm{\Pi }_3^p`$
As a consequence, the complexity of redundancy for consequence-equivalence is exactly characterized for Reiter and justified default logics.
###### Corollary 10
The problem of checking the redundancy of a default or the default redundancy of a theory are $`\mathrm{\Pi }_3^p`$-complete and $`\mathrm{\Sigma }_4^p`$-complete, respectively, for Reiter default logic and consequence-equivalence.
### Acknowledgments
The author thanks Jiang Zhengjun for his comments about on this paper. |
warning/0507/hep-ph0507052.html | ar5iv | text | # Fractional momentum correlations in multiple production of W bosons and of 𝑏𝑏̄ pairs in high energy 𝑝𝑝 collisions
## I Introduction
Multiple parton interactions in high energy hadronic collisions have been discussed long ago by several authors Landshoff ; Paver . Experimentally events with multiparton interactions have been first observed in $`pp`$ collisions by the AFS Collaboration Akesson:1986iv and later, with sizably larger statistics, at Fermilab by the CDF Collaboration CDF .
In multiple parton collisions the hadron is probed in different points contemporarily Paver . The non trivial feature of multiple parton collisions is hence its non-perturbative input, which has a direct relation with the correlations between partons in the hadron structure Calucci:1997ii . As the process is originated by the large population of partons in the initial state, the expectation is nevertheless that correlations should not represent a major feature in the process, with the exception of correlations in transverse space, which are directly measured by the cross section. Indeed the experimental analysis and most of the theoretical estimates have been done with this simplifying assumption and, although the statistics was too low to draw firm conclusions, the CDF Collaboration reported that the cross section is not influenced appreciably when changing the fractional momenta of initial state partons CDF .
The much larger rates of multiple parton collisions expected at the LHC, with the possibility of testing different multiparton processes at different resolution scales, represents however a good motivation for reconsidering the approach to the problem. In particular an interesting process, where multiple parton collisions play an important role and which might be observed at the LHC, is the production of multiple $`W`$ bosons with equal sign Stirling1 , which would allow testing multiple parton interactions at a much larger resolution scale than usually considered. The evolution of the multiparton structure functions will play a non minor role in this case, leading to sizable correlations in fractional momenta.
The purpose of the present note is to give some quantitative indication of the effects in multiparton collisions in a high resolution regime and compare with a case at a lower resolution. After recalling the basic features of the inclusive cross section of double parton collisions, we will hence evolve double parton distributions at high resolution scales. The effect of correlations induced by evolution will be estimated studying the cross sections of multiple production of equal sign $`W`$ bosons and the cross section of multiple production of $`b\overline{b}`$ pairs, in the energy range $`114`$ TeV.
## II Double parton cross section
With the only assumption of factorization of the two hard parton processes A and B, the inclusive cross section of a double parton-scattering process in a hadronic collision is expressed by Paver ; Braun
$$\sigma _{(A,B)}^D=\frac{m}{2}\underset{i,j,k,l}{}\mathrm{\Gamma }_{ij}(x_1,x_2,b)\widehat{\sigma }_{ik}^A(x_1,x_1^{})\widehat{\sigma }_{jl}^B(x_2,x_2^{})\mathrm{\Gamma }_{kl}(x_1^{},x_2^{};b)𝑑x_1𝑑x_1^{}𝑑x_2𝑑x_2^{}d^2b,$$
(1)
where $`\mathrm{\Gamma }_{ij}(x_1,x_2,b)`$ are the double parton distribution functions, depending on the fractional momenta $`x_1,x_2`$ and on the relative transverse distance $`b`$ of the two partons undergoing the hard processes A and B, the indices $`i`$ and $`j`$ refer to the different parton species and $`\widehat{\sigma }_{ik}^A`$ and $`\widehat{\sigma }_{jl}^B`$ are the partonic cross sections. The dependence on the resolution scales is implicit in all quantities. The factor $`m/2`$ is a consequence of the symmetry of the expression for interchanging $`i`$ and $`j`$; specifically $`m=1`$ for indistinguishable parton processes and $`m=2`$ for distinguishable parton processes.
The double distributions $`\mathrm{\Gamma }_{ij}(x_1,x_2,b)`$ are the main reason of interest in multiparton collisions. The distributions $`\mathrm{\Gamma }_{ij}(x_1,x_2,b)`$ contain in fact all the information of probing the hadron in two different points contemporarily, though the hard processes A and B.
The cross section for multiparton process is sizable when the flux of partons is large, namely at small $`x`$, and dies out quickly for larger values. Given the large parton flux one may hence expect that correlations in momentum fraction will not be a major effect and partons to be rather correlated in transverse space (as they must anyhow all belong to the same hadron). Neglecting the effect of parton correlations in $`x`$ one writes
$$\mathrm{\Gamma }_{ij}(x_1,x_2;b)=\mathrm{\Gamma }_i(x_1)\mathrm{\Gamma }_j(x_2)F_j^i(b),$$
(2)
where $`\mathrm{\Gamma }_i(x)`$ are the usual one body parton distribution function and $`F_j^i(b)`$ is a function normalized to one and representing the parton pair density in transverse space. The inclusive cross section hence simplifies to
$$\sigma _{(A,B)}^D=\frac{m}{2}\underset{ijkl}{}\mathrm{\Theta }_{kl}^{ij}\widehat{\sigma }_{ij}(A)\widehat{\sigma }_{kl}(B),$$
(3)
where $`\widehat{\sigma }_{ij}(A)`$ and $`\widehat{\sigma }_{kl}(B)`$ are the hadronic inclusive cross sections for the two partons labelled $`i`$ and $`j`$ to undergo the hard interaction labelled $`A`$ and for two partons $`k`$ and $`l`$ to undergo the hard interaction labelled $`B`$;
$$\mathrm{\Theta }_{kl}^{ij}=d^2bF_k^i(b)F_l^j(b)$$
(4)
are geometrical coefficients with dimension an inverse cross section and depending on the various parton processes. In the simplified scheme above, the coefficients $`\mathrm{\Theta }_{kl}^{ij}`$ are the experimentally accessible quantities carrying the information of the parton correlations in transverse space.
In the experimental search of multiple parton collisions the cross section has been further simplified assuming that the densities $`F_j^i`$ do not depend on the indices $`i`$ and $`j`$, which leads to the expression
$$\sigma _{(A,B)}^D=\frac{m}{2}\frac{\widehat{\sigma }(A)\widehat{\sigma }(B)}{\sigma _{eff}}\sigma _{fact}^D,$$
(5)
where all information on the structure of the hadron in transverse space is summarized in the value of a single the scale factor, $`\sigma _{eff}`$. In the experimental study of double parton collisions CDF quotes $`\sigma _{eff}=14.5mb`$ CDF .
The experimental evidence is not inconsistent with the simplest hypothesis of neglecting correlations in momentum fractions, the resolution scale probed in the CDF experiment is however not very large, the transverse momenta of final state partons being of the order of 5 GeV. We will hence approach the problem in more general terms, focusing on multiple production of equal sign $`W`$ bosons and of $`b\overline{b}b\overline{b}`$ pairs, keeping into account the correlations in fractional momenta induced by evolution.
## III Two-body distribution functions
The evolution of the double parton distribution function has been discussed in refs Kirschner ; Shelest and more recently in Snigirev . The approach is essentially the same used to study particle correlations in the fragmentation functions Puhala , using the jet calculus rules Konishi .
Introducing the dimensionless variable
$$t=\frac{1}{2\pi b}\mathrm{ln}\left[1+\frac{g^2(\mu ^2)}{4\pi }b\mathrm{ln}\left(\frac{Q^2}{\mu ^2}\right)\right],b=\frac{332n_f}{12\pi },$$
where $`g^2(\mu ^2)`$ is the running coupling constant at the reference scale $`\mu ^2`$ and $`n_f`$ the number of active flavors, the probability $`D_h^{j_1j_2}(x_1,x_2;t)`$ to find two partons of types $`j_1`$ and $`j_2`$ with fractional momenta $`x_1`$ and $`x_2`$ satisfy the generalized Lipatov-Altarelli-Parisi-Dokshitzer evolution equation
$`{\displaystyle \frac{dD_h^{j_1j_2}(x_1,x_2;t)}{dt}}`$ $`={\displaystyle \underset{j_1^{}}{}}{\displaystyle _{x_1}^{1x_2}}{\displaystyle \frac{dx_1^{}}{x_1^{}}}D_h^{j_1^{}j_2}(x_1^{},x_2;t)P_{j_1^{}j_1}\left({\displaystyle \frac{x_1}{x_1^{}}}\right)`$
$`+{\displaystyle \underset{j_2^{}}{}}{\displaystyle _{x_2}^{1x_1}}{\displaystyle \frac{dx_2^{}}{x_2^{}}}D_h^{j_1j_2^{}}(x_1,x_2^{};t)P_{j_2^{}j_2}\left({\displaystyle \frac{x_2}{x_2^{}}}\right)`$
$`+{\displaystyle \underset{j^{}}{}}D_h^j^{}(x_1+x_2;t){\displaystyle \frac{1}{x_1+x_2}}P_{j^{}j_1j_2}\left({\displaystyle \frac{x_1}{x_1+x_2}}\right),`$ (6)
where the subtraction terms are included in the evolution kernels $`P`$.
If at the scale $`\mu ^2`$ one assumes the factorized form
$$D_h^{j_1j_2}(z_1,z_2,0)=D_h^{j_1}(z_1;0)D_h^{j_2}(z_2;0)\theta (1z_1z_2),$$
(7)
at a larger scale one obtains a solution which may be expressed as the sum of a factorized and of two non-factorized contributions:
$$D_h^{j_1j_2}(x_1,x_2;t)=D_h^{j_1}(x_1;t)D_h^{j_2}(x_2;t)\theta (1x_1x_2)+D_{h,corr,1}^{j_1j_2}(x_1,x_2;t)+D_{h,corr,2}^{j_1j_2}(x_1,x_2;t),$$
(8)
where the non-factorized contributions are expressed by the convolutions:
$`D_{h,corr,1}^{j_1j_2}(x_1,x_2;t)`$ $`=\theta (1x_1x_2)[{\displaystyle \underset{j_1^{}j_2^{}}{}}{\displaystyle _{x_1}^1}{\displaystyle \frac{dz_1}{z_1}}{\displaystyle _{x_2}^1}{\displaystyle \frac{dz_2}{z_2}}D_h^{j_1^{}}(z_1,0)D_{j_1^{}}^{j_1}({\displaystyle \frac{x_1}{z_1}};t)`$
$`\times D_h^{j_2^{}}(z_2,0)D_{j_2^{}}^{j_2}({\displaystyle \frac{x_2}{z_2}};t)[\theta (1z_1z_2)1]]`$ (9)
$`D_{h,corr,2}^{j_1j_2}(x_1,x_2;t)`$ $`={\displaystyle \underset{j^{}j_1^{}j_2^{}}{}}{\displaystyle _0^t}𝑑t^{}{\displaystyle _{x_1}^1}{\displaystyle \frac{dz_1}{z_1}}{\displaystyle _{x_2}^{1x_1}}{\displaystyle \frac{dz_2}{z_2}}D_h^j^{}(z_1+z_2;t^{}){\displaystyle \frac{1}{z_1+z_2}}`$
$`\times P_{j^{}j_1^{}j_2^{}}\left({\displaystyle \frac{z_1}{z_1+z_2}}\right)D_{j_1^{}}^{j_1}({\displaystyle \frac{x_1}{z_1}};tt^{})D_{j_2^{}}^{j_2}({\displaystyle \frac{x_2}{z_2}};tt^{});`$ (10)
and the distribution functions $`D_i^j(x;t)`$ satisfy the evolution equation
$$\frac{dD_i^j(x;t)}{dt}=\underset{j^{}}{}_x^1\frac{dx^{}}{x^{}}D_i^j^{}(x^{};t)P_{j^{}j}\left(\frac{x}{x^{}}\right).$$
(11)
with initial condition $`D_i^j(x;t=0)=\delta _{ij}\delta (1x)`$
Equations(11) are solved by introducing the Mellin transforms
$$D_i^j(n;t)=_0^1𝑑xx^nD_i^j(x;t),$$
(12)
which lead to a system of ordinary linear-differential equations at the first order. The solution is given by the inverse Mellin transform
$`D_i^j(x;t)`$ $`=`$ $`{\displaystyle \frac{dn}{2\pi ı}x^nD_i^j(n;t)}=^1(D_i^j(n;t),\mathrm{ln}(x)),`$ (13)
where the integration runs along the imaginary axis at the right of all the $`n`$ singularities, while $`^1`$ represents the Inverse Laplace operator.
The double distributions can then be obtained numerically. For inverting the Laplace Transform we have followed two different procedures Abate : the Gaver-Wynn-Rho (GWR) algorithm and the fixed Talbot (FT) method. The first procedure (GWR) is based on a special acceleration sequence of the Gaver functionals and requires to evaluate the transform only on the real axes; the second procedure (FT) is based on the deformation of the contour of the Bromwich inversion integral and requires complex arithmetic. Comparing the two methods we have found more stable results when using the (FT) method. The double distributions have hence been obtained by numerical integration with the Vegas algorithm Vegas , using the MRS99 MRS99 as input parton distribution function at the scale $`\mu ^2`$.
In the kinematical range of interest for the actual case (we never exceed $`x=.1`$) the contribution of the term $`D_{h,corr,1}^{j_1j_2}`$ in eq.(8) is negligible. The first term in eq.(8) represents the factorized contribution usually considered and is the solution of the homogeneous (LAPD) evolution equation, while the third term is a particular solution of the complete equation.
The effect of the correlation terms induced by evolution is shown for gluon-gluon and for quark-quark in Fig., where the ratio
$$R^{j_1j_2}(x_1,x_2;t)=\frac{D_{h,corr,1}^{j_1j_2}(x_1,x_2;t)+D_{h,corr,2}^{j_1j_2}(x_1,x_2;t)}{D_h^{j_1}(x_1;t)D_h^{j_2}(x_2;t)}.$$
(14)
is plotted as a function of $`x`$, with $`x_1=x_2=x`$, with the following choice of parameters: $`\mu =1.2GeV`$, $`n_f=4`$, factorization scale equal to the $`W`$ mass, $`m_w=80.4GeV`$ (solid curves) and factorization scale equal to the bottom quark mass, $`m_b=4.6GeV`$ (dashed curves).
As shown in Fig., the ratio $`R^{gg}`$ is nearly $`35\%`$ for $`x0.1`$ and decreases up to $`810\%`$ for $`x0.01`$ and to $`2\%`$ for $`x0.001`$, when the $`W`$ mass is used as factorization scale. When taking the $`b`$ quark mass as factorization scale, the value of the ratio is of the order of $`1012\%`$ for $`x0.1`$ and decreases up to $`5\%`$ and to $`2\%`$ for $`x0.01`$ and $`x0.001`$ respectively. The ratio would of course be much larger (up to $`60\%`$) if going to larger $`x`$ values.
The ratio $`R^{qq}`$ is shown in Fig. for a few flavor choices. With the $`W`$ mass as factorization scale, the ratios are of the order of $`35,\mathrm{\hspace{0.17em}20},\mathrm{\hspace{0.17em}10}\%`$ for $`x0.1,\mathrm{\hspace{0.17em}0.01},\mathrm{\hspace{0.17em}0.001}`$. With the $`b`$ quark mass as factorization scale the ratios are of order of $`23,\mathrm{\hspace{0.17em}10},\mathrm{\hspace{0.17em}5}\%`$ for $`x0.1,\mathrm{\hspace{0.17em}0.01},\mathrm{\hspace{0.17em}0.001}`$.
Apart from the case of hadron-nucleus collisions, when two different target nucleons take part to the process Strikman:2001gz , the non-perturbative input of the double parton scattering cross section is not represented however by the distribution functions $`D_h^{j_1j_2}(x_1,x_2;t)`$ in Eq.(8), where all transverse variables have been integrated. The double parton scattering cross section, Eq.(1), depends in fact in a direct way also on the relative separation of partons in transverse space, which is of the order of the hadron size and hence outside the control of perturbation theory.
Considering that the longitudinal and the transverse momenta of initial state partons are essentially decoupled in the process, because of the different scales involved, it’s not unreasonable to assume phenomenologically a factorized dependence of the double distribution functions on the longitudinal and transverse degrees of freedom. Given the different origin of the terms in $`D_h^{j_1j_2}`$, it’s also not unnatural to consider the possibility of having different non-perturbative scales, for the transverse separation of the factorized and of the correlated terms. In fact, although in the general case evolution would mix the two scales in the $`D_{h,corr,1}^{j_1j_2}`$ term, the term $`D_{h,corr,1}^{j_1j_2}`$ is very small in the kinematical regime of interest and the hypothesis of two different transverse scales is not inconsistent.
We hence assume that the typical transverse distance between partons in $`D_{h,fact}^{j_1j_2}`$ and in $`D_{h,corr,1}^{j_1j_2}`$ corresponds to the relatively low resolution scale process observed by CDF and, to have an idea on the effects of the presence of two different scales in the double parton densities, we introduce a different transverse distance in the term $`D_{h,corr,2}^{j_1j_2}`$, related to the size of the gluon cloud of a valence quark, and corresponding to a relatively shorter range correlation term. The double parton distributions are hence expressed in the following way:
$$D_h^{j_1j_2}(x_1,x_2;b;t)=\left(D_{h,fact}^{j_1j_2}(x_1,x_2;t)+D_{h,corr,1}^{j_1j_2}(x_1,x_2;t)\right)F_{\sigma _{eff}}(b)+D_{h,corr,2}^{j_1j_2}(x_1,x_2;t)F_{\sigma _r}(b)$$
where the parton pair densities $`F_i(b)`$ satisfy
$$d^2bF_i(b)=1d^2bF_i(b)^2=\frac{1}{\sigma _i}$$
with $`i=\sigma _{eff},\sigma _r`$.
While $`F_{\sigma _{eff}}`$ represents the transverse density of partons at a relatively low resolution scale, relevant in the kinematical conditions of the CDF experiment and leading to the measured value of the scale factor $`\sigma _{eff}=14.5`$ mb, $`F_{\sigma _r}`$ is rather the transverse parton density characterizing partons correlated in fractional momentum, which becomes increasingly important when the resolution scale is large. To study the effect of the two scales we have let the smaller scale vary in the interval $`\sigma _{r_0}\sigma _{eff}`$ assuming $`\sigma _{r_0}=2.8`$ mb Povh , which might represent the size of the gluon cloud of a valence quark in the hadron. To disentangle the effects of the correlation in fractional momenta we have neglected a possible dependence of the parton pair densities $`F_i(b)`$ on the partons flavor.
## IV Multiple production of $`b\overline{b}`$ pairs and of equal sign $`W`$ bosons in $`pp`$ collisions
For the purpose of the present analysis we have hence evaluated the contributions to multiple production of equal sign $`W`$ bosons and to multiple production of $`b\overline{b}`$ pairs, due to multiple (disconnected) parton collision processes, taking into account the correlation terms in fractional momenta induced by evolution.
As a matter of fact higher order corrections in $`\alpha _S`$ are very important in heavy quark production. To the purpose of the present analysis we have evaluated the cross section at the lowest order in perturbation theory, taking higher order corrections into account by rescaling the lowest order results with a constant factor $`K`$, defined as the ratio between the inclusive cross-section for $`b\overline{b}`$ production, $`\sigma (b\overline{b})`$, and the result of the lowest-order calculation in pQCD. Our assumption is hence that higher order corrections in $`b\overline{b}b\overline{b}`$ production may be taken into account by multiplying the cross section of each connected process by the same factor $`K`$, so that higher order corrections are taken into account by multiplying the lowest order cross section by the $`K`$-factor at the second power. In the actual calculation we have used a $`K`$ factor equal to $`5.7`$ (Cattaruzza ; DelFabbro4 ) and the value $`m_b=4.6GeV`$ for the mass of the bottom quark.
The multiparton distributions have been obtained, as described in the previous paragraph, using as input distributions at the scale $`\mu ^2`$ the MRS99 MRS99 parton distribution functions. Factorization and renormalization scale have been set equal to the transverse mass of the produced quarks. As for the dependence on the transverse variables, in addition to the usual factorized contribution, leading to the scale factor $`1/\sigma _{eff}`$, in the present case the cross section includes also non factorized contributions, corresponding to the couplings of $`D_{h,corr,2}^{ik}`$ both with $`D_{h,fact,1}^{jl}`$ and with $`D_{h,corr,2}^{ik}`$. We have assumed a gaussian distribution for $`F_{\sigma _{eff}}(b)`$ and for $`F_{\sigma _r}(b)`$. The scale factors are correspondingly $`2/(\sigma _{eff}+\sigma _r)`$ and $`1/\sigma _r`$.
In Fig. we plot the $`gg`$ correlation (the dominant contribution to $`b\overline{b}`$ is gluon fusion) while the expected rise of the total $`b\overline{b}b\overline{b}`$ cross-section is plotted in Fig. (*left-panel*) as a function of the center of mass energy. The dashed curve refers to the double-parton scattering factorized term ($`\sigma _{fact}^D`$) given by eq.(5); the continuous curves refer to the double-parton scattering correlation contributions ($`\sigma _{corr}^D`$), with geometrical factors determined by setting $`r=r_0`$ (upper curve) and $`r=r_{eff}`$ (lower curve). The ratio between the contribution of the terms with correlations and the factorized term is shown in Fig. (*right panel*) as a function of center of mass energy. The effect of the terms with correlations decreases by increasing the center of mass energy; depending on the values of $`r[r_{eff},r_0]`$, correction effects may vary between $`(1220)\%`$ at $`\sqrt{s}=1TeV`$ and $`(3.56)\%`$ at $`\sqrt{s}=14TeV`$. The decrease is faster as $`\sqrt{s}5TeV`$: for larger c.m. energies the average fractional momentum $`<x>`$ becomes smaller than $`0.01`$, where the fraction $`R^{gg}`$ stabilizes around $`0.030.05`$, consistently with the amount of correction obtained for $`\sqrt{s}>5TeV`$.
In Fig. we plot the $`b\overline{b}b\overline{b}`$ production cross-section at $`\sqrt{s}=14TeV`$ (*left-panel*) and at $`\sqrt{s}=5.5TeV`$ (*right-panel*), as a function of the minimum value of transverse momenta of the outgoing $`b`$ quarks $`p_t^{min}`$, in the pseudorapidity interval $`|\eta |<0.9`$. At $`\sqrt{s}=14TeV`$ with $`p_t^{min}[0,10]GeV`$ one has $`<x>[1.2,3.4]\times 10^3`$, which leads to a contribution of the correlation terms of the order of $`(24)\%`$ and of $`(47)\%`$, respectively for the lower and the higher choices of $`p_t^{min}`$ Fig.(5)(*left-panel*). At $`\sqrt{s}=5.5TeV`$, in the considered range of variability of $`p_t^{min}`$ one has $`<x>[3.5,6.5]\times 10^3`$ and the contribution of the correlation terms can become of the order of $`12\%`$, Fig.(5)(*right-panel*).
The cross sections of like-sign W pair production are evaluated at the leading order, hence including only quark initiated processes in the elementary interaction ($`q\overline{q}^{}W`$). Higher order corrections are taken into account multiplying the lowest order cross section by the factor $`K1+(8\pi /9)\alpha _s(M_W^2)`$ Barger . We plot in fig. (*left-panel*) the $`W^+W^+`$ cross-section as a function of the $`pp`$ center of mass energy. As in the case of $`b\overline{b}b\overline{b}`$ production, the dashed curve refers to the double-parton scattering factorized term ($`\sigma _{fact}^D`$), while the solid curves to the contribution of the terms with correlations ($`\sigma _{corr}^D`$), for the two different choices $`r=r_0`$ and $`r=r_{eff}`$. As one may infer from the behavior of the qq-correlation ratio, for $`<x>[0.2,\mathrm{\hspace{0.17em}6}]\times 10^2`$, which corresponds to the energy interval considered, the corrections due to the correlation terms range from $`(2745)\%`$ at $`\sqrt{s}=1TeV`$ to $`(7.513)\%`$ at $`\sqrt{s}=14TeV`$, depending on the choice of $`\sigma _r`$. The results for $`W^{}W^{}`$ production are presented in fig.(7). As shown in the right panel the correlation terms can give contributions ranging from $`(2340)\%`$ at $`1TeV`$ to $`(1220)\%`$ at $`14TeV.`$
## V Conclusions
As an effect of evolution, the multiparton distributions functions are expected to become strongly correlated in momentum fraction at large $`Q^2`$ and finite $`x`$ Kirschner ; Shelest ; Snigirev . On the other hand, the indications from the experimental observation of multiparton collisions at Fermilab CDF are not in favor of strong correlation effects in fractional momenta. The most likely reason being that the kinematical domain observed, relatively low $`x`$ values and limited resolution scale, is far from the limiting case considered in QCD.
The possibility of testing multiparton collisions at high resolution scales at the LHC will open the opportunity of testing the correlations predicted by evolution. To have an indication on the importance of the effects to be expected, we have considered a high resolution scale multiparton process (equal sign $`W`$ pair production) and, for comparison, a sizably smaller resolution scale process ($`b\overline{b}b\overline{b}`$ production) in $`pp`$ collisions in the energy range $`1\mathrm{T}\mathrm{e}\mathrm{V}\sqrt{s}14\mathrm{T}\mathrm{e}\mathrm{v}`$. In both cases the production process may take place either by single (connected) or by multiple (disconnected) hard parton collisions, while the two contributions may be disentangled applying proper cuts in the final state CDF ; Stirling1 ; DelFabbro4 . To study the effects of correlations we have hence worked out the disconnected contributions to the cross sections after evolving the multiparton distribution functions at high resolution scales.
Our result is that the contribution of the terms with correlations, in equal sign $`W`$ pairs production, might be almost 40% of the cross section at 1 TeV and might still be a 20% effect at the LHC. The effect is much smaller in $`b\overline{b}b\overline{b}`$ production, where corrections to the usually considered factorized distribution are typically between 5 and 10%.
###### Acknowledgements.
This work was partially supported by the Italian Ministry of University and of Scientific and Technological Researches (MIUR) by the Grant COFIN2003. |
warning/0507/astro-ph0507383.html | ar5iv | text | # I Introduction
## I Introduction
Hydrodynamics simulation in general relativity is probably the best theoretical approach for investigating dynamical phenomena in relativistic astrophysics such as stellar core collapse to a neutron star and a black hole, and the merger of binary neutron stars. In the past several years, this field has been extensively developed (e.g., ) and, as a result, now it is feasible to perform accurate simulations of such general relativistic phenomena for yielding scientific results (e.g., for our latest results). For example, with the current implementation, radiation reaction of gravitational waves in the merger of binary neutron stars can be taken into account within $`1\%`$ error in an appropriate computational setting . This fact illustrates that the numerical relativity is a robust approach for detailed theoretical study of astrophysical phenomena and gravitational waves emitted.
However, so far, most of the scientific simulations in full general relativity have been performed without taking into account detailed effects except for general relativistic gravity and pure hydrodynamics. For example, simplified ideal equations of state have been adopted instead of realistic ones (but see ). Also, the effect of magnetic fields has been neglected although it could often play an important role in the astrophysical phenomena (but see ). In the next stage of numerical relativity, it is necessary to incorporate these effects for more realistic simulations. As a step toward a more realistic simulation, we have incorporated an implementation for ideal magnetohydrodynamics (MHD) equations in fully general relativistic manner. In this paper, we describe our approach for these equations and then present numerical results for test problems computed by the new implementation.
Magnetic fields indeed play an important role in determining the evolution of a number of relativistic objects. In the astrophysical context, the plasma is usually highly conducting, and hence, the magnetic fields are frozen in the matter. This implies that a small seed field can wind up and grow in the complex motion of the matter, resulting in a significant effect in the dynamics of the matter such as magnetically driven wind or jet and angular momentum redistribution. Specifically, in the context of the general relativistic astrophysics, the magnetic fields will play a role in the following phenomena and objects: Stellar core collapse of magnetized massive stars to a protoneutron star or a black hole, stability of accretion disks (which are either non-self-gravitating or self-gravitating) around black holes and neutron stars, magnetic braking of differentially rotating neutron stars which are formed after merger of binary neutron stars and stellar core collapse , and magnetically induced jet around the compact objects (e.g., ). To clarify these phenomena, fully general relativistic MHD (GRMHD) simulation (involving dynamical spacetimes) is probably the best theoretical approach.
In the past decade, numerical implementations for GRMHD simulation in the fixed gravitational field have been extensively developed (e.g., ). In particular, it is worth to mention that Refs. have recently presented implementations for which detailed tests have been carried out for confirmation of the reliability of their computation, in contrast with the attitude in an early work . They are applied for simulating magnetorotational instability (MRI) of accretion disks and subsequently induced winds and jets around black holes and neutron stars. On the other hand, little effort has been paid to numerical implementations of fully GRMHD (in the dynamical gravitational field). About 30 years ago, Wilson performed a simulation for collapse of a magnetized star in the presence of poloidal magnetic fields in general relativity. However, he assumes that the three-metric is conformally flat , and hence, the simulation is not fully general relativistic, although recent works have indicated that the conformally flat approximation works well in the axisymmetric collapse (e.g., compare results among , , and ). The first fully GRMHD simulation for stellar collapse was performed by Nakamura about 20 years ago . He simulated collapse of nonrotating stars with poloidal magnetic fields to investigate the criteria for formation of black holes and naked singularities. Very recently, Duez et al. have presented a new implementation capable of evolution for the Einstein-Maxwell-MHD equations for general cases . They report successful results for test simulations. Valencia group has also developed a GRMHD implementation very recently .
In this paper, we present our new implementation for fully GRMHD which is similar to but in part different from that in <sup>*</sup><sup>*</sup>*For instance, our formulation for Einstein’s evolution equations, gauge conditions, and our numerical scheme for GRMHD equations are different from those in as mentioned in Secs. II and III.. As a first step toward scientific simulations, we have performed simulations in standard test problems including special relativistic magnetized shocks, general relativistic Bondi flow in stationary spacetime, and long term evolution of fully general relativistic stars with magnetic fields. We here report the successful results for these test problems.
Before proceeding, we emphasize that it is important to develop new GRMHD implementations. In the presence of magnetic fields, matter motion often becomes turbulence-like due to MRIs in which a small scale structure often grows most effectively . Furthermore in the presence of general relativistic self-gravity which has a nonlinear nature, the matter motion may be even complicated. Perhaps, the outputs from the simulations will contain results which have not been well understood yet, and thus, are rich in new physics. Obviously high accuracy is required for such frontier simulation to confirm novel numerical results. However, because of the restriction of computational resources, it is often very difficult to get a well-resolved and completely convergent numerical result in fully general relativistic simulation. In such case, comparison among various results obtained by different numerical implementations is crucial for checking the reliability of the numerical results. From this point of view, it is important to develop several numerical implementations in the community of numerical relativity. By comparing several results computed by different implementations, reliability of the numerical results will be improved each other. Our implementation presented here will be useful not only for finding new physics but also for checking numerical results by other implementations such as that very recently presented in .
In Sec. II, we present formulations for Einstein, Maxwell, and GRMHD equations. In Sec. III, numerical methods for solving GRMHD equations are described. In Sec. IV, methods for a solution of initial value problem in general relativity is presented. In Secs. V and VI, numerical results for special and general relativistic test simulations are shown. In the final subsection of Sec. VI, we illustrate that our implementation can follow growth of magnetic fields of accretion disks in fully general relativistic simulation. Sec. VII is devoted to a summary and a discussion. Throughout this paper, we adopt the geometrical units in which $`G=c=1`$ where $`G`$ and $`c`$ are the gravitational constant and the speed of light. Latin and Greek indices denote spatial components and spacetime components, respectively. $`\eta _{\mu \nu }`$ and $`\delta _{ij}(=\delta ^{ij})`$ denote the flat spacetime metric and the Kronecker delta, respectively.
## II Basic equations
### A Definition of variables
Basic equations consist of the Einstein equations, general relativistic hydrodynamic equations, and Maxwell equations. In this subsection, we define the variables used in these equations. The fundamental variables for geometry are $`\alpha `$: lapse function, $`\beta ^k`$: shift vector, $`\gamma _{ij}`$: metric in three-dimensional spatial hypersurface, and $`K_{ij}`$: extrinsic curvature. The spacetime metric $`g_{\mu \nu }`$ is written as
$`g_{\mu \nu }=\gamma _{\mu \nu }n_\mu n_\nu ,`$ (1)
where $`n^\mu `$ is a unit normal to a spacelike spatial hypersurface $`\mathrm{\Sigma }`$ and is written as
$`n^\mu =({\displaystyle \frac{1}{\alpha }},{\displaystyle \frac{\beta ^i}{\alpha }}),\mathrm{or}n_\mu =(\alpha ,0).`$ (2)
In the BSSN formalism , one defines $`\gamma \eta e^{12\varphi }=\mathrm{det}(\gamma _{ij})`$: determinant of $`\gamma _{ij}`$, $`\stackrel{~}{\gamma }_{ij}=e^{4\varphi }\gamma _{ij}`$: conformal three-metric, $`K=K_k^k`$: trace of the extrinsic curvature, and $`\stackrel{~}{A}_{ij}e^{4\varphi }(K_{ij}K\gamma _{ij}/3)`$: a tracefree part of the extrinsic curvature. Here, $`\eta `$ denotes the determinant of flat metric; in the Cartesian coordinates, $`\eta =1`$, and in the cylindrical coordinates $`(\varpi ,\phi ,z)`$, $`\eta =\varpi ^2`$. In the following, $`_\mu `$, $`D_i`$, and $`\stackrel{~}{D}_i`$ denote the covariant derivatives with respect to $`g_{\mu \nu }`$, $`\gamma _{ij}`$, and $`\stackrel{~}{\gamma }_{ij}`$, respectively. $`\mathrm{\Delta }`$ and $`\stackrel{~}{\mathrm{\Delta }}`$ denote the Laplacians with respect to $`\gamma _{ij}`$ and $`\stackrel{~}{\gamma }_{ij}`$. $`R_{ij}`$ and $`\stackrel{~}{R}_{ij}`$ denote the Ricci tensors with respect to $`\gamma _{ij}`$ and $`\stackrel{~}{\gamma }_{ij}`$, respectively.
The fundamental variables in hydrodynamics are $`\rho `$: rest-mass density, $`\epsilon `$ : specific internal energy, $`P`$ : pressure, and $`u^\mu `$ : four velocity. From these variables, we define the following variables which often appear in the basic equations:
$`\rho _{}\rho we^{6\varphi },`$ (3)
$`v^i{\displaystyle \frac{dx^i}{dt}}={\displaystyle \frac{u^i}{u^t}}=\beta ^i+\gamma ^{ij}{\displaystyle \frac{u_j}{u^t}},`$ (4)
$`h1+\epsilon +{\displaystyle \frac{P}{\rho }},`$ (5)
$`w\alpha u^t.`$ (6)
Here, $`\rho _{}`$ is a weighted baryon rest mass density from which the conserved baryon rest mass can be computed as
$`M_{}={\displaystyle \rho _{}\eta ^{1/2}d^3x}.`$ (7)
The fundamental variable in the ideal MHD is only $`b^\mu `$: magnetic field. The electric field $`E^\mu `$ in the comoving frame $`F^{\mu \nu }u_\nu `$ is assumed to be zero, and electric current $`j^\mu `$ is not explicitly necessary for evolving the field variables. Using the electromagnetic tensor $`F^{\mu \nu }`$, $`b_\mu `$ is defined by
$`b_\mu {\displaystyle \frac{1}{2}}ϵ_{\mu \nu \alpha \beta }u^\nu F^{\alpha \beta },`$ (8)
where $`ϵ_{\mu \nu \alpha \beta }`$ is the Levi-Civita tensor with $`ϵ_{t123}=\sqrt{g}`$ and $`ϵ^{t123}=1/\sqrt{g}`$. Equation (8) implies
$`b^\mu u_\mu =0.`$ (9)
Using Eq. (8), $`F^{\mu \nu }`$ in the ideal MHD is written as
$`F^{\mu \nu }=ϵ^{\mu \nu \alpha \beta }u_\alpha b_\beta ,`$ (10)
and thus, it satisfies the ideal MHD condition
$`F_{\mu \nu }u^\nu =0.`$ (11)
The dual tensor of $`F_{\mu \nu }`$ is defined by
$`F_{\mu \nu }^{}{\displaystyle \frac{1}{2}}ϵ_{\mu \nu \alpha \beta }F^{\alpha \beta }=b_\mu u_\nu b_\nu u_\mu .`$ (12)
For rewriting the induction equation for the magnetic fields into a simple form (see Sec. II D), we define the three-magnetic field as
$`^ie^{6\varphi }\gamma _j^iF^{j\mu }n_\mu =e^{6\varphi }(wb^i\alpha b^tu^i).`$ (13)
Here, we note that $`^t=0`$ (i.e., $`^\mu n_\mu =0`$), and thus, $`_i=\gamma _{ij}^j`$. Equations (13) and (9) lead to
$`b^t={\displaystyle \frac{^\mu u_\mu }{\alpha e^{6\varphi }}}\mathrm{and}b_i={\displaystyle \frac{1}{we^{6\varphi }}}\left(_i+^ju_ju_i\right)`$ (14)
Using the hydrodynamic and electromagnetic variables, energy-momentum tensor is written as
$`T_{\mu \nu }`$ $`=`$ $`T_{\mu \nu }^{\mathrm{Fluid}}+T_{\mu \nu }^{\mathrm{EM}}.`$ (15)
$`T_{\mu \nu }^{\mathrm{Fluid}}`$ and $`T_{\mu \nu }^{\mathrm{EM}}`$ denote the fluid and electromagnetic parts defined by
$`T_{\mu \nu }^{\mathrm{Fluid}}=(\rho +\rho \epsilon +P)u_\mu u_\nu +Pg_{\mu \nu }=\rho hu_\mu u_\nu +Pg_{\mu \nu },`$ (16)
$`T_{\mu \nu }^{\mathrm{EM}}=F_{\mu \sigma }F_\nu ^\sigma {\displaystyle \frac{1}{4}}g_{\mu \nu }F_{\alpha \beta }F^{\alpha \beta }=\left({\displaystyle \frac{1}{2}}g_{\mu \nu }+u_\mu u_\nu \right)b^2b_\mu b_\nu ,`$ (17)
where
$`b^2=b_\mu b^\mu ={\displaystyle \frac{^2+(^iu_i)^2}{w^2e^{12\varphi }}}.`$ (18)
Thus, $`T_{\mu \nu }`$ is written as
$`T_{\mu \nu }=(\rho h+b^2)u_\mu u_\nu +\left(P+{\displaystyle \frac{1}{2}}b^2\right)g_{\mu \nu }b_\mu b_\nu .`$ (19)
For the following, we define magnetic pressure and total pressure as $`P_{\mathrm{mag}}=b^2/2`$ and $`P_{\mathrm{tot}}P+b^2/2`$, respectively.
The (3+1) decomposition of $`T_{\mu \nu }`$ is
$`\rho _\mathrm{H}T_{\mu \nu }n^\mu n^\nu =(\rho h+b^2)w^2P_{\mathrm{tot}}(\alpha b^t)^2,`$ (20)
$`J_iT_{\mu \nu }n^\mu \gamma _i^\nu =(\rho h+b^2)wu_i\alpha b^tb_i,`$ (21)
$`S_{ij}T_{\mu \nu }\gamma _i^\mu \gamma _j^\nu =(\rho h+b^2)u_iu_j+P_{\mathrm{tot}}\gamma _{ij}b_ib_j.`$ (22)
Using these, the energy-momentum tensor is rewritten in the form
$`T_{\mu \nu }=\rho _\mathrm{H}n_\mu n_\nu +J_i\gamma _\mu ^in_\nu +J_i\gamma _\nu ^in_\mu +S_{ij}\gamma _\mu ^i\gamma _\nu ^j.`$ (23)
This form of the energy-momentum tensor is useful for deriving the basic equations for GRMHD presented in Sec. II C. For the following, we define
$`S_0e^{6\varphi }\rho _\mathrm{H},`$ (24)
$`S_ie^{6\varphi }J_i.`$ (25)
These variables together with $`\rho _{}`$ and $`^i`$ are evolved explicitly in the numerical simulation of the ideal MHD (see Sec. II C).
### B Einstein’s equation
Our formulation for Einstein’s equations is the same as in in three spatial dimensions and in in axial symmetry. Here, we briefly review the basic equations in our formulation. Einstein’s equations are split into constraint and evolution equations. The Hamiltonian and momentum constraint equations are written as
$`R_k^k\stackrel{~}{A}_{ij}\stackrel{~}{A}^{ij}+{\displaystyle \frac{2}{3}}K^2=16\pi \rho _\mathrm{H},`$ (26)
$`D_i\stackrel{~}{A}_j^i{\displaystyle \frac{2}{3}}D_jK=8\pi J_j,`$ (27)
or, equivalently
$`\stackrel{~}{\mathrm{\Delta }}\psi ={\displaystyle \frac{\psi }{8}}\stackrel{~}{R}_k^k2\pi \rho _\mathrm{H}\psi ^5{\displaystyle \frac{\psi ^5}{8}}\left(\stackrel{~}{A}_{ij}\stackrel{~}{A}^{ij}{\displaystyle \frac{2}{3}}K^2\right),`$ (28)
$`\stackrel{~}{D}_i(\psi ^6\stackrel{~}{A}_j^i){\displaystyle \frac{2}{3}}\psi ^6\stackrel{~}{D}_jK=8\pi J_j\psi ^6,`$ (29)
where $`\psi e^\varphi `$. These constraint equations are solved to set initial conditions. A method in the case of GRMHD is presented in Sec. IV.
In the following of this subsection, we assume that Einstein’s equations are solved in the Cartesian coordinates $`(x,y,z`$) for simplicity. Although we apply the implementation described here to axisymmetric issues as well as nonaxisymmetric ones, this causes no problem since Einstein’s equations in axial symmetry can be solved using the so-called Cartoon method in which an axisymmetric boundary condition is appropriately imposed in the Cartesian coordinates : In the Cartoon method, the field equations are solved only in the $`y=0`$ plane, and grid points of $`y=\pm \mathrm{\Delta }x`$ ($`\mathrm{\Delta }x`$ denotes the grid spacing in the uniform grid) are used for imposing the axisymmetric boundary conditions.
We solve Einstein’s evolution equations in our latest BSSN formalism . In this formalism, a set of variables ($`\stackrel{~}{\gamma }_{ij},\varphi ,\stackrel{~}{A}_{ij},K,F_i`$) are evolved. Here, we adopt an auxiliary variable $`F_i\delta ^{jl}_l\stackrel{~}{\gamma }_{ij}`$ that is the one originally proposed and different from the variable adopted in in which $`_i\stackrel{~}{\gamma }^{ij}`$ is used. Evolution equations for $`\stackrel{~}{\gamma }_{ij}`$, $`\varphi `$, $`\stackrel{~}{A}_{ij}`$, and $`K`$ are
$`(_t\beta ^l_l)\stackrel{~}{\gamma }_{ij}=2\alpha \stackrel{~}{A}_{ij}+\stackrel{~}{\gamma }_{ik}\beta _{,j}^k+\stackrel{~}{\gamma }_{jk}\beta _{,i}^k{\displaystyle \frac{2}{3}}\stackrel{~}{\gamma }_{ij}\beta _{,k}^k,`$ (30)
$`(_t\beta ^l_l)\stackrel{~}{A}_{ij}=e^{4\varphi }\left[\alpha \left(R_{ij}{\displaystyle \frac{1}{3}}e^{4\varphi }\stackrel{~}{\gamma }_{ij}R_k^k\right)\left(D_iD_j\alpha {\displaystyle \frac{1}{3}}e^{4\varphi }\stackrel{~}{\gamma }_{ij}\mathrm{\Delta }\alpha \right)\right]`$ (31)
$`+\alpha (K\stackrel{~}{A}_{ij}2\stackrel{~}{A}_{ik}\stackrel{~}{A}_j^k)+\beta _{,i}^k\stackrel{~}{A}_{kj}+\beta _{,j}^k\stackrel{~}{A}_{ki}{\displaystyle \frac{2}{3}}\beta _{,k}^k\stackrel{~}{A}_{ij}`$ (32)
$`8\pi \alpha \left(e^{4\varphi }S_{ij}{\displaystyle \frac{1}{3}}\stackrel{~}{\gamma }_{ij}S_k^k\right),`$ (33)
$`(_t\beta ^l_l)\varphi ={\displaystyle \frac{1}{6}}\left(\alpha K+\beta _{,k}^k\right),`$ (34)
$`(_t\beta ^l_l)K=\alpha \left[\stackrel{~}{A}_{ij}\stackrel{~}{A}^{ij}+{\displaystyle \frac{1}{3}}K^2\right]\mathrm{\Delta }\alpha +4\pi \alpha (\rho _\mathrm{H}+S_k^k).`$ (35)
For a solution of $`\varphi `$, the following conservative form may be adopted :
$$_te^{6\varphi }_i(\beta ^ie^{6\varphi })=\alpha Ke^{6\varphi }.$$
(36)
For computation of $`R_{ij}`$ in the evolution equation of $`\stackrel{~}{A}_{ij}`$, we decompose
$$R_{ij}=\stackrel{~}{R}_{ij}+R_{ij}^\varphi ,$$
(37)
where
$`R_{ij}^\varphi =2\stackrel{~}{D}_i\stackrel{~}{D}_j\varphi 2\stackrel{~}{\gamma }_{ij}\stackrel{~}{\mathrm{\Delta }}\varphi +4\stackrel{~}{D}_i\varphi \stackrel{~}{D}_j\varphi 4\stackrel{~}{\gamma }_{ij}\stackrel{~}{D}_k\varphi \stackrel{~}{D}^k\varphi ,`$ (38)
$`\stackrel{~}{R}_{ij}={\displaystyle \frac{1}{2}}\left[\delta ^{kl}(h_{ij,kl}+h_{ik,lj}+h_{jk,li})+2_k(f^{kl}\stackrel{~}{\mathrm{\Gamma }}_{l,ij})2\stackrel{~}{\mathrm{\Gamma }}_{kj}^l\stackrel{~}{\mathrm{\Gamma }}_{il}^k\right].`$ (39)
In Eq. (39), we split $`\stackrel{~}{\gamma }_{ij}`$ and $`\stackrel{~}{\gamma }^{ij}`$ as $`\delta _{ij}+h_{ij}`$ and $`\delta ^{ij}+f^{ij}`$, respectively. $`\stackrel{~}{\mathrm{\Gamma }}_{ij}^k`$ is the Christoffel symbol with respect to $`\stackrel{~}{\gamma }_{ij}`$, and $`\stackrel{~}{\mathrm{\Gamma }}_{k,ij}=\stackrel{~}{\gamma }_{kl}\stackrel{~}{\mathrm{\Gamma }}_{ij}^l`$. Because of the definition det$`(\stackrel{~}{\gamma }_{ij})=1`$ (in the Cartesian coordinates), we use $`\stackrel{~}{\mathrm{\Gamma }}_{ki}^k=0`$.
In addition to a flat Laplacian of $`h_{ij}`$, $`\stackrel{~}{R}_{ij}`$ involves terms linear in $`h_{ij}`$ as $`\delta ^{kl}h_{ik,lj}+\delta ^{kl}h_{jk,li}`$. To perform numerical simulation stably, we replace these terms by $`F_{i,j}+F_{j,i}`$. This is the most important part in the BSSN formalism, pointed out originally by Nakamura . The evolution equation of $`F_i`$ is derived by substituting Eq. (30) into the momentum constraint as
$`(_t\beta ^l_l)F_i`$ $`=`$ $`16\pi \alpha J_i+2\alpha \left\{f^{kj}\stackrel{~}{A}_{ik,j}+f_{,j}^{kj}\stackrel{~}{A}_{ik}{\displaystyle \frac{1}{2}}\stackrel{~}{A}^{jl}h_{jl,i}+6\varphi _{,k}\stackrel{~}{A}_i^k{\displaystyle \frac{2}{3}}K_{,i}\right\}`$ (40)
$`+`$ $`\delta ^{jk}\left\{2\alpha _{,k}\stackrel{~}{A}_{ij}+\beta _{,k}^lh_{ij,l}+\left(\stackrel{~}{\gamma }_{il}\beta _{,j}^l+\stackrel{~}{\gamma }_{jl}\beta _{,i}^l{\displaystyle \frac{2}{3}}\stackrel{~}{\gamma }_{ij}\beta _{,l}^l\right)_{,k}\right\}.`$ (41)
We also have two additional notes for handling the evolution equation of $`\stackrel{~}{A}_{ij}`$. One is on the method for evaluation of $`R_k^k`$ for which there are two options, use of the Hamiltonian constraint and direct calculation by
$`R_{ij}\gamma ^{ij}=e^{4\varphi }(\stackrel{~}{R}_k^k+R_{ij}^\varphi \stackrel{~}{\gamma }^{ij}).`$ (42)
We always adopt the latter one since with this, the conservation of the relation $`\stackrel{~}{A}_{ij}\stackrel{~}{\gamma }^{ij}=0`$ is much better preserved. The other is on the handling of a term of $`\stackrel{~}{\gamma }^{ij}\delta ^{kl}h_{ij,kl}`$ which appears in $`\stackrel{~}{R}_k^k`$. This term is written by
$$\stackrel{~}{\gamma }^{ij}\delta ^{kl}h_{ij,kl}=\delta ^{kl}h_{ij,k}f_{,l}^{ij},$$
(43)
where we use $`\mathrm{det}(\stackrel{~}{\gamma }_{ij})=1`$ (in the Cartesian coordinates).
As the time slicing condition, an approximate maximal slice condition $`K0`$ is adopted following previous papers (e.g., ). As the spatial gauge condition, we adopt a hyperbolic gauge condition as in . Successful numerical results for merger of binary neutron stars and stellar core collapse in these gauge conditions are presented in . We note that these are also different from those in .
### C GRMHD equations
Hydrodynamic equations are composed of
$`_\mu (\rho u^\mu )=0,`$ (44)
$`\gamma _i^\nu _\mu T_\nu ^\mu =0,`$ (45)
$`n^\nu _\mu T_\nu ^\mu =0.`$ (46)
The first, second, and third equations are the continuity, Euler, and energy equations, respectively. In the following, the equations are described for general coordinate systems since the hydrodynamic equations are solved in the cylindrical coordinates as well as in the Cartesian coordinates.
The continuity equation (44) is immediately written to
$`_t\rho _{}+{\displaystyle \frac{1}{\sqrt{\eta }}}_i(\rho _{}\sqrt{\eta }v^i)=0.`$ (47)
Equations (45) and (46) are rewritten as
$`_\mu (\sqrt{g}T_i^\mu ){\displaystyle \frac{\sqrt{g}}{2}}T^{\mu \nu }_ig_{\mu \nu }=0,`$ (48)
$`_\mu (\sqrt{g}T_\nu ^\mu n^\nu )\sqrt{g}T^{\mu \nu }_\mu n_\nu =0.`$ (49)
Then, using Eq. (23), they are written to
$`_\mu [\sqrt{g}(n^\mu J_j+\gamma ^{\mu i}S_{ij})]=\sqrt{\gamma }\left(\rho _\mathrm{H}_j\alpha +J_i_j\beta ^i{\displaystyle \frac{\alpha }{2}}S_{ik}_j\gamma ^{ik}\right),`$ (50)
$`_\mu [\sqrt{g}(\rho _\mathrm{H}n^\mu +\gamma ^{\mu i}J_i)]=\sqrt{\gamma }\left(\alpha K^{ij}S_{ij}J_iD^i\alpha \right),`$ (51)
where we use
$`n^\mu n^\nu _jg_{\mu \nu }=2_j\mathrm{ln}\alpha ,`$ (52)
$`n^\mu \gamma _i^\nu _jg_{\mu \nu }=\alpha ^1\gamma _{ik}_j\beta ^k,`$ (53)
$`\gamma _i^\mu \gamma _k^\nu _jg_{\mu \nu }=_j\gamma _{ik},`$ (54)
$`_\mu n_\nu =K_{\mu \nu }n_\mu D_\nu \mathrm{ln}\alpha .`$ (55)
The explicit forms of Eqs. (50) and (51) are
$`_tS_j+{\displaystyle \frac{1}{\sqrt{\eta }}}_i\left[\sqrt{\eta }\left\{S_jv^i+\alpha e^{6\varphi }P_{\mathrm{tot}}\delta _j^i{\displaystyle \frac{\alpha }{w^2e^{6\varphi }}}^i(_j+u_j^ku_k)\right\}\right]`$ (56)
$`=S_0_j\alpha +S_k_j\beta ^k+\alpha e^{6\varphi }\left[2S_k^k_j\varphi +P_{\mathrm{tot}}_j\mathrm{ln}\sqrt{\eta }\right]{\displaystyle \frac{1}{2}}\alpha e^{2\varphi }\widehat{S}_{ik}_j\stackrel{~}{\gamma }^{ik},`$ (57)
$`_tS_0+{\displaystyle \frac{1}{\sqrt{\eta }}}_i\left[\sqrt{\eta }\left\{S_0v^i+e^{6\varphi }P_{\mathrm{tot}}(v^i+\beta ^i){\displaystyle \frac{\alpha }{we^{6\varphi }}}(^ku_k)^i\right\}\right]`$ (58)
$`={\displaystyle \frac{1}{3}}\alpha e^{6\varphi }KS_k^k+\alpha e^{2\varphi }\widehat{S}_{ij}\stackrel{~}{A}^{ij}S_kD^k\alpha ,`$ (59)
where
$`\widehat{S}_{ij}=S_{ij}P_{\mathrm{tot}}\gamma _{ij}.`$ (60)
In the axisymmetric case, the equations for $`(\rho _{},S_i,S_0)`$ should be written in the cylindrical coordinates $`(\varpi ,\phi ,z)`$ when we adopt the Cartoon method for solving Einstein’s evolution equations . On the other hand, in the standard Cartoon method, Einstein’s equations are solved in the $`y=0`$ plane for which $`x=\varpi `$, $`S_\varpi =S_x`$, $`S_\phi =xS_y`$, and other similar relations hold for vector and tensor quantities. Taking into this fact, the hydrodynamic equations in axisymmetric spacetimes may be written using the Cartesian coordinates replacing $`(\varpi ,\phi )`$ by $`(x,y)`$. Then,
$`_t\rho _{}+{\displaystyle \frac{1}{x}}_x(\rho _{}v^xx)+_z(\rho _{}v^z)=0,`$ (61)
$`_tS_A+{\displaystyle \frac{1}{x}}_x\left[x\left\{S_Av^x+\alpha e^{6\varphi }P_{\mathrm{tot}}\delta _A^x{\displaystyle \frac{\alpha }{w^2e^{6\varphi }}}^x\left(_A+u_A^iu_i\right)\right\}\right]`$ (62)
$`+_z\left[S_Av^z+\alpha e^{6\varphi }P_{\mathrm{tot}}\delta _A^z{\displaystyle \frac{\alpha }{w^2e^{6\varphi }}}^z(_A+u_A^iu_i)\right]`$ (63)
$`=S_0_A\alpha +S_k_A\beta ^k+\alpha e^{6\varphi }\left[2S_k^k_A\varphi +{\displaystyle \frac{P_{\mathrm{tot}}}{x}}\delta _A^x\right]{\displaystyle \frac{1}{2}}\alpha e^{2\varphi }\widehat{S}_{ik}_A\stackrel{~}{\gamma }^{ik}`$ (64)
$`+\left[{\displaystyle \frac{S_yv^y}{x}}{\displaystyle \frac{\alpha }{xw^2e^{6\varphi }}}^y(_y+^iu_iu_y)\right]\delta _A^x,`$ (65)
$`_tS_y+{\displaystyle \frac{1}{x^2}}_x\left[x^2\left\{S_yv^x{\displaystyle \frac{\alpha }{w^2e^{6\varphi }}}^x\left(_y+u_y^iu_i\right)\right\}\right]`$ (66)
$`+_z\left[S_yv^z{\displaystyle \frac{\alpha }{w^2e^{6\varphi }}}^z\left(_y+u_y^iu_i\right)\right]=0,`$ (67)
$`_tS_0+{\displaystyle \frac{1}{x}}_x\left[x\left\{S_0v^x+e^{6\varphi }P_{\mathrm{tot}}(v^x+\beta ^x){\displaystyle \frac{\alpha }{we^{6\varphi }}}^iu_i^x\right\}\right]+_z\left[S_0v^z+e^{6\varphi }P_{\mathrm{tot}}(v^z+\beta ^z){\displaystyle \frac{\alpha }{we^{6\varphi }}}^iu_i^z\right]`$ (68)
$`={\displaystyle \frac{1}{3}}\alpha e^{6\varphi }KS_k^k+\alpha e^{2\varphi }\widehat{S}_{ij}\stackrel{~}{A}^{ij}S_kD^k\alpha ,`$ (69)
where $`A`$ denotes $`x`$ or $`z`$, while $`i,j,k,\mathrm{}`$ are $`x`$ or $`y`$ or $`z`$.
After evolving $`\rho _{}`$, $`S_i`$, and $`S_0`$ together with $`^i`$ (see next subsection for the equations), we have to determine the primitive variables such as $`\rho `$, $`\epsilon `$, $`u_i`$, and $`u^t`$ (or $`w=\alpha u^t`$). For this procedure, we make an equation from the definition of $`S_i`$ as
$`s^2\rho _{}^2\gamma ^{ij}S_iS_j=\left(h+B^2w^1\right)^2(w^21)D^2(hw)^2(B^2+2hw),`$ (70)
where $`B^2`$ and $`D^2`$ are determined from the evolved variables $`(\rho _{},S_i,^i,\varphi )`$ as
$`B^2={\displaystyle \frac{^2}{\rho _{}e^{6\varphi }}}\mathrm{and}D^2={\displaystyle \frac{(^iS_i)^2}{\rho _{}^3e^{6\varphi }}},`$ (71)
and for getting Eq. (70), we use the relation $`S_i^i=\rho _{}h^iu_i`$. Equation (70) is regarded as a function of $`h`$ and $`w`$ for given data sets of $`s^2`$, $`B^2`$, and $`D^2`$. From the definition of $`S_0`$, we also make a function of $`h`$ and $`w`$ as
$`{\displaystyle \frac{S_0}{\rho _{}}}=hw{\displaystyle \frac{P}{\rho w}}+B^2{\displaystyle \frac{1}{2w^2}}(B^2+D^2h^2).`$ (72)
Here, $`P/\rho `$ may be regarded as a function of $`h`$ and $`w`$ for a given data sets of $`\rho _{}`$ and $`S_0`$. This is indeed the case for frequently used equations of state such as $`\mathrm{\Gamma }`$-law equations of state $`P=(\mathrm{\Gamma }1)\rho \epsilon `$ where $`\mathrm{\Gamma }`$ is an adiabatic constant and hybrid equations of state for which $`P`$ is written in the form $`P(\rho ,h)`$ (e.g., ). Thus, Eqs. (70) and (72) constitute simultaneous equations for $`h`$ and $`w`$ for given values of $`\rho _{}`$, $`S_i`$, $`S_0`$, $`^i`$, and geometric variables. The solutions for $`h`$ and $`w`$ are numerically computed by the Newton-Raphson method very simply. Typically, a convergent solution is obtained with four iterations according to our numerical experiments.
### D Maxwell equations
The Maxwell equations are
$`_\mu F^{\mu \nu }=4\pi j^\nu ,`$ (73)
$`_\mu F_{\alpha \beta }+_\alpha F_{\beta \mu }+_\beta F_{\mu \alpha }=0.`$ (74)
In the ideal MHD, Eq. (73) is not necessary, and only Eq. (74) has to be solved. Using the dual tensor, Eq. (74) is rewritten to
$`_\mu F_\nu ^\mu =0.`$ (75)
This immediately leads to
$`_k(\eta ^{1/2}^k)=0,`$ (76)
$`_t^k={\displaystyle \frac{1}{\eta ^{1/2}}}_i\left[\eta ^{1/2}(^iv^k^kv^i)\right].`$ (77)
Equation (76) is the no-monopoles constraint, and Eq. (77) is the induction equation. The constraint equation (76) is solved in giving initial conditions, and the induction equation is solved for the evolution.
In the axisymmetric case, these equations in the $`y=0`$ plane are written as
$`{\displaystyle \frac{1}{x}}_x(x^x)+_z^z=0,`$ (78)
$`_t^x=_z(^xv^z^zv^x),`$ (79)
$`_t^z={\displaystyle \frac{1}{x}}_x\left[x(^xv^z^zv^x)\right],`$ (80)
$`_t^y=_x(^xv^y^yv^x)+_z(^zv^y^yv^z).`$ (81)
Equations (79)–(81) together with Eqs. (61)–(69) constitute basic equations for ideal MHD in the axisymmetric case.
### E Definition of global quantities
In numerical simulations for self-gravitating system, in addition to the total baryon rest mass $`M_{}`$, we refer to the ADM mass and the angular momentum of the system, which are given by
$`M`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _r\mathrm{}}_i\psi dS_i`$ (83)
$`={\displaystyle \left[\rho _\mathrm{H}e^{5\varphi }+\frac{e^{5\varphi }}{16\pi }\left(\stackrel{~}{A}_{ij}\stackrel{~}{A}^{ij}\frac{2}{3}K^2\stackrel{~}{R}_k^ke^{4\varphi }\right)\right]d^3x},`$
$`J`$ $`{\displaystyle \frac{1}{8\pi }}{\displaystyle _r\mathrm{}}\phi ^i\stackrel{~}{A}_i^je^{6\varphi }𝑑S_j`$ (85)
$`={\displaystyle e^{6\varphi }\left[S_i\phi ^i+\frac{1}{8\pi }\left(\stackrel{~}{A}_i^j_j\phi ^i\frac{1}{2}\stackrel{~}{A}_{ij}\phi ^k_k\stackrel{~}{\gamma }^{ij}+\frac{2}{3}\phi ^j_jK\right)\right]d^3x},`$
where $`dS_j=r^2_jrd(\mathrm{cos}\theta )d\phi `$, $`\phi ^j=y(_x)^j+x(_y)^j`$, and $`\psi =e^\varphi `$. In this paper, simulations are performed in axial symmetry, and hence, $`J`$ is conserved. $`M`$ is approximately conserved since the emission of gravitational waves is negligible. Thus, conservation of these quantities is checked during numerical simulations.
The violation of the Hamiltonian constraint is locally measured by the equation as
$`f_\psi `$ $`\left|\stackrel{~}{\mathrm{\Delta }}\psi {\displaystyle \frac{\psi }{8}}\stackrel{~}{R}_k^k+2\pi \rho _\mathrm{H}\psi ^5+{\displaystyle \frac{\psi ^5}{8}}\left(\stackrel{~}{A}_{ij}\stackrel{~}{A}^{ij}{\displaystyle \frac{2}{3}}K^2\right)\right|`$ (87)
$`\times \left[|\stackrel{~}{\mathrm{\Delta }}\psi |+|{\displaystyle \frac{\psi }{8}}\stackrel{~}{R}_k^k|+|2\pi \rho _\mathrm{H}\psi ^5|+{\displaystyle \frac{\psi ^5}{8}}\left(\stackrel{~}{A}_{ij}\stackrel{~}{A}^{ij}+{\displaystyle \frac{2}{3}}K^2\right)\right]^1.`$
Following , we define and monitor a global quantity as
$$H\frac{1}{M_{}}\rho _{}f_\psi d^3x.$$
(88)
Hereafter, this quantity will be referred to as the averaged violation of the Hamiltonian constraint.
## III Numerical scheme for solving GRMHD equations
### A GRMHD equations
As described in Sec. II C, we write the GRMHD equations in the conservative form. In this case, roughly speaking, there are two options for numerically handling the transport terms . One is to use the Godunov-type, approximate Riemann solver , and the other is to use the high-resolution central (HRC) scheme . We adopt a HRC scheme proposed by Kurganov and Tadmor and very recently used in special relativistic simulations by Lucas-Sarrano et al. . Thus our numerical scheme for a solution of GRMHD equations is slightly different from that in , in which the HLL scheme is basically adopted.
The basic equations can be schematically written as
$`{\displaystyle \frac{𝐔}{t}}+{\displaystyle \frac{𝐅^i}{x^i}}+{\displaystyle \frac{\mathrm{ln}\sqrt{\eta }}{x^i}}𝐅^i=𝐒,`$ (89)
where
$`𝐔=(\rho _{},S_i,S_0,^i),`$ (90)
$`𝐅^j=(\rho _{}v^j,S_iv^j+\alpha e^{6\varphi }P_{\mathrm{tot}}\delta _i^j\tau _i^{Bj},S_0v^j+e^{6\varphi }P_{\mathrm{tot}}(v^j+\beta ^j)\tau _0^{Bj},^iv^j^jv^i),`$ (91)
and $`𝐒`$ denotes the terms associated with the gravitational force. Here, $`\tau _\mu ^{Bj}`$ denotes a magnetic stress defined by
$`\tau _i^{Bj}={\displaystyle \frac{\alpha }{w^2e^{6\varphi }}}^j[_i+u_i(^ku_k)],`$ (92)
$`\tau _0^{Bj}={\displaystyle \frac{\alpha }{we^{6\varphi }}}(^ku_k)^j.`$ (93)
In addition to $`𝐔`$, we define a set of variables as
$`𝐏=(\rho _{},\widehat{u}_i,\epsilon ,^i).`$ (94)
$`\widehat{u}_i`$ and $`\epsilon `$ are computed at each time step from Eqs. (70) and (72). We use $`𝐏`$ for the reconstruction of $`𝐅`$ at cell interfaces. In standard method, one often uses a set of primitive variables $`(\rho ,v^i,\epsilon ,^i)`$ instead of $`𝐏`$ for reconstruction of $`𝐅`$. We have found in the test problems that even using $`\rho _{}`$ and $`\widehat{u}_i`$ instead of $`\rho `$ and $`v^i`$, it is possible to guarantee the similar accuracy and stability.
To evaluate $`𝐅`$, we use a HRC scheme . The fluxes are defined at cell faces. A piece-wise parabolic interpolation from the cell centers gives $`𝐏_R`$ and $`𝐏_L`$, the primitive variables at the right- and left-hand side of each cell interface, as
$`Q_L=Q_i+{\displaystyle \frac{\mathrm{\Phi }(r_{i1}^+)\mathrm{\Delta }_{i1}}{6}}+{\displaystyle \frac{\mathrm{\Phi }(r_i^{})\mathrm{\Delta }_i}{3}},`$ (95)
$`Q_R=Q_{i+1}{\displaystyle \frac{\mathrm{\Phi }(r_i^+)\mathrm{\Delta }_i}{3}}{\displaystyle \frac{\mathrm{\Phi }(r_{i+1}^{})\mathrm{\Delta }_{i+1}}{6}}.`$ (96)
Here, $`Q`$ denotes a component of $`𝐏`$ and $`\mathrm{\Delta }_{i+1}=Q_{i+1}Q_i`$. $`\mathrm{\Phi }`$ denotes a limiter function defined by
$`\mathrm{\Phi }(r)=\mathrm{minmod}(1,br)(1b4\mathrm{for}\mathrm{TVD}\mathrm{condition}),`$ (97)
where $`r_i^\pm =\mathrm{\Delta }_{i\pm 1}/\mathrm{\Delta }_i`$, and
$`\mathrm{minmod}(1,x)=\{\begin{array}{cc}1\hfill & \mathrm{if}x>1\hfill \\ x\hfill & \mathrm{if}1>x>0\hfill \\ 0\hfill & \mathrm{if}x<0\hfill \end{array}.`$ (101)
For the simulations presented in Secs. V and VI, we choose $`b=2`$ unless otherwise stated. We have found that the dissipation is relatively large for $`b=1`$ with which it is difficult to evolve isolated neutron stars for a long time scale accurately. On the other hand, for $`b3`$, the dissipation is so small that instabilities often occur around strong discontinuities, and around the region for which $`P_{\mathrm{tot}}P`$.
From $`𝐏_L`$ and $`𝐏_R`$, we calculate the maximum wave speed $`c_L`$ and $`c_R`$, and the fluxes $`𝐅_L`$ and $`𝐅_R`$ at the right- and left-hand sides of each cell interface. Then, we define $`c_{\mathrm{max}}=\mathrm{max}(c_L,c_R)`$, and the flux
$`𝐅={\displaystyle \frac{1}{2}}\left[𝐅_L+𝐅_Rc_{\mathrm{max}}(𝐔_R𝐔_L)\right].`$ (102)
In adopting the central schemes, the eigen vectors for the Jacobi matrix $`𝐅/𝐔`$ are not required in contrast to the case with the Godunov-type scheme . However, the eigen values for each direction are still necessary to evaluate characteristic wave speeds $`c_L`$ and $`c_R`$. The equation for the seven eigen values $`\lambda `$ is derived by Anile and Pennisi : Three of the seven solutions for $`\lambda `$ in $`x^i`$ direction are described by
$`\lambda =v^i,{\displaystyle \frac{b^i\pm u^i\sqrt{\rho h+b^2}}{b^t\pm u^t\sqrt{\rho h+b^2}}},`$ (103)
and rest four are given by the solutions for the following fourth order equation
$`(u^t)^4(\lambda v^i)^4(1\zeta )+\left[c_s^2{\displaystyle \frac{(b^i\lambda b^t)^2}{\rho h+b^2}}(u^t)^2(\lambda v^i)^2\left(\gamma ^{ii}{\displaystyle \frac{\beta ^i+\lambda }{\alpha ^2}}\right)\zeta \right]=0(\mathrm{no}\mathrm{summation}\mathrm{for}i).`$ (104)
Here, $`\zeta `$, the sound velocity $`c_s`$, and the Alfvén velocity $`v_A`$ are defined, respectively, by
$`\zeta v_A^2+c_s^2v_A^2c_s^2,`$ (105)
$`c_s^2{\displaystyle \frac{1}{h}}\left[{\displaystyle \frac{P}{\rho }}|_\epsilon +{\displaystyle \frac{P}{\rho ^2}}{\displaystyle \frac{P}{\epsilon }}|_\rho \right],`$ (106)
$`v_A^2{\displaystyle \frac{b^2}{\rho h+b^2}}.`$ (107)
In the central schemes, we only need the maximum characteristic speed, and thus, only the solutions for Eq. (104), which contain the fast mode, are relevant. The solutions for the fourth-order equation are determined either analytically or by standard numerical methods. However, for simplicity and for saving computational time, we use the prescription proposed by Gammie et al. , who have found it convenient to replace the fourth-order equation approximately by a second-order one:
$`(u^t)^2(\lambda v^i)^2(1\zeta )\zeta \left(\gamma ^{ii}{\displaystyle \frac{\beta ^i+\lambda }{\alpha ^2}}\right)=0(\mathrm{no}\mathrm{summation}\mathrm{for}i).`$ (108)
The solution of Eq. (108) for an arbitrary direction $`x^i`$ is written as
$`\lambda ^i={\displaystyle \frac{1}{\alpha ^2V_kV^k\zeta }}`$ $`[v^i\alpha ^2(1\zeta )\beta ^i\zeta (\alpha ^2V^2)`$ (110)
$`\pm \alpha \sqrt{\zeta }\sqrt{(\alpha ^2V^2)\{\gamma ^{ii}(\alpha ^2V^2\zeta )(1\zeta )V^iV^i\}}](\mathrm{no}\mathrm{summation}\mathrm{for}i),`$
where $`V^i=v^i+\beta ^i`$ and $`V^2=\gamma _{ij}V^iV^j`$. This is equivalent to that obtained by replacing $`c_s`$ by $`\sqrt{\zeta }`$ in the solution for the pure hydrodynamic case .
### B Induction equation
The induction equation may be solved using the same scheme as in solving the hydrodynamic equations described above. However, with such a scheme, the violation of the constraint equation (76) is often accumulated with time, resulting in a nonreliable solution. Thus, we adopt a constraint transport scheme . Namely, we put the components of the magnetic field at the cell-face centers. Here, we specifically consider the axisymmetric case with the cylindrical coordinates $`(x,\phi ,z)`$ ($`\varpi `$ is replaced by $`x`$). Extension to the nonaxisymmetric case is straightforward, and the description below can be used with slight modification.
In the axisymmetric case with the cylindrical coordinates, the numerical computation is performed in a discretized cell for $`(x,z)`$. Here, we denote the cell center for $`(x,z)`$ by $`(i,j)`$. Then, we put $`^x`$ at $`(i+1/2,j)`$, and $`^z`$ at $`(i,j+1/2)`$ while components of the gravitational field and fluid variables as well as $`^yx^\phi `$ are put at the cell center $`(i,j)`$. In this case, the induction equations for $`^x`$ and $`^z`$ are solved in a constraint transport scheme , while that for $`^y`$ is solved in the same method as that for the continuity equation of $`\rho _{}`$.
Computing the flux at cell edges for the induction equation is different from that for the fluid equation. This is because numerical fluxes have to be defined so that the constraint equation (76) is satisfied. For example, for the $`x`$ component of the induction equation, the flux in the $`z`$ direction is written as $`v^z^xv^x^zF_1`$. On the other hand, for the $`z`$ component of the induction equation, the flux in the $`x`$ direction is written as $`v^x^zv^z^xF_2`$. Both $`F_1`$ and $`F_2`$ have to be defined at cell edges $`(i+1/2,j+1/2)`$, and for the constraint equation (76) to be satisfied, we have to require $`F_1=F_2=F`$ at each cell edge. In addition, an upwind scheme should be adopted for numerical stability: For the induction equation of $`^x`$, the upwind prescription should be applied for the $`z`$ component of the flux. On the other hand, for the induction equation of $`^z`$, the upwind prescription should be applied for the $`x`$ component of the flux. $`F`$ has to be determined taking into account these requirements.
We here adopt a scheme proposed by Del Zanna et al. , which satisfies such requirements. In this scheme, the flux is written as
$`F={\displaystyle \frac{1}{4}}(F^{LL}+F^{LR}+F^{RL}+F^{RR}){\displaystyle \frac{c_{\mathrm{max}}^z}{2}}[(^x)^R(^x)^L]+{\displaystyle \frac{c_{\mathrm{max}}^x}{2}}[(^z)^R(^z)^L],`$ (111)
where, e.g., $`F^{LR}`$ is the flux defined at the left-hand side in the $`x`$ direction and at right-hand side in the $`z`$ direction. These fluxes are computed by a piece-wise parabolic interpolation. $`c_{\mathrm{max}}^i`$ is the characteristic speed for the prescription of an upwind flux-construction and calculated at cell edges using the interpolated variables. For simplicity, we set $`c_{\mathrm{max}}^z=\mathrm{max}(v_L^z,v_R^z)`$ and $`c_{\mathrm{max}}^x=\mathrm{max}(v_L^x,v_R^x)`$.
For solving other equations, it is necessary to define the magnetic field at the cell center. Since the $`x`$ and $`z`$ components of the magnetic field are defined at the cell face centers (i.e., $`^x`$ at $`(i+1/2,j)`$ and $`^z`$ at $`(i,j+1/2)`$), this is done by a simple averaging as
$`_{i,j}^x={\displaystyle \frac{1}{2}}(_{i+1/2,j}^x+_{i1/2,j}^x),`$ (112)
$`_{i,j}^z={\displaystyle \frac{1}{2}}(_{i,j+1/2}^z+_{i,j1/2}^z).`$ (113)
Also, $`v^i`$ at the cell face center is necessary for computing $`c_{\mathrm{max}}^x`$ and $`c_{\mathrm{max}}^z`$. To compute them, we also use a simple averaging. For the definition of $`v_{i+1/2,j}^k`$ and $`v_{i,j+1/2}^k`$, we have also tried the Roe-type averaging in terms of $`\rho _{}^{1/2}`$, but any significant modification in the results has not been found.
Before closing this section, we note that our scheme for the induction equation is different from that adopted in , in which a Tóth’s method is used .
## IV Initial value problem
In the fully general relativistic and dynamical simulations, we have to solve the constraint equations of general relativity for preparing the initial condition. One solid method is to give an equilibrium state. For rigidly rotating stars of poloidal magnetic fields in axial symmetry, such equilibrium has been already computed . However, for differentially rotating stars or nonaxisymmetric cases, the method has not been established. Thus, we here present a simple method for preparing an initial condition which is similar to that in . In the following, we assume that axisymmetric matter fields $`\rho _{}`$, $`\widehat{e}hwP/\rho w`$, $`h`$, and $`\widehat{u}_ihu_i`$ are a priori given (e.g., those for rotating stars of no magnetic field in equilibrium are given). Although we assume the axial symmetry, the same method can be applied for the nonaxisymmetric case.
Initial conditions for magnetic fields have to satisfy Eq. (76). A solution of Eq. (76) is written as
$`^k=e^{kij}_iA_j,`$ (114)
where $`A_j`$ is an arbitrary vector potential and $`e^{kij}`$ is a Levi-Civita tensor of flat three-space. If we choose
$`A_x=A_z=0,\mathrm{and}A_\phi 0,`$ (115)
the magnetic fields are poloidal. Here, we assume to use the cylindrical coordinates $`(x,\phi ,z)`$ ($`\varpi `$ is replaced by $`x`$).
In the axisymmetric case, we can also choose pure toroidal magnetic fields as
$`^x=^z=0,\mathrm{and}^\phi 0,`$ (116)
where $`^\phi `$ may be an arbitrary function. In the following, we give a nonzero function either for $`A_\phi `$ or for $`^\phi `$.
Initial conditions also have to satisfy Eqs. (28) and (29). In the following, we assume that $`\stackrel{~}{\gamma }_{ij}`$ and $`K`$ are given functions in these equations. Remind that $`\rho _\mathrm{H}`$ and $`J_i`$ are written as
$`\rho _\mathrm{H}=\rho _{}\widehat{e}\psi ^6+\psi ^{12}\left(^2{\displaystyle \frac{^2+(^iu_i)^2}{2w^2}}\right),`$ (117)
$`J_i=\rho _{}\widehat{u}_i\psi ^6+{\displaystyle \frac{1}{w\psi ^{12}}}\left(^2u_i_i^ju_j\right),`$ (118)
where $`\psi `$ denotes the conformal factor $`(=e^\varphi )`$, and $`w=\sqrt{1+\psi ^4\stackrel{~}{\gamma }^{ij}u_iu_j}`$. Thus, if $`\rho _{}`$, $`\widehat{e}`$, $`h`$, $`\widehat{u}_i`$, $`A_\phi `$, $`^\phi `$, $`K`$, and $`\stackrel{~}{\gamma }_{ij}`$ are given, the remaining unknown functions are $`\psi `$ and $`\stackrel{~}{A}_{ij}`$. This implies that the constraint equations are solved for these variables using the technique developed by York .
First, we decompose the tracefree part of the extrinsic curvature as
$`\widehat{A}_{ij}\psi ^6\stackrel{~}{A}_{ij}=\stackrel{~}{D}_iW_j+\stackrel{~}{D}_jW_i{\displaystyle \frac{2}{3}}\stackrel{~}{\gamma }_{ij}\stackrel{~}{D}_kW^k+K_{ij}^{\mathrm{TT}},`$ (119)
where $`W_i`$ is a three vector, and $`K_{ij}^{\mathrm{TT}}`$ is a transverse-tracefree tensor which satisfies
$`\stackrel{~}{D}^iK_{ij}^{\mathrm{TT}}=0=K_{ij}^{\mathrm{TT}}\stackrel{~}{\gamma }^{ij}.`$ (120)
$`K_{ij}^{\mathrm{TT}}`$ would be composed mainly of gravitational waves. Hereafter, we set $`K_{ij}^{\mathrm{TT}}=0`$ for simplicity. Using Eq. (119), Eq. (29) is rewritten to
$`\stackrel{~}{\mathrm{\Delta }}W_j+{\displaystyle \frac{1}{3}}\stackrel{~}{D}_j\stackrel{~}{D}_iW^i+\stackrel{~}{R}_{ji}W^i{\displaystyle \frac{2}{3}}\psi ^6\stackrel{~}{D}_jK=8\pi J_i\psi ^6.`$ (121)
This equation can be solved for an initial trial function of $`\psi `$. Then, $`\widehat{A}_{ij}`$ is computed from Eq. (119). Substituting $`\widehat{A}_{ij}`$, the Hamiltonian constraint (28) is solved in the next step. Then we solve the momentum constraint again, and repeat these procedures until a sufficient convergence is achieved.
## V Special relativistic tests
In this section, we present numerical results for a number of special relativistic tests. In the tests, we adopt the $`\mathrm{\Gamma }`$-law equations of state as
$`P=(\mathrm{\Gamma }1)\rho \epsilon ,`$ (122)
with $`\mathrm{\Gamma }=4/3`$ or $`5/3`$. Simulations are always performed using the uniform grid in all the axis directions.
### A One dimensional tests
Any numerical implementation of the MHD equations has to be checked if it can produce the basic waves such as shock and rarefaction waves accurately. Komissarov has proposed a suite of one-dimensional test problems in special relativity: Propagation of fast and slow shocks, fast and slow rarefaction waves, Alfvén waves, compound waves, shock tube tests, and collision of two flows. We have performed all the tests except for the compound wave following . Our implementation can integrate each of remaining eight tests although in some cases we have to reduce the Courant number significantly to avoid numerical instabilities as reported by Gammie . On the other hand, we adopt the same limiter, $`b=2`$, for all the simulations. Numerical results are shown in Figs. 15. Grid size $`N`$ and spacing $`\mathrm{\Delta }x`$ we adopt for each of test simulations are approximately the same as those by Komissarov, and described in the figure captions.
Figure 1 shows the results for fast and slow shocks. In these problems, the system is stationary with respect to the frame comoving with the shock front. The velocity of the shocks is 0.2 and 0.5 for the fast and slow shocks, respectively. As the previous works illustrate , the fast shock can be computed accurately with a relatively large grid spacing. On the other hand, in the numerical solution of the slow shock, a spurious modulation is found for $`\rho `$ in the region of $`1<x<1.3`$ as in the previous works . This is always generated soon after the onset of the simulation irrespective of grid resolutions. Thus, it is impossible to avoid such small error in our implementation. Although the modulation is always present, its wavelength and amplitude gradually decrease with improving the grid resolution. We computed an L1 norm defined for the difference between the numerical and exact solutions, and found that it decreases as the grid spacing is smaller. In this case, the convergence is achieved at first order since discontinuities are present, around which the transport terms of hydrodynamic equations are computed with the first-order accuracy.
Figure 2 shows the results for switch-off and switch-on rarefaction waves. Although we have not compared the results precisely with those by other authors , the accuracy of our results is similar to that reported by others. For the switch-off waves, a spurious bump is found at $`x0.6`$ as in the previous works . As in the slow shock problem, this bump is generated at $`t=0`$ irrespective of grid resolutions, and with improving the grid resolution, the magnitude of the L1 norm decreases at first order. On the other hand, the numerical solution for the switch-on waves, spurious bumps are not present, and with $`\mathrm{\Delta }x=0.005`$, a good convergent result appears to be obtained with our implementation.
Figure 3(a) shows the results for an Alfvén wave test, demonstrating that the Alfvén wave can be computed accurately with our implementation as in . In this problem, the density and pressure should be unchanged. In our results, this is achieved within $`1\%`$ error for $`\mathrm{\Delta }x=0.0025.`$ Since no discontinuities are present in this problem, the convergence of the numerical solution to the exact one should be achieved approximately at second order . To check if this is the case, we compute an L1 norm defined by the difference between the numerical and exact solutions for $`\rho `$ and $`P`$. The results are shown in Fig. 3(b), which illustrates that the convergence is achieved approximately at second order (slightly better than second order).
In Fig. 4, numerical results for shock-tube problems are presented. For the problem of Fig. 4(a), shocks are very strong since the ratio of the pressure in the left- and right-hand sides at $`t=0`$ is $`10^3`$. However, since the magnetic field lines are normal to the discontinuities, the effects of the magnetic field for the formation and propagation of shocks are absent. In this case, a large spurious overshooting is found around the shock for $`u^x`$. This is partly due to our limiter ($`b=2`$) which is not very dissipative. If we use the minmod limiter ($`b=1`$), height of the overshooting decreases although the shocks are less sharply computed. For the problem of Fig. 4(b), shocks are not as strong as those in 4(a). However, the magnetic fields affect the formation and propagation of shocks since they are parallel to the shocks. The results shown in Fig. 4(b) are very similar to those in , and hence, are likely to be as accurate as them. This indicates that our implementation can compute magnetized shocks as accurately as the previous ones.
In Fig. 5, numerical results for collision of two magnetized flows are presented. It shows that four separate discontinuities generated at $`t=0`$ are computed accurately. As found in previous papers , a small dip spuriously appears in $`\rho `$ around $`x=0`$. As in the case of the slow shock and switch-off rarefaction wave, this is spuriously generated at $`t=0`$ irrespective of grid resolution, and with improving the resolution, the magnitude of the error is decreased at first order.
### B Multi dimensional tests
For multidimensional tests, following Del Zanna et al. , we performed simulations for (i) a cylindrical blast explosion, (ii) a rotating cylinder in two-dimensional Cartesian coordinates with a uniform magnetized medium, and (iii) propagation of a jet in cylindrical coordinates in a magnetized background. The parameters for the initial conditions adopted here are the same as those in . On the other hand, we varied the grid spacing to see the convergence in contrast to the previous works.
In the test (i), the Cartesian grid of $`(x,y)`$ is adopted with the range $`[0.6,0.6]`$ for both directions. The grid spacing chosen is 0.004, 0.005, and 0.006. The initial condition is
$`(\rho ,P,^x,^y)=\{\begin{array}{cc}(1,10^3,4,0)\hfill & \mathrm{for}\sqrt{x^2+y^2}0.08,\hfill \\ (1,10^2,4,0)\hfill & \mathrm{for}\sqrt{x^2+y^2}>0.08,\hfill \end{array}`$ (125)
with $`u^i=0`$ and $`\mathrm{\Gamma }=4/3`$. Because of the large internal energy in the central region, the outward explosion occurs. In this problem, the shocks generated at $`t=0`$ are very strong, and hence, the minmod limiter with $`b=1`$ is adopted to avoid numerical instability. With the limiter of $`b=2`$, the computation crashes because of the appearance of negative internal energy (or $`h<1`$) irrespective of the grid resolution. We have found that the computation first crashes along the line of $`x=y`$ and $`x=y`$ for which the accuracy is likely to be worst.
In Fig. 6, we display the snapshot of the numerical results at $`t=0.4`$. In Fig. 7, configurations of the density, pressure, magnetic pressure, and Lorentz factor along $`x`$ and $`y`$ axes are shown for three levels of grid resolution. The expansion velocity of the blast wave is largest along the $`x`$ axis because of the confinement by the magnetic pressure. The maximum value of the Lorentz factor is about $`4`$ at $`t=0.4`$ with the best resolved case. Along the $`y`$ axis, the magnetic field lines are squeezed yielding the highest magnetic pressure. These features agree with those found in . As mentioned in , the total energy is completely conserved since we solve the MHD equations in the conservative form and do not add any dissipative terms in contrast with the treatment in .
One point to be mentioned is that convergence around the density peak along the $`x`$ axis is not achieved well within the adopted resolution although for other region, convergence is achieved well. The likely reason is that the discontinuities around the peak is very thin for which it is very difficult to resolve with the chosen grid resolutions. Thus, it is difficult to accurately derive the maximum values of the density, pressure, and Lorentz factor which are underestimated in this test problem.
In the test (ii), the Cartesian grid of $`(x,y)`$ is also adopted with the range $`[0.6,0.6]`$ for both directions. The grid spacing chosen is 0.0025, 0.003, and 0.004. The initial condition is
$`(\rho ,P,^x,^y)=\{\begin{array}{cc}(10,1,1,0)\hfill & \mathrm{for}\sqrt{x^2+y^2}0.1\hfill \\ (1,1,1,0)\hfill & \mathrm{for}\sqrt{x^2+y^2}>0.1,\hfill \end{array}`$ (128)
with $`v^i=[\omega y,\omega x]`$ for $`\sqrt{x^2+y^2}0.1`$ where $`\omega =0.995`$ and thus the Lorentz factor at the surface of the rotating cylinder is initially about 10. $`\mathrm{\Gamma }`$ is chosen to be 5/3 following .
In Fig. 8, we display the snapshot of the numerical results at $`t=0.4`$. In Fig. 9, configurations of the density, pressure, magnetic pressure, and Lorentz factor along $`x`$ and $`y`$ axes are shown for three levels of grid resolution. In this problem, the magnetic field lines keep winding-up, and at $`t=0.4`$, the central field lines are rotated by an angle of $`90`$ degrees. Because of magnetic braking, the rotational speed is decreased monotonically, and at $`t=0.4`$ the maximum Lorentz factor is decreased to $`1.7`$. Due to the outward explosion induced by the rotation, the density in the central region becomes an uniformly low value of $`0.44`$ while an ellipsoidal density peak is formed around the central region. As in the test (i), it is difficult to obtain a convergent value for the peak density with the chosen grid resolutions. The likely reason is that the thickness of the density peak is so small that the grid resolutions are not sufficient. However, for the other region, convergent results are obtained.
In the test (iii), the cylindrical grid of $`(x,z)`$ is adopted with the range $`[0,8]`$ and $`[0,20]`$, respectively. The grid spacing is 0.06, 0.08, and 0.1. The initial condition is
$`(\rho ,P,v^z,^z)=\{\begin{array}{cc}(10,10^2,0.99,0.1)\hfill & \mathrm{for}0x1\mathrm{and}0z1\hfill \\ (0.1,10^2,0,0.1)\hfill & \mathrm{otherwise},\hfill \end{array}`$ (131)
with $`v^x=0`$, $`^x=0`$, and $`\mathrm{\Gamma }=5/3`$. The region with $`0x1`$ and $`0z1`$ is defined to be a jet-inlet zone, and the stationary condition is artificially imposed. In the simulation, the regularity condition is imposed along the symmetric axis $`x=0`$. For the boundary conditions at $`z=0`$, extrapolation is assumed following . In this test, we adopt the minmod limiter with $`b=1`$ since with $`b=2`$, the computation soon crashes irrespective of the grid resolutions.
In Fig. 10, we show the snapshot of the density contour curves and magnetic field lines at $`t=15`$ and $`30`$ with $`\mathrm{\Delta }x=0.06`$. The contour curves and field lines are similar to those in . The maximum Lorentz factor at $`t=0`$ is $`7.09`$. At the head of the jet, the density becomes maximum and shocks are formed, inducing back flows at the shocks. These flows make a cocoon which is to expand in the direction of the cylindrical radius, squeezing the magnetic field lines. A part of the matter is back-scattered toward the $`z=0`$ plane dragging the magnetic field lines together. As a result, the magnetic field lines are highly deformed. The deformation is computed more accurately with finer grid resolutions.
However, we found that precise computation for the deformation of the magnetic field lines increases the risk for crash of the computation. For $`\mathrm{\Delta }x=0.1`$ and 0.08, computations can be continued until the shock front of the jets reaches the outer boundary. However, for $`\mathrm{\Delta }x=0.06`$, the computation crashes at $`t35`$ in spite of the fact that the motion of the jet head is still stably computed. If we adopt a better resolution with $`\mathrm{\Delta }x<0.05`$, the computation crashes before $`t`$ reaches 30. The instabilities always occur near the boundary region of the jet-inlet zone around which the magnetic field configuration is deformed to be highly complicated. This seems to be due to the fact that we impose the stationary condition inside the jet-inlet zone. This artificial handling makes the field configuration near the boundary of the jet-inlet zone nonsmooth (i.e., the derivative of the magnetic field variables can be artificially larger for better grid resolutions). Here, we note that this problem happens only in the presence of magnetic fields. Thus, for continuing the computation for a longer time, probably, it is necessary to include a resistivity for inducing reconnections of magnetic fields near the jet-inlet zone for stabilization. The other method is to change the stationary condition we adopt here to other appropriate boundary conditions near the jet-inlet zone .
## VI General relativistic tests
### A Relativistic Bondi accretion
As the first test for general relativistic implementation, we perform a simulation for spherical accretion onto the fixed background of a Schwarzschild black hole. The relativistic Bondi solution is known to describe a stationary flow, and thus, by comparing the numerical solution with the analytical one, it is possible to check the suitability of the numerical implementation for general relativistic hydrodynamics problems . Furthermore, it has been shown that the relativistic Bondi solution is unchanged even in the presence of a divergence-free pure radial magnetic field . Thus, it can be also used for checking the GRMHD implementations. The advantage of this test is that the exact solution can be obtained very easily while it involves strong gravitational fields, relativistic flows, and strong magnetic fields all together.
Following previous authors , we write the metric in Kerr-Schild coordinates in which all the variables are well behaved at the event horizon ($`r=2M`$; where $`r`$ and $`M`$ are the radial coordinate and the mass of the black hole). Nevertheless, the hydrostatic equations for the stationary flow are the same forms as those in the Schwarzschild coordinates, and thus, the stationary solution is determined from an algebraic equation which can be easily solved by standard numerical methods .
For this test, we adopt the same solution used in . Namely, the sonic radius is set at $`r=8M`$, the accretion rate $`\dot{M}=4\pi \rho r^2u^r`$ is set to be $`1`$, and the adiabatic index for the equation of state is $`4/3`$. The simulation is performed in an axisymmetric implementation with the cylindrical coordinates $`(x,z)`$. The computational domain is set to be $`[0,18M]`$ for $`x`$ and $`z`$, and the radius of $`r=1.9M`$ is chosen as the excision radius. The uniform grid is adopted. At the excision radius and outer boundaries, we impose the condition that the system is stationary. The (semi) analytic solution for the stationary Bondi flow is put as the initial condition, and we evolve for $`100M`$ following previous authors .
The simulations are performed changing the grid spacing $`\mathrm{\Delta }x`$. Irrespective of the grid resolution, the system relaxes to a stationary state long before $`100M`$. When evolved with a finite-difference implementation, discretization errors will cause small deviations in the flow from the exact stationary configuration. These deviations should converge to zero at second order with improving the grid resolution. To diagnose the behavior of our numerical solution, we measure an L1 norm for $`\rho _{}\rho _{}^{\mathrm{exact}}`$ where $`\rho _{}^{\mathrm{exact}}`$ denotes the exact stationary value of $`\rho _{}`$. Specifically, the L1 norm is here defined by
$`{\displaystyle _{r2M}}|\rho _{}\rho _{}^{\mathrm{exact}}|d^3x/{\displaystyle _{r2M}}\rho _{}^{\mathrm{exact}}d^3x.`$ (132)
For the convergence test, the grid spacing is changed from $`0.06M`$ to $`0.4M`$.
The radial magnetic field strength is also changed for a wide range. Following , we denote the magnetic field strength by
$`\widehat{\beta }{\displaystyle \frac{b^2}{\rho }}|_{r=2M}.`$ (133)
We note that the ratio of the magnetic pressure to the gas pressure $`b^2/2P`$ is $`3.85\widehat{\beta }`$ at $`r=2M`$ for the solution chosen in this test problem.
In Fig. 11, we show the L1 norm as a function of the grid spacing for $`0\widehat{\beta }63`$. Irrespective of the magnetic field strength, the numerical solution converges approximately at second order to the exact solution for $`\mathrm{\Delta }x0`$. The L1 norm is larger for the stronger magnetic fields, implying that the relaxed state deviates more from the true stationary solution for the larger value of $`\widehat{\beta }`$. Specifically, the velocity field configuration deviates significantly from the exact solution with increasing the value of $`\widehat{\beta }`$, although the deviation for the density configuration is not very outstanding.
In this test simulation, we have found several interesting behaviors of the numerical solutions. First, for a given value of the grid spacing with $`\mathrm{\Delta }x>0.1M`$, there is the maximum allowed value of $`\widehat{\beta }`$ above which the computation crashes. The maximum value is larger for better grid resolution; e.g. for $`\mathrm{\Delta }x=0.1M`$, $`0.2M`$, and $`0.3M`$, the maximum allowed values of $`\widehat{\beta }`$ are $`45`$, $`25`$, and $`10`$, respectively. For $`\mathrm{\Delta }x<0.1M`$, on the other hand, the maximum allowed value is $`\widehat{\beta }70`$ irrespective of the grid resolution. The limitation is due to the well-known weak point in the conservative scheme that the small error in the magnetic energy density in the magnetically dominated region with $`\widehat{\beta }1`$ leads to fractionally large errors in other components of the total energy density, by which the computation crashes (typically, the internal energy density becomes negative). For the poorer grid resolutions, the numerical error is larger, and hence, the computation crashes for the smaller value of the magnetic field. Second, the maximum allowed value of $`\widehat{\beta }`$ found here $`(70)`$ is by about one order of magnitude smaller than that found in . This is probably due to the difference of the coordinate system adopted; we use the cylindrical coordinates while the authors in use the spherical polar coordinates which obviously have advantage for handling the spherically symmetric problem. However, we note that even in the cylindrical coordinates, it is possible to handle the flow with a very high value of $`\widehat{\beta }60`$ if a sufficient grid resolution is guaranteed. In , the authors suggest that in the cylindrical coordinates, the maximum allowed value of $`\widehat{\beta }`$ is at most $`10`$. We have not found such severe limitation in our numerical experiment. Their failure for simulating the flow with high values of $`\widehat{\beta }`$ is probably due to the fact that they use an excision boundary which may be applicable for more general problems (e.g., for simulation of dynamical spacetimes). Even in the cylindrical coordinates, a high value of $`\beta `$ will be achieved if a stationary inner boundary condition is imposed.
### B Longterm evolution for system of a rotating star and a disk with no magnetic field
Next, we illustrate that with our implementation (for axisymmetric systems), self-gravitating objects can be simulated accurately. In a previous paper , we have already illustrated that our implementation with a HRC scheme can simulate rapidly rotating compact neutron stars for more than 20 rotational periods accurately. Thus, we here choose a more complicated system; an equilibrium system composed of a rapidly rotating neutron star and a massive disk. By this test, it is possible to check that our implementation is applicable to a longterm evolution not only for an isolated rotating star but also for a self-gravitating disk rotating around a compact object.
The equilibrium configuration is determined by solving equations for the gravitational field and hydrostatic equations self consistently. For simplicity, we here adopt a conformally flat formalism for the spatial metric . As shown in , a good approximate solution for axisymmetric rotating stars can be obtained even in this approximation. Thus, the initial condition presented here can be regarded as a slightly perturbed equilibrium state. At the start of the simulations, we further add a slight perturbation by reducing the pressure by 0.1% to investigate if a quasiradial oscillation is followed stably and accurately. The magnitude of the perturbation in association with the conformally flat approximation is much smaller than this pressure perturbation.
The Euler equation for axisymmetric stars in equilibrium can be integrated to give the first integral, which is written as
$`\mathrm{ln}{\displaystyle \frac{h}{u^t}}+{\displaystyle u^tu_\phi 𝑑\mathrm{\Omega }}=C,`$ (134)
or
$`{\displaystyle \frac{h}{u^t}}+{\displaystyle hu_\phi 𝑑\mathrm{\Omega }}=C^{},`$ (135)
where $`C`$ and $`C^{}`$ are integral constants. Equation (134) is a well-known form . However here, we adopt Eq. (135), and set that the specific angular momentum $`hu_\phi `$ is constant $`(=j)`$ for the disk and $`\mathrm{\Omega }=`$const for the central star.
A hybrid, parametric equation of state is used in this simulation following previous papers . In this equation of state, one assumes that the pressure consists of the sum of polytropic and thermal parts as
$$P=P_\mathrm{P}+P_{\mathrm{th}}.$$
(136)
The polytropic part, which denotes the cold part of the equations of state, is given by
$`P_\mathrm{P}=\{\begin{array}{cc}K_1\rho ^{\mathrm{\Gamma }_1},\hfill & \rho \rho _{\mathrm{nuc}},\hfill \\ K_2\rho ^{\mathrm{\Gamma }_2},\hfill & \rho \rho _{\mathrm{nuc}},\hfill \end{array}`$ (139)
where $`K_1`$ and $`K_2`$ are polytropic constants. $`\rho _{\mathrm{nuc}}`$ denotes the nuclear density and is set to be $`2\times 10^{14}\mathrm{g}/\mathrm{cm}^3`$. In this paper, we choose $`\mathrm{\Gamma }_1=4/3`$ and $`\mathrm{\Gamma }_2=2.5`$. Since $`P_\mathrm{P}`$ should be continuous, the relation, $`K_2=K_1\rho _{\mathrm{nuc}}^{\mathrm{\Gamma }_1\mathrm{\Gamma }_2}`$, is required. Here, the value of $`K_1`$ is chosen to be $`2.5534\times 10^{14}`$ in the cgs unit. With this value, the maximum ADM (baryon rest) mass for the cold and spherical neutron star becomes about $`1.84M_{}`$ ($`2.05M_{}`$) which is a close to that derived in realistic equations of state .
Since the specific internal energy should be also continuous at $`\rho =\rho _{\mathrm{nuc}}`$, the polytropic specific internal energy $`\epsilon _\mathrm{P}`$ is defined as
$`\epsilon _\mathrm{P}=\{\begin{array}{cc}{\displaystyle \frac{K_1}{\mathrm{\Gamma }_11}}\rho ^{\mathrm{\Gamma }_11},\hfill & \rho \rho _{\mathrm{nuc}},\hfill \\ {\displaystyle \frac{K_2}{\mathrm{\Gamma }_21}}\rho ^{\mathrm{\Gamma }_21}+{\displaystyle \frac{(\mathrm{\Gamma }_2\mathrm{\Gamma }_1)K_1\rho _{\mathrm{nuc}}^{\mathrm{\Gamma }_11}}{(\mathrm{\Gamma }_11)(\mathrm{\Gamma }_21)}},\hfill & \rho \rho _{\mathrm{nuc}}.\hfill \end{array}`$ (142)
The thermal part of the pressure $`P_{\mathrm{th}}`$ plays an important role in the case that shocks are generated. $`P_{\mathrm{th}}`$ is related to the thermal energy density $`\epsilon _{\mathrm{th}}\epsilon \epsilon _\mathrm{P}`$ as
$$P_{\mathrm{th}}=(\mathrm{\Gamma }_{\mathrm{th}}1)\rho \epsilon _{\mathrm{th}}.$$
(143)
For simplicity, the value of $`\mathrm{\Gamma }_{\mathrm{th}}`$, which determines the strength of shocks, is chosen to be equal to $`\mathrm{\Gamma }_1`$. For computing initial equilibria, we set $`\epsilon =\epsilon _\mathrm{P}`$ and $`P=P_\mathrm{P}`$.
For the simulation, we choose a sufficiently deformed star with the axial ratio of the minor axis to major axis $`0.6`$. The ADM mass is $`1.888M_{}`$, total baryon rest mass $`2.074M_{}`$, the central density $`1.3\times 10^{15}\mathrm{g}/\mathrm{cm}^3`$, the circumferential radius at equator $`16.2`$ km, the rotational period $`P_c=0.841`$ ms, and $`J/M^2=0.545`$. Thus, the neutron star considered is massive and rapidly rotating. The baryon rest mass of the disk is much smaller than the central star as $`4.9\times 10^5M_{}`$ with the maximum density $`2\times 10^{10}\mathrm{g}/\mathrm{cm}^3`$. Since it is of low density, the disk is composed of $`\mathrm{\Gamma }=4/3`$ polytropic equation of state. Orbital radius of inner edges of the disk is $`20`$ km, and thus, the uniform specific angular momentum is small as $`j3.45M`$ which is very close to the value for a particle orbiting an innermost stable circular orbit. The rotational periods of the disk at the inner and outer edges in the equatorial plane are 1.03 ms($`=1.2P_c`$) and 2.58 ms($`=3.1P_c`$), respectively. The simulations are performed in axial symmetry with (241,241), (193, 193), and (161,161) grid sizes for which the grid spacing is 0.165, 0.202, and 0.248 km, respectively. The reflection symmetry with respect to the equatorial plane is assumed. The outer boundaries along each axis are located at 39.6 km.
An atmosphere of small density $`\rho =2\times 10^4\mathrm{g}/\mathrm{cm}^3`$ is added uniformly outside the neutron star and disk at $`t=0`$, since the vacuum is not allowed in grid-based hydrodynamics implementations. We note that the density of atmosphere can be chosen to be much smaller than the nuclear density $`\rho _{\mathrm{nuc}}`$. This is the advantage of HRC schemes in which such low density can be handled in contrast with high-resolution shock-capturing schemes . Since the atmosphere is added as well as a small pressure perturbation is imposed, the Hamiltonian and momentum constraints are enforced at $`t=0`$ using the method described in Sec. IV.
In Fig. 12(a), we show the evolution of the central density of the neutron star, and mass and angular momentum of the disk which are defined by
$`M_{\mathrm{disk}}{\displaystyle _{xx_{\mathrm{in}}}}\rho _{}d^3x,`$ (144)
$`J_{\mathrm{disk}}{\displaystyle _{xx_{\mathrm{in}}}}S_\phi d^3x,`$ (145)
where $`x_{\mathrm{in}}`$ denotes the initial coordinate radius of the inner edges of the disk. The figure shows that our implementation keeps the equilibrium system to be in equilibrium for more than $`20P_c`$. With the grid of size (241,241), increase of the density, which is perhaps associated with the outward transport of the angular momentum, is at most $`1\%`$ at $`t=20P_c`$. The change in the baryon rest mass and angular momentum of the disk, which is caused spuriously by the mass transfer from the central star and mass accretion to the central star due to a numerical error, is smaller than $`0.1\%`$. Also, the numerical results converge at better than second order with improving the grid resolution.
In Fig. 12(b), we also show the evolution of the ADM mass, angular momentum, and averaged violation of the Hamiltonian constraint. It is found that the ADM mass is conserved within $`1\%`$ error for $`t20P_c`$ with $`(241,241)`$ grid resolution. An outstanding feature is that the angular momentum is conserved with much better accuracy than the ADM mass. This is a feature when a HRC scheme is adopted . The averaged violation of the Hamiltonian constraint also remains to be a small magnitude for $`t20P_c`$ and converges at better than second order. All these results illustrate that our implementation can compute self-gravitating equilibrium systems accurately.
### C Winding-up of magnetic field lines in a disk around a neutron star
Next, we add magnetic fields confined only in the disk around the neutron star. For this test, we use the same system of a neutron star and a disk which is described in Sec. VI B. Similar test in a fixed background spacetime of a black hole has been performed in . Here, we perform the test in full general relativity replacing the black hole by a neutron star. The purpose of this subsection is to illustrate that our implementation can follow the growth of magnetic fields by winding-up due to differential rotation of the disk. Subsequent papers will focus on detailed scientific aspect of this issue .
Following , the $`\phi `$ component of the vector potential $`A_\phi `$ is chosen as
$`A_\phi =\{\begin{array}{cc}A(\rho \rho _0)\hfill & \mathrm{for}\rho \rho _0,\hfill \\ 0\hfill & \mathrm{for}\rho <\rho _0,\hfill \end{array}`$ (148)
where $`A`$ is a constant which determines the magnetic field strength. Then the magnetic fields are given by $`^z=x^1_xA_\phi `$ and $`^x=x^1_zA_\phi `$. This choice of $`A_\phi `$ produces poloidal field loops that coincide with isodensity contours. Here, $`\rho _0`$ is chosen as $`0.3\rho _{\mathrm{max}:\mathrm{disk}}`$ where $`\rho _{\mathrm{max}:\mathrm{disk}}`$ is the maximum density inside the disk. In the following, all the simulations are performed in axial symmetry with (301, 301) grid size and with the grid spacing of 0.165 km. The reflection symmetry with respect to the equatorial plane is assumed. We note that the boundary condition for the magnetic field is $`^x=^y=0`$ and $`_z^z=0`$ at the equatorial plane (in contrast to those for velocity fields $`v^i`$ and $`u_i`$ for which, e.g., $`_zv^x=_zv^y=0`$ and $`v^z=0`$ at the equatorial plane). Outer boundary conditions are not necessary for the magnetic field in the present simulations since the location of the outer boundary is far enough from the center that the magnetic field lines do not reach the outer boundaries. The Hamiltonian and momentum constraints are enforced at $`t=0`$ using the method described in Sec. IV. Since the magnetic field strength we choose is very weak initially, the obtained initial condition is approximately the same as that of no magnetic fields.
Simulations are performed for various values of $`A`$ which is chosen so that the magnetic pressure is initially much smaller than the gas pressure. In the following we specify the model in terms of the initial ratio of the energy of magnetic fields to the internal energy of the disk (hereafter $`R_B`$) instead of $`A`$. Here, the energy of magnetic fields and the internal energy of the disk is simply defined by
$`U_{\mathrm{mag}}{\displaystyle _{\mathrm{disk}}}b^2d^3x,`$ (149)
$`U_{\mathrm{disk}}{\displaystyle _{\mathrm{disk}}}\rho _{}\epsilon d^3x,`$ (150)
and thus, $`R_BU_{\mathrm{mag}}/U_{\mathrm{disk}}`$ at $`t=0`$. We note that the precise definition of the magnetic energy is unknown in general relativity, but the present definition is likely to give a guideline for the magnitude within an error of a factor of $`2`$–3.
In Fig. 13, we show the evolution of $`U_{\mathrm{mag}}`$ for three values of $`R_B`$. Here, the magnetic energy is plotted in units of the initial value of $`U_{\mathrm{disk}}`$ (hereafter $`U_{\mathrm{disk0}}`$) which is about $`1.8\times 10^4M_{\mathrm{disk}}`$. It is found that $`U_{\mathrm{mag}}`$ grows monotonically until the growth is saturated irrespective of the value of $`R_B`$. The growth rate is in proportional to $`R_B^{1/2}`$ in the early phase before the saturation is reached. This indicates that differential rotation winds up the magnetic field lines for amplifying the field strength . After the saturation occurs, $`U_{\mathrm{mag}}/U_{\mathrm{disk0}}`$ relaxes to $`0.02`$–0.2. These values indicate that the magnetic $`\beta `$ parameter often referred in is of order $`10`$. These relaxed values are in good agreement with previous results obtained in the simulation with a fixed background . The magnetic energy reached after the saturation depends on $`R_B`$, indicating that not only the winding-up of the field lines but also other mechanisms (which may be MRI or other instabilities associated with the magnetic fields) are likely to determine the final value.
To check that the growth of the magnetic fields occurs irrespective of grid resolution, we performed additional simulations for $`R_B=5\times 10^5`$ with grid sizes of (241, 241) and (201, 201) without changing the location of the outer boundaries. In the small panel of Fig. 13, evolution of the magnetic energy for these cases as well as for (301, 301) grid size is displayed. It is shown that the growth rate depends very weakly on the grid resolution. This confirms that our simulation can follow the winding-up of the magnetic field lines well. On the other hand, it should be mentioned that the fastest growing mode of the MRI cannot be resolved in the present computational setting since the characteristic wavelength for this mode $`2\pi v_A/\mathrm{\Omega }`$, where $`v_A`$ denotes the characteristic Alfvén speed, is approximately as small as the grid size (for $`R_B=5\times 10^5`$, $`2\pi v_A/\mathrm{\Omega }\mathrm{\Delta }x`$) in the current setting. To resolve the fastest growing mode, the grid size should be at least one tenth of the present one. Performing such a simulation of high resolution is beyond scope of this paper and an issue for the next step.
In Fig. 14, snapshots of the density profile of disks are displayed for which $`R_B=2\times 10^4`$. It shows that with the growth of magnetic fields due to winding-up of the field lines, a wind is induced to blow the matter in the outer part of the disk off. Also, the matter in the inner part of the disk gradually falls into the neutron star because of the angular momentum transport by the magnetic fields from the inner to the outer parts (see Fig. 15). After the nonlinear development of the turbulence, the disk settles down to a quasi stationary state. As explained in , this is probably due to the imposition of axial symmetry which precludes the development of the azimuthal unstable modes. Also, in the present numerical simulation, MRI which could induce turbulence is not well resolved. This may be also a reason.
In Fig. 15, we show the evolution of mass and angular momentum of the disk. It shows that after the saturation of the nonlinear growth of the magnetic fields, these quantities decrease. Decrease rates of the mass and angular momentum take maximum values soon after the growth of the magnetic pressure is saturated (e.g., at $`t5P_c`$ for $`R_B=2\times 10^4`$; cf. the dashed curves). Then, the mass and angular momentum relax to approximately constants (cf. the dashed curves). This indicates that the disk settles down to a quasistationary state. An interesting feature is that $`J_{\mathrm{disk}}`$ is approximately proportional to $`M_{\mathrm{disk}}`$ throughout the evolution. This is reasonable because the specific angular momentum $`j`$ is constant in the disk at $`t=0`$, and approximately so is the matter fallen to the neutron star as long as the magnetic pressure is much smaller than the gas pressure. However, in the case of $`R_B=2\times 10^4`$, at $`t20P_c`$ for which growth of the magnetic field has already saturated enough, $`J_{\mathrm{disk}}/M_{\mathrm{disk}}`$ slightly deviates from the initial value. This indicates that angular momentum is transported by the effect of magnetic fields.
We also performed a simulation for a toroidal magnetic field $`B^\phi `$. For $`B^\phi `$, we gave
$`B^\phi =\{\begin{array}{cc}C(\rho \rho _0)z/(z+z_0)\hfill & \mathrm{for}\rho \rho _0,\hfill \\ 0\hfill & \mathrm{for}\rho <\rho _0,\hfill \end{array}`$ (153)
where $`\rho _0=0.3\rho _{\mathrm{max}:\mathrm{disk}}`$ and $`z_0`$ is a constant much smaller than the scale hight of the disk. We note that $`B^\phi `$ has to be zero in the equatorial plane because we impose a reflection symmetry for the matter field with respect to this plane. In this simulation, magnetic energy decreases monotonically due to a small expansion of the disk induced by the magnetic pressure. In this case, no instability sets in. This is a natural consequence since the field lines are parallel to the rotational motion, and hence, they are not wound by the differential rotation. Obviously, the assumption of the axial symmetry prohibits deformation of the magnetic field lines and plays a crucial role for stabilization. If a nonaxisymmetric simulation is performed, MRI could set in .
## VII Summary and discussion
In this paper, we describe our new implementation for ideal GRMHD simulations. In this implementation, Einstein’s evolution equations are evolved by a latest version of BSSN formalism, the MHD equations by a HRC scheme, and the induction equation by a constraint transport method. We performed a number of simulations for standard test problems in relativistic MHD including special relativistic magnetized shocks, general relativistic magnetized Bondi flow in the stationary spacetime, and fully general relativistic simulation for a self-gravitating system composed of a neutron star and a disk. Our implementation yields accurate and convergent results for all these test problems. In addition, we performed simulations for a magnetized accretion disk around a neutron star in full general relativity. It is shown that magnetic fields in the differentially rotating disk are wound, and as a result, the magnetic field strength increases monotonically until a saturation is achieved. This illustrates that our implementation can be applied for investigation of growth of magnetic fields in self-gravitating systems.
In the future, we will perform a wide variety of simulations including magnetized stellar core collapse, MRI for self-gravitating neutron stars and disks, magnetic braking of differentially rotating neutron stars, and merger of binary magnetized neutron stars. Currently, we consider that the primary target is stellar core collapse of a strongly magnetized star to a black hole and a neutron star which could be a central engine of gamma-ray bursts. Recently, simulations aiming at clarifying these high energy phenomena have been performed . In such simulations, however, one neglects self-gravity and also assumes the configuration of the disks around the central compact object and magnetic fields without physical reasons. On the other hand, Newtonian MHD simulations including self-gravity consistently have recently performed in . However, stellar core collapse to a black hole and gamma-ray bursts are relativistic phenomena. For a self consistent study, it is obviously necessary to perform a general relativistic simulation from the onset of stellar core collapse throughout formation of a neutron star or a black hole with surrounding disks. Subsequent phenomena such as ejection of jets and onset of MRI of disks should be investigated using the output of the collapse simulation. In previous papers , we performed fully general relativistic simulations of stellar core collapse to formation of a neutron star and a black hole in the absence of magnetic fields. As an extension of the previous work, simulation for stellar core collapse with a strongly magnetized massive star should be a natural next target.
It is also important and interesting to clarify how MRI sets in and how long the time scale for the angular momentum transport after the onset of the MRI is in differentially rotating neutron stars. Recent numerical simulations for merger of binary neutron stars in full general relativity have clarified that if total mass of the system is smaller than a critical value, the outcome after the merger will be a hypermassive neutron star for which the self-gravity is supported by strong centrifugal force generated by rapid and differential rotation. Furthermore, the latest simulations have clarified that the hypermassive neutron star is likely to have an ellipsoidal shape with a large ellipticity , implying that it can be a strong emitter of high-frequency gravitational waves which may be detected by advanced laser interferometric gravitational wave detectors . In our estimation of amplitude of gravitational waves , we assume that there is no magnetic field in the neutron stars. However, the neutron stars in nature are magnetized, and hence, the hypermassive neutron stars should also be. If the differential rotation of the hypermassive neutron stars amplifies the seed magnetic field via winding-up of magnetic fields or MRI very rapidly, the angular momentum may be redistributed and hence the structure of the hypermassive neutron stars may be significantly changed. In , we evaluate the emission time scale of gravitational waves for the hypermassive neutron stars is typically $`50`$$`100`$ ms for the mass $`M2.4`$$`2.7M_{}`$ assuming the absence of the magnetic effects. Here, the time scale of $`50`$$`100`$ ms is an approximate dissipation time scale of angular momentum via gravitational radiation, and hence in this case, after $`50`$$`100`$ ms, the hypermassive neutron stars collapse to a black hole because the centrifugal force is weaken. Thus, it is interesting to ask if the dissipation and/or transport time scale of angular momentum by magnetic fields is shorter than $`50`$$`100`$ ms so that they can turn on before collapsing to a black hole. Rotational periods of the hypermassive neutron stars are 0.5–1 ms. Thus, if the magnetic fields grow in the dynamical time scale associated with the rotational motion via MRI, the amplitude and frequency of gravitational waves may be significantly affected. According to a theory of MRI , the wavelength of the fastest growing mode is $`10(B/10^{12}\mathrm{gauss})(\rho /10^{15}\mathrm{g}/\mathrm{cm}^3)^{1/2}(P/1\mathrm{ms})`$ cm where $`B`$, $`\rho `$, and $`P`$ denotes a typical magnetic field strength, density, and rotational period, respectively. This indicates that a turbulence composed of small eddies (for which the typical scale is much smaller than the stellar radius) will set in. Subsequently, it will contribute to a secular angular momentum transport for which the time scale is likely to be longer than the growth time scale of MRI $``$ a few ms although it is not clear if it is longer than $`100`$ ms. On the other hand, if the transport time scale is not as short as $`100`$ ms, other effects associated with magnetic fields will not affect the evolution of the hypermassive neutron stars. Indeed, Ref. indicates that the typical time scale associated with magnetic braking (winding-up of magnetic field lines) depends on the initial strength of the magnetic fields, and it is much longer than the dynamical time scale as $`100(10^{12}\mathrm{gauss}/B)`$ s. In this case, the hypermassive neutron stars can be strong emitters of gravitational waves as indicated in . As is clear from this discussion, it is important to clarify the growth time scale of magnetic fields in differentially rotating neutron stars. This is also the subject in our subsequent papers .
###### Acknowledgements.
We are grateful to Stu Shapiro for many valuable discussions and to Yuk-Tung Liu for providing solutions for Alfvén wave tests presented in Sec. V and valuable discussions. We also thank Miguel Aloy, Matt Duez, Toni Font, S. Inutsuka, A. Mizuta, Branson Stephens, and R. Takahashi for helpful discussions. Numerical computations were performed on the FACOM VPP5000 machines at the data processing center of NAOJ and on the NEC SX6 machine in the data processing center of ISAS in JAXA. This work was in part supported by Monbukagakusho Grant (Nos. 15037204, 15740142, 16029202, 17030004, and 17540232). |
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